id
stringlengths 27
33
| source
stringclasses 1
value | format
stringclasses 1
value | text
stringlengths 13
1.81M
|
|---|---|---|---|
warning/0005/hep-ph0005242.html
|
ar5iv
|
text
|
# CURRENT ISSUES IN HEAVY QUARK PRODUCTION
## Acknowledgements
This work was supported in part by NSF PHY-9722101. I would like to thank B. Harris, E. Laenen and W.L. van Neerven for comments on this report.
## References
|
warning/0005/cond-mat0005105.html
|
ar5iv
|
text
|
# Critical Exponents of the KPZ Equation via Multi-Surface Coding Numerical Simulations
## 1 Introduction
The KPZ equation , in its apparent simplicity, involves many issues that need clarification. The continuum equation (that describes the local growth of an interface profile) is
$$\frac{h}{t}=\nu \stackrel{}{}^2h+\frac{\lambda }{2}\left(\stackrel{}{}h\right)^2+\eta (\stackrel{}{r},t).$$
(1)
The real behavior of this equation is not well understood. The absence of a complete mean field theory does not help, and the fact that we have do understand the effects of a strong coupling fixed point makes very difficult the development of a perturbative renormalization approach.
A practical approach for numerical studied of the problem is to consider lattice Restricted Solid On Solid (RSOS) surfaces (see for example for a review): one considers a discretized height field on a $`D`$ dimensional lattice, and imposes the constrain that the absolute value of the distance among two neighboring surface elements can only take the values zero and one.
Old extensive numerical simulations do not clarify the situation completely. After a number of interesting Field Theoretical results , the introduction of new Renormalization Group based techniques is a potentially promising direction of development .
The main theoretical points which still deserve proper explanation are two: it is still not clear whether equation (1) has a finite upper critical dimension ($`D_>`$) or not, and the exact quantification of the related critical exponents. Regarding arguments in favor of $`D_>=4`$ we address the reader to , while arguments supporting $`D_>=\mathrm{}`$ can be found in .
In this light we introduce here a new numerical technique and present some precise numerical simulations that allow us to estimate critical exponents. Thanks to the new, precise technique we succeed to clarify some questions: we show for example that a conjecture of is not founded, and we give quantitative estimates about the behavior of sub-leading corrections.
## 2 The Multi-Surface Coding
The precise results on which this paper is based are due to a new technique we introduce to simulate RSOS surfaces, where adjacent elements of the surface can only be at a distance of $`1`$, $`0`$ or $`1`$ lattice spacings. The technique is a generalization of the so called multi-spin coding technique (well discussed by Rieger in the context of the study of disordered systems), where, by using the fact that the $`\pm 1`$ spins can be represented by Boolean logical variables, one stores $`64`$ copies of the system in one single computer word (we will assume in the following we are using a $`64`$ bit computer). The main idea is that we can simulate, basically at the cost of one single usual simulation, $`64`$ copies of the system, by rewriting the basic operations (like summing spins for computing the effective force) as boolean operations, and by exploiting the fact that when, for example, the computer is calculating a logical bit $`AND`$ it is indeed doing that $`64`$ times at once. The gain of such an approach over an usual single spin simulation is of order $`30`$: one gains a factor $`64`$ for the number of systems that updates at once, and looses a factor of order $`2`$ in computational complexity (adding the spin in the Boolean form takes some more time ).
We generalize here the Ising case to the case of a field of differences, that can have three values, since a given element of the surface can be at the same height of a given neighbor, or one step behind it, or one step ahead of it. The method is new as applied to such a system, where one does not have to compute the value of an energy, but uses the boolean operations to determine if the element can be moved without violating the geometrical constraint.
Let us thinks for simplicity about $`2D`$, where we have a $`2D`$ support where a height field (the surface height) $`h_i=h_{x,y}`$ is defined. One surface element $`h_i`$ has $`2D=4`$ first neighbors: in the RSOS model the difference $`\mathrm{\Delta }_{i,\mu }h_ih_{i\pm \widehat{\mu }}`$, where $`\mu =x,y`$, can only take the three values $`1`$, $`0`$ and $`+1`$. We store the $`DL^2`$ values of the $`\mathrm{\Delta }_{i,\mu }`$ (and, as we will discuss better, each $`\mathrm{\Delta }_{i,\mu }`$ needs two bit to be stored. This is a redundant way to store the information, since a factor two in memory could be easily saved, but it is very convenient from the point of view of computer time: as usual one trades memory for time, and avoids a cumbersome reconstruction by storing more information).
If a given element is behind its neighbor the difference is $`1`$, if it has the same height than the neighbor the difference is $`0`$, while if it is ahead of the neighbor the difference is $`+1`$. Since we have three allowed values we can represent each of these difference with two bits (that we will call respectively $`H`$, high bit, and $`L`$, low bit): we can for example code with $`00`$ the situation where the given element is behind its neighbor, with $`01`$ the situation where they have the same height, and with $`10`$ the situation where it is ahead (the value $`11`$ is forbidden).
If the element is ahead of even a single one of the $`2D`$ neighbors the move is forbidden (it would violate the RSOS constrain). It is clear this is easily implemented in our coding: one just needs to do a logical $`OR`$ of the $`2D`$ $`H`$ bits related to the site $`i`$, and if it is $`1`$ at least one of the neighbors is behind and the move is not allowed. Let us see better why. We start comparing our element $`i`$ to the first of its $`2D=4`$ neighbors (let us say the one in the positive direction $`+1`$): if $`i`$ is ahead the neighbor the relevant $`H`$ bit (that we have stored in $`\mathrm{\Delta }_{i,+1}`$) is $`1`$, and we already know we cannot move. If $`i`$ is on the contrary at the same height of the neighbor or behind it we have that, as far as this neighbor is concerned, the element $`i`$ can advance of one unit. Now we look at the second neighbor, where the same reasoning holds: if we look at the logical $`OR`$ of the two relevant high bit we will find that we cannot move if this quantity is one. Looking at all the $`2D`$ neighbor we see that the move is forbidden if the logical $`OR`$ of the $`2D`$ $`H`$ bit is $`1`$: clearly this operation, as all the other present in the core of the code, updates with a single computer cycle the $`64`$ copies of the system. After doing that, if the element $`i`$ cannot be moved there is nothing that has be done: if it gets moved we have to update the $`\mathrm{\Delta }_{i,\mu }`$ related to the site $`i`$ and to the neighbors, to describe the new situation.
The codes simulating systems in a different number of spatial dimensions (in our case from $`D=2`$ to $`D=4`$) are simple generalization one of the other (when adding a dimension one has to add checks on the new neighbors and their update).
As usual in these kind of simulations the random noise is implemented with a random choice of the site to be updated. We also implement a fifty percent probability of really updating a surface element that according to the RSOS constrain could be updated. This is very important, since starting from random independent surfaces is not enough: because of our parallel scheme the sites of our $`64`$ copies have to be updated in the same order, and such an updating algorithm with updating probability one is attractive, and the $`64`$ configurations asymptotically at large times become equal . The probability of not accepting an allowing change, that depends on each of the $`64`$ configurations, solve this problem.
We are aware of a parallel algorithm to update surfaces : it is very different in spirit from our algorithm (since it parallelizes on different sites of the lattice). We have not compared in detail the performances of the two algorithms, but we believe that on one side the algorithm of is more general, and not limited to RSOS models, but on the other side our algorithms is more regular (there are no exceptional loops), and is probably better performing for the model we study.
## 3 The Numerical Simulation
We have based this work on numerical simulations of lattices of volume $`VL^D`$. The spatial dimensionality $`D`$ is the spatial support, where a $`1`$ dimensional surface take values. We study the $`D=2`$, $`D=3`$ and $`D=4`$ cases. At a given time $`t`$ of the dynamical evolution the position of the surface can be expressed by the values $`h_i(t)`$ (that we reconstruct by the differences we store in our code, see the former section). $`i`$ represents, in lexicografic order, a $`D`$-ple labeling the spatial sites.
We consider a dynamics which generates a Restricted Solid on Solid (RSOS) growth: distances of first neighboring elements of the surface cannot be larger than one. At each trial step we move elements of the surface that are not constrained not to do so because of the RSOS restriction with probability $`\frac{1}{2}`$: we have to use a probability different from one in order to keep independent the different surfaces we simulate in the same computer word.
We consider a large number of different lattice sizes. In $`D=2`$ we take $`L`$ going from $`5`$ to $`641`$, in $`D=3`$ we consider $`L`$ going from $`5`$ to $`103`$ while in $`D=4`$ $`L`$ goes from $`5`$ to $`28`$. All our data have been averaged over $`64`$ different dynamical runs.
Because of the way we use to determine critical exponents, by trying to measure precisely the asymptotic time behavior, we use very long runs, and we always try to check that we have reached the asymptotic plateau in a clear way (see the discussion and the figures in the next section). Let $`t`$ be the time labeling sweeps of our simulation (we define a sweep as the trial update of $`V`$ random sites). In any of our simulations we run $`T`$ updating sweeps: we give in table 1 the number of full lattice sweeps for each run on different lattice sizes and number of dimensions. We measure the observable every $`1000`$ lattice sweeps: when analyzing the large time asymptotic behavior of the system we discard the first half sweeps (we call $`T_0`$ the time of our first measurement): this is a very conservative attitude, but we prefer to be safe on not having any systematic bias by paying a price of making maybe the statistical error ten percent larger of the best we could do.
We define the time dependent observables
$$\overline{h(t)}\frac{1}{V}\underset{i=1}{\overset{V}{}}\left[h_i(t)\right],$$
(2)
and
$$w_k(t)\frac{1}{V}\underset{i=1}{\overset{V}{}}\left[\left(h_i(t)\overline{h(t)}\right)^k\right],$$
(3)
that we compute for $`k=2`$, $`3`$ and $`4`$ and for different $`D`$ and $`L`$ values (when needed we will label $`w`$ with the upper script $`(L)`$, to make clear to which lattice size we are referring). We define the large time asymptotic limit of $`w_k(t)`$ as
$$w_k^{(L)}\frac{1}{TT_0+1}\underset{t=T_0}{\overset{T}{}}w_k(t),$$
(4)
for simulation on a lattice of linear size $`L`$. We always check (and this is one of the crucial points of this note, that we will discuss in better detail in the next section) that $`T_0`$ and $`T`$ are large enough to make our result unbiased in the precision of our statistical error.
## 4 Analysis
We will discuss here the analysis of our numerical data. We want to determine the critical exponents of the asymptotic behavior of the $`w_k`$ we have defined in the former section. For example we have that the asymptotic infinite time value of $`w_2`$ has a leading scaling behavior
$$w_2^{(L)}=L^{2\chi },$$
(5)
while at intermediate times (large enough for being in the scaling region but small enough not to feel the finite size of the lattice)
$$w_2(t)=t^{\frac{2\chi }{z}}.$$
(6)
The first behavior is obtained by taking large times on different lattice sizes, and by studying the time asymptotic value as a function of $`L`$. The second behavior is studied by simulating large lattices, and analyzing the behavior of the systems for times larger than one, but very smaller than the thermalization time (at the given value of the lattice size). An exact (Galilean) invariance of the KPZ equation implies that $`z+\chi =2`$. Deciding if it is better to measure accurately $`z`$ or $`\chi `$ is a practical matter.
As opposed to the choice, for example, of reference , here we have mainly based our analysis on fitting the behavior of the large time asymptotic value $`w_k^{(L)}`$, and we have only used fits to the intermediate time behavior to substantiate our results. We believe in fact that it is very difficult to determine a precise quantitative estimate of the exponent of the time scaling. The problem with the time dependent behavior of $`w_k(t)`$ is that, in order to get an unbiased value, one needs a double cutoff, both at small and at large time. At small times the behavior of $`w_k(t)`$ is not a pure power, and one has to discard small lattice, and/or to use corrections to scaling, in order to remove this effect. At large times one starts to feel the finiteness of the lattice system, and a new crossover (toward the asymptotic, constant time behavior) intervenes. In other words a careful analysis of the time exponent needs a double sliding window, moving both at small and at large time. We also find that the crossover effects at large time are very important: even on large lattices one can soon see systematic effects on the exponent estimate due to the finiteness of the lattice.
On the contrary the time asymptotic behavior only needs one cutoff, that excludes small times, where the asymptotic value has not yet been reached. This is easy, and we do it by using a logarithmic division of our data. We are in this way able to check with very high precision that we are computing an unbiased (effective, $`L`$-dependent) exponent. Again, we will also show results obtained from a direct fit of $`z`$, to show they are consistent with the $`\chi `$ values we determine.
Our main analysis is done by fitting at the same time the three sets of data :
$`w_2`$ $``$ $`A_2L^{2\chi }\left(1+B_2L^\omega \right),`$
$`w_3`$ $``$ $`A_3L^{3\chi }\left(1+B_3L^\omega \right),`$ (7)
$`w_4`$ $``$ $`A_4L^{4\chi }\left(1+B_4L^\omega \right).`$
We always compare this fit to the best fit of $`w_2`$ alone (typically by only including the large volumes and by ignoring scaling corrections) and check that things are coherent. We have also check independently that, for example, $`w_3`$ really scales as $`w_2^{\frac{3}{2}}`$: this is clearly true for our data. Without the use of all our set of data that we have defined in (7) we would not have been able to determine $`\omega `$ with a reasonable statistical precision.
With this definition the exponent $`\chi `$ is the same of $`\chi `$ of and of $`\alpha `$ of (our definition of dimensionality of the system excludes the time dimension, and is always only the dimension of the space).
The error analysis is done by using a jack-knife approach : we divide our statistical sample in $`10`$ parts all including all of data but one tenth (each part excludes a different tenth of the data), we fit ten time the behavior of, for example, (7), and compute the error on $`A_k`$, $`B_k`$, $`\chi `$ and $`\omega `$ by the fluctuations of the ten results (multiplying the error times a factor $`10`$, due to the fact that the individual parts we have formed are correlated ). In short, we present a reliable estimate of the statistical errors over the quantities we determine.
The first point we have to discuss is about the exponents of finite size corrections that appear in equation (7). Do the corrections to even and odd momenta really scale with the same exponent? This has been questioned in , and we provide here accurate evidence that $`\omega _2=\omega _3=\omega _4\omega `$ in our model.
Let us start from arguing that we are not in a situation in which the scaling exponent of the odd moments, $`\omega _{odd}`$ is smaller than the scaling exponent of the even moments, $`\omega _{even}`$ (this is the opposite scenario of the one proposed in ).
Following we define
$$R_3\frac{w_3}{w_2^{\frac{3}{2}}},R_4\frac{w_4}{w_2^2}.$$
(8)
The exponent $`\chi `$ disappears from these ratios, and, on general ground, if $`\omega _{odd}<\omega _{even}`$, one has that asymptotically for large $`L`$, ignoring sub-leading corrections,
$$R_3c_3+d_3t^{\omega _{odd}},R_4c_4+d_4t^{\omega _{even}},$$
(9)
i.e.
$$R_4c+d\left(R_3+c\right)^{\omega _{even}/\omega _{odd}},$$
(10)
and $`R_3`$ is a linear function of $`R_4`$ if and only if $`\omega _{odd}\omega _{even}`$ (in which case the two ratios asymptotically scale with the same exponent). We plot in figure 1 $`R_4`$ versus $`R_3`$, and notice that the linearity is impressive. In these data there is no sign of a discrepancy among scaling exponents of odd and even momenta, and they surely exclude that $`\omega _{odd}<\omega _{even}`$. The last point on the left in the figure is our asymptotic extrapolation, with attached the (small) statistical error (the best estimate of is $`R_3=0.27\pm 0.01`$ and $`R_4=3.15\pm 0.02`$, well compatible with our data but with a very larger error).
We use now figure 2 to also exclude the case $`\omega _{odd}>\omega _{even}`$, establishing in this way that for our model $`\omega _{odd}=\omega _{even}`$: in figure 2 we plot the effective scaling exponent obtained by using separately the data for $`w_2(L)`$, $`w_3(L)`$ and $`w_4(L)`$. The three quantities do all depend linearly, with very good approximation, on $`\frac{1}{L}`$, showing that we are having the same exponent of the scaling corrections. Again, the impressive linearity of the data implies that we are measuring precisely a single exponent of finite size corrections, and also that we are not mislead by finite size corrections.
Let us now discuss the determination of the exponent $`\chi `$.
Let us start by discussing the $`D=2`$ case. A simple analysis of $`w_2^{(L)}`$ without corrections to scaling shows that a fit to lattices of linear size from $`L=19`$ give a good value of the chi squared and $`\chi =0.393`$. A systematic analysis to the form (7) by including lattice with $`L11`$ gives our final best value of
$$\chi _{D=2}=0.393\pm 0.003,\omega _{D=2}=1.1\pm 0.3.$$
(11)
We plot the rescaled $`w_2`$ versus the rescaled time in figure 3: it is clear that the asymptotic plateau is exposed with good accuracy, and that the scaling is very good. Also the best fit to the form (7) is very good: we plot the numerical data for $`w_2`$, $`w_3`$ and $`w_4`$ versus $`L`$ and the best fit in 4.
We can compare to the rational guess of that would give here $`\chi _R=0.4`$. Indeed Kim and Kosterlitz in conjecture that $`\chi (d)=\frac{2}{d+3}`$ (that seemed to fit reasonably the numerical results available at the time). Here Lässig , by using an operator product expansion, also find $`\chi (d=2)=0.4`$. Our result is at three standard deviations from the rational guess, that is a safe distance. Still, since we are dealing with a very complex situation, with many corrections that can possibly affect the result (sub-sub-leading corrections, short time, small volume,…), we perform a further check to determine if $`\chi _R=0.4`$ is a plausible result. We fix $`\chi =0.4`$, and fit our data with now seven and not eight free parameters ($`A_2`$, $`A_3`$, $`A_4`$, $`B_2`$, $`B_3`$, $`B_4`$, and $`\omega `$). Now we get a very small value of $`\omega .28`$, and a chi squared that increases of a factor $`10`$ from our previous best fit (where it was of order one per degree of freedom).
In order to show that the fact that we have been able to exclude that the exponent takes the value $`0.4`$ is not due to the hypothesis that the exponent of the sub-leading correction is the same for all momenta we can look again at figure 2. Already from these data it is clear that the value $`0.4`$ is excluded. The more sophisticated analysis which we have presented before is crucial in obtaining a controlled extrapolation to $`L=\mathrm{}`$, keeping the statistical errors under control. We believe that this is a very strong evidence against the validity of the rational guess. We will see that for higher $`D`$ values we get an even clearer discrepancy.
In $`3D`$ the same analysis of the three cumulant allows to establish that
$$\chi _{D=3}=0.3135\pm 0.0015,\omega _{D=3}=0.98\pm 0.08.$$
(12)
We are including here sizes from $`L=11`$ up to $`L=103`$. We plot the rescaled $`w_2`$ versus the rescaled time in figure 5 and the numerical data for $`w_2`$, $`w_3`$ and $`w_4`$ versus $`L`$ and the best fit in 6. Here the rational guess would give $`\chi _R=\frac{1}{3}0.333`$, and we are sitting at more than ten standard deviations from it (Lässig gives here $`\frac{2}{7}`$ that is far from our estimate). The same check we have done in $`D=2`$ leads to a strong evidence: when fixing $`\chi =\frac{1}{3}`$ we find again a very small value $`\omega 0.2`$, and the chi squared increases of a factor large than $`20`$. Here the evidence against the validity of a rational guess is even stronger than in $`D=2`$.
In $`4D`$ we use data for $`L`$ starting from $`10`$. Again the same three cumulant analysis gives us
$$\chi _{D=4}=0.255\pm 0.003,\omega _{D=4}=0.98\pm 0.09.$$
(13)
We plot the rescaled $`w_2`$ versus the rescaled time in figure 7 and the numerical data for $`w_2`$, $`w_3`$ and $`w_4`$ versus $`L`$ and the best fit in 8.
Here $`\chi _R=\frac{2}{7}0.286`$, and, again, we are sitting a ten standard deviations away from the rational guess.
We summarize our findings in table 2 where we give all the best fit values of the parameters entering equation (7).
Two more observations are interesting. In first we find $`\omega 1`$ for all the $`D`$ values we investigate. In second the pre-factor of the scaling corrections increases with $`D`$: we find $`B_2(D=2)0.08`$, $`B_2(D=3)0.25`$, $`B_2(D=4)0.37`$.
Our recent are not incompatible with the recent ones of , but our small error bars allow us to reach precise conclusions. Also the comparison with the exponents found in is fair: we have stressed the length of our runs, in order to be able to give a clean estimate of the time asymptotic behavior, so that our result is hopefully unbiased.
## 5 Conclusions
The numerical technique we have introduced works well, and has allowed us to run very precise numerical simulations with a limited amount of computer time (a few months of Pentium II processor).
We have been able to determine critical exponents of the KPZ universality class with high accuracy. We have falsified the guess that the exponents are simple rational numbers. It is also unambiguous from our data that the upper critical dimension is larger than $`4`$ (as opposed to the claims of ).
Thanks to our precise measurements (and fitting together the first $`3`$ non trivial moments of $`h`$) we have also been able to determine the exponent $`\omega `$ of the first non-leading scaling corrections. It is interesting to notice that the estimated value of $`\omega `$ is always very close to $`1`$, independent from the dimensionality of the system.
The next interesting step, following the approach of , would be to try and implement a systematic Monte Carlo Renormalization Group: the Multi-Surface Coding technique we have discussed could be a very important ingredient of such a development.
## Acknowledgments
For the numerical simulations described here we have used the Kalix2 parallel computer (built on Pentium II chips), funded by Italian MURST 1998 COFIN. We thank Claudio Castellano and Matteo Marsili for interesting discussions, and Marcel den Nijs and Kay Wiese for a relevant correspondence.
|
warning/0005/cs0005016.html
|
ar5iv
|
text
|
# In: Proc. of ANLP–NAACL, Apr 29 – May 4, 2000, pp.ANLP-325-330 Improving Testsuites via Instrumentation
## 1 Introduction
Computational Linguistics (CL) has moved towards the marketplace: One finds programs employing CL-techniques in every software shop: Speech Recognition, Grammar and Style Checking, and even Machine Translation are available as products. While this demonstrates the applicability of the research done, it also calls for a rigorous development methodology of such CL application products.
In this paper,<sup>1</sup><sup>1</sup>1The work reported here was conducted during my time at the Institut für Maschinelle Sprachverarbeitung (IMS), Stuttgart University, Germany. I describe the adaptation of a technique from Software Engineering, namely code instrumentation, to grammar development. Instrumentation is based on the simple idea of marking any piece of code used in processing, and evaluating this usage information afterwards. The application I present here is the evaluation and improvement of grammar and testsuites; other applications are possible.
### 1.1 Software Engineering vs. Grammar Engineering
Both software and grammar development are similar processes: They result in a system transforming some input into some output, based on a functional specification (e.g., cf. \[Ciravegna et al., 1998\] for the application of a particular software design methodology to linguistic engineering). Although Grammar Engineering usually is not based on concrete specifications, research from linguistics provides an informal specification.
Software Engineering developed many methods to assess the quality of a program, ranging from static analysis of the program code to dynamic testing of the program’s behavior. Here, we adapt dynamic testing, which means running the implemented program against a set of test cases. The test cases are designed to maximize the probability of detecting errors in the program, i.e., incorrect conditions, incompatible assumptions on subsequent branches, etc. (for overviews, cf. \[Hetzel, 1988, Liggesmeyer, 1990\]).
### 1.2 Instrumentation in Grammar Engineering
How can we fruitfully apply the idea of measuring the coverage of a set of test cases to grammar development? I argue that by exploring the relation between grammar and testsuite, one can improve both of them. Even the traditional usage of testsuites to indicate grammar gaps or overgeneration can profit from a precise indication of the grammar rules used to parse the sentences (cf. Sec.4). Conversely, one may use the grammar to improve the testsuite, both in terms of its coverage (cf. Sec.3.1) and its economy (cf. Sec.3.2).
Viewed this way, testsuite writing can benefit from grammar development because both describe the syntactic constructions of a natural language. Testsuites systematically list these constructions, while grammars give generative procedures to construct them. Since there are currently many more grammars than testsuites, we may re-use the work that has gone into the grammars for the improvement of testsuites.
The work reported here is situated in a large cooperative project aiming at the development of large-coverage grammars for three languages. The grammars have been developed over years by different people, which makes the existence of tools for navigation, testing, and documentation mandatory. Although the sample rules given below are in the format of LFG, nothing of the methodology relies on the choice of linguistic or computational paradigm.
## 2 Grammar Instrumentation
Measures from Software Engineering cannot be simply transferred to Grammar Engineering, because the structure of programs is different from that of unification grammars. Nevertheless, the *structure* of a grammar allows the derivation of suitable measures, similar to the structure of programs; this is discussed in Sec.2.1. The actual instrumentation of the grammar depends on the formalism used, and is discussed in Sec.2.2.
### 2.1 Coverage Criteria
Consider the LFG grammar rule in Fig. 1.<sup>2</sup><sup>2</sup>2 Notation: ?/*/+ represent optionality/iteration including/excluding zero occurrences on categories. Annotations to a category specify equality (=) or set membership ($``$) of feature values, or non-existence of features ($`\neg `$); they are terminated by a semicolon (;). Disjunctions are given in braces ({...$``$...}). $``$ ($``$) are metavariables representing the feature structure corresponding to the mother (daughter) of the rule. Comments are enclosed in quotation marks ("..."). Cf. \[Kaplan and Bresnan, 1982\] for an introduction to LFG notation. On first view, one could require of a testsuite that each such rule is exercised at least once. Further thought will indicate that there are hidden alternatives, namely the optionality of the NP and the PP. The rule can only be said to be thoroughly tested if test cases exist which test both presence and absence of optional constituents (requiring 4 test cases for this rule).
In addition to context-free rules, unification grammars contain equations of various sorts, as illustrated in Fig.1. Since these annotations may also contain disjunctions, a testsuite with complete rule coverage is not guaranteed to exercise all equation alternatives. The phrase-structure-based criterion defined above must be refined to cover all equation alternatives in the rule (requiring two test cases for the PP annotation). Even if we assume that (as, e.g., in LFG) there is at least one equation associated with each constituent, equation coverage does not subsume rule coverage: Optional constituents introduce a rule disjunct (without the constituent) that is not characterizable by an equation. A measure might thus be defined as follows:
The disjunct coverage of a testsuite is the quotient
$$T_{\text{dis}}=\frac{\text{number of disjuncts tested}}{\text{number of disjuncts in grammar}}$$
where a disjunct is either a phrase-structure alternative, or an annotation alternative. Optional constituents (and equations, if the formalism allows them) have to be treated as a disjunction of the constituent and an empty category (cf. the instrumented rule in Fig.2 for an example).
Instead of considering disjuncts in isolation, one might take their interaction into account. The most complete test criterion, doing this to the fullest extent possible, can be defined as follows:
The interaction coverage of a testsuite is the quotient
$$T_{\text{inter}}=\frac{\text{number of disjunct combinations tested}}{\text{number of legal disjunct combinations}}$$
There are methodological problems in this criterion, however. First, the set of legal combinations may not be easily definable, due to far-reaching dependencies between disjuncts in different rules, and second, recursion leads to infinitely many legal disjunct combinations as soon as we take the number of usages of a disjunct into account. Requiring complete interaction coverage is infeasible in practice, similar to the path coverage criterion in Software Engineering.
We will say that an analysis (and the sentence receiving this analysis) *relies on* a grammar disjunct if this disjunct was used in constructing the analysis.
### 2.2 Instrumentation
Basically, grammar instrumentation is identical to program instrumentation: For each disjunct in a given source grammar, we add grammar code that will identify this disjunct in the solution produced, iff that disjunct has been used in constructing the solution.
Assuming a unique numbering of disjuncts, an annotation of the form DISJUNCT-nn = + can be used for marking. To determine whether a certain disjunct was used in constructing a solution, one only needs to check whether the associated feature occurs (at some level of embedding) in the solution. Alternatively, if set-valued features are available, one can use a set-valued feature DISJUNCTS to collect atomic symbols representing one disjunct each: DISJUNCT-nn $``$ DISJUNCTS.
One restriction is imposed by using the unification formalism, though: One occurrence of the mark cannot be distinguished from two occurrences, since the second application of the equation introduces no new information. The markers merely unify, and there is no way of counting.
Therefore, we have used a special feature of our grammar development environment: Following the LFG spirit of different representation levels associated with each solution (so-called projections), it provides for a multiset of symbols associated with the complete solution, where structural embedding plays no role (so-called optimality projection; see \[Frank et al., 1998\]). In this way, from the root node of each solution the set of all disjuncts used can be collected, together with a usage count.
Fig. 2 shows the rule from Fig.1 with such an instrumentation; equations of the form DISJUNCT-nn$`o`$ express membership of the disjunct-specific atom DISJUNCT-nn in the sentence’s multiset of disjunct markers.
### 2.3 Processing Tools
Tool support is mandatory for a scenario such as instrumentation: Nobody will manually add equations such as those in Fig. 2 to several hundred rules. Based on the format of the grammar rules, an algorithm instrumenting a grammar can be written down easily.
Given a grammar and a testsuite or corpus to compare, first an instrumented grammar must be constructed using such an algorithm. This instrumented grammar is then used to parse the testsuite, yielding a set of solutions associated with information about usage of grammar disjuncts. Up to this point, the process is completely automatic. The following two sections discuss two possibilities to evaluate this information.
## 3 Quality of Testsuites
This section addresses the aspects of completeness (“does the testsuite exercise all disjuncts in the grammar?”) and economy of a testsuite (“is it minimal?”).
Complementing other work on testsuite construction (cf. Sec.5), we will assume that a grammar is already available, and that a testsuite has to be constructed or extended. While one may argue that grammar and testsuite should be developed in parallel, such that the coding of a new grammar disjunct is accompanied by the addition of suitable test cases, and vice versa, this is seldom the case. Apart from the existence of grammars which lack a testsuite, and to which this procedure could be usefully applied, there is the more principled obstacle of the evolution of the grammar, leading to states where previously necessary rules silently loose their usefulness, because their function is taken over by some other rules, structured differently. This is detectable by instrumentation, as discussed in Sec.3.1.
On the other hand, once there is a testsuite, you want to use it in the most economic way, avoiding redundant tests. Sec.3.2 shows that there are different levels of redundancy in a testsuite, dependent on the specific grammar used. Reduction of this redundancy can speed up the test activity, and give a clearer picture of the grammar’s performance.
### 3.1 Testsuite Completeness
If the disjunct coverage of a testsuite is 1 for some grammar, the testsuite is *complete* w.r.t. this grammar. Such a testsuite can reliably be used to monitor changes in the grammar: Any reduction in the grammar’s coverage will show up in the failure of some test case (for negative test cases, cf. Sec.4).
If there is no complete testsuite, one can – via instrumentation – identify disjuncts in the grammar for which no test case exists. There might be either (i) appropriate, but untested, disjuncts calling for the addition of a test case, or (ii) inappropriate disjuncts, for which one cannot construct a grammatical test case relying on them (e.g., left-overs from rearranging the grammar). Grammar instrumentation singles out all untested disjuncts automatically, but cases (i) and (ii) have to be distinguished manually.
Checking completeness of our local testsuite of 1787 items, we found that only 1456 out of 3730 grammar disjuncts in our German grammar were tested, yielding $`T_{dis}=0.39`$ (the TSNLP testsuite containing 1093 items tests only 1081 disjuncts, yielding $`T_{dis}=0.28`$).<sup>3</sup><sup>3</sup>3There are, of course, unparsed but grammatical test cases in both testsuites, which have not been taken into account in these figures. This explains the difference to the overall number of 1582 items in the German TSNLP testsuite. Fig.3 shows an example of a gap in our testsuite (there are no examples of circumpositions), while Fig.4 shows an inapproppriate disjunct thus discovered (the category ADVadj has been eliminated in the lexicon, but not in all rules). Another error class is illustrated by Fig.5, which shows a rule that can never be used due to an LFG coherence violation; the grammar is inconsistent here.<sup>4</sup><sup>4</sup>4Test cases using a free dative pronoun may be in the testsuite, but receive no analysis since the grammatical function FREEDAT is not defined as such in the configuration section.
### 3.2 Testsuite Economy
Besides being complete, a testsuite must be economical, i.e., contain as few items as possible without sacrificing its diagnostic capabilities. Instrumentation can identify redundant test cases. Three criteria can be applied in determining whether a test case is redundant:
There is a set of other test cases which jointly rely on all disjunct on which the test case under consideration relies.
There is a single test case which relies on exactly the same combination(s) of disjuncts.
There is a single test case which is equivalent to and, additionally, relies on the disjuncts exactly as often as, the test case under consideration.
For all criteria, lexical and structural ambiguities must be taken into account. Fig.6 shows some equivalent test cases derived from our testsuite: Example 1 illustrates the distinction between equivalence and strict equivalence; the test cases contain different numbers of attributive adjectives, but are nevertheless considered equivalent. Example 2 shows that our grammar does not make any distinction between adverbial usage and secondary (subject or object) predication. Example 3 shows test cases which should not be considered equivalent, and is discussed below.
The reduction we achieved in size and processing time is shown in Table 1, which contains measurements for a test run containing only the parseable test cases, one without equivalent test cases (for every set of equivalent test cases, one was arbitrarily selected), and one without similar test cases. The last was constructed using a simple heuristic: Starting with the sentence relying on the most disjuncts, working towards sentences relying on fewer disjuncts, a sentence was selected only if it relied on a disjunct on which no previously selected sentence relied. Assuming that a disjunct working correctly once will work correctly more than once, we did not consider strict equivalence.
We envisage the following use of this redundancy detection: There clearly are linguistic reasons to distinguish all test cases in example 2, so they cannot simply be deleted from the testsuite. Rather, their equivalence indicates that the grammar is not yet perfect (or never will be, if it remains purely syntactic). Such equivalences could be interpreted as a reminder which linguistic distinctions need to be incorporated into the grammar. Thus, this level of redundancy may drive your grammar development agenda. The level of equivalence can be taken as a limited interaction test: These test cases represent one complete selection of grammar disjuncts, and (given the grammar) there is nothing we can gain by checking a test case if an equivalent one was tested. Thus, this level of redundancy may be used for ensuring the quality of grammar changes prior to their incorporation into the production version of the grammar. The level of similarity contains much less test cases, and does not test any (systematic) interaction between disjuncts. Thus, it may be used during development as a quick rule-of-thumb procedure detecting serious errors only.
Coming back to example 3 in Fig.6, building equivalence classes also helps in detecting grammar errors: If, according to the grammar, two cases are equivalent which actually aren’t, the grammar is incorrect. Example 3 shows two test cases which are syntactically different in that the first contains the adverbial *oft*, while the other doesn’t. The reason why they are equivalent is an incorrect rule that assigns an incorrect reading to the second test case, where the infinitival particle “zu” functions as an adverbial.
## 4 Negative Test Cases
To control overgeneration, appropriately marked ungrammatical sentences are important in every testsuite. Instrumentation as proposed here only looks at successful parses, but can still be applied in this context: If an ungrammatical test case receives an analysis, instrumentation informs us about the disjuncts used in the incorrect analysis. One (or more) of these disjuncts must be incorrect, or the sentence would not have received a solution. We exploit this information by accumulation across the entire test suite, looking for disjuncts that appear in unusually high proportion in parseable ungrammatical test cases.
In this manner, six grammar disjuncts are singled out by the parseable ungrammatical test cases in the TSNLP testsuite. The most prominent disjunct appears in 26 sentences (listed in Fig.7), of which group 1 is really grammatical and the rest fall into two groups: A partial VP with object NP, interpreted as an imperative sentence (group 2), and a weird interaction with the tokenizer incorrectly handling capitalization (group 3).
Far from being conclusive, the similarity of these sentences derived from a suspicious grammar disjunct, and the clear relation of the sentences to only two exactly specifiable grammar errors make it plausible that this approach is very promising in reducing overgeneration.
## 5 Other Approaches to Testsuite Construction
Although there are a number of efforts to construct reusable large-coverage testsuites, none has to my knowledge explored how existing grammars could be used for this purpose.
Starting with \[Flickinger et al., 1987\], testsuites have been drawn up from a linguistic viewpoint, *“informed by \[the\] study of linguistics and \[reflecting\] the grammatical issues that linguists have concerned themselves with”* \[Flickinger et al., 1987, , p.4\]. Although the question is not explicitly addressed in \[Balkan et al., 1994\], all the testsuites reviewed there also seem to follow the same methodology. The TSNLP project \[Lehmann and Oepen, 1996\] and its successor DiET \[Netter et al., 1998\], which built large multilingual testsuites, likewise fall into this category.
The use of corpora (with various levels of annotation) has been studied, but even here the recommendations are that much manual work is required to turn corpus examples into test cases (e.g., \[Balkan and Fouvry, 1995\]). The reason given is that corpus sentences neither contain linguistic phenomena in isolation, nor do they contain systematic variation. Corpora thus are used only as an inspiration.
\[Oepen and Flickinger, 1998\] stress the interdependence between application and testsuite, but don’t comment on the relation between grammar and testsuite.
## 6 Conclusion
The approach presented tries to make available the linguistic knowledge that went into the grammar for development of testsuites. Grammar development and testsuite compilation are seen as complementary and interacting processes, not as isolated modules. We have seen that even large testsuites cover only a fraction of existing large-coverage grammars, and presented evidence that there is a considerable amount of redundancy within existing testsuites.
To empirically validate that the procedures outlined above improve grammar and testsuite, careful grammar development is required. Based on the information derived from parsing with instrumented grammars, the changes and their effects need to be evaluated. In addition to this empirical work, instrumentation can be applied to other areas in Grammar Engineering, e.g., to detect sources of spurious ambiguities, to select sample sentences relying on a disjunct for documentation, or to assist in the construction of additional test cases. Methodological work is also required for the definition of a practical and intuitive criterion to measure limited interaction coverage.
Each existing grammar development environment undoubtely offers at least some basic tools for comparing the grammar’s coverage with a testsuite. Regrettably, these tools are seldomly presented publicly (which accounts for the short list of such references). It is my belief that the thorough discussion of such infrastructure items (tools and methods) is of more immediate importance to the quality of the lingware than the discussion of open linguistic problems.
|
warning/0005/hep-th0005119.html
|
ar5iv
|
text
|
# 1 Introduction
## 1 Introduction
The simplest theory where saturated strings exist in the weak coupling regime is supersymmetric electrodynamics with the Fayet-Iliopoulos term . Topologically stable solutions in this model and its modifications were considered more than once in the past . If we consider the $`𝒩=2`$ supersymmetric Yang-Mills theory softly broken near the monopole or dyon singularities, the Abrikosov strings develop . They were discussed in the literature previously. We can see them in the effective Lagrangians near the singularities , where the superpotential for the monoploes (or dyons) can be written as
$$𝒲=\mu u(a_D)+\sqrt{2}\stackrel{~}{M}a_DM$$
(1)
where $`M`$ is the monopole field. Minimization of the potential yields the monopole condensation, as a result, the standard Abrikosov-Nielsen-Olesen string appears with the tension proportional to the mass of the $`𝒩=1`$ chiral field $`\mu `$.
We can expand the $`u(a_D)`$ term in Eq. (1) as
$$\mu u(a_D)=\mu a_D+\eta a_D^2+\mathrm{},$$
where the zero-th order approximation corresponds to the linear term in $`a_D`$ which preserves $`𝒩=2`$ supersymmetry . Our task is to find the correction due to the quadratic term in $`a_D`$. In the Sec. 2 we investigate the supersymmetric electrodynamics (SQED) with only linear term in $`a_D`$. The Abrikosov strings in this case are found to be $`1/2`$-BPS saturated. In the Sec. 3, we consider the additional quadratic term in the superpotential which destroys the BPS property. And the correction in the string tension in this case is calculated numerically to the first order of the perturbation parameter.
## 2 $`\frac{1}{2}`$-BPS Saturated Abrikosov Strings
Consider $`𝒩=2`$ supersymmetric electrodynamics . The “photon” $`A_\mu `$ is accompanied by its $`𝒩=2`$ superpartners(photinos)–two neutral Weyl spinors $`\lambda `$ and $`\psi `$, and a complex neutral scalar $`a`$. They form an irreducible $`𝒩=2`$ representation that can be decomposed as a sum of two $`𝒩=1`$ representations: $`a`$ and $`\psi `$ are in a chiral representation, $`\mathrm{\Phi }`$, while $`A_\mu `$ and $`\lambda `$ are in a vector representation, $`W_\alpha `$. Matter sector consists of two $`𝒩=1`$ chiral multiplets $`M`$ and $`\stackrel{~}{M}`$ with opposite electric charge. The renormalizable $`𝒩=2`$ invariant Lagrangian is described in an $`𝒩=1`$ language by canonical kinetic terms and minimal gauge couplings for all the fields as well as a superpotential
$$𝒲=\sqrt{2}\mathrm{\Phi }M\stackrel{~}{M}\mu \mathrm{\Phi },$$
(2)
here we replace the $`a_D`$ in Eq. (1) by $`\mathrm{\Phi }`$ for simplicity.
The Lagrangian in component fields is given by
$``$ $`=`$ $`{\displaystyle \frac{1}{4}}F_{\mu \nu }F^{\mu \nu }+^\mu \overline{a}_\mu a+i\lambda \sigma ^\mu _\mu \overline{\lambda }+i\psi \sigma ^\mu _\mu \overline{\psi }`$ (3)
$`+D_\mu \overline{M}D^\mu M+D_\mu \overline{\stackrel{~}{M}}D^\mu \stackrel{~}{M}+i\psi _M\sigma ^\mu D_\mu \overline{\psi }_M+i\psi _{\stackrel{~}{M}}\sigma ^\mu D_\mu \overline{\psi }_{\stackrel{~}{M}}`$
$`+(\sqrt{2}i\psi _M\lambda \overline{M}+h.c.)+(\sqrt{2}i\psi _{\stackrel{~}{M}}\lambda \overline{\stackrel{~}{M}}+h.c.)`$
$`+(\sqrt{2}aMF_{\stackrel{~}{M}}\sqrt{2}a\psi _M\psi _{\stackrel{~}{M}}+\sqrt{2}aF_M\stackrel{~}{M}\sqrt{2}\psi \psi _{\stackrel{~}{M}}M`$
$`\sqrt{2}\psi \psi _M\stackrel{~}{M}+\sqrt{2}FM\stackrel{~}{M}\mu F+h.c.)`$
$`+F\overline{F}+{\displaystyle \frac{1}{2}}D^2+D\left(\overline{M}M\overline{\stackrel{~}{M}}\stackrel{~}{M}\right)+F_M\overline{F}_M+F_{\stackrel{~}{M}}\overline{F}_{\stackrel{~}{M}}.`$
Here we set the coupling constant $`e=1`$ for the convenience, which is not important in our later results. We stick to this convention in what follows. $`F,F_M,F_{\stackrel{~}{M}}`$ and $`D`$ are the auxiliary fields which are given by
$`F`$ $`=`$ $`\mu \sqrt{2}\overline{M}\overline{\stackrel{~}{M}},`$
$`F_M`$ $`=`$ $`\sqrt{2}\overline{a}\overline{\stackrel{~}{M}},`$
$`F_{\stackrel{~}{M}}`$ $`=`$ $`\sqrt{2}\overline{a}\overline{M},`$
$`D`$ $`=`$ $`\left(\overline{M}M\overline{\stackrel{~}{M}}\stackrel{~}{M}\right).`$ (4)
Here $`M,\stackrel{~}{M}`$ are the lowest components of the corresponding superfields, respectively, with the electric charges $`\pm 1`$, e.g.
$$D_\mu M=_\mu MiA_\mu M,D_\mu \stackrel{~}{M}=_\mu \stackrel{~}{M}+iA_\mu \stackrel{~}{M}.$$
The scalar potential is minimized at
$$F=F_M=F_{\stackrel{~}{M}}=0,D=0,$$
(5)
which occurs when
$$a=0,M\stackrel{~}{M}=\frac{\mu }{\sqrt{2}},\text{and}\left|M\right|=|\stackrel{~}{M}|.$$
Then the supersymmetric transformations which preserve $`𝒩=2`$ supersymmetry are given by
$`\delta a`$ $`=`$ $`\sqrt{2}\xi \lambda +\sqrt{2}\epsilon \psi ,`$
$`\delta \psi `$ $`=`$ $`i\xi Di\xi \sigma ^{\mu \nu }F_{\mu \nu }+\sqrt{2}\epsilon F+\sqrt{2}i\sigma ^\mu \overline{\epsilon }_\mu a,`$
$`\delta F`$ $`=`$ $`\sqrt{2}i\xi \sigma ^\mu _\mu \overline{\lambda }+\sqrt{2}i_\mu \psi \sigma ^\mu \overline{\epsilon },`$
$`\delta A_\mu `$ $`=`$ $`i\xi \sigma _\mu \overline{\psi }+i\psi \sigma _\mu \overline{\xi }i\epsilon \sigma _\mu \overline{\lambda }+i\lambda \sigma _\mu \overline{\epsilon },`$
$`\delta \lambda `$ $`=`$ $`\sqrt{2}\xi \overline{F}+\sqrt{2}i\sigma ^\mu \overline{\xi }_\mu a+iD\epsilon +i\sigma ^{\mu \nu }F_{\mu \nu }\epsilon ,`$
$`\delta D`$ $`=`$ $`_\mu \psi \sigma ^\mu \overline{\xi }+\xi \sigma ^\mu _\mu \overline{\psi }+_\mu \lambda \sigma ^\mu \overline{\epsilon }+\epsilon \sigma ^\mu _\mu \overline{\lambda },`$
$`\delta M`$ $`=`$ $`\sqrt{2}\epsilon \psi _M+\sqrt{2}\overline{\xi }\overline{\psi }_{\stackrel{~}{M}},`$ (6)
$`\delta \stackrel{~}{M}`$ $`=`$ $`\sqrt{2}\overline{\xi }\overline{\psi }_M+\sqrt{2}\epsilon \psi _{\stackrel{~}{M}},`$
$`\delta \psi _M`$ $`=`$ $`\sqrt{2}\epsilon F_M+\sqrt{2}i\sigma ^\mu \overline{\epsilon }D_\mu M+\sqrt{2}\sigma ^\mu \overline{\xi }D_\mu \overline{\stackrel{~}{M}}2i\overline{a}\xi M,`$
$`\delta \psi _{\stackrel{~}{M}}`$ $`=`$ $`\sqrt{2}\sigma ^\mu \overline{\xi }D_\mu \overline{M}+\sqrt{2}i\sigma ^\mu \overline{\epsilon }D_\mu \stackrel{~}{M}+\sqrt{2}\epsilon F_{\stackrel{~}{M}}2i\overline{a}\xi \stackrel{~}{M},`$
$`\delta F_M`$ $`=`$ $`\sqrt{2}iD_\mu \psi _M\sigma ^\mu \overline{\epsilon }2\overline{\xi }\overline{\lambda }\overline{\stackrel{~}{M}}2i\overline{a}\xi \psi _M,`$
$`\delta F_{\stackrel{~}{M}}`$ $`=`$ $`\sqrt{2}iD_\mu \psi _{\stackrel{~}{M}}\sigma ^\mu \overline{\epsilon }2\overline{\xi }\overline{\lambda }\overline{M}2i\overline{a}\xi \psi _{\stackrel{~}{M}},`$
where the spinorial indices are suppressed.
Without loss of generality, we can assume the Abrikosov string axis lies along the $`z`$ axis, while the string profile depends only on $`x,y`$. Then we obtain the saturation equations by requiring the fermionic fields transformations in Eq. (6) to vanish as follows
$`F_{12}`$ $`=`$ $`\sqrt{2}\left(\sqrt{2}\overline{M}\overline{\stackrel{~}{M}}\mu \right),`$
$`\left(D_1+iD_2\right)M`$ $`=`$ $`0,`$ (7)
$`\left(D_1iD_2\right)\stackrel{~}{M}`$ $`=`$ $`0,`$
with the constraint determining the parameter of the residual supersymmetry,
$$i\tau _3\xi =\epsilon ,$$
(8)
which reduces the number of supersymmetries from eight to four. The Abrikosov-Nielsen-Olesen string is $`1/2`$-BPS saturated.
The Ansatz which goes through Eq. (7) is
$`M`$ $`=`$ $`\left({\displaystyle \frac{\mu }{\sqrt{2}}}\right)^{\frac{1}{2}}e^{i\varphi }f(r),`$
$`\stackrel{~}{M}`$ $`=`$ $`\left({\displaystyle \frac{\mu }{\sqrt{2}}}\right)^{\frac{1}{2}}e^{i\varphi }f(r),`$ (9)
$`A_\varphi `$ $`=`$ $`2{\displaystyle \frac{g(r)}{r}},`$
with the boundary conditions
$`f(0)`$ $`=`$ $`g(0)=0,`$
$`\underset{r\mathrm{}}{lim}f(r)`$ $`=`$ $`1,`$
$`\underset{r\mathrm{}}{lim}g(r)`$ $`=`$ $`{\displaystyle \frac{1}{2}}.`$
The profile functions, $`f(r)`$ and $`g(r)`$ satisfy the first-order differential equations
$`f^{}`$ $`=`$ $`{\displaystyle \frac{f}{r}}(1+2g),`$
$`g^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}r(1f^2),`$ (10)
where the prime denotes differentiation over $`r`$.
One can calculate the string tension as follows
$`𝒯`$ $`=`$ $`{\displaystyle }d^2x\{{\displaystyle \frac{1}{2}}F_{12}^2+D_1\overline{M}D_1M+D_2\overline{M}D_2M+D_1\overline{\stackrel{~}{M}}D_1\stackrel{~}{M}+D_2\overline{\stackrel{~}{M}}D_2\stackrel{~}{M}`$
$`+(\mu \sqrt{2}\overline{M}\overline{\stackrel{~}{M}})(\mu \sqrt{2}M\stackrel{~}{M})+{\displaystyle \frac{1}{2}}(\overline{M}M\overline{\stackrel{~}{M}}\stackrel{~}{M})^2\}`$
$`=`$ $`{\displaystyle }d^2x\left\{\right|{\displaystyle \frac{1}{\sqrt{2}}}F_{12}(\sqrt{2}\overline{M}\overline{\stackrel{~}{M}}\mu )|^2+(D_1+iD_2)\overline{M}(D_1+iD_2)M`$
$`+\left(D_1iD_2\right)\overline{\stackrel{~}{M}}\left(D_1iD_2\right)\stackrel{~}{M}+{\displaystyle \frac{1}{2}}\left(\overline{M}M\overline{\stackrel{~}{M}}\stackrel{~}{M}\right)^2\sqrt{2}\mu F_{12}`$
$`i[_1\left(\overline{M}D_2M\right)_2\left(\overline{M}D_1M\right)]+i[_1\left(\overline{\stackrel{~}{M}}D_2\stackrel{~}{M}\right)_2\left(\overline{\stackrel{~}{M}}D_1\stackrel{~}{M}\right)]\}.`$
Applying Eq. (7) and neglecting the total derivative terms, we get
$$𝒯=\sqrt{2}\mu d^2xF_{12}.$$
(12)
## 3 Small Perturbation in Superpotential
We can add a small perturbation in the superpotential Eq. (2)
$$𝒲=\sqrt{2}\mathrm{\Phi }M\stackrel{~}{M}\mu \mathrm{\Phi }+\eta \mathrm{\Phi }^2,$$
(13)
where $`\eta `$ is a real small perturbation parameter.
Then we can go over the analysis in Sec. 2 in the similar way. But the small perturbation will break $`𝒩=2`$ supersymmetry and the resultant Abrikosov string is no longer BPS saturated which can be seen clearly in the string tension
$`𝒯`$ $`=`$ $`{\displaystyle }d^2x\{{\displaystyle \frac{1}{2}}F_{12}^2+D_1\overline{M}D_1M+D_2\overline{M}D_2M+D_1\overline{\stackrel{~}{M}}D_1\stackrel{~}{M}+D_2\overline{\stackrel{~}{M}}D_2\stackrel{~}{M}`$ (14)
$`+\left(\mu \sqrt{2}\overline{M}\overline{\stackrel{~}{M}}2\eta \overline{a}\right)\left(\mu \sqrt{2}M\stackrel{~}{M}2\eta a\right)+{\displaystyle \frac{1}{2}}\left(\overline{M}M\overline{\stackrel{~}{M}}\stackrel{~}{M}\right)^2`$
$`+_1\overline{a}_1a+_2\overline{a}_2a+2a\overline{a}M\overline{M}+2a\overline{a}\stackrel{~}{M}\overline{\stackrel{~}{M}}\}`$
$`=`$ $`{\displaystyle }d^2x\left\{\right|{\displaystyle \frac{1}{\sqrt{2}}}F_{12}(\sqrt{2}\overline{M}\overline{\stackrel{~}{M}}\mu )|^2+(D_1+iD_2)\overline{M}(D_1+iD_2)M`$
$`+\left(D_1iD_2\right)\overline{\stackrel{~}{M}}\left(D_1iD_2\right)\stackrel{~}{M}+{\displaystyle \frac{1}{2}}\left(\overline{M}M\overline{\stackrel{~}{M}}\stackrel{~}{M}\right)^2\sqrt{2}\mu F_{12}`$
$`i\left[_1\left(\overline{M}D_2M\right)_2\left(\overline{M}D_1M\right)\right]+i\left[_1\left(\overline{\stackrel{~}{M}}D_2\stackrel{~}{M}\right)_2\left(\overline{\stackrel{~}{M}}D_1\stackrel{~}{M}\right)\right]`$
$`+_1\left(\overline{a}_1a\right)+_2\left(\overline{a}_2a\right)2\eta a(\mu \sqrt{2}\overline{M}\overline{\stackrel{~}{M}})\},`$
where we have used the equation of motion for field $`a`$
$$_1^2a_2^2a=2\eta \left(\mu \sqrt{2}M\stackrel{~}{M}2\eta a\right)2aM\overline{M}2a\stackrel{~}{M}\overline{\stackrel{~}{M}}.$$
(15)
To the first order of $`\eta `$, we can still use the Ansatz (9) for the fields $`M,\stackrel{~}{M}`$ and $`A_\mu `$. After applying Ansatz (9), Eq. (15) becomes to the first order in $`\eta `$
$$_1^2a+_2^2a=2\eta \mu \left(1f^2\right)+2\sqrt{2}\mu af^2.$$
(16)
Then the string tension turns out to be
$$𝒯=d^2x\left\{\left(\sqrt{2}\mu F_{12}\right)2\eta a\left(\mu \sqrt{2}\overline{M}\overline{\stackrel{~}{M}}\right)\right\}.$$
(17)
Then from Eq. (12),(17), we can find the difference of the string tensions between non-BPS and BPS saturated situation to be
$`\mathrm{\Delta }𝒯`$ $`=`$ $`{\displaystyle d^2x\left\{2\eta a\left(\mu \sqrt{2}\overline{M}\overline{\stackrel{~}{M}}\right)\right\}}`$ (18)
$`=`$ $`2\eta \mu {\displaystyle d^2xa\left(1f^2\right)},`$
where we have used the Ansatz (9).
To see Eq. (18) more clearly, one can switch to dimensionless quantities
$`x`$ $``$ $`{\displaystyle \frac{1}{(\sqrt{2}\mu )^{\frac{1}{2}}}}x,`$
$`y`$ $``$ $`{\displaystyle \frac{1}{(\sqrt{2}\mu )^{\frac{1}{2}}}}y,`$
$`a`$ $``$ $`\eta a.`$ (19)
Then we get
$$\mathrm{\Delta }𝒯=\sqrt{2}\eta ^2a\left(1f^2\right)d^2x,$$
(20)
where $`a,f`$, and $`x`$ here are dimensionless.
We can solve Eq. (10),(16), and calculate $`\mathrm{\Delta }𝒯`$ in Eq. (20). The result is
$$\mathrm{\Delta }𝒯=2\sqrt{2}\pi \eta ^2\mathrm{\hspace{0.17em}0.68}<\mathrm{\hspace{0.17em}0}.$$
(21)
## 4 Conclusions
We investigated the Abrikosov-Nielsen-Olesen string solution in $`𝒩=2`$ supersymmetric electrodynamics with some $`𝒩=2`$-preserving superpotential. The string solution is due to the superpotential rather than due to the Fayet-Iliopoulos term. The Abrikosov string was found to be $`1/2`$-BPS saturated which follows directly from the $`𝒩=2`$ supersymmetric transformations. After the $`𝒩=2`$ supersymmetry is broken to $`𝒩=1`$ by the perturbation in the superpotential, the Abrikosov string is no longer BPS saturated. And the string tension in this case was found to be less than that of the BPS case.
## 5 Acknowledgments
I am grateful to M. Shifman for suggesting the problem to me and for numerous useful discussions. I would also like to thank A. Vainshtein and A. Yung for useful discussions. This work was supported in part by DOE under the grant number DE-FG02-94ER408.
|
warning/0005/math0005166.html
|
ar5iv
|
text
|
# Generalization of Wigner’s Unitary-Antiunitary Theorem for Indefinite Inner Product Spaces
## Abstract.
We present a generalization of Wigner’s unitary-antiunitary theorem for pairs of ray transformations. As a particular case, we get a new Wigner-type theorem for non-Hermitian indefinite inner product spaces.
The classical Wigner unitary-antiunitary theorem plays a fundamental role in the foundations of quantum mechanics and it also has deep connections with the theory of projective spaces. It states that every ray transformation (see below) on a Hilbert space which preserves the transition probabilities can be lifted to a (linear) unitary or a (conjugate-linear) antiunitary operator on $`H`$ (see ). So, Wigner’s result concerns definite inner product spaces. On the other hand, it has become quite clear by now that the indefinite inner product spaces might be even more useful for the discussion of several physical problems. For example, this is the case in relation to the divergence problem in quantum field theory, or when one wants to preserve some basic properties of the field like relativistic covariance and locality (see the introduction of ). This raises the need to study Wigner’s theorem in the ”indefinite” setting as well. Previous results in this direction were presented in . The aim of this paper is to contribute to this study by giving a very general Wigner-type theorem which involves not one but two ray transformations and then apply it to get a generalization of Wigner’s theorem for indefinite inner product spaces. The main difference which distiguishes our result from the previous ones is that we do not assume even that the indefinite inner product under consideration is Hermitian. What allows us to reach this result is that we refine our algebraic approach to Wigner’s theorem which has already been proved to be fruitful in our recent papers . The main feature of this approach is that instead of manipulating in the underlying space, we push the problem to an operator algebra over our space and apply some classical results from pure ring theory. Hence, our method is completely different from those used previously in the papers dealing with Wigner’s theorem in indefinite inner product spaces.
Let us fix the definitions and notation that we shall use throughout. In what follows, let $`H`$ be a Hilbert space. Given a vector $`xH`$, the set of all vectors of the form $`\lambda x`$ with $`\lambda `$, $`|\lambda |=1`$ is called the ray associated to $`x`$ and it is denoted by $`\underset{¯}{x}`$. For any $`x,yH`$ we define
$$\underset{¯}{x}\underset{¯}{y}=|x,y|.$$
The notation $`\underset{¯}{H}`$ stands for the set of all rays in $`H`$. The algebra of all bounded linear operators on $`H`$ is denoted by $`B(H)`$, and $`F(H)`$ stands for the ideal of all finite rank operators in $`B(H)`$. If $`x,yH`$ are arbitrary vectors, then $`xy`$ is an element of $`F(H)`$ which is defined by $`(xy)z=z,yx`$ $`(zH)`$. A linear map $`\varphi :𝒜`$ between the algebras $`𝒜`$ and $``$ is called a Jordan homomorphism if
$$\varphi (x^2)=\varphi (x)^2(x𝒜).$$
Our main result which follows presents a Wigner-type result for pairs of ray transformations.
###### Theorem.
Let $`H`$ be a complex Hilbert space of dimension at least 3. Let $`T,S:\underset{¯}{H}\underset{¯}{H}`$ be bijective transformations with the property that
$$T\underset{¯}{x}S\underset{¯}{y}=\underset{¯}{x}\underset{¯}{y}(\underset{¯}{x},\underset{¯}{y}\underset{¯}{H}).$$
Then there are bounded invertible either both linear or both conjugate-linear operators $`U,V:HH`$ such that $`V=U_{}^{}{}_{}{}^{1}`$ and
$$T\underset{¯}{x}=\underset{¯}{Ux},S\underset{¯}{x}=\underset{¯}{Vx}(xH).$$
###### Proof.
For every $`xH`$ pick a vector from $`T\underset{¯}{x}`$. In that way we get a function, which will be denoted by the same symbol $`T`$, from $`H`$ into itself with the property that for every vector $`yH`$, there exists a vector $`xH`$ such that $`y=\lambda Tx`$ for some $`\lambda `$ of modulus 1. Let us do the same with the other transformation $`S`$. Clearly, we have
$$|Tx,Sy|=|x,y|(x,yH).$$
Obviously, for every unit vector $`xH`$ we can choose a scalar $`\lambda _x`$ with $`|\lambda _x|=1`$ such that $`\lambda _xTx,Sx=1`$. By the properties of our original transformation $`T`$, we can clearly suppose that here in fact we have $`Tx,Sx=1`$. We define a function $`\mu `$ on the set $`P_f(H)`$ of all finite rank projections (self-adjoint idempotents) on $`H`$ as follows. If $`PP_f(H)`$, then there are pairwise orthogonal unit vectors $`x_1,\mathrm{},x_nH`$ such that $`P=x_1x_1+\mathrm{}+x_nx_n`$. We set
$$\mu (P)=Tx_1Sx_1+\mathrm{}+Tx_nSx_n.$$
Apparently, the operators $`Tx_1Sx_1,\mathrm{},Tx_nSx_n`$ are pairwise orthogonal rank-one idempotents (two idempotents $`P,Q`$ are said to be orthogonal if $`PQ=QP=0`$). Hence, $`\mu (P)`$ is a rank-$`n`$ idempotent. We have to check that $`\mu `$ is well-defined. This follows from the following observation. We have
$$\text{rng }(\underset{k=1}{\overset{n}{}}Tx_kSx_k)=[Tx_1,\mathrm{},Tx_n]$$
and
$$\text{ker }(\underset{k=1}{\overset{n}{}}Tx_kSx_k)=[Sx_1,\mathrm{},Sx_n]^{},$$
where $`[.]`$ denotes generated subspace. Now, suppose that the pairwise orthogonal unit vectors $`x_1^{},\mathrm{},x_n^{}`$ generate the same subspace as $`x_1,\mathrm{},x_n`$ do. Let $`yH`$. Then there exist a vector $`xH`$ and a scalar $`\lambda `$ of modulus 1 such that $`y=\lambda Sx`$. We have
$$y[Tx_1,\mathrm{},Tx_n]Sx[Tx_1,\mathrm{},Tx_n]$$
(1)
$$x[x_1,\mathrm{},x_n]x[x_1^{},\mathrm{},x_n^{}]$$
$$Sx[Tx_1^{},\mathrm{},Tx_n^{}]y[Tx_1^{},\mathrm{},Tx_n^{}].$$
This shows that the range of $`_{k=1}^nTx_kSx_k`$ is the same as that of $`_{k=1}^nTx_k^{}Sx_k^{}`$. The same applies for the kernels. Since the idempotents are determined by their ranges and kernels, this proves that $`\mu `$ is well-defined. It is now clear that $`\mu `$ is an orthoadditive measure on $`P_f(H)`$. We show that $`\mu `$ is bounded on the set $`P_1(H)`$ of all rank-one projections which is equivalent to
$$\underset{x=1}{sup}TxSx<\mathrm{}.$$
Suppose, on the contrary, that there is a sequence $`(u_n)`$ of unit vectors in $`H`$ for which $`Tu_nSu_n\mathrm{}`$. Since $`(u_n)`$ is bounded, it has a subsequence $`(u_{k_n})`$ weakly converging to a vector, say, $`uH`$. We have
$$|Tu_{k_n},Sv|=|u_{k_n},v||u,v|.$$
Since this holds for every $`vH`$, we deduce that $`(Tu_{k_n})`$ is weakly bounded which implies that it is in fact norm-bounded. The same argument applies in relation to $`S`$. Hence, we obtain that $`(u_n)`$ has a subsequence $`(u_{l_n})`$ such that $`Tu_{l_n},Su_{l_n}`$ are bounded which is a contradiction. Consequently, $`\mu `$ is bounded on $`P_1(H)`$.
By Gleason’s theorem $`\mu `$ can be extended to a Jordan homomorphism of $`F(H)`$. In fact, if $`AF(H)`$ is self-adjoint, then there are finite-rank projections $`P_1,\mathrm{},P_n`$ (here, we do not require that they are pairwise orthogonal) and scalars $`\lambda _1,\mathrm{},\lambda _n`$ such that $`A=\lambda _1P_1+\mathrm{}+\lambda _nP_n`$. Let
$$\varphi (A)=\lambda _1\mu (P_1)+\mathrm{}+\lambda _n\mu (P_n).$$
Consider a finite dimensional subspace $`H_0`$ of $`H`$ with dimension at least 3 which contains all the subspaces $`\text{rng }A,\text{ker }A^{},\text{rng }P_1,\mathrm{},\text{rng }P_n`$. Since $`\mu `$ is bounded on $`P_1(H_0)`$, by the variation \[5, Theorem 3.2.16\] of Gleason’s theorem, for every $`x,yH`$ there is an operator $`T_{xy}`$ on $`H_0`$ such that
$$\lambda _1\mu (P_1)+\mathrm{}+\lambda _n\mu (P_n)x,y=\lambda _1\mu (P_1)x,y+\mathrm{}+\lambda _n\mu (P_n)x,y=$$
$$\lambda _1\text{tr }(P_1T_{xy})+\mathrm{}+\lambda _n\text{tr }(P_nT_{xy})=\text{tr }(AT_{xy}).$$
We now easily obtain that $`\varphi `$ is well-defined and real-linear on the set of all self-adjoint finite rank operators. If $`AF(H)`$ is arbitrary, then there exist self-adjoint finite rank operators $`A_1,A_2`$ such that $`A=A_1+iA_2`$. Define $`\varphi (A)=\varphi (A_1)+i\varphi (A_2)`$. Clearly, $`\varphi `$ is a linear map on $`F(H)`$ which sends projections to idempotents. It is a standard algebraic argument to verify that $`\varphi `$ is then a Jordan homomorphism (see, for example, the proof of \[9, Theorem 2\]). Since $`F(H)`$ is a locally matrix ring, we can apply a classical theorem of Jacobson and Rickart. By \[8, Theorem 8\] we obtain that $`\varphi `$ can be written as $`\varphi =\varphi _1+\varphi _2`$, where $`\varphi _1`$ is a homomorphism and $`\varphi _2`$ is an antihomomorphism. Since $`\varphi (P)`$ is a rank-one idempotent and $`\varphi _1(P)`$, $`\varphi _2(P)`$ are idempotents, we infer from $`\varphi (P)=\varphi _1(P)+\varphi _2(P)`$ that either $`\varphi _1(P)=0`$ or $`\varphi _2(P)=0`$. Since the ring $`F(H)`$ is simple, we obtain that either $`\varphi _1=0`$ or $`\varphi _2=0`$. Therefore, $`\varphi `$ is either a homomorphism or an antihomomorphism. Without loss of generality we can assume that $`\varphi `$ is a homomorphism. We assert that $`\varphi `$ is rank-preserving. Let $`AF(H)`$ be a rank-$`n`$ operator. Then there is a rank-$`n`$ projection $`P`$ such that $`PA=A`$. The rank of $`\varphi (P)`$ is also $`n`$. We have $`\varphi (A)=\varphi (PA)=\varphi (P)\varphi (A)`$ which proves that $`\varphi (A)`$ is of rank at most $`n`$. If $`Q`$ is any rank-$`n`$ projection, then there are finite rank operators $`U,V`$ such that $`Q=UAV`$. Since $`\varphi (Q)=\varphi (U)\varphi (A)\varphi (V)`$ and the rank of $`\varphi (Q)`$ is $`n`$, it follows that the rank of $`\varphi (A)`$ is at least $`n`$. Therefore, $`\varphi `$ is rank-preserving. We now refer to Hou’s work on the form of linear rank preservers on operator algebras. It follows from the argument leading to \[6, Theorem 1.2\] that there are linear operators $`U,V`$ on $`H`$ such that $`\varphi `$ is of the form
(2)
$$\varphi (xy)=(Ux)(Vy)(x,yH)$$
(recall that we have assumed that $`\varphi `$ is a homomorphism). If $`xH`$ is a unit vector, then we have $`TxSx=\varphi (xx)=UxVx`$. Taking traces, we obtain $`1=Tx,Sx=Ux,Vx`$. Since this holds for every unit vector $`x`$, by the linearity of $`U,V`$, using polarization we get that
(3)
$$Ux,Vy=x,y(x,yH).$$
We assert that $`U,V`$ are surjective. Consider, for example, the case of $`U`$. Let $`0xH`$ be any vector and let $`0\lambda `$ be any scalar. It is easy to see that $`[Tx]^{}=[T(\lambda x)]^{}`$ (see (1)). Therefore, $`T(\lambda x)=\lambda ^{}Tx`$ with some scalar $`\lambda ^{}`$. Denote $`x_e=x/x`$. We compute
$$UxVx=x^2Ux_eVx_e=x^2\varphi (x_ex_e)=x^2Tx_eSx_e.$$
This gives us that $`Tx_e[Ux]`$. But $`Tx`$ is in the one-dimensional subspace generated by $`Tx_e`$. So, we have
(4)
$$Tx[Ux].$$
Since $`\text{rng }U`$ is a linear subspace of $`H`$ and $`T`$ is ”almost” surjective, we obtain the surjectivity of $`U`$. Similar argument applies to $`V`$. We next show that $`U,V`$ are bounded. Let $`(x_n)`$ be a sequence converging to 0 and let $`yH`$ be such that $`Ux_ny`$. If $`xH`$ is arbitrary, then we have
$$Ux_n,Vx=x_n,x0.$$
Since $`V`$ is surjective, we obtain that $`(Ux_n)`$ weakly converges to 0. It follows that $`y=0`$. By the closed graph theorem we deduce that $`U`$ is bounded. Similar argument proves the boundedness of $`V`$. It follows from (3) that $`V^{}U=I`$. This gives us that $`U`$ is injective. Therefore, $`U`$ and $`V`$ are invertible and $`V=U_{}^{}{}_{}{}^{1}`$.
By (4) and the similar relation $`Sx[Vx]`$ $`(xH)`$, there are functions $`\phi ,\psi :H`$ such that
$$Tx=\phi (x)Ux,Sx=\psi (x)Vx(xH).$$
We have
$$|\phi (x)||\psi (y)||x,y|=|\phi (x)||\psi (y)||Ux,Vy|=|Tx,Sy|=|x,y|,$$
that is, $`|\phi (x)||\psi (y)|=1`$ if $`x,y0`$. This easily implies that $`|\phi |`$ and $`|\psi |`$ are both constant. Multiplying $`U,V`$, $`\phi ,\psi `$ by suitable constants, we obtain the statement of the theorem. The proof is complete. ∎
In the following corollary of our theorem we give a generalization of Wigner’s theorem for the indefinite inner product space generated by any invertible operator $`AB(H)`$. Since we do not assume that $`A`$ is self-adjoint, this result can, in some sense, be considered as a generalization of the results in .
###### Corollary 1.
Let $`H`$ be a complex Hilbert space with $`dimH3`$ and let $`AB(H)`$ be invertible. For any $`x,yH`$ define $`\underset{¯}{x}_A\underset{¯}{y}=|Ax,y|`$. Let $`T:\underset{¯}{H}\underset{¯}{H}`$ be a bijective transformation such that
$$T\underset{¯}{x}_ATy=\underset{¯}{x}_A\underset{¯}{y}(x,yH).$$
Then there is a bounded invertible either linear or conjugate-linear operator $`U`$ on $`H`$ with $`U^{}AU=ϵA`$ for some scalar $`ϵ`$ of modulus 1 such that
$$T\underset{¯}{x}=\underset{¯}{Ux}(xH).$$
###### Proof.
Just as in the proof of our theorem above, we can define an ”almost” surjective map (that is, which has values in every ray) on the underlying Hilbert space $`H`$ denoted by the same symbol $`T`$ such that
$$|ATx,Ty|=|Ax,y|(x,yH).$$
Set $`S=ATA^1`$. The proof of our theorem now applies and we find that there is a bounded invertible either linear or conjugate-linear operator $`U`$ on $`H`$ and a scalar function $`\phi :H`$ such that $`Tx=\phi (x)Ux`$ $`(xH)`$. Since
(5)
$$|\phi (x)||\phi (y)||AUx,Uy|=|ATx,Ty|=|Ax,y|(x,yH),$$
it follows that $`[U^{}AUx]^{}=[Ax]^{}`$ for every $`xH`$. Therefore, the linear operators $`U^{}AU`$ and $`A`$ are locally linearly dependent which means that $`U^{}AUx`$ and $`Ax`$ are linearly dependent for every $`xH`$. Since none of the operators $`U^{}AU`$ and $`A`$ is of rank 1, by \[7, Lemma 3\] we obtain that there is a scalar $`c`$ such that $`U^{}AU=cA`$. Let $`x,yH`$ be arbitrary nonzero vectors. Pick $`zH`$ such that $`Ax,z,Ay,z0`$. From (5) we now infer that
$$|\phi (x)||\phi (z)||c|=1,|\phi (y)||\phi (z)||c|=1.$$
This shows that $`|\phi |`$ is constant. If $`d`$ denotes this constant, then we have $`d^2|c|=1`$. Let $`ϵ=d^2c`$. Then $`ϵ`$ is of modulus 1 and we have
$$(dϵU)^{}A(dϵU)=d^2U^{}AU=d^2cA=ϵA.$$
Consider the factorization
$$Tx=\left(\frac{1}{dϵ}\phi (x)\right)(dϵU).$$
Since $`\frac{1}{dϵ}\phi (x)`$ is of modulus 1, the proof is complete. ∎
In the finite dimensional case, Corollary 1 can be reformulated in the following way.
###### Corollary 2.
Let $`H`$ be a finite dimensional complex Hilbert space with $`dimH3`$. Let $`B:H\times H`$ be a sesquiliner form which is non-degenerate in the sense that $`B(x,y)=0`$ $`(yH)`$ implies $`x=0`$. Define $`\underset{¯}{x}_B\underset{¯}{y}=|B(x,y)|`$ $`(x,yH)`$. Let $`T:\underset{¯}{H}\underset{¯}{H}`$ be a bijective transformation such that
$$T\underset{¯}{x}_BTy=\underset{¯}{x}_B\underset{¯}{y}(x,yH).$$
Then either there is an invertible linear operator $`U`$ on $`H`$ such that $`B(Ux,Uy)=ϵB(x,y)`$ $`(x,yH)`$ for some scalar $`ϵ`$ of modulus 1 and
$$T\underset{¯}{x}=\underset{¯}{Ux}(xH),$$
or there is an invertible conjugate-linear operator $`U^{}`$ on $`H`$ such that $`\overline{B(U^{}x,U^{}y)}=ϵ^{}B(x,y)`$ $`(x,yH)`$ for some scalar $`ϵ^{}`$ of modulus 1 and
$$T\underset{¯}{x}=\underset{¯}{U^{}x}(xH).$$
###### Proof.
Since $`H`$ is finite dimensional, it is easy to see that there exists an invertible linear operator $`A`$ on $`H`$ such that $`B(x,y)=Ax,y`$ $`(x,yH)`$. Now, Corollary 1 applies. ∎
###### Remark 1.
Our results are valid in real Hilbert spaces as well. In order to see it, we must refine the argument we have presented in the complex case. Namely, one can follow the argument that has been applied in the proof of \[10, Theorem 3\]. Observe that in the papers the authors considered only complex spaces.
Acknowledgements
This research was supported from the following sources: (1) Hungarian National Foundation for Scientific Research (OTKA), Grant No. T–030082 F–019322, (2) A grant from the Ministry of Education, Hungary, Reg. No. FKFP 0304/1997.
|
warning/0005/hep-ph0005180.html
|
ar5iv
|
text
|
# The Chiral Fermion Meson Model at Finite Temperature
## I Introduction
Several models have been proposed to describe hadron properties in the regime of low energies. Among these models, we adopt the linear $`\sigma `$ model of Gell-Mann and Levy which is a phenomenological model of Quantum Chromodynamics-QCD that incorporates two important features of QCD: chiral symmetry and partial conservation of axial vector current. The model was originally proposed as a model for strong interactions , but nowadays it serves as an effective model for the low energy (low temperature) phase of QCD. It has the advantage of being renormalizable at zero and finite temperature. Altough the linear sigma model lagrangian exhibits chiral symmetry, quantum effects break this symmetry spontaneously. Both from theoretical and experimental point of view, there exist a great amount of interest in the study of chiral symmetry restoration at finite temperature.
However, quantum field theory at high temperature has a well known problem that is the breakdown of the perturbative expansion. This happens in theories with spontaneous symmetry breaking (SSB) or in massless field theories because powers of the temperature can compensate for powers of the coupling constant. Resummation techniques which try consistently to take into account higher-loops are required.
A systematic self-consistent approximation approaches based on the meson sector of the linear $`\sigma `$ model was previously studied by Baym and Grinstein. After, Banerjee and Mallik proposed a modified perturbation expansion with the objective of calculating the two-point functions up to second order in the $`\lambda \varphi ^4`$ theory. A resummed perturbative expansion was proposed by Parwani in order to go beyond leading order in the same model. More recently, Chiku and Hatsuda in the study of the $`O(N)`$ $`\varphi ^4`$ model presented a novel resummation adding a mass parameter determined later by the fastest apparent convergence (FAC) condition. We employ imaginary-time formulation (ITF) whereas in real-time formulation (RTF) is used in the development of the optimized perturbation theory (OPT). In this paper we develop a modified self-consistent resummation at finite temperature and apply it to the investigation of the chiral fermion meson model. We study the temperature dependence of the chiral condensate and the effective meson and fermion masses by this self-consistent non-perturbative approximation up to one-loop order in the perturbative expansion. In the application of the MSCR to the study of the chiral fermion meson model at finite temperature, we divided the problem into three physical regions: low, intermediate and high temperatures. This is essential to identify the regions where resummation is crucial. In each region renormalization and satisfaction of Goldstone’s theorem are discussed in detail. Our study addresses problems found in the context of the well studied $`O(4)`$ linear sigma model and deals with a usually avoided point: the inclusion of the fermions. Also, we re-examine the chiral phase transition in static equilibrium in terms of the linear sigma model with our MSCR. Instead of demanding a infinite gap-equation, as has been done often in the recent literature, we perform the renormalization in stages in order to get finite gap-equations.
We also treat an explicit chiral symmetry breaking term in the Lagrangian which generates the realistic finite pion mass. Symmetry is never restored in this case. It is shown that in the limit of vanishing pion mass, namely when the chiral symmetry is exact, the inclusion of fermions does not change the order (nature) of the phase transition but only lowers the value of the critical temperature.
This paper is organized as follows. In Section II we discuss the chiral fermion meson model and some of its features at zero temperature. In Section III the temperature is introduced via the partition function of the model which lead to the thermodynamical potential. The inclusion of loop corrections and the thermal gap equations are addressed to Section IV. In section V we apply the MSCR to the study of the massless $`\lambda \varphi ^4`$ model in the weak coupling limit at high temperature. The renormalization of the self-energy is studied in Section VI. The numerical results are presented in Section VII. Section VIII is devoted to conclusions.
## II The Chiral Fermion Meson Model at Zero Temperature
The Lagrangian density of the chiral fermion meson model which provides an explicit realization of chiral symmetry is given by
$$_{sym}=\overline{\psi }\left[i\gamma ^\mu _\mu g(\sigma ^{}+i\gamma ^5\stackrel{}{\pi }\stackrel{}{\tau })\right]\psi +\frac{1}{2}\left[(\sigma ^{})^2+(\stackrel{}{\pi })^2\right]\frac{\lambda }{4}(\sigma ^2+\stackrel{}{\pi }^2f_\pi ^2)^2,$$
(1)
where $`\psi `$, $`\sigma ^{}`$, and $`\pi `$ represent the quark, sigma and pion fields, respectively, $`\lambda `$ and $`g`$ are positive coupling constants and $`f_\pi `$ is the pion decay constant in vacuum.
If the up and down quark masses were zero, QCD would have a chiral $`SU(2)_L\times SU(2)_R`$ symmetry. In the vacuum this symmetry is spontaneously broken by quantum effects, with the result that there exists a triplet of Goldstone bosons. In reality the quark masses are very small but nonzero, so that chiral symmetry is only approximate and the pion has a small mass . An explicit chiral symmetry breaking term is added to the Lagrangian which generates the realistic finite pion mass so that
$$^{}=_{sym}+_{symb}$$
(2)
with
$$_{symb}=c\sigma ^{},$$
(3)
where $`c`$ is small and positive.
The term $`_{sym}`$ is symmetric and invariant under an $`SU(2)_L\times SU(2)_R`$ chiral group and $`_{symb}`$ is the symmetry breaking term. Two Noether currents associated with (1), namely the vector current and the axial vector current, are given by
$`\stackrel{}{V}_\mu =\overline{\psi }\gamma _\mu {\displaystyle \frac{\stackrel{}{\tau }}{2}}\psi +\stackrel{}{\pi }\times _\mu \stackrel{}{\pi },`$
$`\stackrel{}{A}_\mu =\overline{\psi }\gamma _\mu \gamma _5{\displaystyle \frac{\stackrel{}{\tau }}{2}}\psi +\sigma ^{}_\mu \stackrel{}{\pi }\stackrel{}{\pi }_\mu \sigma ^{}`$
respectively. The equations of motion for the fields derived from the Lagrangian density (1) give the PCAC relations
$$_\mu \stackrel{}{A}^\mu =c\stackrel{}{\pi }.$$
(4)
The effect of the term $`c\sigma ^{}`$ on the classical fundamental state, can be found by looking at the minimum of the potential
$$V_0(\sigma ^{},\stackrel{}{\pi })=\frac{\lambda }{4}(\sigma ^2+\stackrel{}{\pi }^2f_\pi ^2)^2c\sigma ^{}$$
(5)
$$\frac{V_0(\sigma ^{},\stackrel{}{\pi })}{\sigma ^{}}=\lambda (\sigma ^2+\stackrel{}{\pi }^2f_\pi ^2)\sigma ^{}c=0$$
(6)
$$\frac{V_0(\sigma ^{},\stackrel{}{\pi })}{\pi ^a}=\lambda (\sigma ^2+\stackrel{}{\pi }^2f_\pi ^2)\pi ^a=0$$
(7)
whose (unique) solutions are
$`\stackrel{}{\pi }_0=0,`$ (8)
$`\lambda (\sigma _0^2f_\pi ^2)\sigma _0^{}=c`$ (9)
To first order in $`c`$, we have
$$\sigma _0^{}=f_\pi +\frac{c}{2\lambda f_\pi ^2}\nu $$
(10)
From (10), we see that $`\sigma _0^{}`$ has a non-zero vacuum expectation value. It is convenient to redefine the sigma field as $`\sigma ^{}\sigma +\nu `$ such that $`\sigma `$ has zero expectation value. As an effect of this shift the fermion field acquires a mass given by
$$m_\psi =g\nu $$
(11)
The shifted Lagrangian, $`_s`$, of the new quantum theory reads
$`_s={\displaystyle \frac{\lambda }{4}}(f_\pi ^2\nu ^2)^2+c\nu \lambda (\nu ^3\nu f_\pi ^2{\displaystyle \frac{c}{\lambda }})\sigma +`$ (12)
$`\overline{\psi }[i\gamma ^\mu _\mu m_\psi ]\psi +{\displaystyle \frac{1}{2}}[(\stackrel{}{\pi })^2m_\pi ^2\stackrel{}{\pi }^2+(\sigma )^2m_\sigma ^2\sigma ^2]+`$ (13)
$`g\overline{\psi }[\sigma +i\gamma ^5\stackrel{}{\pi }\stackrel{}{\tau })]\psi {\displaystyle \frac{\lambda }{4}}[(\stackrel{}{\pi }^2+\sigma ^2)^2+4\nu \sigma (\stackrel{}{\pi }^2+\sigma ^2)]=`$ (14)
$`U(\nu )+_0+_I`$ (15)
where $`U(\nu )`$ is the mean field energy density, $`_0`$ is the free Lagrangian and $`_I`$ is the interaction Lagrangian, defined by
$$U(\nu )\frac{\lambda }{4}(f_\pi ^2\nu ^2)^2c\nu ,$$
(16)
$$_0\overline{\psi }[i\gamma ^\mu _\mu m_\psi ]\psi +\frac{1}{2}[(\stackrel{}{\pi })^2m_\pi ^2\stackrel{}{\pi }^2+(\sigma )^2m_\sigma ^2\sigma ^2],$$
(17)
$$_Ig\overline{\psi }[\sigma +i\gamma ^5\stackrel{}{\pi }\stackrel{}{\tau })]\psi \frac{\lambda }{4}[(\stackrel{}{\pi }^2+\sigma ^2)^2+4\nu \sigma (\stackrel{}{\pi }^2+\sigma ^2)],$$
(18)
respectively.
The meson masses read out of the shifted Lagrangian (15) are
$$m_\pi ^2=m^2+\lambda \nu ^2,$$
(19)
$$m_\sigma ^2=m^2+3\lambda \nu ^2$$
(20)
where $`m^2=\lambda f_\pi ^2<0`$.
It is easy to see that the coefficient of the linear term in the sigma field, $`\lambda (\nu ^3\nu f_\pi ^2\frac{c}{\lambda })`$, in Lagrangian (15) is identically zero by the minimal condition (8). This is due to the fact that the vacuum expectation value of the sigma field, $`\sigma `$, should vanish at any order of perturbation theory, even if we include thermal corrections. The one-loop thermal tadpole corrections will modify this relation which will become temperature dependent. If $`\nu `$ is allowed to be temperature dependent, the masses are temperature dependent as well. At any temperature, $`\nu `$ is such that $`\sigma =0`$. At zero temperature, when $`c`$ continuously approaches zero, we have the solutions $`\stackrel{}{\pi }=0`$ and $`\sigma ^{}=f_\pi `$ which minimize the potential satisfying the Goldstone’s theorem.
The contact with phenomenology is made by fixing the parameters of the model to agree with the observable value of the particle masses in vacuum. Then, the tree level parameters of the Lagrangian are
$$\lambda =\frac{m_{\sigma ;0}^23m_{\pi ;0}^2}{2f_\pi ^2},$$
(21)
$$c=\nu m_{\pi ;0}^2f_\pi m_{\pi ;0}^2,$$
(22)
$$g=\frac{m_{\psi ;0}(m_{\sigma ;0}^23m_{\pi ;0}^2)}{f_\pi (m_{\sigma ;0}^22m_{\pi ;0}^2)},$$
(23)
where $`m_{\pi ;0}=139MeV`$ , $`m_{\sigma ;0}=600MeV`$ , $`m_{\psi ;0}=340MeV`$ and $`f_\pi =93MeV`$.
As we have mentioned earlier, our goal in this work is to study the chiral phase transition in the chiral fermion meson model and to analyze the thermal behavior of the temperature dependent meson condensate $`\nu `$ and the meson and fermion masses. So it will be necessary to compute all the one-loop self-energies for the particles present in the model. Such self-energy diagrams have divergent pieces which must be renormalized if we want reliable results. In most of the approximations found in the literature several difficulties have been found in the tentative of renormalizing the divergent gap-equations. Sometimes the undesirable parts have been ignored. The renormalization of the self-energy is studied in sectionVI whereas the effective potential renormalization is performed in appendix A. For the purpose of renormalization it is necessary to add to $`_s`$ a counterterm Lagrangian, $`L_{C.T.}`$, needed to render the theory finite ,
$$=_s+L_{C.T.}$$
(24)
where
$$L_{C.T.}=C.T.+D_1m_\psi ^4+D_2m_\pi ^4+D_3m_\sigma ^4$$
(25)
In (25) $`C.T.`$ contains the appropriate counterterms to be used in the renormalization of the masses while $`D_{1,2,3}m_{\psi ,\pi ,\sigma }^4`$ are necessary to keep the thermodynamical potential finite, as we will see further. As we are interested only in the study of the thermal effective masses at one loop order in the perturbative expansion, other counterterms necessary to renormalize the coupling constants are not explicitly shown.
## III The Partition Function of the Model and the link to Statistical Mechanics
One of the most fundamental objects in thermodynamics is the partition function, defined by,
$$Z=\mathrm{Tr}e^{\beta \widehat{H}},$$
(26)
where $`\widehat{H}`$ is the Hamiltonian of the system, $`\beta =1/k_BT=T^1`$ with the Boltzmann constant, $`k_B`$, set equal to one and the trace, $`\mathrm{Tr}`$, in eq. (26) meaning the sum of the elements of the matrix $`e^{\beta \widehat{H}}`$ in all independent states of the system. All information concerning the equilibrium thermodynamic macroscopic properties of the system are obtained from $`Z`$.
In relativistic quantum field theory, the partition function can be derived from the Feynman’s functional formalism . The bridge between quantum mechanics and statistical mechanics is achieved by the heuristic introduction of a variable defined as $`\tau =it`$. Also, the fields are constrained to obey periodic(anti) boundary conditions: $`\varphi (𝐫,0)=\varphi (𝐫,\beta )`$ for bosons and $`\psi (𝐫,0)=\psi (𝐫,\beta )`$ for fermions. Following these prescriptions, we get
$`Z[\overline{\psi },\psi ,\sigma ,\pi ]=N^{}{\displaystyle 𝒟[\varphi ]\mathrm{exp}[_0^\beta 𝑑\tau d^3x(\varphi ,\varphi )]}=`$ (27)
$`\mathrm{exp}{\displaystyle _\beta }d^4x\left[{\displaystyle \frac{\lambda }{4}}(f_\pi ^2\nu ^2)^2+c\nu +D_1m_\psi ^4+D_2m_\pi ^4+D_3m_\sigma ^4\right]N^{}{\displaystyle 𝒟[\varphi ]e^{S_0}e^{S_I}}`$ (28)
Here we have introduced a short hand notation for the Euclidean space-time integral: $`S=_\beta d^4x_0^\beta 𝑑\tau d^3x`$, $`𝒟[\varphi ]`$ is an abbreviation for the integral over $`\overline{\psi }`$, $`\psi `$, $`\sigma `$ and $`\pi `$, $`N^{}`$ is an unimportant infinit constant and $``$ is given by (24).
Next we introduce the thermodynamical potential, $`\mathrm{\Omega }`$, difined by
$$\mathrm{\Omega }(T,\nu )=\frac{T}{V}\mathrm{ln}Z$$
(29)
where $`lnZ=\frac{V}{T}(U(\nu )+D_{1,2,3}m_{\psi ,\pi ,\sigma }^4)+lnZ_0+lnZ_I`$. Since the interaction action $`S_I`$ contains terms which are more than quadratic in the fields, it is not possible to carry out the functional integration above in closed form. For a while we will neglect $`S_I`$ in our calculations. This amounts to considering only the tadpole contributions. Thus,
$`\mathrm{\Omega }_1(T,\nu ){\displaystyle \frac{\lambda }{4}}(f_\pi ^2\nu ^2)^2c\nu \left[D_1m_\psi ^4+D_2m_\pi ^4+D_3m_\sigma ^4\right]{\displaystyle \frac{T}{V}}lnZ_0=`$ (30)
$`{\displaystyle \frac{\lambda }{4}}(f_\pi ^2\nu ^2)^2c\nu \left[D_1m_\psi ^4+D_2m_\pi ^4+D_3m_\sigma ^4\right]+`$ (31)
$`{\displaystyle \frac{d^3p}{(2\pi )^3}\left\{\frac{1}{2}\omega _\sigma +T\mathrm{ln}(1e^{\beta \omega _\sigma })+\frac{3}{2}\omega _\pi +3T\mathrm{ln}(1e^{\beta \omega _\pi })22\left[\omega _\psi +2T\mathrm{ln}(1+e^{\beta \omega _\psi })\right]\right\}}`$ (32)
with $`\omega _\sigma ^2𝐩^2+m_\sigma ^2`$, $`\omega _\pi ^2𝐩^2+m_\pi ^2`$ and $`\omega _\psi ^2𝐩^2+m_\psi ^2`$. In the third line of (30) the first factor 2 multiplying the bracket which contains the fermion contribution, comes from the spin degrees of freedom, whereas the other factor two is due the isospin degrees of freedom. Inside this same bracket there is another factor 2 corresponding to the particle and antiparticle contributions. The thermodynamical potential $`\mathrm{\Omega }_1`$ is precisely the one-loop effective potential of the linear sigma model, and it can be expressed as,
$$\mathrm{\Omega }_1(T,\nu )=U(\nu )+D_{1,2,3}m_{\psi ,\pi ,\sigma }^4+\mathrm{\Omega }_1^0(\nu )+\mathrm{\Omega }_1^\beta (T,\nu )$$
(33)
The equation of state of the noninteracting system composed by a (free) relativistic boson and fermion gas is
$`P_0={\displaystyle \frac{T}{V}}\mathrm{ln}Z_0=\mathrm{\Omega }_1^\beta (T,\nu )=`$ (34)
$`T{\displaystyle \frac{d^3p}{(2\pi )^3}\left[\mathrm{ln}(1e^{\beta \omega _\sigma })+3\mathrm{ln}(1e^{\beta \omega _\pi })8\mathrm{ln}(1+e^{\beta \omega _\psi })\right]}`$ (35)
where $`P_0`$ is the thermal pressure.
The integration over the temperature independent terms $`\mathrm{\Omega }_1^0(\nu )\frac{d^3p}{(2\pi )^3}[\frac{1}{2}\omega _\sigma +\frac{3}{2}\omega _\pi 4\omega _\psi ]`$ actually diverges. It is exactly the counterterms $`D_{1,2,3}m_{\psi ,\pi ,\sigma }^4`$ which will take care of these divergences. We discuss this further in section VI. It can be known, from thermodynamical considerations, that in thermal equilibrium $`\mathrm{\Omega }`$ is a minimum with respect to variations in $`\nu `$. Applying this extremum condition, we have
$$\frac{\mathrm{\Omega }_1(T,\nu )}{\nu }=0$$
(36)
Since $`\mathrm{\Omega }_1^\beta (T=0,\nu )=0`$, the divergent quantity $`\frac{\mathrm{\Omega }_1^0(\nu )}{\nu }=\frac{d^3p}{(2\pi )^3}\left[3\lambda (\frac{1}{2\omega _\sigma }+\frac{1}{2\omega _\pi })8g^2\frac{1}{2\omega _\psi }\right]\nu `$ represents the sum of the tadpoles at $`T=0`$. On the other hand, $`\frac{\mathrm{\Omega }_1^\beta (T,\nu )}{\nu }`$ is finite at $`T0`$ because of the natural cutoff $`\frac{1}{(e^{\beta \omega _{\sigma \pi ,(\psi )}}(+)1)}`$ present in the integrals. We could let the counterterms absorb the finite parts of $`\frac{\mathrm{\Omega }_1^0(\nu )}{\nu }`$ together with the infinities since the vacuum contribution is irrelevant for the discussions of thermodynamics. But, as we will see in section IV C, these terms are important in the verification of the Goldstone theorem in our self-consistent treatment. After these considerations, equation (36) is written as
$$\left[F(\omega _{\pi ,\sigma ,\psi }(m_{\pi ,\sigma ,\psi }),T)+G\left(\omega _{\pi ,\sigma ,\psi }(m_{\pi ,\sigma ,\psi })\right)\lambda (f_\pi ^2\nu ^2)\right]\nu c=0$$
(37)
where the functions $`F`$ and $`G`$ are defined as
$$F(\omega _{\pi ,\sigma ,\psi }(m_{\pi ,\sigma ,\psi }),T)\frac{d^3p}{(2\pi )^3}\left[3\lambda \left(\frac{1}{\omega _\sigma (e^{\beta \omega _\sigma }1)}+\frac{1}{\omega _\pi (e^{\beta \omega _\pi }1)}\right)+8g^2\frac{1}{\omega _\psi (e^{\beta \omega _\psi }+1)}\right],$$
(38)
$$G\left(\omega _{\pi ,\sigma ,\psi }(m_{\pi ,\sigma ,\psi })\right)\frac{3\lambda }{(4\pi )^2}\left(m_\pi ^2\mathrm{ln}(\frac{m_\pi ^2}{\mu ^2})+m_\sigma ^2\mathrm{ln}(\frac{m_\sigma ^2}{\mu ^2})\right)\frac{8g^2}{(4\pi )^2}m_\psi ^2\mathrm{ln}(\frac{m_\psi ^2}{\mu ^2}),$$
(39)
respectively. In eq.(39), $`\mu `$ is the renormalization scale.
As a first approximation, we consider only the thermal loop corrections to the effective potential. This approximation allows us to get an analytic expression for the approximate critical temperature.
$$\frac{d^3p}{(2\pi )^3}\left[3\lambda \left(\frac{n_\pi (\omega _\pi )}{\omega _\pi }+\frac{n_\sigma (\omega _\sigma )}{\omega _\sigma }\right)+8g^2\frac{n_\psi (\omega _\psi )}{\omega _\psi }\right]\nu \lambda (f_\pi ^2\nu ^2)\nu c=0$$
(40)
where $`n_{\pi ,\sigma }`$ and $`n_\psi `$ are the usual distribution functions for bosons and fermions given by
$$n_{\pi ,\sigma }(\omega _{\pi ,\sigma };T)=\frac{1}{e^{\beta \omega _{\pi ,\sigma }}1}$$
(41)
$$n_\psi (\omega _\psi ;T)=\frac{1}{e^{\beta \omega _\psi }+1},$$
(42)
respectively. In the above expression, when one minimizes the effective potential, one is summing the thermal tadpole contributions to the usually called mean field equation. The chiral condensate, $`\nu `$, which is a non trivial solution of this integral equation now depends on $`T`$. This equation can be solved with an explicit analytic form in the high temperature limit. The leading terms in the high temperature approximation for this integral equation are
$$\nu ^3+\left[\frac{1}{2}(1+\frac{2g^2}{3\lambda })T^2f_\pi ^2\right]\nu \frac{c}{\lambda }=0.$$
(43)
The above equation has a real solution that is a slowly decreasing function of temperature, but does not vanish. Thus, when $`c0`$ the symmetry is never restored. On the other hand, when $`c=0`$ the non trivial solution of (43) is
$$\nu ^2=f_\pi ^2\frac{1}{2}(1+\frac{2g^2}{3\lambda })T^2$$
(44)
The critical temperature is defined as the temperature where the condensate goes to zero. It is given by
$$T_c^2=\frac{2f_\pi ^2}{(1+\frac{2g^2}{3\lambda })}$$
(45)
It shows that the inclusion of fermions does not change the order of the phase transition, but only lowers the value of $`T_c`$. Note that the interactions of the mesons with the fermions forces the “critical” temperature to depend on the coupling constants $`\lambda `$ and $`g`$. If $`g=0`$ we recover the result of .
## IV Inclusion of Loop Corrections
### A The first necessity: Beyond the mean field approximation
Let us analyze the finite temperature behavior of the tree-level meson masses (19) and (20) as functions of the thermal expectation value of the sigma field, $`\nu (T)`$. Since $`\nu ^2(T)`$ decreases as $`T`$ increases and $`m^2<0`$, the particle masses becomes tachyonic. Another problem which arises is the fact that Goldstone’s theorem is not satisfied in the ordered phase (when $`c=0`$), i.e., replacing (44) on (19), we obtain a non zero pion mass given by $`m_\pi ^2=\frac{\lambda }{2}(1+\frac{2g^2}{3\lambda })T^2`$. This pathological behavior is due to the fact that in our approximation we have neglect the interaction action $`S_I`$ in the thermodynamical potential (29). The result is that the mean field approximation can be trusted only in the approximate prediction of a phase transition at $`T_c`$ given by (45). It is incorrect in what concern the description of the finite temperature behavior of the meson and fermion masses. So it is necessary to include all one-loop corrections from all 1PI diagrams present in $`S_I`$ to the masses.
Following the program of we expand the partition function in powers of the interaction, in order to get the one-loop self-energy corrections.
$$\mathrm{ln}Z_I=\mathrm{ln}(1+\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n!}\frac{[d\varphi ]e^{S_0}S_I^n}{[d\varphi ]e^{S_0}})$$
(46)
The one-loop 1PI graphs come from $`\mathrm{ln}Z_1+\mathrm{ln}Z_2`$, which are given by
$$\mathrm{ln}Z_1=\frac{[d\varphi ]e^{S_0}S_I}{[d\varphi ]e^{S_0}}$$
(47)
$$\mathrm{ln}Z_2=\frac{1}{2}(\frac{[d\varphi ]e^{S_0}S_I}{[d\varphi ]e^{S_0}})^2+\frac{1}{2}\frac{[d\varphi ]e^{S_0}S_I^2}{[d\varphi ]e^{S_0}}$$
(48)
where the disconnected diagrams cancel in $`lnZ_2`$ and the diagrams which gives rise to tadpoles in the self-energy are not to be included, since their effect is already considered in the mean field equation. The terms which “survive” come from:
$`\mathrm{ln}Z_I^{2loop}={\displaystyle \frac{\lambda }{4}}{\displaystyle \frac{[d\varphi ]e^{S_0}𝑑\tau d^3x(\pi ^2+\sigma ^2)^2}{[d\varphi ]e^{S_o}}}+`$ (49)
$`{\displaystyle \frac{1}{2}}{\displaystyle _0^\beta }𝑑\tau _1𝑑\tau _2{\displaystyle d^3x_1d^3x_2\frac{[d\varphi ]e^{S_0}\left\{\lambda ^2\nu ^2[(\sigma \pi ^2)^2+(\sigma ^3)^2]+g^2[(\overline{\psi }\sigma \psi )^2+(\overline{\psi }i\gamma ^5\stackrel{}{\pi }\stackrel{}{\tau }\psi )^2]\right\}}{[d\varphi ]e^{S_0}}}`$ (50)
The 1PI graphs from this expression can be represented diagrammatically as shown in Fig.1.
The self-energy for bosons and fermions are defined, respectively by
$$𝐃(\omega _n,𝐤)_{\sigma ,\pi }^1=𝐃_{0\sigma ,\pi }(\omega _n,𝐤)^1+\mathrm{\Pi }_{\sigma ,\pi }(\omega _n,𝐤)$$
(51)
$$𝒟(\omega _n,𝐤)^1=𝒟_0(\omega _n,𝐤)^1+\mathrm{\Sigma }(\omega _n,𝐤)$$
(52)
where $`𝐃_{0\sigma ,\pi }(\omega _n,𝐩)`$ and $`𝒟_0(\omega _n,𝐩)`$ are the tree-level boson and fermion propagators, expressed respectively as
$$𝐃_{0\sigma ,\pi }(\omega _n,𝐤)^1=\omega _n^2+𝐤^2+m_{\sigma ,\pi }^2$$
(53)
$$𝒟_0(\omega _n,𝐤)^1=\gamma _\mu k^\mu m_\psi $$
(54)
Here, $`\omega _n`$ are the Matsubara frequencies, defined as $`\omega _n=2n\pi T`$ for bosons and $`\omega _n=(2n+1)\pi T\omega _{nf}`$ for fermions.
To one-loop order the self-energy expressions are given by
$$\mathrm{\Pi }_{\sigma ,\pi }=2(\frac{\delta \mathrm{ln}Z_I^{2loop}}{\delta 𝐃_{0\sigma ,\pi }})_{1PI}$$
(55)
$$\mathrm{\Sigma }=(\frac{\delta \mathrm{ln}Z_I^{2loop}}{\delta 𝒟_{0\psi }})_{1PI}$$
(56)
The self-energy graphs to each particle can be pictorially represented as by cutting one of the corresponding loops in the diagrams representing $`(\mathrm{ln}Z_I^{2loop})_{1PI}`$. After the integration in $`x`$, $`\tau `$ and in the fields in $`(\mathrm{ln}Z_I^{2loop})_{1PI}`$ and the differentiations above, we obtain the following expressions for the self-energies at one loop order
$`\mathrm{\Pi }_\pi (k_0,𝐤)={\displaystyle \underset{i=1}{\overset{4}{}}}\mathrm{\Pi }_{\pi i}=\lambda 5T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^3p}{(2\pi )^3}𝐃_{0\pi }(\omega _n,𝐩)}+\lambda T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^3p}{(2\pi )^3}𝐃_{0\sigma }(\omega _n,𝐩)}+`$ (57)
$`4\lambda ^2\nu ^2T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^3p}{(2\pi )^3}𝐃_{0\sigma }(\omega _{n+l},𝐩+𝐤)𝐃_{0\pi }(\omega _n,𝐩)}+g^2T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^3p}{(2\pi )^3}\mathrm{Tr}[\gamma ^5𝒟_{0\psi }(\omega _{n+l},𝐩+𝐤)\gamma ^5𝒟_{0\psi }(\omega _n,𝐩)]}`$ (58)
$`\mathrm{\Pi }_\sigma (k_0,𝐤)={\displaystyle \underset{i=1}{\overset{5}{}}}\mathrm{\Pi }_{\sigma i}=3\lambda T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^3p}{(2\pi )^3}𝐃_{0\sigma }(\omega _n,𝐩)}+3\lambda T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^3p}{(2\pi )^3}𝐃_{0\pi }(\omega _n,𝐩)}+`$ (59)
$`6\lambda ^2\nu ^2T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^3p}{(2\pi )^3}𝐃_{0\pi }(\omega _{n+l},𝐩+𝐤)𝐃_{0\pi }(\omega _n,𝐩)}18\lambda ^2\nu ^2T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^3p}{(2\pi )^3}𝐃_{0\sigma }(\omega _{n+l},𝐩+𝐤)𝐃_{0\sigma }(\omega _n,𝐩)}+`$ (60)
$`g^2T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^3p}{(2\pi )^3}\mathrm{Tr}[𝒟_{0\psi }(\omega _{n+l},𝐩+𝐤)𝒟_{0\psi }(\omega _n,𝐩)]}`$ (61)
$`\mathrm{\Sigma }(k_0,𝐤)={\displaystyle \underset{i=1}{\overset{2}{}}}\mathrm{\Sigma }_i=g^2T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^3p}{(2\pi )^3}𝐃_{0\sigma }(\omega _{n+l},𝐩+𝐤)𝒟_{0\psi }(\omega _n,𝐩)}+`$ (62)
$`3g^2T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^3p}{(2\pi )^3}𝐃_{0\pi }(\omega _{n+l},𝐩+𝐤)𝒟_{0\psi }(\omega _n,𝐩)}`$ (63)
The diagrams representing the pion, sigma and nucleon one-loop self-energies are drawn in figures 2,3 and 4.
We note here that we could get the same results for the self-energies directly applying Feynman rules to construct the diagrams with the appropriate substitutions: the $`\delta `$-function at each vertex is replaced with a Kronecker delta which imposes conservation of the discrete energy ($`k_0=i\omega _n`$), and round each loop of a thermal graph with $`\frac{d^4k}{(2\pi )^4}\frac{i}{\beta }_n\frac{d^3k}{(2\pi )^3}`$,which are the finite-temperature Feynman rules. When the summation over $`n`$ is performed, each graph of the self-energies is separated into two parts, namely a temperature independent part(at this stage), which is divergent, and a temperature dependent part containing the Bose-Einstein distribution in the case of bosons or the Fermi-Dirac distribution for fermions.
We will adopt the definition of mass at finite temperature as the real part of the pole of the corrected propagator at zero momentum $`(𝐤=0)`$. Thus, from eqs. (51) and (52) we have
$`𝐃_\pi (\omega _n,\left|𝐤\right|=0)^1=\omega _n^2+m_\pi ^2+\mathrm{\Pi }_\pi (\omega _n,\left|𝐤\right|=0)=0`$ (64)
$`k_{0,\pi }^2+m_\pi ^2+\mathrm{\Pi }_\pi (T,k_{0,\pi },\left|𝐤\right|=0)=0`$ (65)
$`𝐃_\sigma (\omega _n,\left|𝐤\right|=0)^1=\omega _n^2+m_\sigma ^2+\mathrm{\Pi }_\sigma (\omega _n,\left|𝐤\right|=0)=0`$ (66)
$`k_{0,\sigma }^2+m_\sigma ^2+\mathrm{\Pi }_\sigma (T,k_{0,\sigma },\left|𝐤\right|=0)=0`$ (67)
$`𝒟_\psi (\omega _nf,\left|𝐤\right|=0)^1=\gamma ^\mu k_\mu m_\psi +\mathrm{\Sigma }_s+\gamma ^\mu \mathrm{\Sigma }_\mu =0`$ (68)
$`k_{0,\psi }+\mathrm{\Sigma }_0(T,k_{0,\psi },\left|𝐤\right|=0)+\mathrm{\Sigma }_s(T,k_{0,\psi }\left|𝐤\right|=0)m_\psi =0,`$ (69)
where the arrow indicates an analytical continuation from discrete to continuous energies in Minkowski space. Hence the physical masses are the values of the $`k_{0;\pi ,\sigma ,\psi }`$ which are the zeros of the functions (64), (66) and (68) above, i.e., the location of the poles in the limit $`𝐤=0`$. The full self-energy expressions $`\mathrm{\Pi }_\pi `$, $`\mathrm{\Pi }_\sigma `$, $`\mathrm{\Sigma }_0`$ and $`\mathrm{\Sigma }_s`$ are shown explicitly in appendix B. The renormalization of the self-energy is studied in section VI. Through out this paper, we will use dimensional regularization, but omitting, for notational simplicity the factor $`\mu ^{2ϵ}`$ which multiplies $`\lambda `$. Since our calculations does not require traces involving an odd number of $`\gamma ^5`$ matrices, we use the definition of $`\gamma ^5`$ as in.
Since these expressions are self-consistent they have to be solved numerically. For each fixed temperature one finds a value of $`M`$ which satisfies the equations above. On the other hand, if one is interested only in the meson sector of the linear sigma model, the integrals in the self-energies could be evaluated exactly in the high temperature limit and at low frequency where the boson diagrams involving three-point vertices which are proportional to $`\lambda ^2\nu ^2`$ may be neglected. This is not consistent if one wants to study the behavior of the condensate $`\nu `$ and the particle masses in all ranges of temperatures. It is important to note that when $`c=0`$ the three-point vertex boson diagrams are significant in the region $`T_i(0)<T<T_f(\nu (T))`$ and when $`c0`$, $`\nu (T)0`$ for any finite value of T.
### B The second necessity: The resummation
The expressions for the self-energies appearing in eq. (64), (66) and (68) are functions of $`\omega _{\pi ,\sigma ,\psi }`$ which are expressed in terms of the mean-field masses. As we discussed, the meson masses become negative as the temperature increases. Thus, in the computation of the one-loop corrections, the masses running in the loops become tachyonic. A proper resummation of higher order loops is naturally necessary . Various resummation methods have been proposed in a tentative of curing the problem of the breaking down of the perturbative expansion at high temperature. In effective models, when a phase transition occurs, one can find tachyonic masses even below $`T_c`$. The O(N) linear $`\sigma `$ model which is one of the laboratory effective models employed to study QCD has been investigated by different authors using different techniques. One of these methods is the CJT formalism which provides for a consistent loop expansion of the effective potential in terms of the full propagator. The CJT formalism elegantly provides for a gap equation from stationarity conditions for the daisy and super-daisy effective potential. However some authors use this non-perturbative approach with an ansatz for the full (corrected) propagator in which the thermal corrections are momentum independent. These corrections are the finite piece of the divergent integral (which is temperature dependent through the gap equation for $`M`$) plus the finite explicit temperature dependent piece. This is the Hartree approximation, and means resuming only the “bubble diagrams” that are dominant at high temperatures. Another non-perturbative approach widely found in the literature is the large-$`N`$ approximation. The $`N\mathrm{}`$ limit facilitates the calculations, but can lead to inaccuracies . One must be careful in taking the large-$`N`$ limit since its truncation depends on the problem to be studied and the relevant value of $`N`$ . In this case the three-point vertex diagrams are omitted which in principle makes sense only in the $`N\mathrm{}`$ limit since these sunset diagrams are of order $`1/N`$. In these two kinds of treatment one can not study the bosons interacting with fermions (with the interactions of the linear sigma model) since the self-energy diagrams are momentum dependent which invalidates the ansatz cited above. It is worth to remember that the Hartree approximation does not satisfies the Goldstone theorem. This fact may be attributed to the non inclusion of these diagrams. Although the finite temperature mass in these approaches is the pole of the corrected propagator, it is not the true mass, since their $`\mathrm{\Pi }_{\pi ,\sigma }`$ are not the true one-loop self-energy functions (see discussion below). The corrections included only shift the masses. Once we are interested in the study of the masses behavior also in the range $`0<T<T_c`$ (if $`c=0`$) the three-particle vertex diagrams will not be neglected. The inclusion of these diagrams brings an additional complication since the self-energy now depends on the momentum.
### C A non-perturbative resummation method: The MSCR
Let us now introduce our procedure which resumms higher loop diagrams in the mean-field (tree-level) propagators. The method consist in recalculating the self-energy, in steps, using in each step the masses obtained in the previous one such that $`M_n^2=(A_n+1)M_{n1}^2+\mathrm{\Pi }(M_{n1})`$, where $`n`$ is the order of the non-perturbative correction and $`A_n`$ is the coefficient of the appropriate counterterm. With this procedure it is easier to identify and absorb the divergent parts of the self-energy in order to have finite gap-equations. The goal is to make renormalization possible since the masses which multiply the divergences are necessarily the same as in counterterms.
Application of MSCR
Here we apply the MSCR in the study of the chiral fermion meson model at finite temperature. The analysis of the problem has to be done carefully which will be divided into three regions.
Region I: The low temperature region
The first region is for $`0T<T^{}`$, where $`T^{}`$ is the temperature where $`\nu (T^{})=f_\pi `$. This implies that $`m_\pi ^2=0`$ and consequently the appearance of infrared divergences in the self-energy. So, in this region
Step 1:
Start with the mean-field effective Lagrangian where the condensate and the masses are given by:
$$\left[F(\omega _{\pi ,\sigma ,\psi }(m_{\pi ,\sigma ,\psi ;0}),T)+G\left(\omega _{\pi ,\sigma ,\psi }(m_{\pi ,\sigma ,\psi ;0})\right)\lambda (f_\pi ^2\nu ^2)\right]\nu c=0,$$
(70)
$$M_{\pi ,0}^2=m_\pi ^2=m^2+\lambda \nu ^2,$$
(71)
$$M_{\sigma ,0}^2=m_\sigma ^2=m^2+3\lambda \nu ^2,$$
(72)
$$M_{\psi ,0}=m_\psi =g\nu .$$
(73)
Step 2:
Evaluate the one-loop self-energy corrections to these masses from the equations presented in appendix B and define the condensate and the first order corrected masses as
$$\left[F(\omega _{\pi ,\sigma ,\psi }(M_{\pi ,\sigma ,\psi ,0}),T)+G\left(\omega _{\pi ,\sigma ,\psi }(M_{\pi ,\sigma ,\psi ,0})\right)\lambda (f_\pi ^2\nu ^2)\right]\nu c=0,$$
(74)
$`M_{\pi ,1}^2=M_{\pi ,0}^2+\mathrm{\Pi }_\pi (M_{\pi ,0}^2,M_{\sigma ,0}^2,M_{\psi ,0})=`$ (75)
$`(A_1+E_1+1)m^2+(\overline{A}_1+E_1+1)\lambda \nu ^2+\mathrm{\Pi }_\pi (M_{\pi ,0}^2,M_{\sigma ,0}^2,M_{\psi ,0})=`$ (76)
$`m^2+\lambda \nu ^2+\mathrm{\Pi }_\pi ^{Ren}(M_{\pi ,0}^2,M_{\sigma ,0}^2,M_{\psi ,0}),`$ (77)
$`M_{\sigma ,1}^2=M_{\sigma ,0}^2+\mathrm{\Pi }_\sigma (M_{\pi ,0}^2,M_{\sigma ,0}^2,M_{\psi ,0})=`$ (78)
$`(B_1+F_1+1)m^2+(F_1+1)\lambda \nu ^2+(\overline{B}_1+1)2\lambda \nu ^2+\mathrm{\Pi }_\sigma (M_{\pi ,0}^2,M_{\sigma ,0}^2,M_{\psi ,0})=`$ (79)
$`m^2+3\lambda \nu ^2+\mathrm{\Pi }_\sigma ^{Ren}(M_{\pi ,0}^2,M_{\sigma ,0}^2,M_{\psi ,0}),`$ (80)
$`M_{\psi ,1}=M_{\psi ,0}\mathrm{\Sigma }(M_{\pi ,0}^2,M_{\sigma ,0}^2,M_{\psi ,0})=`$ (81)
$`(C_1+1)g\nu \mathrm{\Sigma }(M_{\pi ,0}^2,M_{\sigma ,0}^2,M_{\psi ,0})=`$ (82)
$`g\nu \mathrm{\Sigma }^{Ren}(M_{\pi ,0}^2,M_{\sigma ,0}^2,M_{\psi ,0}).`$ (83)
where $`\mathrm{\Pi }_{\pi ,\sigma }=\mathrm{\Pi }_{\pi ,\sigma }(k_{0,\pi ,\sigma }=M_{\pi ,\sigma };𝐤=0)^0+\mathrm{\Pi }_{\pi ,\sigma }(T;k_{0,\pi ,\sigma }=0;𝐤=0)^\beta `$, $`A_1`$, $`\overline{A}_1`$, $`B_1`$, $`\overline{B}_1`$, $`C_1`$, $`E_1`$ and $`F_1`$ are the appropriate coefficients of the counterterms added to the mean-field effective Lagrangian needed to render the model finite up to this order, which are shown in section VI. The requirement that $`k_{0,\pi ,\sigma }=0`$ in $`\mathrm{\Pi }_{\pi ,\sigma }^\beta `$ excludes the possibility of tachyonic tree-level masses since thermal effects provides for the pion a non-zero width due to the Landau damping process. For the fermions we adopt the requirement that $`\mathrm{\Sigma }=\mathrm{\Sigma }(k_{0,\psi }=M_\psi ;𝐤=0)^0+\mathrm{\Sigma }(T;k_{0,\psi }=0;𝐤=0)^\beta `$ to prevent a similar consequence. The resummation has to be done exactly to avoid this problem. In this range there is no necessity of resummation since the masses running in the loops are positive, i.e., $`M_{\pi ,0}^2>0`$, $`M_{\sigma ,0}^2>0`$ and $`M_{\psi ,0}>0`$.
Renormalization:
The renormalization is done normally since the masses multiplying the divergences are the same as in the counterterms. The coefficients of the counterterms are found in section VI.
Goldstone’s Theorem:
In the exact chiral limit ($`c=0`$) and low temperature phase (where $`\nu 0`$) from eq. (64) at $`(k_00,\left|𝐤\right|=0)`$, we have
$`M_{\pi ,1}^2=M_{\pi ,0}^2+\mathrm{\Pi }_\pi (k_00,\left|𝐤\right|=0)=m^2+\lambda \nu ^2+\mathrm{\Pi }_\pi ^{Ren}(k_00,\left|𝐤\right|=0)=`$ (84)
$`\{{\displaystyle \frac{3\lambda }{(4\pi )^2}}[M_{\pi ,0}^2\mathrm{ln}\left({\displaystyle \frac{M_{\pi ,0}^2}{e\mu ^2}}\right)+M_{\sigma ,0}^2\mathrm{ln}\left({\displaystyle \frac{M_{\sigma ,0}^2}{e\mu ^2}}\right)]{\displaystyle \frac{8g^2}{(4\pi )^2}}M_{\psi ,0}^2\mathrm{ln}\left({\displaystyle \frac{M_{\psi ,0}^2}{e\mu ^2}}\right)+`$ (85)
$`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dpp^2}{2\pi ^2}}[3\lambda ({\displaystyle \frac{n_\sigma (M_{\sigma ,0})}{\omega _\sigma (M_{\sigma ,0})}}+{\displaystyle \frac{n_\pi (M_{\pi ,0})}{\omega _\pi (M_{\pi ,0})}})+8g^2{\displaystyle \frac{n_\psi (M_{\psi ,0})}{\omega _\psi (M_{\psi ,0})}}]\}+`$ (86)
$`{\displaystyle \frac{5\lambda }{(4\pi )^2}}M_{\pi ,0}^2\mathrm{ln}\left({\displaystyle \frac{M_{\pi ,0}^2}{e\mu ^2}}\right)+{\displaystyle \frac{5\lambda }{2}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dpp^2}{\pi ^2}}{\displaystyle \frac{n_\pi }{\omega _\pi }}+{\displaystyle \frac{\lambda }{(4\pi )^2}}M_{\sigma ,0}^2\mathrm{ln}\left({\displaystyle \frac{M_{\sigma ,0}^2}{e\mu ^2}}\right)+{\displaystyle \frac{\lambda }{2}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dpp^2}{\pi ^2}}{\displaystyle \frac{n_\sigma }{\omega _\sigma }}+`$ (87)
$`4{\displaystyle \frac{\lambda ^2\nu ^2}{(4\pi )^2}}{\displaystyle \frac{M_{\pi ,0}^2\mathrm{ln}\left(\frac{M_{\pi ,0}^2}{e\mu ^2}\right)M_{\sigma ,0}^2\mathrm{ln}\left(\frac{M_{\sigma ,0}^2}{e\mu ^2}\right)}{M_{\sigma ,0}^2M_{\pi ,0}^2}}2\lambda ^2\nu ^2{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dpp^2}{\pi ^2}}\left[{\displaystyle \frac{n_\pi }{\omega _\pi }}{\displaystyle \frac{1}{M_{\sigma ,0}^2M_{\pi ,0}^2}}{\displaystyle \frac{n_\sigma }{\omega _\sigma }}{\displaystyle \frac{1}{M_{\sigma ,0}^2M_{\pi ,0}^2}}\right]+`$ (88)
$`{\displaystyle \frac{8g^2}{(4\pi )^2}}M_{\psi ,0}^2\mathrm{ln}\left({\displaystyle \frac{M_{\psi ,0}^2}{e\mu ^2}}\right)+4g^2{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dpp^2}{\pi ^2}}{\displaystyle \frac{n_\psi }{\omega _\psi }}=0`$ (89)
since $`M_{\sigma ,0}^2M_{\pi ,0}^2=2\lambda \nu ^2`$ (in deriving eq.(84) we have used eq.(37)).
Region II: The intermediate temperature region
This is the region for $`T^{}TT_c`$ where the resummation is really necessary, since the masses in the loops are zero or tachyonic. The problem here is more complicated than in the other temperature regions since phase transition takes place. As is well known, around the critical temperature quantum fluctuations become essential and one loop corrections may not be enough. In fact, as we proceed to show here, this scheme (i.e., allowing only one loop corrections in the perturbative expansion) prevents us from adequately renormalizing the gap equations for the tree-level resummed masses and satisfying Goldstone’s theorem.
We show next that this is indeed the case and that a possible (but inconsistent !) way of achieving renormalization, would be e.g. to consider only the diagrams for which $`\mathrm{\Pi }_\pi =\mathrm{\Pi }_\sigma `$. This is very frequently implemented in the literature .
We will adopt the point of view that the MSCR fails in this region and higher order loop corrections in the perturbative expansion will be necessary at this stage. This is subject of current investigation. However we explicitly show were the problem appears.
Step 1:
Start with the mean-field effective Lagrangian where the condensate and the masses are given by:
$$\left[F(\omega _{\pi ,\sigma ,\psi }(m_{\pi ,\sigma ,\psi ;0}),T)+G\left(\omega _{\pi ,\sigma ,\psi }(m_{\pi ,\sigma ,\psi ;0})\right)\lambda (f_\pi ^2\nu ^2)\right]\nu c=0,$$
(90)
$$M_{\pi ,0}^2=m_\pi ^2=m^2+\lambda \nu ^2,$$
(91)
$$M_{\sigma ,0}^2=m_\sigma ^2=m^2+3\lambda \nu ^2,$$
(92)
$$M_{\psi ,0}=m_\psi =g\nu .$$
(93)
Step 2:
Evaluate the one-loop self-energy corrections to these masses from the equations presented in appendix B and define the condensate and the first order corrected masses as
$$\left[F(\omega _{\pi ,\sigma ,\psi }(M_{\pi ,\sigma ,\psi ,0}),T)+G\left(\omega _{\pi ,\sigma ,\psi }(M_{\pi ,\sigma ,\psi ,0})\right)\lambda (f_\pi ^2\nu ^2)\right]\nu c=0,$$
(94)
$`M_{\pi ,1}^2=M_{\pi ,0}^2+\mathrm{\Pi }_\pi (M_{\pi ,0}^2,M_{\sigma ,0}^2,M_{\psi ,0})=`$ (95)
$`(A_1+E_1+1)m^2+(\overline{A}_1+E_1+1)\lambda \nu ^2+\mathrm{\Pi }_\pi (M_{\pi ,0}^2,M_{\sigma ,0}^2,M_{\psi ,0})=`$ (96)
$`m^2+\lambda \nu ^2+\mathrm{\Pi }_\pi ^{Ren}(M_{\pi ,0}^2,M_{\sigma ,0}^2,M_{\psi ,0}),`$ (97)
$`M_{\sigma ,1}^2=M_{\sigma ,0}^2+\mathrm{\Pi }_\sigma (M_{\pi ,0}^2,M_{\sigma ,0}^2,M_{\psi ,0})=`$ (98)
$`(B_1+F_1+1)m^2+(F_1+1)\lambda \nu ^2+(\overline{B}_1+1)2\lambda \nu ^2+\mathrm{\Pi }_\sigma (M_{\pi ,0}^2,M_{\sigma ,0}^2,M_{\psi ,0})=`$ (99)
$`m^2+3\lambda \nu ^2+\mathrm{\Pi }_\sigma ^{Ren}(M_{\pi ,0}^2,M_{\sigma ,0}^2,M_{\psi ,0}),`$ (100)
$`M_{\psi ,1}=M_{\psi ,0}\mathrm{\Sigma }(M_{\pi ,0}^2,M_{\sigma ,0}^2,M_{\psi ,0})=`$ (101)
$`(C_1+1)g\nu \mathrm{\Sigma }(M_{\pi ,0}^2,M_{\sigma ,0}^2,M_{\psi ,0})=`$ (102)
$`g\nu \mathrm{\Sigma }^{Ren}(M_{\pi ,0}^2,M_{\sigma ,0}^2,M_{\psi ,0}).`$ (103)
At this stage of the application of the method, the renormalization is possible, Goldstone’s theorem is verified, but the tree-level masses are zero or tachyonic.
Step 3:
Now we take the masses computed in the previous step and improve the results defining a next-order non-perturbative correction. With this we get a new effective Lagrangian where the condensate and masses are given by
$$\left[F(\omega _{\pi ,\sigma ,\psi }(M_{\pi ,\sigma ,\psi ,1}),T)+G\left(\omega _{\pi ,\sigma ,\psi }(M_{\pi ,\sigma ,\psi ,1})\right)\lambda (f_\pi ^2\nu ^2)\right]\nu c=0,$$
(104)
$`M_{\pi ,2}^2=M_{\pi ,1}^2+\mathrm{\Pi }_\pi (M_{\pi ,1}^2,M_{\sigma ,1}^2,M_{\psi ,1})=`$ (105)
$`(A_2+E_2+1)m^2+(\overline{A}_2+E_2+1)\lambda \nu ^2+`$ (106)
$`(\overline{\overline{A}}_2+E_2+1)\mathrm{\Pi }_\pi ^{Ren}(M_{\pi ,0}^2,M_{\sigma ,0}^2,M_{\psi ,0})+\mathrm{\Pi }_\pi (M_{\pi ,1}^2,M_{\sigma ,1}^2,M_{\psi ,1})=`$ (107)
$`m^2+\lambda \nu ^2+\mathrm{\Pi }_\pi ^{Ren}(M_{\pi ,1}^2,M_{\sigma ,1}^2,M_{\psi ,1}){\displaystyle \frac{\lambda }{(4\pi )^2}}\mathrm{\Pi }_\sigma ^{Ren}(M_{\pi ,1}^2,M_{\sigma ,1}^2,M_{\psi ,1}){\displaystyle \frac{1}{\stackrel{~}{ϵ}}},`$ (108)
$`M_{\sigma ,2}^2=M_{\sigma ,1}^2+\mathrm{\Pi }_\sigma (M_{\pi ,1}^2,M_{\sigma ,1}^2,M_{\psi ,1})=`$ (109)
$`(B_2+F_2+1)m^2+(F_2+1)\lambda \nu ^2+(\overline{B}_2+1)2\lambda \nu ^2+`$ (110)
$`(\overline{\overline{B}}_2+F_2+1)\mathrm{\Pi }_\sigma ^{Ren}(M_{\pi ,0}^2,M_{\sigma ,0}^2,M_{\psi ,0})+\mathrm{\Pi }_\sigma (M_{\pi ,1}^2,M_{\sigma ,1}^2,M_{\psi ,1})=`$ (111)
$`m^2+3\lambda \nu ^2+\mathrm{\Pi }_\sigma ^{Ren}(M_{\pi ,1}^2,M_{\sigma ,1}^2,M_{\psi ,1}){\displaystyle \frac{3\lambda }{(4\pi )^2}}\mathrm{\Pi }_\pi ^{Ren}(M_{\pi ,1}^2,M_{\sigma ,1}^2,M_{\psi ,1}){\displaystyle \frac{1}{\stackrel{~}{ϵ}}},`$ (112)
$`M_{\psi ,2}=M_{\psi ,1}\mathrm{\Sigma }(M_{\pi ,1}^2,M_{\sigma ,1}^2,M_{\psi ,1})=`$ (113)
$`(C_2+1)g\nu (\overline{C}_2+1)\mathrm{\Sigma }^{Ren}(M_{\pi ,0}^2,M_{\sigma ,0}^2,M_{\psi ,o})\mathrm{\Sigma }(M_{\pi ,1}^2,M_{\sigma ,1}^2,M_{\psi ,1})=`$ (114)
$`g\nu \mathrm{\Sigma }^{Ren}(M_{\pi ,1}^2,M_{\sigma ,1}^2,M_{\psi ,1}).`$ (115)
It is shown that in the $`\lambda \varphi ^4`$ model, this (first) recalculation is equivalent to the sum of an infinite set of diagrams, namely the “daisy” sum or the set of ring . The coefficients of the temperature dependent mass counterterms $`\overline{\overline{A}}_2`$, $`\overline{\overline{B}}_2`$ and $`\overline{C}_2`$ are fixed in a manner to cancel not only divergences proportional to $`\mathrm{\Pi }_\pi ^{Ren}(M_{\pi ,0}^2,M_{\sigma ,0}^2,M_{\psi ,0})`$, $`\mathrm{\Pi }_\sigma ^{Ren}(M_{\pi ,0}^2,M_{\sigma ,0}^2,M_{\psi ,0})`$ and $`\mathrm{\Sigma }^{Ren}(M_{\pi ,0}^2,M_{\sigma ,0}^2,M_{\psi ,o})`$ respectively, but also these terms together. That is, at each stage of the procedure, for $`n>1`$, in the expressions for $`M_{\pi ,\sigma ,\psi ,n}`$, the self-energy $`\mathrm{\Pi }_{\pi ,\sigma }(M_{\pi ,\sigma ,\psi ,n2})`$ (or $`\mathrm{\Sigma }(M_{\pi ,\sigma ,\psi ,n2})`$) have to be cancelled to avoid overcounting of diagrams.
This shows explicitly that renormalization can not be performed within this approximation scheme. This is not surprising since in this temperature region quantum fluctuation may need a more thorough description. So, there is no reason to believe that only the “daisy” diagrams should be resummed at low and intermediate temperatures. In fact, the “daisy” graphs contributions are dominant at high temperature .
Renormalization:
Since in this region $`M_{\sigma ,n}^2M_{\pi ,n}^2=2\lambda \nu ^2+\mathrm{\Delta }\mathrm{\Pi }`$, where $`\mathrm{\Delta }\mathrm{\Pi }\mathrm{\Pi }_\sigma \mathrm{\Pi }_\pi `$, there is the presence of undesirable non-renormalizible terms. These terms are the last ones on the r.h.s. of equations (105) and (109) which come from equations (B2) and (B10) respectively and can not be absorbed in the counterterms.
Goldstone’s Theorem:
In this region, Goldstone’s theorem is satisfied only if $`\mathrm{\Delta }\mathrm{\Pi }0`$ and the contribution in (B3) is decoupled into integrals proportional to $`\lambda `$. This would assure the cancellation of $`M_\pi `$ at $`(k_00,\left|𝐤\right|=0)`$. The reason for this frustration is the same as for the lack of renormalizability.
Step 4: (applicable only in the case where $`\mathrm{\Delta }\mathrm{\Pi }=0`$. This guarantees that renormalization and Goldstone theorem can be satisfactorily implemented at each step. In particular this will be the case for the high temperature region as will be shown next.)
Proceeding with the iteration, in the limit $`n\mathrm{}`$ the masses $`M_n`$ have formally the same expressions as the masses $`M_{n1}`$ which are already renormalized. Thus, in this limit we will have,
$$\left[F(\omega _{\pi ,\sigma ,\psi }(M_{\pi ,\sigma ,\psi ,n}),T)+G\left(\omega _{\pi ,\sigma ,\psi }(M_{\pi ,\sigma ,\psi ,n})\right)\lambda (f_\pi ^2\nu ^2)\right]\nu c=0,$$
(116)
$$M_{\pi ,n}^2=m^2+\lambda \nu ^2+\mathrm{\Pi }_\pi ^{Ren}(M_{\pi ,n}^2,M_{\sigma ,n}^2,M_{\psi ,n}),$$
(117)
$$M_{\sigma ,n}^2=m^2+3\lambda \nu ^2+\mathrm{\Pi }_\pi ^{Ren}(M_{\pi ,n}^2,M_{\sigma ,n}^2,M_{\psi ,n}),$$
(118)
$$M_{\psi ,n}=g\nu \mathrm{\Sigma }^{Ren}(M_{\pi ,n}^2,M_{\sigma ,n}^2,M_{\psi ,n}).$$
(119)
At each intermediate step, in the loops we set $`k_{0\pi ,\psi }^2=M_{\pi ,\psi ,n1}^2`$ in the computation of $`M_{\pi ,\sigma ,\psi ,n}^2`$. This ensures the cancellation of the divergences in all stages of the process, since the masses in the counterterms will necessarily be the same as in the divergences. In the end, in the resulting equations of interest (to be solved nummericaly), $`K_{0\pi ,\psi }^2=M_{\pi ,\psi }^2`$ as it should. By our MSCR we have gotten a set of four coupled non-linear integral equations to be solved self-consistently, with finite gap equations for the tree-level masses, which read
$$M_\pi ^2=m^2+\lambda \nu ^2+\mathrm{\Pi }_\pi ^{Ren}(k_{0\pi }=M_\pi ,\left|𝐤\right|=0)$$
(120)
$$M_\sigma ^2=m^2+3\lambda \nu ^2+\mathrm{\Pi }_\pi ^{Ren}(k_{0\pi }=M_\pi ,\left|𝐤\right|=0)$$
(121)
$$M_\psi =g\nu \mathrm{\Sigma }_0^{Ren}(k_{0\psi }=M_\psi ,\left|𝐤\right|=0)\mathrm{\Sigma }_s^{Ren}(k_{0\psi }=M_\psi ,\left|𝐤\right|=0)$$
(122)
$`\nu \{m^2+\lambda \nu ^2+{\displaystyle \frac{3\lambda }{(4\pi )^2}}[M_\pi ^2\mathrm{ln}\left({\displaystyle \frac{M_\pi ^2}{e\mu ^2}}\right)+M_\sigma ^2\mathrm{ln}\left({\displaystyle \frac{M_\sigma ^2}{e\mu ^2}}\right)]{\displaystyle \frac{8g^2}{(4\pi )^2}}M_\psi ^2\mathrm{ln}\left({\displaystyle \frac{M_\psi ^2}{e\mu ^2}}\right)+`$ (123)
$`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dpp^2}{2\pi ^2}}[3\lambda ({\displaystyle \frac{n_\sigma (M_\sigma )}{\omega _\sigma (M_\sigma )}}+{\displaystyle \frac{n_\pi (M_\pi )}{\omega _\pi (M_\pi )}})+8g^2{\displaystyle \frac{n_\psi (M_\psi )}{\omega _\psi (M_\psi )}}]\}=c`$ (124)
where the renormalization scale $`\mu `$ can be determined by a physical condition. We choose $`\mu `$ such that the pion mass has the correct value at $`T=0`$.
Now we have to go back to real world encoded by eq. (64 to 68). In order to get finite physical masses from the pole of these equations it is necessary to sum and subtract the finite quantities $`\mathrm{\Pi }_\pi ^{Ren}`$ and $`\mathrm{\Sigma }^{Ren}`$ which will be regarded as mass parameters . This corresponds to the reorganization of the perturbative expansion. Now we rewrite eq. (24) as
$`={\displaystyle \frac{\lambda }{4}}(f_\pi ^2\nu ^2)^2+c\nu +\overline{\psi }[i\gamma ^\mu _\mu M_\psi ]\psi +`$ (125)
$`{\displaystyle \frac{1}{2}}[(\stackrel{}{\pi })^2M_\pi ^2\stackrel{}{\pi }^2+(\sigma )^2M_\sigma ^2\sigma ^2]g\overline{\psi }[\sigma +i\gamma ^5\stackrel{}{\pi }\stackrel{}{\tau })]\psi +`$ (126)
$`{\displaystyle \frac{\lambda }{4}}[(\stackrel{}{\pi }^2+\sigma ^2)^2+4\nu \sigma (\stackrel{}{\pi }^2+\sigma ^2)]+\mathrm{\Sigma }^{Ren}\overline{\psi }\psi +{\displaystyle \frac{1}{2}}\mathrm{\Pi }_\pi ^{Ren}(\stackrel{}{\pi }^2+\sigma ^2)+C.T.`$ (127)
The last two terms on the third line of eq.(127) must be considered as extra “interaction” terms and will naturally be present in $`\mathrm{ln}Z_I`$, eq.(49). The counterterm structure of eq.(127) is the same as the one present in the method for the pion and fermion, differing only by numerical factors in the case of the sigma mass renormalization. By eqs.(55) and (56) the extra contribution to the self-energy read
$$\mathrm{\Pi }_\pi ^{extra}=\mathrm{\Pi }_\sigma ^{extra}=\mathrm{\Pi }_\pi ^{Ren}$$
(128)
$$\mathrm{\Sigma }^{extra}=\mathrm{\Sigma }^{Ren}$$
(129)
where we have defined
$$\mathrm{\Pi }_\pi ^{Ren}\mathrm{\Pi }_\pi (k_{0,\pi }=M_\pi ,𝐤=0)^0+\mathrm{\Pi }_\pi (T,k_0=0,𝐤=0)^\beta $$
(130)
$$\mathrm{\Sigma }^{Ren}\mathrm{\Sigma }(k_0=M_\psi ,𝐤=0)^0+\mathrm{\Sigma }(T,k_0=0,𝐤=0)^\beta $$
(131)
As a result, the final resummmed tree-level masses (eqs.(120), (121) and (122)), may be used in eqs. (64), (66) and (68) if one wants to study for instance spectral functions as the authors of, or decay width as in .
Goldstone’s Theorem:
If the sunset type graphs are neglected, the self-energy function at one-loop is not complete, despite of the fact that the definition of masses as poles of propagators at zero momentum is still valid. Now, having the result for the full one-loop (and higher-order loops contributions from the resumation) self-energy function we can test algebraically the fulfillment of Goldstone’s theorem in the exact chiral limit ($`c=0`$) and low temperature phase (where $`\nu 0`$). From eq. (64) at $`(k_00,\left|𝐤\right|=0)`$, we have
$`M_\pi ^2+\mathrm{\Pi }_\pi ^{total}=M_\pi ^2+\mathrm{\Pi }_\pi (k_00,\left|𝐤\right|=0)+\mathrm{\Pi }_\pi ^{extra}=m^2+\lambda \nu ^2+\mathrm{\Pi }_\pi ^{Ren}(k_00,\left|𝐤\right|=0)=`$ (132)
$`\{{\displaystyle \frac{3\lambda }{(4\pi )^2}}[M_\pi ^2\mathrm{ln}\left({\displaystyle \frac{M_\pi ^2}{e\mu ^2}}\right)+M_\sigma ^2\mathrm{ln}\left({\displaystyle \frac{M_\sigma ^2}{e\mu ^2}}\right)]{\displaystyle \frac{8g^2}{(4\pi )^2}}M_\psi ^2\mathrm{ln}\left({\displaystyle \frac{M_\psi ^2}{e\mu ^2}}\right)+`$ (133)
$`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dpp^2}{2\pi ^2}}[3\lambda ({\displaystyle \frac{n_\sigma (M_\sigma )}{\omega _\sigma (M_\sigma )}}+{\displaystyle \frac{n_\pi (M_\pi )}{\omega _\pi (M_\pi )}})+8g^2{\displaystyle \frac{n_\psi (M_\psi )}{\omega _\psi (M_\psi )}}]\}+`$ (134)
$`{\displaystyle \frac{5\lambda }{(4\pi )^2}}M_\pi ^2\mathrm{ln}\left({\displaystyle \frac{M_\pi ^2}{e\mu ^2}}\right)+{\displaystyle \frac{5\lambda }{2}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dpp^2}{\pi ^2}}{\displaystyle \frac{n_\pi }{\omega _\pi }}+{\displaystyle \frac{\lambda }{(4\pi )^2}}M_\sigma ^2\mathrm{ln}\left({\displaystyle \frac{M_\sigma ^2}{e\mu ^2}}\right)+{\displaystyle \frac{\lambda }{2}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dpp^2}{\pi ^2}}{\displaystyle \frac{n_\sigma }{\omega _\sigma }}+`$ (135)
$`4{\displaystyle \frac{\lambda ^2\nu ^2}{(4\pi )^2}}{\displaystyle \frac{M_\pi ^2\mathrm{ln}\left(\frac{M_\pi ^2}{e\mu ^2}\right)M_\sigma ^2\mathrm{ln}\left(\frac{M_\sigma ^2}{e\mu ^2}\right)}{M_\sigma ^2M_\pi ^2}}2\lambda ^2\nu ^2{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dpp^2}{\pi ^2}}\left[{\displaystyle \frac{n_\pi }{\omega _\pi }}{\displaystyle \frac{1}{M_\sigma ^2M_\pi ^2}}{\displaystyle \frac{n_\sigma }{\omega _\sigma }}{\displaystyle \frac{1}{M_\sigma ^2M_\pi ^2}}\right]+`$ (136)
$`{\displaystyle \frac{8g^2}{(4\pi )^2}}M_\psi ^2\mathrm{ln}\left({\displaystyle \frac{M_\psi ^2}{e\mu ^2}}\right)+4g^2{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dpp^2}{\pi ^2}}{\displaystyle \frac{n_\psi }{\omega _\psi }}=0`$ (137)
since $`M_\sigma ^2M_\pi ^2=2\lambda \nu ^2`$ (in deriving eq.(132) we have used eq.(123)).
To our opinion the fullfilment of Goldstone’s theorem is ultimately related to the preservation of the relation imposed by chiral symmetry to the tree-level masses. Moreover, it is crucial to keep all diagrams of a given order. This is due to the fact that, strictly speaking, a loop expansion is an expansion in powers of the Lagrangian. As discussed in in order to respect the symmetries of the Lagrangian, one must retain all diagrams to the given number of loops.
Then, in the absence of the explicit chiral symmetry breaking term, one has,
for $`0<T<T_c`$
$`M_\pi ^2=0`$ (138)
$`M_\sigma ^2=2\lambda \nu ^2`$ (139)
for $`T=T_c`$
$$M_\pi ^2=M_\sigma ^2=0$$
(140)
and for $`T>T_c`$
$$M_\pi ^2=M_\sigma ^2=m^2+3\lambda _0^{\mathrm{}}\frac{dpp^2}{\pi ^2}\frac{n_b}{\omega _b}+4g^2_0^{\mathrm{}}\frac{dpp^2}{\pi ^2}\frac{n_f}{\omega _f}$$
(141)
which shows chiral symmetry restoration. Here b stands for bosons and f for fermions. We could interpret the result in the r.h.s. of eq. (141) as if each independent pion effectively “sees” one sigma and the other two pions and four fermions (since $`\mu `$, the chemical potential, here is zero). On the other hand the sigma “sees” the three pions and four fermions. This equation serves to define the critical temperature in which the common masses of the particles vanish. In the high-temperature limit of these integrals, we find that $`T_c^2=\frac{2f_\pi ^2}{(1+\frac{2g^2}{3\lambda })}`$, as predicted by the mean-field analysis in eq.(45).
Region III: The high temperature region
This is the region of high temperatures, $`TT_c`$, if $`c=0`$ and $`\nu =0`$ or $`TT_i`$, where $`T_i`$ is defined as a “inflexion” temperature, for the case $`c0`$ and $`\nu <<f_\pi `$ such that $`M_{\pi ,0}^2M_{\sigma ,0}^2=m^2`$.
$`M_{\pi ,1}^2=M_{\pi ,0}^2+\mathrm{\Pi }_\pi (M_{\pi ,0}^2,M_{\sigma ,0}^2,M_{\psi ,0})=`$ (142)
$`m^2+\mathrm{\Pi }_\pi (M_{\pi ,0}^2,M_{\sigma ,0}^2,M_{\psi ,0}),`$ (143)
$`M_{\sigma ,1}^2=M_{\sigma ,0}^2+\mathrm{\Pi }_\sigma (M_{\pi ,0}^2,M_{\sigma ,0}^2,M_{\psi ,0})=`$ (144)
$`m^2+\mathrm{\Pi }_\sigma (M_{\pi ,0}^2,M_{\sigma ,0}^2,M_{\psi ,0})=`$ (145)
$`m^2+\mathrm{\Pi }_\pi (M_{\pi ,0}^2,M_{\sigma ,0}^2,M_{\psi ,0})=M_{\pi ,1}^2M_1^2,`$ (146)
$`M_{\psi ,1}=M_{\psi ,0}\mathrm{\Sigma }(M_{\pi ,0}^2,M_{\sigma ,0}^2,M_{\psi ,0})=`$ (147)
$`g\nu \mathrm{\Sigma }(M_{\pi ,0}^2,M_{\sigma ,0}^2,M_{\psi ,0}).`$ (148)
Note that the pion and sigma masses become degenerate and the problem encountered in the previous region ($`T^{}TT_c`$) is no longer here since $`\mathrm{\Delta }\mathrm{\Pi }=0`$ in this region of temperatures. In this case, the masses in the loops can be neglected, and we have
$`M_1^2=(A_1+1)M_0^2+\mathrm{\Pi }(M_0)=m^2+{\displaystyle \frac{\lambda }{2}}\left(1+{\displaystyle \frac{2g^2}{3\lambda }}\right)T^2=`$ (149)
$`\lambda f_\pi ^2\left[{\displaystyle \frac{T^2}{T_c^2}}1\right]`$ (150)
If we set $`g=0`$ these results agrees with the ones obtained by Bochkarev and Kapusta.
Following the iterations, we find for the $`n`$-th iterated mass
$`M_n^2=(A_n+1)M_{n1}^2+\mathrm{\Pi }(M_{n1})=`$ (151)
$`m^2+{\displaystyle \frac{\lambda }{2}}\left(1+{\displaystyle \frac{2g^2}{3\lambda }}\right)T^2\left[1{\displaystyle \frac{3}{\pi T}}M_{n1}\right]`$ (152)
In the limit $`n\mathrm{}`$, we get
$$M^2=m^2+\frac{\lambda }{2}\left(1+\frac{2g^2}{3\lambda }\right)T^2\left[1\frac{3}{\pi T}M\right]$$
(153)
which can be easily solved for $`M`$,
$$M=\left[\left(\frac{3\lambda f_\pi ^2}{2\pi }\frac{T}{T_c^2}\right)^2+\lambda f_\pi ^2\left(\frac{T^2}{T_c^2}1\right)\right]^{\frac{1}{2}}\frac{3\lambda f_\pi ^2}{2\pi }\frac{T}{T_c^2}.$$
(154)
For $`T>>T_c`$
$$M^2=\frac{\lambda }{2}\left(1+\frac{2g^2}{3\lambda }\right)T^2$$
(155)
## V The Massless $`\lambda \varphi ^4`$ At High-Temperature
Now we apply the MSCR to study a very popular model: the massless $`\lambda \varphi ^4`$ model in the weak coupling limit
$$=\frac{1}{2}(_\mu \varphi )^2\frac{\lambda }{4!}\varphi ^4$$
(156)
$`M_0=0`$ (157)
$`M_1^2=M_0^2+\mathrm{\Pi }(M_0)={\displaystyle \frac{\lambda T^2}{24}}`$ (158)
at this stage of the procedure there is no necessity of adding counterterms since up to this order there are no ultraviolet divergences in dimensional regularization. Here $`\mathrm{\Pi }`$ is the 1PI one-loop self-energy to lowest order, namely the “bubble” of Fig. 5(a).
$`M_2^2=M_1^2+\mathrm{\Pi }(M_1)=(A_2+1)\mathrm{\Pi }^{Ren}(M_0)+\mathrm{\Pi }(M_1)=`$ (159)
$`{\displaystyle \frac{\lambda T^2}{24}}\left(1{\displaystyle \frac{3M_1}{\pi T}}\right)+O(\lambda ^2\mathrm{ln}\lambda )=`$ (160)
$`M_1^2\left[13\left({\displaystyle \frac{\lambda }{24\pi ^2}}\right)^{\frac{1}{2}}\right]+O(\lambda ^2\mathrm{ln}\lambda )`$ (161)
with the result that this correction to the mass is of order $`\lambda ^{\frac{3}{2}}`$, which is an signature of the non-perturbative resummation. The temperature dependent counterterm is fixed so as to cancel the divergence and avoid overcounting of diagrams, as explained before. So, $`A_2=1+\frac{\lambda }{2(4\pi )^2}\frac{1}{ϵ}`$. The diagrams used, in a given number of loops, in any resummation method must be the same in all stages of the process. What changes is the masses running in the loops at each iteration. This is because one must keep the same fundamental theory in the recalculation of the self-energy. The result shown in eq. (159) is in agreement with the one obtained by Parwani’s resummed perturbative expansion (see eq.(2.12) of his paper). The second iteration corrected mass, $`M_2`$, which was obtained in our method evaluating Fig. 5(a) with $`M_1`$ in that loop can equivalently be achieved calculating the “daisy” sum, that is a summation of the infinite set of “daisy” diagrams of Fig. 5(b) with $`M_0`$ in the loops. In this case all “daisy” types diagrams are IR-divergent since $`M_0=0`$, but their sum is IR-finite.
Continuing the iterations, we find for the next correction
$`M_3^2=M_2^2+\mathrm{\Pi }(M_2)=`$ (162)
$`{\displaystyle \frac{\lambda T^2}{24}}\left[1{\displaystyle \frac{3M_1}{\pi T}}\left(1{\displaystyle \frac{3M_1}{\pi T}}\right)^{\frac{1}{2}}\right].`$ (163)
When $`\lambda <<1`$ we get
$$M_3^2=\frac{\lambda T^2}{24}\left[13\left(\frac{\lambda }{24\pi ^2}\right)^{\frac{1}{2}}+\frac{9}{2}\left(\frac{\lambda }{24\pi ^2}\right)\right]$$
(164)
and for the n-th iteration, we obtain
$$M_n^2=\frac{\lambda T^2}{24}\left\{1+\underset{j=1}{\overset{n}{}}\frac{1}{2^{j1}}\left[3\left(\frac{\lambda }{24\pi ^2}\right)^{\frac{1}{2}}\right]^j\right\}$$
(165)
The “superdaisy” sum corresponds to the limit $`n\mathrm{}`$ of equation (165) and it can be summed up (for $`\lambda <<1`$) to give
$$M^2=\frac{\lambda T^2}{24}\left[\frac{1\frac{3}{2}\left(\frac{\lambda }{24\pi ^2}\right)^{\frac{1}{2}}}{1+\frac{3}{2}\left(\frac{\lambda }{24\pi ^2}\right)^{\frac{1}{2}}}\right]$$
(166)
## VI Renormalization
### A Determination of the counter-terms
The divergences are regulated via dimensional rugularization. To renormalize the divergences, we use the Minimal Subtraction scheme where only the poles are eliminated by the appropriate counterterms. The first-order parameters of the temperature-dependent counterterms read
$$A_1=\frac{6\lambda }{(4\pi )^2}\frac{1}{\stackrel{~}{ϵ}},\overline{A}_1=\frac{12\lambda }{(4\pi )^2}\frac{1}{\stackrel{~}{ϵ}},E_1=\frac{4g^2}{(4\pi )^2}\frac{1}{\stackrel{~}{ϵ}},$$
(167)
with $`\frac{1}{\stackrel{~}{ϵ}}\frac{2}{4d}\gamma +\mathrm{log}(4\pi )`$, where $`\gamma `$ is the Euler constant.
$$B_1=\frac{6\lambda }{(4\pi )^2}\frac{1}{\stackrel{~}{ϵ}},\overline{B}_1=\frac{6\lambda }{(4\pi )^2}\frac{1}{\stackrel{~}{ϵ}},F_1=\frac{4g^2}{(4\pi )^2}\frac{1}{\stackrel{~}{ϵ}}$$
(168)
$$C_1=\frac{8g^2}{(4\pi )^2}\frac{1}{\stackrel{~}{ϵ}}$$
(169)
$$D_{1,1}=8\left[\frac{1}{64\pi ^2}\frac{1}{\stackrel{~}{ϵ}}\right],D_{2,1}=3\left[\frac{1}{64\pi ^2}\frac{1}{\stackrel{~}{ϵ}}\right],D_{3,1}=\frac{1}{64\pi ^2}\frac{1}{\stackrel{~}{ϵ}}.$$
(170)
For all steps we always have
$$A_n=A_1,\overline{A}_n=\overline{A}_1,E_n=E_1,B_n=B_1,\overline{B}_n=\overline{B}_1,F_n=F_1,C_n=C_1,D_{1,2,3,n}=D_{1,2,3,1}$$
(171)
and for $`n>1`$
$$\overline{\overline{A}}_n=1+A_1,\overline{\overline{B}}_n=1+B_1,\overline{\overline{C}}_n=1C_1$$
(172)
### B Comments related with the presence of the fermions in the game
We must remark that, on eqs. (B6), (B15) and (B18) of appendix B, the following terms: $`\frac{8g^2}{(4\pi )^2}m_\psi ^2\frac{1}{\stackrel{~}{ϵ}}`$, $`\frac{8g^2}{(4\pi )^2}3m_\psi ^2\frac{1}{\stackrel{~}{ϵ}}`$, $`\frac{g^2}{(4\pi )^2}\left[\frac{m_\sigma ^2}{2k_0}\right]\frac{1}{\stackrel{~}{ϵ}}`$ and $`\frac{g^2}{(4\pi )^2}\left[3\frac{m_\pi ^2}{2k_0}\right]\frac{1}{\stackrel{~}{ϵ}}`$ should be neglected. As stated by the authors of and remarked by the authors of , these terms will be canceled by contributions from higher order loops. Since we are concerned only about the one-loop approximation, we do not have to worry about them. Nevertheless, in the $`O(4)`$ linear sigma model, i.e., when $`g=0`$, none of the above terms will be present, and our model will be order by order renormalizable in the regions of validity of the MSCR. This occurs because our tree-level resummed masses are related by a symmetry relation that always guarantees the cancellation of the UV divergences.
## VII Numerical Analysis
In this section, we present numerical solutions of the gap equations for the tree-level meson and fermion masses and the condensate derived in section IV C including all diagrams which belong to the one loop order.
As an approximation, only for the sake of obtaining continuous curves, in the numerical evaluation we considered $`\mathrm{\Delta }\mathrm{\Pi }=0`$ also in the intermediate temperature region. Rigorously speaking, the curves should only be trusted in the low and high temperature regions. Figure 6 shows the tree-level resummed meson masses, eqs. (120) and (121), as functions of the temperature. We show in Figure 7 the tree-level fermion resummed mass, eq. (122), as a function of temperature. The tree-level masses behavior exhibit the fact that the MSCR has solved the problem of tachyonic masses. In Figure 8 the chiral condensate $`\nu `$, eq.(123), as a function of temperature is shown whereas in Figure 9 the condensate is ploted in the case $`M_\pi =0`$.
Since at low temperatures the condensate dominates, the mesons masses suffers its influence in this region. The sigma mass decreases and they approach each other to become degenerate in a temperature of about $`300MeV`$. This confirms the results we found in a phenomenological approach to the linear sigma model.
The condensate is a slowly decreasing function of the temperature, which is a signature of the order parameter when the symmetry breaking term is present. The qualitative behavior of the results shown in Figs. 6 and 8 can be compared with the ones obtained by Chiku and Hatsuda since OPT also sums three-point vertex diagrams, as our method does. Some differences may be attributed to the incorporation of the fermions, as performed in our method. In the absence of the chiral symmetry breaking term, i.e., when $`c=0`$, the non-vanishing solutions of the extremum condition, eq. (36), are obtained numerically by equation (123) with $`M_\pi ^2=0`$, $`M_\sigma ^2=2\lambda \nu ^2`$ and $`M_\psi =M_\psi (M_\pi ^2=0,M_\sigma ^2=2\lambda \nu ^2)`$. The solution is depicted in Fig. 9 and gives an indication of first order phase transition. This result agrees with the predictions of first order phase transition found in previous analysis by Roh and Matsui, Petropoulos, Chiku and Hatsuda, Randrup and Bilic. Of course we have to bear in mind that our result is at one loop order in the perturbative expansion. It may well be that near the critical temperature higher order corrections become crucial and may change the order of the phase transition.
The tree fermion mass, Fig. 7, does not become zero when chiral symmetry is restored and $`\nu 0`$ since we considered contributions from the mesons, given by eq. (62). On the contrary, when the temperature is $`200MeV`$ these contributions dominate the variable $`\nu `$ and the fermion mass increases with temperature. The behavior of the fermion mass is in agreement with the results found by Panda in for the quark meson coupling model.
## VIII Concluding Remarks
In this paper, we presented a modified self-consistent resummation (MSCR) at finite temperature. Results for the chiral fermion meson model and the massless $`\lambda \varphi ^4`$ model in the weak coupling limit were obtained and analyzed. We have shown that our procedure properly resumes higher order terms which cures the problem of the breakdown of the perturbative expansion.
We have also shown that the MSCR, when applied to the study of the chiral fermion meson model, has the essential features which leads to the satisfaction of Goldstone’s theorem and renormalization of the UV divergences, in the low and high temperature regions. We have explicitly shown that the scheme breaks down around $`T_c`$ i.e., in the region of intermediate temperatures. The application of the MSCR in these three physically different regions (low, intermediate and high temperatures) revealed a source of mistakes usually found in the literature, that is to treat all ranges of temperatures in the same way. It is naive to expect that the same approximations which is valid, e.g., for high temperatures would be enough in the intermediate temperature region, since quantum fluctuations are known to play a major role there.
This division was essential to identify the regions where higher order terms and resummation are crucial. It is valid to remember that even when higher order loops are taken into account, the resummation is still necessary since the tree-level masses will become tachyonic even below the critical temperature (in theories with spontaneous symmetry breaking) and break the perturbative expansion. This breakdown of the perturbative expansion can also happen in massless field theories, like QCD, due to the appearance of infrared divergences. As we discussed, the breakdown of perturbative expansion in finite temperature field theory requires resummation techniques as the MSCR to recover the reliability of perturbative expansion.
In each region renormalization and satisfaction of Goldstone’s theorem were discussed in detail. In our study, we have also addressed a usually avoided point: the inclusion of the fermions. Finally, we have re-examined the chiral phase transition in static equilibrium in terms of the linear sigma model with our MSCR.
The gap equations for the tree-level masses, are constructed by our method and in the effective Lagrangian they are renormalized. For the particular case of intermediate temperatures region, the gap equations would be renormalized in the (reorganized) effective Lagrangian only if $`\mathrm{\Delta }\mathrm{\Pi }=\mathrm{\Pi }_\sigma \mathrm{\Pi }_\pi =0`$. In most of the approximations found in the literature, the gap equations are reached by some technique or via some ad-hoc procedure but the Lagrangian is yet the original one. This makes the renormalization process non-trivial, unless a finite cut-off is used and the theory is treated as an effective model . As pointed out by Chiku and Hatsuda , the resummation must be done also in the counterterms, which is essential to show the renormalization.
At this point, it is extremely worth emphasizing that, although one has the freedom of adding and subtracting mass parameters to the Lagrangian, in this case they can not be completely arbitrary. If the mass parameters introduced were different for the pion and sigma fields (i.e., $`\frac{1}{2}M_1\sigma ^2+\frac{1}{2}M_2\stackrel{}{\pi }^2`$ and, of course, the same quantities subtracted, with $`M_1M_2`$), neither the $`O(4)`$ linear $`\sigma `$ model is renormalizable in any given order nor Goldstone s theorem is satisfied. This will happen even if the mass parameters are determined by some physical condition as FAC or principle of minimal sensitivity (PMS). So, the most important fact behind the fulfillment of Goldstone’s theorem and renormalizability of theories with SSB is the chiral symmetry that must dictate which mass parameter should be introduced to the Lagrangian.
## Acknowledgements
One of the authors (H. C. G. Caldas) thanks the hospitality given by the Nuclear Theory group during his visit at the University of Minnesota were part of this work was done. He is gratefully indebted to Professors J. I. Kapusta and P. J. Ellis for various helpful advices and enlightening discussions. He is also grateful to Dr. M. Hott for valuables conversations about this problem, Professor A. Das for comments concerning the effective potential and finally Drs. J. Lenaghan and S. Chiku for useful e-conversations. The authors would like to thank B. Hiller and A. Blin for useful comments on the subject of the paper. H. C. G. Caldas thanks the generous support provided by the Faculty of UFMG and FUNREI.
## A Renormalization of the effective potential
As mentioned earlier the vacuum contribution to $`\mathrm{\Omega }_1(T,\nu )`$ is divergent and requires renormalization. In this subsection we also use dimensional regularization in the computation of the effective potential.
$$\mathrm{\Omega }_1^0(m)\frac{d^3p}{(2\pi )^3}\frac{\omega }{2}$$
(A1)
$$\frac{\mathrm{\Omega }_1^0(m)}{m}=m\frac{d^3p}{(2\pi )^3}\frac{1}{2\omega }=mL(m)$$
(A2)
where $`L(m)`$ is the usual zero temperature loop integral
$$L(m)\frac{d^4p}{(2\pi )^4}\frac{1}{p^2+m^2}=\frac{d^4p}{(2\pi )^4}\frac{1}{p_4^2+𝐩^2+m^2},$$
(A3)
with $`d^4p=dp_4d^3p`$ being the four Euclidean momentum.
The divergent integral $`L(m)`$ can be evaluated in the standard manner
$$\frac{m^2}{(4\pi )^2}\left[\frac{1}{\stackrel{~}{ϵ}}1+\mathrm{ln}\left(\frac{m^2}{\mu ^2}\right)\right].$$
(A4)
The quantity $`\mathrm{\Omega }_1^0(m)`$ is then obtained with the integration of $`mL\left(m\right)`$
$$\mathrm{\Omega }_1^0(m)=\frac{m^4}{64\pi ^2}(\mathrm{ln}\frac{m^2}{\mu ^2}\frac{3}{2}\frac{1}{\stackrel{~}{ϵ}})$$
(A5)
With this expression we can find the zero temperature effective potential,
$$\mathrm{\Omega }_1^0(\nu )=\frac{m_\sigma ^4}{64\pi ^2}\frac{1}{\stackrel{~}{ϵ}}+\frac{m_\sigma ^4}{64\pi ^2}(\mathrm{ln}\frac{m_\sigma ^2}{\mu ^2}\frac{3}{2})+3(m_\sigma m_\pi )8(m_\sigma m_\psi )$$
(A6)
The renormalization of the thermodynamical potential at the end amounts to the determination of the parameters $`D_{1,2,3}`$,
$$D_1=8\left[\frac{1}{64\pi ^2}\frac{1}{\stackrel{~}{ϵ}}\right],$$
(A7)
$$D_2=3\left[\frac{1}{64\pi ^2}\frac{1}{\stackrel{~}{ϵ}}\right],$$
(A8)
$$D_3=\frac{1}{64\pi ^2}\frac{1}{\stackrel{~}{ϵ}}.$$
(A9)
## B One-loop self-energy at finite temperature
At zero momentum the expressions for the self-energies are given by
$$\mathrm{\Pi }_{\pi 1}(k_0,\left|𝐤\right|=0)=\mathrm{\Pi }_{\pi 1}^0+\mathrm{\Pi }_{\pi 1}^\beta =\frac{5\lambda }{(4\pi )^2}m_\pi ^2\left[\frac{1}{\stackrel{~}{ϵ}}1+\mathrm{ln}\left(\frac{m_\pi ^2}{\mu ^2}\right)\right]+\frac{5\lambda }{2}_0^{\mathrm{}}\frac{dpp^2}{\pi ^2}\frac{n_\pi }{\omega _\pi }$$
(B1)
$$\mathrm{\Pi }_{\pi 2}(k_0,\left|𝐤\right|=0)=\mathrm{\Pi }_{\pi 2}^0+\mathrm{\Pi }_{\pi 2}^\beta =\frac{\lambda }{(4\pi )^2}m_\sigma ^2\left[\frac{1}{\stackrel{~}{ϵ}}1+\mathrm{ln}\left(\frac{m_\sigma ^2}{\mu ^2}\right)\right]+\frac{\lambda }{2}_0^{\mathrm{}}\frac{dpp^2}{\pi ^2}\frac{n_\sigma }{\omega _\sigma }$$
(B2)
$`\mathrm{\Pi }_{\pi 3}(k_0,\left|𝐤\right|=0)=\mathrm{\Pi }_{\pi 3}^0+\mathrm{\Pi }_{\pi 3}^\beta =`$ (B3)
$`{\displaystyle \frac{4\lambda ^2\nu ^2}{(4\pi )^2}}\left[{\displaystyle \frac{1}{\stackrel{~}{ϵ}}}1+\mathrm{ln}\left({\displaystyle \frac{m_\pi ^2}{\mu ^2}}\right)+{\displaystyle \frac{k_0^2+m_\sigma ^2m_\pi ^2}{2k_0^2}}\mathrm{ln}\left({\displaystyle \frac{m_\pi ^2}{m_\sigma ^2}}\right)+{\displaystyle \frac{\sqrt{\mathrm{\Delta }_3}}{k_0^2}}(\mathrm{\Delta }_4+\mathrm{\Delta }_5)\right]+`$ (B4)
$`2\lambda ^2\nu ^2{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dpp^2}{\pi ^2}}\left[{\displaystyle \frac{n_\pi }{\omega _\pi }}{\displaystyle \frac{k_0^2+m_\sigma ^2m_\pi ^2}{(k_0^2+m_\sigma ^2m_\pi ^2)^22k_0^2(\omega _\pi ^2+\omega _\sigma ^2)}}+{\displaystyle \frac{n_\sigma }{\omega _\sigma }}{\displaystyle \frac{k_0^2+m_\pi ^2m_\sigma ^2}{(k_0^2+m_\pi ^2m_\sigma ^2)^22k_0^2(\omega _\pi ^2+\omega _\sigma ^2)}}\right]`$ (B5)
$`\mathrm{\Pi }_{\pi 4}(k_0,\left|𝐤\right|=0)=\mathrm{\Pi }_{\pi 4}^0+\mathrm{\Pi }_{\pi 4}^\beta ={\displaystyle \frac{8g^2}{(4\pi )^2}}\left[m_\psi ^2{\displaystyle \frac{k_0^2}{2}}\right]{\displaystyle \frac{1}{\stackrel{~}{ϵ}}}+`$ (B6)
$`{\displaystyle \frac{8g^2}{(4\pi )^2}}\left[\left(m_\psi ^2{\displaystyle \frac{k_0^2}{2}}\right)\mathrm{ln}\left({\displaystyle \frac{m_\psi ^2}{\mu ^2}}\right)+\mathrm{\Delta }_2\sqrt{\mathrm{\Delta }_1}K_0^2\right]+`$ (B7)
$`4g^2{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dpp^2}{\pi ^2}}{\displaystyle \frac{n_\psi }{\omega _\psi }}\left[1+{\displaystyle \frac{k_0^2}{4\omega _{\psi ^2}k_0^2}}\right]`$ (B8)
$$\mathrm{\Pi }_{\sigma 1}(k_0,\left|𝐤\right|=0)=\mathrm{\Pi }_{\sigma 1}^0+\mathrm{\Pi }_{\sigma 1}^\beta =\frac{3\lambda }{(4\pi )^2}m_\sigma ^2\left[\frac{1}{\stackrel{~}{ϵ}}1+\mathrm{ln}\left(\frac{m_\sigma ^2}{\mu ^2}\right)\right]+\frac{3\lambda }{2}_0^{\mathrm{}}\frac{dpp^2}{\pi ^2}\frac{n_\sigma }{\omega _\sigma }$$
(B9)
$$\mathrm{\Pi }_{\sigma 2}(k_0,\left|𝐤\right|=0)=\mathrm{\Pi }_{\sigma 2}^0+\mathrm{\Pi }_{\sigma 2}^\beta =\frac{3\lambda }{(4\pi )^2}m_\pi ^2\left[\frac{1}{\stackrel{~}{ϵ}}1+\mathrm{ln}\left(\frac{m_\pi ^2}{\mu ^2}\right)\right]+\frac{3\lambda }{2}_0^{\mathrm{}}\frac{dpp^2}{\pi ^2}\frac{n_\pi }{\omega _\pi }$$
(B10)
$`\mathrm{\Pi }_{\sigma 3}(k_0,\left|𝐤\right|=0)=\mathrm{\Pi }_{\sigma 3}^0+\mathrm{\Pi }_{\sigma 3}^\beta =18{\displaystyle \frac{\lambda ^2\nu ^2}{(4\pi )^2}}\left[{\displaystyle \frac{1}{\stackrel{~}{ϵ}}}2+\mathrm{ln}\left({\displaystyle \frac{m_\sigma ^2}{\mu ^2}}\right)+2f_1(k_0)arctan\left({\displaystyle \frac{1}{f_1(k_0)}}\right)\right]+`$ (B11)
$`18\lambda ^2\nu ^2{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dpp^2}{\pi ^2}}{\displaystyle \frac{n_\sigma }{\omega _\sigma }}{\displaystyle \frac{1}{4\omega _\sigma ^2k_0^2}}`$ (B12)
$`\mathrm{\Pi }_{\sigma 4}(k_0,\left|𝐤\right|=0)=\mathrm{\Pi }_{\sigma 4}^0+\mathrm{\Pi }_{\sigma 4}^\beta =6{\displaystyle \frac{\lambda ^2\nu ^2}{(4\pi )^2}}\left[{\displaystyle \frac{1}{\stackrel{~}{ϵ}}}2+\mathrm{ln}\left({\displaystyle \frac{m_\pi ^2}{\mu ^2}}\right)+2f_2(k_0)arctan\left({\displaystyle \frac{1}{f_2(k_0)}}\right)\right]+`$ (B13)
$`6\lambda ^2\nu ^2{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dpp^2}{\pi ^2}}{\displaystyle \frac{n_\pi }{\omega _\pi }}{\displaystyle \frac{1}{4\omega _\pi ^2k_0^2}}`$ (B14)
$`\mathrm{\Pi }_{\sigma 5}(k_0,\left|𝐤\right|=0)=\mathrm{\Pi }_{\sigma 5}^0+\mathrm{\Pi }_{\sigma 5}^\beta ={\displaystyle \frac{8g^2}{(4\pi )^2}}\left[3m_\psi ^2{\displaystyle \frac{k_0^2}{2}}\right]{\displaystyle \frac{1}{\stackrel{~}{ϵ}}}+`$ (B15)
$`{\displaystyle \frac{8g^2}{(4\pi )^2}}\left[{\displaystyle \frac{1}{2}}(6m_\psi ^2k_0^2)\mathrm{ln}\left({\displaystyle \frac{m_\psi ^2}{\mu ^2}}\right)+(4m_\psi ^2k_0^2)\left({\displaystyle \frac{\mathrm{\Delta }_2\sqrt{\mathrm{\Delta }_1}}{k_0^2}}1\right)\right]+`$ (B16)
$`4g^2{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dpp^2}{\pi ^2}}{\displaystyle \frac{n_\psi }{\omega _\psi }}[1+{\displaystyle \frac{4m_{\psi ^2}k_0^2}{k_0^24\omega _{\psi ^2}}}],`$ (B17)
where $`f_1(k_0)=\sqrt{\frac{4m_\sigma ^2}{k_0^2}1}`$, $`f_2(k_0)=\sqrt{\frac{4m_\pi ^2}{k_0^2}1}`$, $`\mathrm{\Delta }_1=k_0^2(k_0^24m_\psi ^2)`$, $`\mathrm{\Delta }_2=arctan\left(\frac{1}{\sqrt{1\frac{4m_\psi ^2}{k_0^2}}}\right)`$, $`\mathrm{\Delta }_3=k_0^42k_0^2(m_\pi ^2+m_\sigma ^2)+(m_\pi ^2m_\sigma ^2)^2`$, $`\mathrm{\Delta }_4=arctanh\left(\frac{k_0^2+(m_\pi ^2+m_\sigma ^2)}{\sqrt{\mathrm{\Delta }_3}}\right)`$ and $`\mathrm{\Delta }_5=arctanh\left(\frac{k_0^2(m_\pi ^2+m_\sigma ^2)}{\sqrt{\mathrm{\Delta }_3}}\right)`$.
$`\mathrm{\Sigma }(k_0,\left|𝐤\right|=0)=(\mathrm{\Sigma }_0+\mathrm{\Sigma }_s)_\sigma +3(\mathrm{\Sigma }_0+\mathrm{\Sigma }_s)_\pi =(\mathrm{\Sigma }_0^0+\mathrm{\Sigma }_0^\beta +\mathrm{\Sigma }_s^0+\mathrm{\Sigma }_s^\beta )_\sigma +3(\mathrm{\Sigma }_0^0+\mathrm{\Sigma }_0^\beta +\mathrm{\Sigma }_s^0+\mathrm{\Sigma }_s^\beta )_\pi =`$ (B18)
$`{\displaystyle \frac{g^2}{(4\pi )^2}}\left[m_\psi +{\displaystyle \frac{1}{2k_0}}(k_0^2+m_\psi ^2m_\sigma ^2)\right]{\displaystyle \frac{1}{\stackrel{~}{ϵ}}}+`$ (B19)
$`{\displaystyle \frac{g^2}{(4\pi )^2}}m_\psi \left[\mathrm{ln}\left({\displaystyle \frac{m_\psi ^2}{\mu ^2}}\right)+Z\right]{\displaystyle \frac{1}{(4\pi )^2}}{\displaystyle \frac{g^2}{2k_0}}(k_0^2+m_\psi ^2m_\sigma ^2)\left[\mathrm{ln}\left({\displaystyle \frac{m_\psi ^2}{\mu ^2}}\right)+Z\right]+`$ (B20)
$`{\displaystyle \frac{g^2}{2}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dpp^2}{\pi ^2}}{\displaystyle \frac{n_\sigma }{\omega _\sigma }}{\displaystyle \frac{[k_0(k_0^2+\omega _\sigma ^2+\omega _\psi ^2)+m_\psi (k_0^2\omega _\sigma ^2+\omega _\psi ^2)]}{[k_0^2(\omega _\psi \omega _\sigma )^2][k_0^2(\omega _\psi +\omega _\sigma )^2]}}+`$ (B21)
$`{\displaystyle \frac{g^2}{2}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dpp^2}{\pi ^2}}{\displaystyle \frac{n_\psi }{\omega _\psi }}{\displaystyle \frac{[2k_0\omega _\psi ^2+m_\psi (k_0^2\omega _\sigma ^2+\omega _\psi ^2)]}{[k_0^2(\omega _\psi \omega _\sigma )^2][k_0^2(\omega _\psi +\omega _\sigma )^2]}}+3(m_\sigma m_\pi )`$ (B22)
where $`\mathrm{\Sigma }_s(T,k_0,\left|𝐤\right|=0)`$ is the scalar contribution, proportional to the unit matrix and $`\mathrm{\Sigma }_0(T,k_0,\left|𝐤\right|=0)`$ is the contribution proportional to the matrix $`\gamma ^0`$, and
$`\left(\mathrm{\Sigma }_s^\beta \right)_\sigma ={\displaystyle \frac{g^2m_\psi }{2\pi ^2}}{\displaystyle _0^{\mathrm{}}}dpp^2[{\displaystyle \frac{n_\sigma }{\omega _\sigma }}{\displaystyle \frac{m_\psi ^2m_\sigma ^2k_0^2}{[k_0^2(\omega _\psi \omega _\sigma )^2][k_0^2(\omega _\psi +\omega _\sigma )^2]}}+`$ (B23)
$`{\displaystyle \frac{n_\psi }{\omega _\psi }}{\displaystyle \frac{m_\psi ^2m_\sigma ^2+k_0^2}{[k_0^2(\omega _\psi \omega _\sigma )^2][k_0^2(\omega _\psi +\omega _\sigma )^2]}}]`$ (B24)
$`\left(\mathrm{\Sigma }_0^\beta \right)_\sigma ={\displaystyle \frac{g^2k_0}{2\pi ^2}}{\displaystyle _0^{\mathrm{}}}dpp^2[{\displaystyle \frac{n_\sigma }{\omega _\sigma }}{\displaystyle \frac{\omega _\psi ^2+\omega _\sigma ^2k_0^2}{[k_0^2(\omega _\psi \omega _\sigma )^2][k_0^2(\omega _\psi +\omega _\sigma )^2]}}+`$ (B25)
$`{\displaystyle \frac{n_\psi }{\omega _\psi }}{\displaystyle \frac{2\omega _\psi ^2}{[k_0^2(\omega _\psi \omega _\sigma )^2][k_0^2(\omega _\psi +\omega _\sigma )^2]}}]`$ (B26)
$$\left(\mathrm{\Sigma }_s^0\right)_\sigma =\frac{g^2}{(4\pi )^2}m_\psi \left[\mathrm{ln}\left(\frac{m_\psi ^2}{\mu ^2}\right)+Z\right]$$
(B27)
$$\left(\mathrm{\Sigma }_0^0\right)_\sigma =\frac{1}{(4\pi )^2}\frac{g^2}{2k_0}(k_0^2+m_\psi ^2m_\sigma ^2)\left[\mathrm{ln}\left(\frac{m_\psi ^2}{\mu ^2}\right)+Z\right],$$
(B28)
with $`Z`$ defined as
$$Z\frac{1}{(4\pi )^2}\left[\frac{k_0^2+m_\sigma ^2m_\psi ^2}{2k_0^2}\mathrm{ln}\left(\frac{m_\psi ^2}{m_\sigma ^2}\right)+\frac{\sqrt{\mathrm{\Delta }_6}}{k_0^2}(\mathrm{\Delta }_7+\mathrm{\Delta }_8)\right],$$
(B29)
with $`\mathrm{\Delta }_6=k_0^42k_0^2(m_\psi ^2+m_{\sigma ;\pi }^2)+(m_\psi ^2m_{\sigma ;\pi }^2)^2`$, $`\mathrm{\Delta }_7=arctanh\left(\frac{k_0^2+(m_\psi ^2+m_{\sigma ;\pi }^2)}{\sqrt{\mathrm{\Delta }_6}}\right)`$ and $`\mathrm{\Delta }_8=arctanh\left(\frac{k_0^2(m_\psi ^2+m_{\sigma ;\pi }^2)}{\sqrt{\mathrm{\Delta }_6}}\right)`$.
|
warning/0005/cond-mat0005326.html
|
ar5iv
|
text
|
# Anelastic relaxation and 139LaNQR in La2-xSrxCuO4around the critical Sr content x=0.02
## 1 Introduction
From a variety of recent experiments and theoretical descriptions (mostly motivated by the search of the microscopic mechanism underlying high-temperature superconductivity), it has been realized that the electron system in doped two-dimensional (2D) quantum Heisenberg antiferromagnets (AF) exhibits complicated ordering phenomena. On cooling from high temperatures, first a kind of phase separation is expected to occur, causing the formation of charged stripes separating mesoscopic AF domains stripes ; zaanen ; emery . In cuprates in general the stripes should exist only dynamically, with slowing down of their fluctuations on cooling. At lower temperatures the spin degrees of freedom associated to the AF patches gooding between the stripes are known to freeze, generating a cluster spin-glass state weid ; chou ; nieder . Charge and spin freezing both involve a complex spin dynamics. While the stripe dynamic is slow, the motion of the holes along the stripe is much faster than the fluctuation of the stripe itself.
There is also evidence of unusual coupling of the lattice to charge and spin excitations egami . For instance, the <sup>139</sup>LaNMR line broadening, for T $``$ 40 K, in La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub>
(LSCO) for x = 0.12 (a signature of modulated magnetic order) is accompanied by softening of sound velocity suzuki . Neutron diffraction, for 0$``$ x $``$ 0.3, indicates local tilts of octahedra, interpreted as evidence of charged stripes bozin , the local tilt decreasing with increasing x. The local tilts give also rise to tunneling systems, observed by acoustic experiments for x $``$ 0.03, and relaxation rate strongly depends on doping cordero .
Charge localization along stripes and spin freezing have been studied, in La<sub>2</sub>CuO<sub>4</sub>-based compounds, mostly by means of NMR-NQR and $`\mu `$SR spectroscopies, which probe the low frequency excitations through the relaxation times and the modifications in the spectra borsa . In particular, it has been argued hunt that, in the underdoped regime of LSCO, when diffraction experiments indicate complete ordering, the stripes are still fluctuating at low frequencies. LSCO at Sr content around 0.02 is interesting, being at the boundary between the 3D-AF and the spin-glass phase (see Ref. john for a review). A recent neutron scattering study has shown that quasi-3D magnetic ordering occurs below about 40 K in the spin-glass state, with a spin structure related to the diagonal stripe structure matsuda . This new intermediate magnetic state is believed to result from partial freezing of 2D spin fluctuations existing at high temperatures, the spin-glass state being described as a random freezing of quasi-3D spin clusters with anisotropic spin correlations.
Motivated by this scenario of interrelated lattice and spin fluctuation effects, we have undertaken a comparative study of LSCO at x = 0.02 and x = 0.03 based on anelastic relaxation, <sup>139</sup>LaNQR relaxation and NQR spectra.
We first qualitatively recall how slowing down of spin fluctuations and ordering are expected to affect nuclear and mechanical relaxation. Single holes or charged stripes motions cause a time dependence in the hyperfine field h(t) =$`_i`$ A<sub>i</sub>S<sub>i</sub>(t) at the nucleus (S<sub>i</sub> spin operator at the i-th ion, A<sub>i</sub> hyperfine coupling tensor). When a characteristic frequency $`\omega _s`$ of the fluctuating stripes becomes of the order of the measuring quadrupole frequency $`\omega _m`$, a maximum in the spin lattice relaxation rate T$`{}_{}{}^{1}{}_{1}{}^{}`$ = $`W`$ driven by the local time dependence of h(t) is expected. Below this temperature the stripes move very slowly, or are ”pinned”, and an ”anomalous” magnetic moment is induced, associated to the 2D patches of AF correlated ions in between stripes gooding . The randomly distributed magnetic moments $`\mu _i`$ experience cooperative slowing down and a second relaxation mechanism sets in, related to the field at the nuclear sites due to $`\mu _i`$’s. In the correspondent relaxation rate a correlation function of the form $`_{i,j}\mu _i(0)\mu _j(t)=_i\mu _i(0)\mu _i(t)`$ is involved.
One can empirically write this correlation function as $`\mu ^2`$ exp\[-t/$`\tau _f(x,T)`$\], with an average correlation time $`\tau _f`$ which increases on decreasing temperature. At the temperature T<sub>g</sub> where $`\tau _f`$ becomes of the order of $`\omega _m^1`$ another peak appears in T$`{}_{}{}^{1}{}_{1}{}^{}`$, with non-exponential recovery law, a signature of disordered systems. Below T<sub>g</sub> one speaks of spin freezing. In a long-range ordered AF matrix T<sub>g</sub> should increase about linearly with x, due to the increased strength of the interaction among the $`\mu _i`$’s. On the contrary, for an amount of doping x which destroys the long range AF order, with the onset of a cluster spin-glass phase, T<sub>g</sub> decreases with increasing x borsa ; cho ; carretta .
In particular, in LSCO, for x $``$ 0.02 the peaks in <sup>139</sup>LaNQR W have been shown chou to follow the law T<sub>g</sub> = T<sub>f</sub> = bx, indicating spin freezing in an AF matrix. For x $``$ 0.02 (cluster spin-glass phase) the peaks at T<sub>g</sub>(x) have been attributed to the spin freezing of the magnetic moments in short-range AF correlated islands cho .
The stripe localization has been detected indirectly from the wipe-out effect on the <sup>63</sup>CuNQR signal in Nd and Eu doped LSCO at x = 0.12, at T<sub>charge</sub> = 65 K (and in the underdoped regime of LSCO), with a wipe out fraction having a temperature behavior similar to the charge and spin order determined by neutron scattering hunt ; singer . It should be remarked that in principle a wipe out effect should be accompanied, at higher temperature, by a marked enhancement in the spin lattice relaxation rate, when $`\omega _s`$ = $`\omega _m`$. Finally, $`\mu `$SR and <sup>139</sup>LaNQR measurements nieder ; julien pointed out a magnetic transition to a spin-glass like phase well extending into the superconducting regime.
As regards the effects expected in the anelastic relaxation, one notes that the elastic energy loss coefficient Q<sup>-1</sup>, measured by exciting flexural vibrations, is directly proportional to the imaginary part of the mechanical susceptibility $`\chi \mathrm{"}(\omega )`$. Thus Q<sup>-1</sup> is related to the spectral density J$`{}_{latt}{}^{}(\omega )`$ of the motions causing dissipation, according to the law Q$`{}_{}{}^{1}\chi \mathrm{"}\omega J_{latt}`$.
Since the motions of the stripes involve sizeable lattice effects, when a characteristic frequency decreases down into the kHz range they can be detected as maxima in $`\chi \mathrm{"}/\omega `$. Thus, from a combination of anelastic relaxation and of magnetic NQR relaxation, one can in principle probe the lattice and the spin fluctuations associated to the stripe motions. The investigation reported here was aimed at this purpose.
## 2 Experimentals results and discussion
Two LSCO ceramic samples grown by standard solid state reaction ferr have been investigated. According to x-ray diffraction the final amounts were x = 0.022 and x = 0.032 respectively. A more precise estimate of x was derived by detecting the orthorhombic-tetragonal transition through anelastic relaxation. The relationship of the transition temperature T<sub>0</sub> to the amount of Sr was taken as T<sub>0</sub> (x) = 535 \[1 - (1/0.235) x\] (Refs. john and cho2 ). From the step in the Young’s modulus (Fig.1) at the transition normalized in temperature and amplitude, one as T<sub>0</sub> = 495 K and T<sub>0</sub> = 468 K, corresponding to the Sr amounts x = 0.019$`\pm `$0.0015 and x=0.030$`\pm `$0.001 (hereafter called samples 2 and 3 percent). As it is noted the structural transition appears even sharper than in pure La<sub>2</sub>CuO<sub>4</sub>(where some broadening may be attributed to non-perfect oxygen stoichiometry).
SQUID magnetization measurements yield a magnetic susceptibility as a function of temperature qualitatively typical of a spin-glass phase. Small differences between the field cooled and zero- field cooled data below about 140 K have been attributed to magnetic impurities present in the powders used for the preparation. These impurities do not affect the <sup>139</sup>LaNQR and anelastic relaxation measurements.
In Fig.2 the quantity J<sub>latt</sub> = T Q$`{}_{}{}^{1}/\omega `$ is reported in the temperature range of interest, in correspondence to the measuring frequencies $`\omega /2\pi `$ = 1.29, 6.9 and 17.2 kHz, in LSCO at x = 2 percent. The thermal depinning of the stripes should have cooperative character. For the moment we neglect their frequency distribution. According to the data in Fig. 2, the spectral density of the motion responsible of the dissipation has a diffusive character:
$$J_{latt}(\omega _m)=\frac{2\omega _s}{\omega _s^2+\omega _m^2}$$
(1)
From the temperature where the maxima are observed one can deduce the values of the characteristic frequency $`\omega _s`$. A good fit of the data (Fig.3) is obtained on the basis of a temperature behavior of $`\omega _s`$ of the form
$$\omega _s=\omega _0exp(E/T)]$$
(2)
with $`\omega _0`$ = 4.5$`\times `$10<sup>12</sup> s<sup>-1</sup> and E = 1650 K, consistent with the idea of thermal activation of the stripe fluctuations.
As regards <sup>139</sup>LaNQR relaxation, the recovery laws for the +5/2 - +7/2 line at 3$`\nu _Q`$ and the +3/2 - +5/2 at 2$`\nu _Q`$, are multiexponential. However in the first decade they differ only little from a single exponential. The correspondent effective decay rate $`\tau _e^1`$ can be related to the magnetic relaxation rate $`W_M`$ or to the quadrupolar relaxation rate $`W_Q`$ due to the time dependence of the electric field gradients at the La site, in the following way rega :
$`3\nu _Q\tau _e^1`$ $`=`$ $`(67/21)W_Q=23W_M`$
$`2\nu _Q\tau _e^1`$ $`=`$ $`(64.5/21)W_Q=41.3W_M`$ (3)
For T $``$ 250 K the NQR relaxation rate (Fig.4) are frequency independent and follow the law $`\tau _e^1T^2`$, features characteristic of the relaxation process driven by underdamped phonons adv . Also an order of magnitude estimate corroborates this conclusion, since for this process one expects birge $`W_Q`$ = 5$`\times `$ 10<sup>-4</sup> T<sup>2</sup> s<sup>-1</sup>. Around 280 K some evidence of a quadrupole contribution due to overdamped phonon modes (tilting of the oxygen octahedra in a double well potential prb ) is present rr . Below 250 K the relaxation rates depart from the behavior described above. For temperatures lower than about 180 K the comparison of the data for the 3$`\nu _Q`$ and 2$`\nu _Q`$ lines indicates the insurgence of a magnetic relaxation mechanism. The temperature behavior of $`\tau _e^1`$ for T $``$ 230 K is analyzed in detail in Fig. 5, after subtraction of the background of quadrupole character. Assuming that the modulation of the hyperfine magnetic field h(t) is due to the same motion causing the mechanical relaxation, then for
$$2W_M=\frac{1}{2}\gamma ^2h_+(0)h_{}(t)e^{i\omega t}𝑑t$$
(4)
one writes
$$(\tau _e^1)_{2nQ/3nQ}=aW_M=a\frac{1}{2}\gamma ^2h^|[2\omega _s/(\omega _s^2+\omega _m^2)]$$
(5)
with a = 23 for the 3$`\nu _Q`$ line and a = 41.3 for the 2$`\nu _Q`$ line, $`\omega _m`$ = 3$`\omega _Q`$ and $`\omega _m`$ = 2$`\omega _Q`$ respectively. In Fig. 5 the experimental data for $`aW_M`$ are compared with the theoretical behaviors for the relaxation rates according to Eq.s 2 and 5, having used for $`\omega _s`$ the expression derived from the anelastic relaxation (Eq. 2). The maxima in $`W_M`$ are well reproduced. The departures of the experimental data from the theoretical expressions in the temperature range corresponding to slow motions, i.e. $`\omega _s\omega _m`$, are likely to be due to the simplifying assumption of a monodispersive process. A distribution in $`\omega _s`$ implies a flattening in the relaxation rate around the maximum and a departure from the behavior for monodispersive process more marked in the low temperature range, as it is observed in the Figure. The relevant fact is that the temperature dependence of $`\omega _s`$ deduced from anelastic relaxation in the kHz range seems to justify quantitatively the magnetic NQR relaxation rate in the MHz range.
On the other hand, the maxima in $`W_M`$ around 150 K could reflect the slowing down of the spin dynamics on approaching the transition to an ordered state. Information in this regard can be obtained from the NQR spectra (Fig.6), since the insurgence of a static field $`h`$ at the La site is signaled by a splitting $`\mathrm{\Delta }`$ of the resonance line, proportional to the sublattice magnetization $`S`$. A clear splitting of the line is noticeable only below about 50 K, close to the temperature indicated by neutron scattering matsuda . If the width $`\delta `$ of a single component is kept temperature independent the fitting of the spectra with two gaussian lines yields an ordering temperature T<sub>N</sub>, where $`\mathrm{\Delta }`$ goes to zero, of about 140 K. The temperature dependence of the order parameter $`hS`$ would turn out quite different from the one of a canonical phase transition to the AF phase, experimentally observed macl in pure La<sub>2</sub>CuO<sub>4</sub>.
On the other hand, if the maxima in $`W_M`$ at 150 K are taken as an indication of stripe motion at frequency around $`\omega _Q`$, then an NQR line broadening must be expected below the temperature at which the frequency becomes of the order of the line width itself, about 200 kHz, namely around 110 K according to Fig. 3. An experimental support to the hypothesis that the intrinsic linewidth $`\delta `$ is temperature dependent comes from the comparison of the spectra at 2$`\nu _Q`$ and at 3$`\nu _Q`$ (Fig.6). The ratio $`\delta _{2\nu _Q}`$(T = 77 K) / $`\delta _{2\nu _Q}`$(T = 177 K) = (210 kHz ) / (183 kHz) = 1.15 is the same of the one for the 3$`\nu _Q`$ line: $`\delta _{3\nu _Q}`$(T = 77 K) / $`\delta _{3\nu _Q}`$(T = 177 K) = (320 kHz ) / (280 kHz) = 1.15 One also has $`\delta _{3\nu _Q}`$ = (3/2) $`\delta _{2\nu _Q}`$. These data are not compatible with an effect due to a magnetic field $`h`$, that would cause an extra broadening of the same amount for both the 2$`\nu _Q`$ and the 3$`\nu _Q`$ lines. Thus, if a moderate (of the order of 15-20%) temperature dependence for the single component linewidth $`\delta `$ is allowed, then the NQR spectra indicate T$`{}_{N}{}^{}`$ 50 K. This value is in substantial agreement with the phase diagram commonly accepted in literature john . One could speculate that at this temperature a small bump is observed in the relaxation rates (Fig.4), consistent with the slowing-down of the spin dynamics, just above the temperature range where the drastic increase of $`W_M`$ due to the spin freezing occurs. On cooling, the relaxation rate exhibits a maximum at a temperature around T = 9 K. Also the recovery plots, showing evidence of departure from an exponential recovery towards the t<sup>1/2</sup> law, support the conclusion that the spin-glass quasi-freezing temperature T<sub>f</sub> has been reached.
We compare now the experimental findings in the sample at the boundary between the AF and the spin-glass phase with the one at x = 3 percent, well within the latter phase. The anelastic relaxation shows a temperature behavior similar to the one for x = 2 percent (Fig.7). Both the peaks attributed to the stripes motion and to collective tilting of the octahedra (not shown) are attenuated by a factor of about two. The <sup>139</sup>LaNQR relaxation rates indicate the typical freezing of the spin fluctuations in a cluster spin-glass (Fig.8). The maximum in $`W`$ occurs at T = 8 K, in good agreement with the phase diagram by Cho et al cho ; cho2 .
The relaxation rate measured at 2$`\nu _Q`$ reaches a value about twice the one reported in previous measurements cho ; borsa2 ; riga at 3$`\nu _Q`$ , consistent with a magnetic relaxation mechanism . It is noted that for x = 0.03 magnetization measurements provide direct evidence of the occurrence of a canonical spin-glass state waki .
The temperature dependence of $`W`$ can be discussed in terms of the behavior expected for the effective correlation time. For magnetic moments coupled to a Fermi gas of carriers macf , from Eq. 5 in the fast fluctuations regime one would have $`W1/\omega _s=\tau _f=h/\pi (\rho J)^2kT`$, where $`J`$ is the exchange coupling to the band and $`\rho `$ the density of states at the Fermi level. From Fig. 8 one sees that $`W`$ diverges, on decreasing temperature, more rapidly than T<sup>-1</sup>, in a way close to the law $`\tau exp[E/T]`$ at least for the temperature range where the fast motions condition, namely $`\tau \omega _Q1`$, holds (see inset in Fig. 8). Below T$``$ 10K one notes the behaviour of the relaxation rate expected for a glassy system chou ; cho .
## 3 Conclusions
Anelastic and NQR relaxation and NQR spectra have been combined in the attempt to derive insights on spin and lattice excitations, possibly stripes motions, driving the ordering processes in LSCO. The experimental findings have been discussed within two interpretative frameworks. On one side it could be possible that the phase diagram around the Sr content x = 0.02 separating the AF and the spin-glass phases is more complex that previously assumed, with a quasi-long range ordering temperature as high as 150 K, corroborating recent neutron scattering measurements birge and qualitatively agreeing with the extrapolation at x = 0.02 of a magnetization study waki carried out in samples at x = 0.03, 0.04 and 0.05. The order parameter of such a transition would be characterized by an unconventional temperature dependence.
On the other hand, the NQR spectra can be interpreted as indicating a conventional transition to the AF phase around T = 50 K, in substantial agreement with the phase diagram commonly accepted. In this case the anelastic and La NQR relaxation rates around T = 80 K and T = 160 K respectively, are the first direct experimental evidence of low frequency motions of stripes simultaneously involving spin and lattice excitations. The thermal depinning barriers and the characteristic d̀iffusivef́requencies are then derived, and they do not differ much in the sample at Sr content 0.03.
In the low temperature range a spin freezing process is detected, with a dramatic increase of the <sup>139</sup>LaNQR relaxation rate on cooling.
These experimental findings could be, at least in part, explained with a kind of distribution of transition temperatures T<sub>N</sub>, T<sub>f</sub> and T<sub>g</sub> resulting from a spread in the Sr content around the critical amount x = 0.02. However the tetragonal-orthorhombic transition appears sharp and hence the Sr content seems to be well defined and close to that value.
## 4 Acknowledgments
Alessandro Lascialfari is gratefully thanked for his SQUID measurements and for helpful discussions. Stimulating discussions with F. Borsa, P. Carretta, R. Gooding and M.H. Julien are also acknowledged.
The research has been carried out in the framework of the PRA project SPIS (1998-2000), financed by INFM (Italy).
|
warning/0005/cond-mat0005495.html
|
ar5iv
|
text
|
# Dynamically Induced Multi-Channel Kondo Effect
## I Introduction
The multi-channel Kondo effect has been the subject of intensive theoretical and experimental studies, which is characterized by unusual non-Fermi liquid behaviors. Its applications are now extended not only to standard dilute magnetic alloys, but also to quantum dots, etc. Thus far, theoretical and experimental studies on the multi-channel Kondo effect have been focused on a static Kondo impurity, which has been related to the measurements of the specific heat, the spin susceptibility, the resistivity, etc. This naturally motivates us to address a question whether such a nontrivial phenomenon can be observed in dynamically generated situations.
The photoemission and the inverse photoemission may be one of the key experiments to study non-Fermi liquid behaviors, which reveal the dynamics of a single hole or electron suddenly created in the system. We here propose the dynamically induced multi-channel Kondo effect, when an electron is emitted from (or added to) the Kondo impurity by the photoemission (inverse photoemission). A remarkable point is that the ground state of the system is assumed to be a completely screened Kondo singlet, and non-Fermi liquid properties are generated by an electron or hole suddenly created. We study low-energy critical properties of the spectrum by using the exact solution of the multi-channel Kondo model combined with boundary conformal field theory (CFT). We analyze the one-particle Green function for the impurity to show its typical non-Fermi liquid behavior. It is further demonstrated that this effect can be observed even in a homogeneous system without impurities. To show this explicitly, we apply the analysis to the photoemission spectrum in a quantum spin chain with spin $`S>1/2`$.
This paper is organized as follows. In §2 we briefly illustrate the idea of the dynamically induced multi-channel Kondo effect, and derive low-energy scaling forms of the one-particle Green function. We discuss non-Fermi liquid properties in the spectrum by exactly evaluating the critical exponents. In §3 the analysis is then applied to the photoemission spectrum for a quantum spin chain. Brief summary is given in §4. We note that preliminary results on this issue have been reported in Ref. 11.
## II Low-Energy Dynamics in the Photoemission Spectrum
### A dynamically induced Kondo effect
Let us consider the spin-$`S`$ Kondo impurity which is completely screened by conduction electrons with $`n(=2S)`$ channels. The impurity spin is assumed to be composed of $`n`$ electrons by the strong Hund coupling. To study the core-electron photoemission spectrum, we start with spectral properties of the impurity Green function,
$`G(t)=\mathrm{i}<\mathrm{T}[d(t)d^{}(0)]>`$ (1)
$`{\displaystyle \genfrac{}{}{0pt}{}{}{=G^>(t)+G^<(t),}}`$ (2)
where $`d`$ is the annihilation operator for one of core electrons which compose the impurity spin and T is the conventional time-ordered product. Here, $`G^>(t)`$ ($`G^<(t)`$) is the Green function, which is restricted in $`t>0(t<0)`$. For the photoemission, we consider $`G^<(t)`$. To be specific, we discuss the case that a core electron is emitted as depicted in Fig. 1 (a), for which the binding energy $`\omega _\alpha `$ (measured from the Fermi energy) is assumed to be larger than the band width $`D`$. Then in the excited state the overscreening system is generated, which is referred to as the dynamically induced overscreening Kondo effect. At the low-energy regime around $`\omega _\alpha `$, we may express the operator as $`d(t)\mathrm{e}^{\mathrm{i}\omega _\alpha t}\varphi (t)`$ where $`\varphi (t)`$ is the corresponding boundary operator in boundary CFT, which characterizes the boundary critical phenomena. It is known that the Fermi-edge singularity is reformulated by the boundary operator, in which nontrivial effects for the overscreening Kondo effect are incorporated in $`\varphi _\alpha (t)`$. We write down the one-particle Green function $`G^<(t)`$ as,
$`\mathrm{Im}G^<(\omega )={\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}<\varphi _\alpha ^{}(0)\varphi _\alpha (t)>\mathrm{e}^{\mathrm{i}\omega _\alpha \mathrm{t}}\mathrm{e}^{\mathrm{i}\omega \mathrm{t}}dt.`$ (3)
On the other hand, for the inverse photoemission, an added electron is combined with the local spin $`S`$ to form higher spin $`S+1/2`$ by the strong Hund-coupling, as shown in Fig. 1 (b). Then we may write $`d(t)\mathrm{e}^{\mathrm{i}\omega _\beta t}\varphi _\beta (t)`$, where $`\omega _\beta >D`$ is the energy cost to make $`S+1/2`$ spin, and $`\varphi _\beta (t)`$ is another boundary operator which controls the undersreening Kondo effect induced by the inverse photoemission. We have
$`\mathrm{Im}G^>(\omega )`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}<\varphi _\beta (t)\varphi _\beta ^{}(0)>\mathrm{e}^{\mathrm{i}\omega _\beta t}\mathrm{e}^{\mathrm{i}\omega \mathrm{t}}dt.`$ (4)
In order to evaluate the critical exponents, we now employ the idea of finite-size scaling in CFT. The scaling form of the correlators $`<\varphi _\alpha ^{}(0)\varphi _\alpha (t)>`$ and $`<\varphi _\beta (t)\varphi _\beta ^{}(0)>`$ are given by
$`<\varphi _\alpha ^{}(0)\varphi _\alpha (t)>`$ $`=`$ $`{\displaystyle \underset{N=0}{\overset{\mathrm{}}{}}}|<0|\varphi _\alpha ^{}(0)|N>|^2\mathrm{e}^{\mathrm{i}\frac{\pi v_F}{l}(\mathrm{\Delta }_\alpha +N)t}`$ (5)
$`=`$ $`\left({\displaystyle \frac{{\displaystyle \frac{\pi }{2l}}}{\mathrm{sinh}{\displaystyle \frac{\pi v_F}{2l}}\mathrm{i}t}}\right)^{2\mathrm{\Delta }_\alpha }{\displaystyle \frac{1}{(\mathrm{i}v_Ft)^{2\mathrm{\Delta }_\alpha }}}`$ (6)
$`<\varphi _\beta (t)\varphi _\beta ^{}(0)>`$ $`=`$ $`{\displaystyle \underset{N=0}{\overset{\mathrm{}}{}}}|<N|\varphi _\beta ^{}(0)|0>|^2\mathrm{e}^{\mathrm{i}\frac{\pi v_F}{l}(\mathrm{\Delta }_\beta +N)t}`$ (8)
$`=`$ $`\left({\displaystyle \frac{{\displaystyle \frac{\pi }{2l}}}{\mathrm{sinh}{\displaystyle \frac{\pi v_F}{2l}}\mathrm{i}t}}\right)^{2\mathrm{\Delta }_\beta }{\displaystyle \frac{1}{(\mathrm{i}v_Ft)^{2\mathrm{\Delta }_\beta }}},`$ (9)
in the long-time asymptotic region. According to the finite-size scaling, the boundary dimensions $`\mathrm{\Delta }_\alpha `$ and $`\mathrm{\Delta }_\beta `$ are read from the lowest excitation energy $`\mathrm{\Delta }E`$,
$`\mathrm{\Delta }E={\displaystyle \frac{\pi v_F}{l}}\mathrm{\Delta }_\gamma ,`$ (10)
with $`\gamma =\alpha ,\beta `$, where $`l`$ corresponds to the system size of one dimension in the radial direction. We thus end up with the relevant scaling forms as
$`\mathrm{Im}G^<(\omega )`$ $`=`$ $`{\displaystyle \frac{\pi }{\mathrm{\Gamma }(2\mathrm{\Delta }_\alpha )v_{F}^{}{}_{}{}^{2\mathrm{\Delta }_\alpha }}}\theta (\omega \omega _\alpha )(\omega \omega _\alpha )^{X_\alpha },`$ (11)
$`\mathrm{Im}G^>(\omega )`$ $`=`$ $`{\displaystyle \frac{\pi }{\mathrm{\Gamma }(2\mathrm{\Delta }_\beta )v_{F}^{}{}_{}{}^{2\mathrm{\Delta }_\beta }}}\theta (\omega \omega _\beta )(\omega \omega _\beta )^{X_\beta },`$ (13)
where $`X_\gamma =2\mathrm{\Delta }_\gamma 1`$ that $`\gamma `$ represents $`\alpha `$ and $`\beta `$.
In both cases, the spectral functions have power-law edge singularity due to the dynamically induced multi-channel Kondo effect, which will be shown to exhibit non-Fermi liquid properties.
### B exact critical properties
We now discuss low-energy critical properties by exactly evaluating $`\mathrm{\Delta }_\alpha `$ and $`\mathrm{\Delta }_\beta `$. To this end, we consider the multi-channel Kondo model,
$``$ $`=`$ $`\mathrm{i}{\displaystyle \underset{a,m}{}}{\displaystyle dx\psi _{am}^{}(x)_x\psi _{am}(x)}`$ (16)
$`+2J{\displaystyle \underset{a,b,m}{}}{\displaystyle \underset{\nu }{}}\psi _{am}^{}(0)\sigma _{ab}^\nu \psi _{bm}(0)S^\nu HS^z,`$
where $`\psi _{am}^{}`$ is the creation operator for conduction electrons with spin $`a=,`$ and orbital indices, $`m=1,\mathrm{},n`$. The exact solution of this model is expressed in terms of the Bethe equations for spin rapidities $`\lambda _\alpha `$ and charge rapidities $`k_j`$,
$`\mathrm{e}^{\mathrm{i}k_jL}={\displaystyle \underset{\alpha =1}{\overset{M}{}}}{\displaystyle \frac{\lambda _\alpha +\mathrm{i}n/2}{\lambda _\alpha \mathrm{i}n/2}}`$ (17)
(18)
$`{\displaystyle \frac{\lambda _\alpha +1/J+\mathrm{i}S}{\lambda _\alpha +1/J\mathrm{i}S}}\left({\displaystyle \frac{\lambda _\alpha +\mathrm{i}n/2}{\lambda _\alpha \mathrm{i}n/2}}\right)^N={\displaystyle \underset{\alpha =1}{\overset{M}{}}}{\displaystyle \frac{\lambda _\alpha \lambda _\beta +\mathrm{i}}{\lambda _\alpha \lambda _\beta \mathrm{i}}},`$ (19)
where $`N`$ is the number of electrons and $`L=2l`$ is the one-dimensional system size. It is assumed that the impurity with spin $`S>1/2`$ is completely screened in the ground state. Then, the core-level photoemission suddenly reduces the impurity spin, thus inducing the overscreening Kondo effect with $`n2S=1`$. As for the underscreening effect induced by the inverse photoemission, the condition is replaced by $`n2S=1`$. At zero temperature the ground-state properties are described by the $`n`$th order string solutions,
$`\lambda _l^{n,\alpha }=\lambda _l^n+{\displaystyle \frac{\mathrm{i}}{2}}(n+12\alpha ),`$ (20)
where $`\alpha =1,\mathrm{},n`$ and $`l=1,\mathrm{},M_n`$ which is restricted by $`M=nM_n`$. Here, $`n`$ represents the number of orbitals. It is well known that a naive application of finite-size techniques based on the string hypothesis turns out to fail for the overscreening case at zero magnetic field. This difficulty comes from an improper treatment of the $`\mathrm{Z}_\mathrm{n}`$ symmetry sector in terms of the string solutions. However, as long as the finite magnetic field is concerned, we can use the Bethe equations to describe its critical properties. We will separately discuss the case of zero-magnetic field, by incorporating $`\mathrm{Z}_\mathrm{n}`$ sector correctly. By applying standard procedures to eq.(19), it is straightforward to exactly evaluate the lowest excitation energy in magnetic fields, for which one of the impurity electrons is assumed to be removed from the system,
$`\mathrm{\Delta }E={\displaystyle \frac{\pi v_F}{l}}\left({\displaystyle \frac{\delta _\alpha ^2}{4n}}+n(n_{\mathrm{imp}})^2\right).`$ (21)
Although the above finite-size correction is apparently similar to that for 1D solvable systems with a static impurity or boundaries, the final-state interaction induced by photoemission is included in the present case. Thus all the features which are governed by the dynamical Kondo effect can be read from this quantity.
By applying the finite-size scaling in eq. (10), the critical exponent $`X_\gamma `$ in eq. (13) is now obtained as,
$`X_\alpha ={\displaystyle \frac{\delta _\alpha ^2}{2n}}+2n(n_{\mathrm{imp}})^21,`$ (22)
where $`\delta _\alpha `$ is the charge scattering phase shift, which is caused by a created hole as in the ordinary Fermi edge singularity. This term depends on the detail of potential scattering. It is mentioned that the second term with the phase shift $`n_{\mathrm{imp}}`$ is caused by the Kondo effect, which is explicitly evaluated as,
$`n_{\mathrm{imp}}={\displaystyle _{\mathrm{}}^{\lambda _0}}\sigma _{\mathrm{imp}}(\lambda )d\lambda ,`$ (23)
where
$`\sigma _{\mathrm{imp}}(\lambda )`$ $`=`$ $`\sigma _{\mathrm{imp}}^0(\lambda +1/J)`$ (26)
$`{\displaystyle _{\mathrm{}}^{\lambda _0}}G_n(\lambda \lambda ^{^{}})\sigma _{\mathrm{imp}}(\lambda ^{^{}})d\lambda ^{^{}}.`$
Here $`\sigma _{\mathrm{imp}}^0(\lambda )`$ and $`G_n(\lambda )`$ are determined by
$`\sigma _{\mathrm{imp}}^0(\lambda )={\displaystyle \frac{1}{\pi }}{\displaystyle \underset{l=1}{\overset{\mathrm{min}(n,2S)}{}}}{\displaystyle \frac{{\displaystyle \frac{1}{2}}(n+2S+12l)}{\lambda ^2+{\displaystyle \frac{1}{4}}(n+2S+12l)^2}},`$ (27)
(28)
$`G_n(\lambda )={\displaystyle \frac{1}{\pi }}{\displaystyle \frac{n}{\lambda ^2+n^2}}+{\displaystyle \frac{2}{\pi }}{\displaystyle \underset{\alpha =1}{\overset{n1}{}}}{\displaystyle \frac{{\displaystyle \frac{1}{2}}(2n2\alpha )}{\lambda ^2+{\displaystyle \frac{1}{4}}(2n2\alpha )^2}}.`$ (29)
The key quantity, $`n_{\mathrm{imp}}`$, is obtained by using Wiener-Hopf method,
$`n_{\mathrm{imp}}={\displaystyle \frac{S}{n}}\left({\displaystyle \frac{S}{n}}{\displaystyle \frac{1}{2}}\right)\theta \left({\displaystyle \frac{S}{n}}{\displaystyle \frac{1}{2}}\right)`$ (30)
(31)
$`+{\displaystyle \frac{\mathrm{i}}{4\pi ^{\frac{3}{2}}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{\mathrm{d}\omega }{\omega \mathrm{i0}}}\mathrm{e}^{\mathrm{i2}\omega \mathrm{log}\frac{H}{T_H}}{\displaystyle \frac{\mathrm{\Gamma }\left(1+\mathrm{i}\omega \right)\mathrm{\Gamma }\left(1/2\mathrm{i}\omega \right)}{\mathrm{\Gamma }(1+\mathrm{i}n\omega )}}`$ (32)
(33)
$`\left({\displaystyle \frac{\mathrm{i}n\omega +0}{\mathrm{e}}}\right)^{\mathrm{i}n\omega }{\displaystyle \frac{\mathrm{e}^{\pi |n2S||\omega |}\mathrm{e}^{\pi (n+2S)|\omega |}}{1\mathrm{e}^{2\pi n|\omega |}}}.`$ (34)
In Fig. 2 we display the critical exponent $`X_\alpha `$ for the overscreening case as a function of the magnetic field. Note that the magnetic-field dependence of $`X_\alpha `$ without $`\delta _\alpha ^2/2n`$ is determined by the dynamical Kondo effect, because the charge scattering phase shift $`\delta _\alpha `$ does not depend on magnetic fields. Particularly in weak magnetic field, $`H<<T_H`$, the obtained exponent behaves as
$`X_\alpha {\displaystyle \frac{\delta ^2}{2n}}+\left({\displaystyle \frac{S}{n}}\mathrm{const}\left({\displaystyle \frac{H}{T_H}}\right)^{{\scriptscriptstyle \frac{2}{n}}}\right)^21.`$ (35)
It is seen that the phase shift, $`n_{\mathrm{imp}}`$, gives rise to the anomalous magnetic-field dependence of the exponent. This non-Fermi liquid behavior is characteristic of the overscreening effect. Another interesting feature in the overscreening effect appears at $`H=0`$. We recall here that the symmetry is enhanced from U(1) to SU(2) at $`H=0`$, for which the boundary dimension $`\mathrm{\Delta }_s`$ for the spin sector is analytically obtained by employing fusion rules hypothesis proposed by Affleck and Ludwig,
$`\mathrm{\Delta }_s={\displaystyle \frac{S(S+1)}{n+2}},`$ (36)
which is a typical conformal dimension for level-$`n`$ SU(2) Kac-Moody algebra. Note that the critical exponent shows a discontinuity at $`H=0`$,
$`X_\alpha (H=0)X_\alpha (H0)=2{\displaystyle \frac{S(S+1)}{n+2}}2{\displaystyle \frac{S^2}{n}},`$ (37)
which is caused by the fact that $`\mathrm{Z}_\mathrm{n}`$ symmetric sector is massless only at $`H=0`$, as already mentioned.
We now move to the underscreening case induced by the inverse photoemission. The calculated critical exponent $`X_\beta `$ is shown as a function of magnetic fields in Fig. 3. For weak magnetic field, the obtained exponent $`X_\beta `$ behaves as
$`X_\beta {\displaystyle \frac{\delta _\beta ^2}{2n}}+n\left({\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{\mathrm{log}(H/T_H)}}\right)^21,`$ (38)
which is characteristic of the underscreening system. In contrast to the overscreening case, there is no discontinuity in the exponent in this case, $`X_\beta (H=0)=X_\beta (H0)`$.
This completes a general description of the dynamically induced Kondo effect. An important point to be emphasized is that this kind of phenomenon may be observed not only for impurity systems but also for other related quantum systems, which will be explicitly discussed in the next section.
## III Application to Spin Chains
We now wish to demonstrate that the dynamically induced Kondo effect proposed here may be observed for gapless quantum spin systems which do not possess impurities. As an example, we consider an integrable antiferromagnetic spin chain with spin $`S>1/2`$, for which the exact solution is available even for the case with doped holes. The photoemission suddenly removes one electron from the spin system, and thus bears an impurity site with spin $`S1/2`$. In the final state, this impurity spin is screened by host spins, and as a result the overscreening Kondo effect may be dynamically induced. It is remarkable that the induced impurity in this case can move through the lattice via the exchange interaction, and the edge singularity is thus governed by a mobile multichannel Kondo impurity.
Let us consider the gapless $`S=1`$ spin chain with a mobile $`S=1/2`$ impurity as an example. Although in more general situations including non-integrable models, the higher spin should be a half-odd integer, the present treatment can be straightforwardly extended to such cases. We here consider an integrable spin chain derived by the quantum inverse scattering method,
$``$ $`={\displaystyle \underset{i=1}{\overset{L}{}}}(1\delta _{S_iS_{i+1},1})𝒫_{i,i+1}(𝐒_i𝐒_{i+1})`$ (40)
$`+`$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{1}{S_iS_{i+1}}}𝐒_i𝐒_{i+1}1+\delta _{S_iS_{i+1},1}[1(𝐒_i𝐒_{i+1})^2]\right)`$ (41)
where the spin $`𝐒_{𝐢}^{}{}_{}{}^{2}=S_i(S_i+1)`$ with $`S_i=1`$ or 1/2, and $`𝒫_{ij}`$ permutes the spin on sites $`i`$ and $`j`$.
In order to deal with the excited states when an electron is emitted from the spin chain, we write down the Bethe equations for the spin-$`S`$ chain with one hole being doped,
$`\left({\displaystyle \frac{\lambda _j+\mathrm{i}S}{\lambda _j\mathrm{i}S}}\right)^L`$ $`=`$ $`{\displaystyle \frac{\lambda _j\nu _{}\mathrm{i}/2}{\lambda _j\nu _+\mathrm{i}/2}}{\displaystyle \underset{kj}{\overset{N_{}+1}{}}}{\displaystyle \frac{\lambda _j\lambda _k+\mathrm{i}}{\lambda _j\lambda _k\mathrm{i}}},`$ (43)
$`1`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{N_{}+1}{}}}{\displaystyle \frac{\nu \lambda _k+\mathrm{i}/2}{\nu \lambda _k\mathrm{i}/2}},`$ (44)
where $`N_{}`$ is the number of down spins and $`\lambda _j(j=1,\mathrm{},N_{}+1)`$ are spin rapidities. Note that the hole rapidity $`\nu `$ appears in the above equation, which characterizes a massive charge excitation suddenly created. Thus, $`\nu `$ specifies the hole momentum $`q`$. It is seen that the above spin-charge scattering term, which describes the final-state interaction, corresponds to the impurity term in eq. (19).
The manipulation illustrated in the previous section enables us to exactly calculate the scaling dimension for the one-particle Green function via the finite-size corrections,
$`x(\nu )={\displaystyle \frac{1}{4\xi _{2S}^2}}(1n_{\mathrm{imp}}(\nu ))+\xi _{2S}^2({\displaystyle \frac{1}{2}}d_{\mathrm{imp}}(\nu ))^2.`$ (45)
The quantity $`\xi _{2S}\xi _{2S}(\lambda _0)`$ often referred to as the dressed charge is given by
$`\xi _{2S}(\lambda )=1{\displaystyle _{\lambda _0}^{\lambda _0}}K_{2S}(\lambda \lambda ^{^{}})\xi _{2S}(\lambda ^{^{}}),`$ (46)
where
$`K_{2S}(\lambda )`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle \frac{\frac{1}{2}(4S)}{\lambda ^2+\frac{1}{4}(4S)^2}}`$ (48)
$`+{\displaystyle \frac{2}{\pi }}{\displaystyle \underset{l=1}{\overset{2S+1}{}}}{\displaystyle \frac{\frac{1}{2}(4S2l)}{\lambda ^2+\frac{1}{4}(4S2l)^2}},`$
and the cut-off parameter $`\lambda _0`$ is related to the magnetization $`m`$,
$`Sm=2S{\displaystyle _{\lambda _0}^{\lambda _0}}\rho _{2S}(\lambda )d\lambda .`$ (49)
The density function $`\rho _{2S}`$ is determined by the following integral equation,
$`\rho _{2S}(\lambda )={\displaystyle \frac{1}{2\pi }}\mathrm{\Theta }_{2S,2S}^{^{}}(\lambda ){\displaystyle _{\lambda _0}^{\lambda _0}}K_{2S}(\lambda \lambda ^{^{}})\rho _{2S}(\lambda ^{^{}}).`$ (50)
We stress that two key quantities $`n_{\mathrm{imp}}(\nu )`$ and $`d_{\mathrm{imp}}(\nu )`$, which contain the effect of a mobile impurity, are introduced in eq. (45),
$`n_{\mathrm{imp}}(\nu )`$ $`=`$ $`{\displaystyle _{\lambda _0}^{\lambda _0}}\rho _{\mathrm{imp}}(\lambda )d\lambda `$ (52)
$`d_{\mathrm{imp}}(\nu )`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle _{\lambda _0}^{\mathrm{}}}{\displaystyle _{\mathrm{}}^{\lambda _0}}\right)\rho _{\mathrm{imp}}(\lambda )\mathrm{d}\lambda `$ (53)
where
$`\rho _{\mathrm{imp}}(\lambda )={\displaystyle \frac{1}{2\pi }}\mathrm{\Theta }_{2S,1}^{^{}}(\lambda \nu ){\displaystyle _{\lambda _0}^{\lambda _0}}K_{2S}(\lambda \lambda ^{^{}})\rho _{\mathrm{imp}}(\lambda ^{^{}}).`$ (54)
Here we have introduced the phase function,
$`{\displaystyle \frac{1}{2\pi }}\mathrm{\Theta }_{n,k}^{^{}}(\lambda )={\displaystyle \frac{1}{\pi }}{\displaystyle \underset{l=1}{\overset{\mathrm{min}(n,k)}{}}}{\displaystyle \frac{{\displaystyle \frac{1}{2}}(n+k+12l)}{\lambda ^2+{\displaystyle \frac{1}{4}}(n+k+12l)^2}}.`$ (55)
These quantities are alternatively represented in terms of the phase shifts $`\delta _L`$ and $`\delta _R`$ at the left and right Fermi points in massless spin excitations: $`n_{\mathrm{imp}}(\nu )=(\delta _L+\delta _R)/2\pi ,d_{\mathrm{imp}}(\nu )=(\delta _L\delta _R)/2\pi `$. We mention that the asymmetric phase shift $`d_{\mathrm{imp}}(\nu )`$ is inherent in a mobile Kondo impurity, different from a localized impurity in §2. Note that the scaling dimension $`x`$ depends on the hole momentum $`q`$ through the asymmetric phase shift.
Let us now discuss low-energy critical properties in the photoemission spectra. We write down the one-particle Green function which depends on the momentum $`q`$,
$`\mathrm{Im}G(q,\omega )(\omega \omega _c(q))^{X(q)},`$ (56)
with $`X(q)=2x1`$, where $`x`$ is the scaling dimension in eq. (45) and $`\omega _c(q)`$ is the dispersion of the charge excitation generated by the photoemission. In Fig. 4 we show the obtained critical exponent as a function of the momentum $`q`$ for the $`S=1`$ case.
This anomalous power-law behavior and the momentum dependence of $`X(q)`$ are caused by a suddenly induced mobile Kondo impurity. As we discussed in §2, the discontinuity of the exponent $`X`$ at $`m=0`$ is caused by the $`\mathrm{Z}_\mathrm{n}`$ symmetric sector..
In this way, the dynamically induced Kondo effect proposed here may be expected to be observed in photoemission experiments for quantum spin chains. In order for this phenomenon to be observed, there may be several problems to be resolved; preparation of a proper 1D sample, experiments with high resolution, etc. Anyway, it is desired to experimentally find or synthesize rather ideal spin chain systems without the magnetic order even at low temperatures.
## IV summary
In summary we have proposed the multi-channel Kondo effect dynamically induced by the photoemission and the inverse photoemission, for which the ground state is a completely screened Kondo singlet. By studying low-energy critical properties in the photoemission spectra, it has been found that the anomalous behavior generated by this effect is indeed characteristic of the multi-channel Kondo system. In particular, we have demonstrated that the idea proposed here can be directly applied to homogeneous quantum spin systems without any impurity. It has been shown in this case that a mobile Kondo impurity suddenly created by the photoemission gives rise to the momentum-dependent anomalous exponent.
Although we have mainly focused on the photoemission spectra for magnetic impurity systems in this paper, it should be noted that the idea can be directly applied to the photoemission spectrum in quantum dot systems, for which a multilevel quantum dot with the Hund coupling plays a role of the Kondo impurity. In this connection, it is also interesting to apply a similar idea to the optical absorption spectra in a quantum dot, as recently demonstrated. In such cases, not only the linear but also the non-linear optics play an important role, which may provide interesting phenomena related to the dynamically induced Kondo effect.
## Acknowledgements
This work was partly supported by a Grant-in-Aid from the Ministry of Education, Science, Sports and Culture, Japan.
|
warning/0005/math0005254.html
|
ar5iv
|
text
|
# Semi-Riemannian submersions with totally geodesic fibres
## 1. Introduction and main results
Riemannian submersions, introduced by O’Neill \[One1\] and Gray \[Gra\], have been used by many authors to construct specific Riemannian metrics. A systematic exposition can be found in Besse’s book \[Bes\]. In this paper, we obtain classification results for semi-Riemannian submersions with totally geodesic fibres.
We first recall briefly some related work on the classification problem of semi-Riemannian submersions. Escobales \[Esc1, Esc2\] and Ranjan \[Ran1\] classified Riemannian submersions with connected totally geodesic fibres from an $`n`$-sphere $`S^n`$, and with connected complex totally geodesic fibres from a complex projective $`n`$-space $`P^n`$, respectively. Ucci \[Ucc\] showed that there are no Riemannian submersions with fibres $`P^3`$ from the complex projective space $`P^7`$ onto $`S^8(4)`$, and with fibres $`P^1`$ from the quaternionic projective space $`P^3`$ onto $`S^8(4)`$. In \[Ran2\], Ranjan obtained a classification theorem for Riemannian submersions with connected totally geodesic fibres from a compact simple Lie group. Gromoll and Grove obtained in \[G-G1\] that, up to equivalence, the only Riemannian submersions of spheres (with connected fibres) are the Hopf fibrations, except possibly for fibrations of the $`15`$-sphere by homotopy $`7`$-spheres. This classification was invoked in the proof of the Diameter Rigidity Theorem (see \[G-G2\]) and of the Radius Rigidity Theorem (see \[Wil\]). Using an approach different from Gromoll and Grove \[G-G1\], Wilking \[Wilk\] proved that a Riemannian submersion $`\pi :S^mB^b`$ is metrically equivalent to the Hopf fibration for $`(m,b)=(15,8)`$ and obtained an improved version of the Diameter Rigidity Theorem as a consequence of his classification theorem.
In comparison, there are few classification results for semi-Riemannian submersions, and the consequences seem to be at least as important as those for Riemannian submersions. In \[Mag\], Magid proved that the only semi-Riemannian submersions with totally geodesic fibres from an anti-de Sitter space onto a Riemannian manifold are the canonical semi-Riemannian submersions $`H_1^{2m+1}H^m`$. In \[Ba-Ia\], the present author and Stere Ianuş classified semi-Riemannian submersions with connected totally geodesic fibres from a pseudo-hyperbolic space onto a Riemannian manifold, and with connected complex totally geodesic fibres from a complex pseudo-hyperbolic space onto a Riemannian manifold.
The aim of this work is to prove new classification results in the theory of semi-Riemannian submersions analogous to those in Riemannian geometry. It is my pleasure to thank Professor Stere Ianuş for useful discussions on this subject.
Now, we list the main results proved in this paper.
###### Theorem 1.1.
Let $`\pi :H_{s+r^{}}^{n+r}B_s^n`$ be a semi-Riemannian submersion with connected totally geodesic fibres from a pseudo-hyperbolic space onto a semi-Riemannian manifold. If the dimension of the fibres is less than or equal to $`3`$ and if the metrics induced on fibres are negative definite, then $`\pi `$ is equivalent to one of the following canonical semi-Riemannian submersions$`:`$
* $`H_{2t+1}^{2m+1}H_t^m,0tm.`$
* $`H_{4t+3}^{4m+3}H_t^m,0tm.`$
###### Theorem 1.2.
Let $`\pi :H_{s+r^{}}^{n+r}B_s^n`$ be a semi-Riemannian submersion with connected totally geodesic fibres from a pseudo-hyperbolic space onto a semi-Riemannian manifold. Assume that one of the following conditions is satisfied$`:`$
* $`B`$ is an isotropic semi-Riemannian manifold, which means that for any $`xB_s^n`$ and any real number $`t`$, the group of isometries $`𝐈(B_s^n,g^{})`$ preserving $`x`$ acts transitively on the set of all nonzero tangent vectors $`X`$ at $`x`$ for which $`g^{}(X,X)=t`$, or
* $`\mathrm{index}(B)\{0,dimB\}`$.
Then $`\pi `$ is equivalent to one of the following canonical semi-Riemannian submersions$`:`$
* $`H_{2t+1}^{2m+1}H_t^m,0tm`$.
* $`H_{4t+3}^{4m+3}H_t^m,0tm`$.
* $`H_{7+8t}^{15}H_{8t}^8(4),t\{0,1\}`$.
###### Theorem 1.3.
Let $`\pi :H_s^nB`$ be a semi-Riemannian submersion from a complex pseudo-hyperbolic space onto a semi-Riemannian manifold. Assume that the fibres are connected complex totally geodesic submanifolds, and one of the following conditions is satisfied$`:`$
* The real dimension of the fibres is $`r2`$ and the fibres are negative definite, or
* $`B`$ is an isotropic semi-Riemannian manifold, or
* $`\mathrm{index}(B)\{0,dimB\}`$.
Then $`\pi `$ is equivalent to the canonical semi-Riemannian submersion
* $`H_{2t+1}^{2m+1}H_t^m,0tm`$.
###### Theorem 1.4.
There exist no semi-Riemannian submersions $`\pi :H_s^nB`$ with connected quaternionic fibres from a quaternionic pseudo-hyperbolic space onto an isotropic semi-Riemannian manifold or onto a semi-Riemannian manifold of $`\mathrm{index}(B)\{0,dim(B)\}`$.
## 2. Preliminaries and examples
In this section we recall several notions and results which will be needed throughout the paper. We also exhibit the construction of canonical semi-Riemannian submersions.
###### Definition 2.1.
Let $`(M,g)`$ be an $`(n+r)`$-dimensional connected semi-Riemannian manifold of index $`s+r^{}`$, and $`(B,g^{})`$ an $`n`$-dimensional connected semi-Riemannian manifold of index $`s`$, where $`0sn`$, $`0r^{}r`$. A semi-Riemannian submersion (see \[One2\]) is a smooth map $`\pi :MB`$ which is surjective and satisfies the following axioms:
* $`\pi _{}|_p`$ is surjective for all $`pM`$;
* the fibres $`\pi ^1(b),bB`$, are semi-Riemannian submanifolds of $`M`$;
* $`\pi _{}`$ preserves scalar products of vectors normal to fibres.
We shall always assume that the fibres are connected, the dimension of the fibres $`dimMdimB>0`$ and $`dimB>0`$. The vectors tangent to fibres are called vertical and those normal to fibres are called horizontal. We denote by $`𝒱`$ the vertical distribution and by $``$ the horizontal distribution.
The geometry of semi-Riemannian submersions is characterized by O’Neill’s tensors $`T`$, $`A`$ (see \[One1\], \[One2\]) defined for vector fields $`E`$, $`F`$ on $`M`$ by
$`A_EF`$ $`=`$ $`h_{hE}vF+v_{hE}hF,`$
$`T_EF`$ $`=`$ $`h_{vE}vF+v_{vE}hF,`$
where $``$ is the Levi-Civita connection of $`g`$, and $`v`$ and $`h`$ denote the orthogonal projections on $`𝒱`$ and $``$, respectively. For basic properties of O’Neill’s tensors see \[One1\], \[One2\], \[Bes\] or \[Ian\].
###### Definition 2.2.
(i) A vector field $`X`$ on $`M`$ is said to be basic if $`X`$ is horizontal and $`\pi `$-related to a vector field $`X^{}`$ on $`B`$.
(ii) A vector field $`X`$ along the fibre $`\pi ^1(x)`$, $`xB`$, is said to be basic along $`\pi ^1(x)`$ if $`X`$ is horizontal and $`\pi _pX(p)=\pi _qX(q)`$ for every $`p`$, $`q\pi ^1(x)`$.
We notice that each vector field $`X^{}`$ on $`B`$ has a unique horizontal lift $`X`$ to $`M`$ which is basic. For a vertical vector field $`V`$ and a basic vector field $`X`$ we have $`h_VX=A_XV`$ (see \[One1\]). We denote by $`R`$, $`R^{}`$ and $`\widehat{R}`$ the Riemann curvature tensors of $`M`$, $`B`$ and of the fibre $`\pi ^1(x)`$, $`xM`$, respectively. We choose the convention for the curvature tensor $`R(E,F)=_E_F_F_E_{[E,F]}`$. The Riemann curvature tensor is defined by
$$R(E,F,G,H)=g(R(G,H)F,E).$$
For O’Neill’s equations of a semi-Riemannian submersion we refer to \[One1\] or \[Bes\].
###### Definition 2.3.
Two semi-Riemannian submersions $`\pi ,\pi ^{}:(M,g)(B,g^{})`$ are said to be equivalent if there exists an isometry $`f`$ of $`M`$ which induces an isometry $`\stackrel{~}{f}`$ of $`B`$ so that $`\pi ^{}f=\stackrel{~}{f}\pi `$. The pair $`(f,\stackrel{~}{f})`$ is called a bundle isometry.
We shall need the following theorem, which is the semi-Riemannian version of Theorem 2.2 in \[Esc1\].
###### Theorem 2.4.
Let $`\pi _1,\pi _2:MB`$ be semi-Riemannian submersions from a complete connected semi-Riemannian manifold $`M`$ onto a semi-Riemannian manifold $`B`$. Assume that the fibres of these submersions are connected and totally geodesic and the metric induced on fibres is negative definite. Let $`f`$ be an isometry of $`M`$ satisfying the following properties at a given point $`pM:`$
* $`f_p:T_pMT_{f(p)}M`$ maps $`_{1p}`$ onto $`_{2f(p)}`$, where $`_i`$ denote the horizontal distributions of $`\pi _i`$ for $`i\{1,2\}`$.
* $`f_{}A_{1E}F=A_{2f_{}E}f_{}F`$ for every $`E`$, $`FT_pM`$, where $`A_i`$ are the integrability tensors associated with $`\pi _i`$.
Then $`f`$ induces an isometry $`\stackrel{~}{f}`$ of $`B`$ so that the pair $`(f,\stackrel{~}{f})`$ is a bundle isometry between $`\pi _1`$ and $`\pi _2`$. In particular, $`\pi _1`$ and $`\pi _2`$ are equivalent.
Escobales’s proof of Theorem 2.2 in \[Esc1\], also works in this semi-Riemannian case. He proves that for any $`bB`$ which can be joined with $`\pi _1(p)`$ by a geodesic we have:
* for every $`x\pi _1^1(b)`$, $`f_x:T_xMT_{f(x)}M`$ maps $`_{1x}`$ onto $`_{2f(x)}`$, and
* $`f`$ maps the fibre $`\pi _1^1(b)`$ into the fibre $`\pi _2^1(\pi _2(f(x)))`$ with $`x\pi _1^1(b)`$.
We notice that for any $`x\pi _1^1(b)`$ with $`bB`$, which can be joined with $`\pi _1(p)`$ by a geodesic, the conditions (1) and (2) are also satisfied for the point $`x`$. Since $`M`$ is connected, $`B`$ is also connected. Therefore, any point $`bB`$ can be joined with $`\pi _1(p)`$ by a broken geodesic. Repeating the argument above, for any corner point of this broken geodesic, we see that for any $`bB`$, $`f`$ maps the fibre $`\pi _1^1(b)`$ into a fibre.
###### Definition 2.5.
Let $`,`$ be the symmetric bilinear form on $`^{m+1}`$ given by
$$x,y=\underset{i=0}{\overset{s}{}}x_iy_i+\underset{i=s+1}{\overset{m}{}}x_iy_i$$
for $`x=(x_0,\mathrm{},x_m),y=(y_0,\mathrm{},y_m)^{m+1}`$. For any $`c<0`$ and any positive integer $`s`$, let $`H_s^m(c)=\{x^{m+1}|x,x=1/c\}`$ be the semi-Riemannian submanifold of
$$_{s+1}^{m+1}=(^{m+1},ds^2=dx^0dx^0\mathrm{}dx^sdx^s+dx^{s+1}dx^{s+1}+\mathrm{}+dx^mdx^m).$$
$`H_s^m(c)`$ is called the $`m`$-dimensional (real) pseudo-hyperbolic space of index $`s`$.
We notice that $`H_s^m(c)`$ has constant sectional curvature $`c`$, whose curvature tensor is given by $`R(X,Y,X,Y)=c(g(X,X)g(Y,Y)g(X,Y)^2)`$. We shall denote simply $`H_s^m=H_s^m(1)`$. It should be remarked that $`H_s^m`$ can be written as a homogeneous space, namely $`H_s^m=SO(s+1,ms)/SO(s,ms)`$, $`H_{2s+1}^{2m+1}=SU(s+1,ms)/SU(s,ms)`$, and $`H_{4s+3}^{4m+3}=Sp(s+1,ms)/Sp(s,ms)`$ (see \[Wol\]).
###### Definition 2.6.
Let $`(,)`$ be the Hermitian form on $`^{m+1}`$ given by
$$(z,w)=\underset{i=0}{\overset{s}{}}z_i\overline{w_i}+\underset{i=s+1}{\overset{m}{}}z_i\overline{w_i}$$
for $`z=(z_0,\mathrm{},z_m),w=(w_0,\mathrm{},w_m)^{m+1}`$. For $`c<0`$, let $`M(c)`$ be the real hypersurface of $`^{m+1}`$ given by $`M(c)=\{z^{m+1}|(z,z)=4/c\}`$, which is endowed with the induced metric of
$$(^{m+1},ds^2=dz^0d\overline{z}^0\mathrm{}dz^sd\overline{z}^s+dz^{s+1}d\overline{z}^{s+1}+\mathrm{}+dz^md\overline{z}^m).$$
The natural action of $`S^1=\{e^{i\theta }|\theta \}`$ on $`^{m+1}`$ induces an action on $`M(c)`$. Let $`H_s^m(c)=M(c)/S^1`$ endowed with the unique indefinite Kähler metric of index $`2s`$ such that the projection $`M(c)M(c)/S^1`$ becomes a semi-Riemannian submersion (see \[Ba-Ro\]). $`H_s^m(c)`$ is called the complex pseudo-hyperbolic space.
Notice that $`H_s^m(c)`$ has constant holomorphic sectional curvature $`c`$, whose curvature tensor is given by $`R(X,Y,X,Y)=(c/4)(g(X,X)g(Y,Y)g(X,Y)^2+3g(I_0X,Y)^2)`$, where $`I_0`$ is the natural complex structure on $`H_s^m(c)`$. We shall denote simply $`H_s^m=H_s^m(4)`$. It is well-known that $`H_s^m`$ is a homogeneous space, namely $`H_s^m=SU(s+1,ms)/S(U(1)U(s,ms))`$ and $`H_{2s+1}^{2m+1}=Sp(s+1,ms)/U(1)Sp(s,ms)`$ (see \[Wol\]).
We shall denote by $`H_s^n`$ the quaternionic pseudo-hyperbolic space of real dimension $`4n`$, and of quaternionic index $`s`$ with quaternionic sectional curvature $`4`$, and by $`S^n`$ and $`S^n(4)`$ the spheres with sectional curvature $`1`$ and $`4`$, respectively.
By a standard construction (see Theorem 9.80 in \[Bes\]), one can obtain many examples of semi-Riemannian submersions with totally geodesic fibres of type $`\pi :G/KG/H`$, where $`G`$ is a Lie group and $`K`$, $`H`$ are closed Lie subgroups of $`G`$ with $`KH`$. In this way the following canonical semi-Riemannian submersions, also called generalized Hopf fibrations, are obtained:
###### Example 1.
Let $`G=SU(t+1,mt)`$, $`H=S(U(1)U(t,mt))`$, $`K=SU(t,mt)`$. For every $`0tm`$, we have the semi-Riemannian submersion
$$H_{2t+1}^{2m+1}=SU(t+1,mt)/SU(t,mt)H_t^m=SU(t+1,mt)/S(U(1)U(t,mt)).$$
###### Example 2.
Let $`G=Sp(t+1,mt)`$, $`H=Sp(1)Sp(t,mt)`$, $`K=Sp(t,mt)`$. For every $`0tm`$, we get the semi-Riemannian submersion
$$H_{4t+3}^{4m+3}=Sp(t+1,mt)/Sp(t,mt)H_t^m=Sp(t+1,mt)/Sp(1)Sp(t,mt).$$
###### Example 3.
a) Let $`G=Spin(1,8)`$, $`H=Spin(8)`$, $`K=Spin(7)`$. Then we have the semi-Riemannian submersion (see \[Ba-Ia\])
$$H_7^{15}=Spin(1,8)/Spin(7)H^8(4)=Spin(1,8)/Spin(8).$$
b) Let $`G=Spin(9)`$, $`H=Spin(8)`$, $`K=Spin(7)`$. Then we have the semi-Riemannian submersion (see \[Bes\])
$$S^{15}=Spin(9)/Spin(7)S^8(4)=Spin(9)/Spin(8).$$
###### Example 4.
Let $`G=Sp(t+1,mt)`$, $`H=Sp(1)Sp(t,mt)`$, $`K=U(1)Sp(t,mt)`$. For every $`0tm`$, we obtain the semi-Riemannian submersion
$$H_{2t+1}^{2m+1}=Sp(t+1,mt)/U(1)Sp(t,mt)H_t^m=Sp(t+1,mt)/Sp(1)Sp(t,mt).$$
In order to prove Theorem 1.2, we need the following nonexistence proposition, which is the semi-Riemannian version of Proposition 5.1 in \[Ran1\].
###### Proposition 2.7.
There exist no semi-Riemannian submersions $`\pi :H_{7+8t}^{23}aH_t^2`$, $`t\{0,1,2\}`$, with totally geodesic fibres from the $`23`$-dimensional pseudo-hyperbolic space of index $`7+8t`$ onto the Cayley pseudo-hyperbolic plane of Cayley index $`t`$ .
We notice that the case $`t=2`$ is Proposition 5.1 in \[Ran1\]. For the case $`t=0`$, see \[Ba-Ia\]. Here we only recall some details of Ranjan’s proof and suggest its modification to the semi-Riemannian case. Ranjan’s argument in \[Ran1\], which leads to a contradiction to the assumption of the existence of such a submersion, is based on finding for every $`X_p`$, $`g(X,X)0`$, an irreducible $`Cl(𝒱_p)`$-submodule $`S`$ of $`_p`$ passing through $`X`$. Here $`Cl(𝒱_p)`$ denotes the Clifford algebra of $`(𝒱_p,\stackrel{~}{g}_p)`$, where $`\stackrel{~}{g}(U,V)=g(U,V)`$ for every $`U`$, $`V𝒱_p`$. $`_p`$ becomes a $`Cl(𝒱_p)`$-module by considering the extension of the map $`𝒰:𝒱_p\mathrm{End}(_p)`$ defined by $`𝒰(V)(X)=A_XV`$ to the Clifford algebra $`Cl(𝒱_p)`$. Since $`\stackrel{~}{g}_p`$ is positive definite, we have $`Cl(𝒱_p)(8)(8)`$. Hence, $`_p`$ splits into two $`8`$-dimensional irreducible $`Cl(𝒱_p)`$-modules. Since the induced metrics on fibres are negative definite, we obtain in a manner similar to Ranjan’s proof that
* for $`g(X,X)>0`$, $`\pi ^1(aH^1)`$ is totally geodesic in $`H_{7+8t}^{23}`$ and is isometric to $`H_7^{15}`$, where $`aH^1`$ denotes the Cayley hyperbolic line through $`\pi _{}X`$, and
* for $`g(X,X)<0`$, $`\pi ^1(aH_1^1)`$ is totally geodesic in $`H_{7+8t}^{23}`$ and is isometric to $`H_{15}^{15}`$, where $`aH_1^1`$ denotes the negative definite Cayley hyperbolic line through $`\pi _{}X`$.
We choose $`S`$ to be the horizontal space of the restricted submersion $`\stackrel{~}{\pi }:H_7^{15}aH^1`$ if $`g(X,X)>0`$ or $`\stackrel{~}{\pi }:H_{15}^{15}aH_1^1`$ if $`g(X,X)<0`$.
## 3. Proof of the main results
The next lemma gives useful properties of O’Neill’s integrability tensor.
###### Lemma 3.1.
Let $`\pi :MB`$ be a semi-Riemannian submersion with connected totally geodesic fibres from a semi-Riemannian manifold $`M`$ with constant curvature $`c0`$. Then the following assertions are true$`:`$
* If $`X`$ is a horizontal vector such that $`g(X,X)0`$, then the map $`A_X:𝒱`$ given by $`A_X(V)=A_XV`$ is injective and the map $`A_X^{}:𝒱`$ given by $`A_X^{}(Y)=A_XY`$ is surjective.
* If $`X`$, $`Y`$ are the horizontal liftings along the fibre $`\pi ^1(\pi (p))`$, $`pM`$, of two vectors $`X^{},Y^{}T_{\pi (p)}B`$ respectively, $`g^{}(X^{},X^{})0`$ and $`(A_XY)(p)=0`$, then $`A_XY=0`$ along the fibre $`\pi ^1(\pi (p))`$.
###### Proof.
(a) By O’Neill’s equations, we get
$$g(A_XV,A_XW)=cg(X,X)g(V,W)$$
for a horizontal vector field $`X`$ and for vertical vector fields $`V`$ and $`W`$. Thus $`A_X^{}A_XV=cg(X,X)V`$ for every vertical vector field $`V`$. Therefore $`A_X:𝒱`$ is injective and $`A_X^{}:𝒱`$ is surjective.
(b) By O’Neill’s equations, we have
$$3g(A_XY,A_XZ)=c[g(X,X)g(Y,Z)g(X,Y)g(X,Z)]R^{}(\pi _{}X,\pi _{}Y,\pi _{}X,\pi _{}Z)$$
for horizontal vector fields $`X`$, $`Y`$ and $`Z`$.
If $`X`$, $`Y`$, $`Z`$ are basic vector fields, then $`g(A_XY,A_XZ)`$ is constant along the fibre $`\pi ^1(\pi (p))`$. Therefore, $`g(A_XA_XY,Z)=0`$ along the fibre $`\pi ^1(\pi (p))`$ for every basic vector field $`Z`$. Hence $`A_XA_XY=0`$ along $`\pi ^1(\pi (p))`$. Since $`A_X:𝒱`$ is injective, it follows that $`A_XY=0`$ along the fibre $`\pi ^1(\pi (p))`$. ∎
###### Lemma 3.2.
If $`\pi :MB`$ is a semi-Riemannian submersion with connected totally geodesic fibres from a semi-Riemannian manifold $`M`$ with constant curvature $`c0`$ onto a semi-Riemannian manifold $`B`$, then the tangent bundle of any fibre is trivial.
###### Proof.
Let $`xB`$ and $`p\pi ^1(x)`$. Let $`\{v_{1p},\mathrm{},v_{rp}\}`$ be an orthonormal basis in $`𝒱_p`$. Let $`Y_1,Y_2,\mathrm{},Y_r`$ be the horizontal liftings along the fibre $`\pi ^1(\pi (p))`$ of $`(1/(cg(X,X)))\pi _{}A_Xv_{1p}`$, $`(1/(cg(X,X)))\pi _{}A_Xv_{2p}`$,…, $`(1/(cg(X,X)))\pi _{}A_Xv_{rp}`$, respectively. Let $`v_i=A_XY_i`$ for each $`i\{1,\mathrm{},r\}`$. Since
$`g(v_j,v_l)`$ $`=`$ $`g(A_XY_j,A_XY_l)`$
$`=`$ $`(1/3)(R^{}(\pi _{}X,\pi _{}Y_j,\pi _{}X,\pi _{}Y_l)cg(X,X)g(Y_j,Y_l)+cg(X,Y_j)g(X,Y_l))`$
is constant along the fibre $`\pi ^1(\pi (p))`$ and
$$g(A_XY_j,A_XY_l)(p)=\frac{1}{c^2}g(A_XA_Xv_{jp},A_XA_Xv_{lp})=g(X,X)^2g(v_{jp},v_{lp})=\epsilon _j\delta _{jl},$$
we see that $`\{v_1,v_2,\mathrm{},v_r\}`$ is a global orthonormal basis of the tangent bundle of the fibre $`\pi ^1(x)`$, which makes the tangent bundle trivial. ∎
We suppose that the curvature of the total space is negative. The case of positive curvature can be reduced to the negative one by changing simultaneously the signs of the metrics on the base and on the total space. We establish relations between the dimensions and the indices of fibres and of base spaces, and see how the geometry of base spaces looks like.
###### Theorem 3.3.
Let $`\pi :MB`$ be a semi-Riemannian submersion with connected totally geodesic fibres from an $`(n+r)`$-dimensional semi-Riemannian manifold $`M`$ of index $`s+r^{}`$ with constant negative curvature $`c`$ onto an $`n`$-dimensional semi-Riemannian manifold $`B`$ of index $`s`$. Then the following hold$`:`$
* $`n=k(r+1)`$ for some positive integer $`k`$ and $`s=q_1(r^{}+1)+q_2(rr^{})`$ for some nonnegative integers $`q_1`$, $`q_2`$ with $`q_1+q_2=k`$.
* If, moreover, $`M`$ is a simply connected complete semi-Riemannian manifold and the dimension of fibres is less than or equal to $`3`$ and the metric induced on fibres is negative definite, then $`B`$ is an isotropic semi-Riemannian manifold and $`r\{1,3\}`$.
###### Proof.
Normalizing the metric on $`M`$, we can suppose $`c=1`$. Let $`pM`$. Since the tangent bundle of the fibre $`\pi ^1(\pi (p))`$ is trivial, we can choose a global orthonormal frame $`\{v_1,v_2,\mathrm{},v_r\}`$ for the tangent bundle of $`\pi ^1(\pi (p))`$. We have $`g(v_i,v_j)=\epsilon _i\delta _{ij}`$, $`\epsilon _i\{1,1\}`$, and card$`\{i|\epsilon _i<0\}=r^{}`$.
(1) Let $`X`$ be the horizontal lifting along the fibre $`\pi ^1(\pi (p))`$ of a vector $`X^{}T_{\pi (p)}B`$, so that $`g(X^{},X^{})\{1,1\}`$. By O’Neill’s equations, we have
$$g(A_YV,A_YV)=g(Y,Y)g(V,V)$$
for a horizontal vector field $`Y`$ and for a vertical vector field $`V`$. Along the fibre $`\pi ^1(\pi (p))`$ we obtain for every $`i,j\{1,\mathrm{},r\}`$
$$g(A_Xv_i,A_Xv_j)=g(X,X)g(v_i,v_j)=g(X,X)\epsilon _i\delta _{ij},$$
$$g(X,A_Xv_i)=g(A_XX,v_i)=0.$$
Thus $`\{X,A_Xv_1,\mathrm{},A_Xv_r\}`$ is an orthonormal system. Hence $`nr+1`$.
Let $`L_0=X`$. For every integer $`\alpha `$ such that $`1\alpha <n/(r+1)`$, let $`L_\alpha `$ be a horizontal vector field along the fibre $`\pi ^1(\pi (p))`$ such that $`L_\alpha `$ is the horizontal lifting of some unit vector (i.e., $`g(L_\alpha ,L_\alpha )\{1,1\}`$), that $`L_\alpha `$ is orthogonal to $`L_0,L_1,\mathrm{},L_{\alpha 1}`$ and that $`L_\alpha (p)\mathrm{ker}A_{L_0(p)}^{}\mathrm{ker}A_{L_1(p)}^{}\mathrm{}\mathrm{ker}A_{L_{\alpha 1}(p)}^{}`$. Then, by Lemma 3.1, $`L_\alpha (q)`$ belongs to $`\mathrm{ker}A_{L_0(q)}^{}\mathrm{ker}A_{L_1(q)}^{}\mathrm{}\mathrm{ker}A_{L_{\alpha 1}(q)}^{}`$ for every $`q\pi ^1(\pi (p))`$. Therefore, for $`j\{1,\mathrm{},r\}`$ and $`\alpha ,\beta 0`$, we get
$$g(A_{L_\alpha }v_j,L_\beta )=g(v_j,A_{L_\alpha }L_\beta )=0$$
along the fibre $`\pi ^1(\pi (p))`$.
By O’Neill’s equations, we obtain
(3.1)
$$\begin{array}{cc}\hfill R(X,U,Y,V)& =g((_UA)_XY,V)+g(A_XU,A_YV)\hfill \\ & =g(_UA_XY,V)g(A_{_UX}Y,V)g(A_X_UY,V)+g(A_XU,A_YV)\hfill \\ & =g(_UA_XY,V)+g(A_YA_XU,V)g(A_XA_YU,V)g(A_YA_XU,V)\hfill \\ & =g(_UA_XY,V)+g(A_YU,A_XV)\hfill \end{array}$$
for basic vector fields $`X`$, $`Y`$ and for vertical vector fields $`U`$, $`V`$. Thus, along the fibre $`\pi ^1(\pi (p))`$ we get for every $`\alpha ,\beta 0`$ and $`j,l\{1,\mathrm{},r\}`$
$`g(A_{L_\alpha }v_j,A_{L_\beta }v_l)`$ $`=`$ $`R(L_\alpha ,v_l,L_\beta ,v_j)g(_{v_l}A_{L_\alpha }L_\beta ,v_j)`$
$`=`$ $`g(L_\alpha ,L_\beta )g(v_l,v_j)v_l(g(A_{L_\alpha }L_\beta ,v_j))+g(A_{L_\alpha }L_\beta ,_{v_l}v_j).`$
Since $`A_{L_\alpha }L_\beta =0`$ along the fibre $`\pi ^1(\pi (p))`$, it follows that
$$g(A_{L_\alpha }v_j,A_{L_\beta }v_l)=g(L_\alpha ,L_\beta )g(v_l,v_j)=g(L_\alpha ,L_\beta )\epsilon _l\delta _{lj}.$$
We proved that for some positive integer $`k`$,
$$=\{L_0,A_{L_0}v_1,\mathrm{},A_{L_0}v_r,\mathrm{},L_{k1},A_{L_{k1}}v_1,\mathrm{},A_{L_{k1}}v_r\}$$
is an orthonormal basis of $``$ along the fibre $`\pi ^1(\pi (p))`$. Thus $`dimB=(1+dim\mathrm{fibre})k`$ for some positive integer $`k`$. Counting the timelike vectors in $``$, we get $`\mathrm{index}(B)=q_1(r^{}+1)+q_2(rr^{})`$ for some nonnegative integers $`q_1`$, $`q_2`$ with $`q_1+q_2=k`$.
(2) Let $`xB`$ and $`X^{}`$, $`Y^{}T_xB`$ such that $`g^{}(X^{},X^{})=g^{}(Y^{},Y^{})0`$. We shall construct an isometry $`\stackrel{~}{f}:BB`$ such that $`\stackrel{~}{f}(x)=x`$ and $`\stackrel{~}{f}_{}X^{}=Y^{}`$. Note that we may assume that $`g^{}(X^{},X^{})=g^{}(Y^{},Y^{})=\pm 1`$. Let $`X`$, $`Y`$ be the horizontal liftings along the fibre $`\pi ^1(x)`$ of $`X^{}`$ and $`Y^{}`$, respectively. Take $`p\pi ^1(x)`$. Let
$$=\{L_0,A_{L_0}v_1,\mathrm{},A_{L_0}v_r,\mathrm{},L_{k1},A_{L_{k1}}v_1,\mathrm{},A_{L_{k1}}v_r\},$$
$$^{}=\{L_0^{},A_{L_0^{}}v_1^{},\mathrm{},A_{L_0^{}}v_r^{},\mathrm{},L_{k1}^{},A_{L_{k1}^{}}v_1^{},\mathrm{},A_{L_{k1}^{}}v_r^{}\}$$
be two orthonormal bases constructed as above such that $`L_0=X`$, $`L_0^{}=Y`$, $`g(L_\alpha ,L_\alpha )=g(L_\alpha ^{},L_\alpha ^{})`$ for $`\alpha \{1,\mathrm{},k1\}`$, and that $`\{v_1=A_XY_1,\mathrm{},v_r=A_XY_r\}`$ and $`\{v_1^{}=A_YY_1^{},\mathrm{},v_r^{}=A_YY_r^{}\}`$ are orthonormal bases of the tangent bundle of the fibre $`\pi ^1(\pi (p))`$, where $`Y_1,\mathrm{},Y_r`$ and $`Y_1^{},\mathrm{},Y_r^{}`$ are the horizontal liftings along $`\pi ^1(\pi (p))`$ of the vectors $`\pi _{}A_Xv_{1p},\mathrm{},\pi _{}A_Xv_{rp}`$ and $`\pi _{}A_Yv_{1p}^{},\mathrm{},\pi _{}A_Yv_{rp}^{}`$, respectively (as in Lemma 3.1), for which $`g(v_i,v_j)=g(v_i^{},v_j^{})`$ for $`i,j\{1,\mathrm{},r\}`$. Let $`\varphi :T_pMT_pM`$ be the linear map given by $`\varphi (L_\alpha )=L_\alpha ^{}`$, $`\varphi (v_j)=v_j^{}`$, $`\varphi (A_{L_\alpha }v_j)=A_{L_\alpha ^{}}v_j^{}`$ for every $`\alpha \{0,\mathrm{},k1\}`$ and $`j\{1,\mathrm{},r\}`$. Since both $``$, $`^{}`$ are orthonormal bases, we see that $`\varphi `$ is a linear isometry.
We shall apply Theorem 2.4. Thus we need to prove that $`\varphi (A_EF)=A_{\varphi (E)}\varphi (F)`$ for every $`E`$, $`FT_pM`$. Indeed, we obtain for $`\alpha ,\beta \{0,\mathrm{},k1\}`$ and $`j,l\{1,\mathrm{},r\}`$,
$`\varphi (A_{L_\alpha }L_\beta )`$ $`=`$ $`\varphi (0)=0=A_{L_\alpha ^{}}L_\beta ^{}=A_{\varphi (L_\alpha )}\varphi (L_\beta ),`$
$`g(v_j,A_{L_\alpha }A_{L_\beta }v_l)`$ $`=`$ $`g(A_{L_\alpha }v_j,A_{L_\beta }v_l)=g(L_\alpha ,L_\beta )g(v_j,v_l)`$
$`=`$ $`g(L_\alpha ^{},L_\beta ^{})g(v_j^{},v_l^{})=g(v_j^{},A_{L_\alpha ^{}}A_{L_\beta ^{}}v_l^{}).`$
Hence $`\varphi (A_{L_\alpha }A_{L_\beta }v_l)=A_{\varphi (L_\alpha )}\varphi (A_{L_\beta }v_l)`$.
###### Lemma 3.4.
$`A_{L_\alpha }v_j`$ is a basic vector field along the fibre $`\pi ^1(\pi (p))`$ for every $`1jr`$ and $`\alpha 0`$.
###### Proof of Lemma 3.4.
We have $`g(A_Xv_j,Z)=g(A_XA_XY_j,Z)=g(A_XY_j,A_XZ)`$. For every basic vector field $`Z`$ along the fibre $`\pi ^1(\pi (p))`$ we know that $`g(A_XY_j,A_XZ)`$ is constant along the fibre $`\pi ^1(\pi (p))`$. Hence $`A_Xv_j`$ is a basic vector field along the fibre $`\pi ^1(\pi (p))`$.
Now we assume $`\alpha 1`$. Since $`dim(\mathrm{ker}A_X^{}\mathrm{ker}A_{L_\alpha }^{})=dim\mathrm{ker}A_X^{}+dim\mathrm{ker}A_{L_\alpha }^{}dim(\mathrm{ker}A_X^{}\mathrm{ker}A_{L_\alpha }^{})=(nr)+(nr)(n2r)=n`$, it follows that $`\mathrm{ker}A_X^{}+\mathrm{ker}A_{L_\alpha }^{}=`$. Hence $`A_{L_\alpha }v_j`$ is a basic vector field along the fibre $`\pi ^1(\pi (p))`$ if and only if the following conditions are satisfied: $`g(A_{L_\alpha }v_j,Z_1)`$ is constant along $`\pi ^1(\pi (p))`$ for every $`Z_1\mathrm{ker}A_X^{}`$, which is a basic vector field along $`\pi ^1(\pi (p))`$, and $`g(A_{L_\alpha }v_j,Z_2)`$ is constant along the fibre $`\pi ^1(\pi (p))`$ for every $`Z_2\mathrm{ker}A_{L_\alpha }^{}`$, which is a basic vector field along $`\pi ^1(\pi (p))`$. If $`Z_2\mathrm{ker}A_{L_\alpha }^{}`$, then $`A_{L_\alpha }^{}Z_2=0`$ along $`\pi ^1(\pi (p))`$. So $`g(A_{L_\alpha }v_j,Z_2)=g(v_j,A_{L_\alpha }Z_2)=0`$ along $`\pi ^1(\pi (p))`$. If $`Z_1\mathrm{ker}A_X^{}`$, then $`A_X^{}Z_1=0`$ along $`\pi ^1(\pi (p))`$. By O’Neill’s equations, we get along the fibre $`\pi ^1(\pi (p))`$
$`R^{}(\pi _{}X,\pi _{}Y_j,\pi _{}L_\alpha ,\pi _{}Z_1)`$ $`=`$ $`R(X,Y_j,L_\alpha ,Z_1)+2g(A_XY_j,A_{L_\alpha }Z_1)`$
$`g(A_{Y_j}L_\alpha ,A_XZ_1)g(A_{L_\alpha }X,A_{Y_j}Z_1)`$
$`=`$ $`g(X,L_\alpha )g(Y_j,Z_1)+g(X,Z_1)g(Y_j,L_\alpha )`$
$`+2g(v_j,A_{L_\alpha }Z_1),`$
since $`A_{L_\alpha }X=A_XL_\alpha =0`$ and $`A_XZ_1=0`$. Hence $`g(v_j,A_{L_\alpha }Z_1)=g(A_{L_\alpha }v_j,Z_1)`$ is constant along $`\pi ^1(\pi (p))`$ for every $`Z_1\mathrm{ker}A_X^{}`$, which is a basic vector field along $`\pi ^1(\pi (p))`$.
We proved that $`A_{L_\alpha }v_j`$ is a basic vector field along $`\pi ^1(\pi (p))`$ for every $`\alpha 0`$ and $`j\{1,\mathrm{},r\}`$. ∎
We denote by $`\widehat{}`$ the induced Levi-Civita connection on the fibre $`\pi ^1(\pi (p))`$.
###### Lemma 3.5.
$`A_{A_{L_\alpha }v_i}A_{L_\beta }v_j=g(L_\alpha ,L_\beta )\widehat{}_{v_i}v_j.`$
###### Proof of Lemma 3.5.
By the relation (3.1) together with Lemma 3.4, we obtain for $`i,j,l\{1,\mathrm{},r\}`$ and $`\alpha ,\beta 0`$ that
$`g(A_{A_{L_\alpha }v_i}A_{L_\beta }v_j,v_l)`$ $`=`$ $`g(A_{A_{L_\alpha }v_i}v_l,A_{L_\beta }v_j)`$
$`=`$ $`R(L_\beta ,v_l,A_{L_\alpha }v_i,v_j)+g(_{v_l}A_{L_\beta }A_{L_\alpha }v_i,v_j)`$
$`=`$ $`g(L_\beta ,A_{L_\alpha }v_i)g(v_l,v_j)+v_lg(A_{L_\beta }A_{L_\alpha }v_i,v_j)`$
$`g(A_{L_\beta }A_{L_\alpha }v_i,_{v_l}v_j)`$
$`=`$ $`v_lg(A_{L_\alpha }v_i,A_{L_\beta }v_j)+g(A_{L_\alpha }v_i,A_{L_\beta }v_t)g(_{v_l}v_j,v_t)\epsilon _t`$
$`=`$ $`g(L_\alpha ,L_\beta )g(\widehat{}_{v_l}v_j,v_i)`$
$`=`$ $`g(L_\alpha ,L_\beta )g(\widehat{}_{v_i}v_j,v_l).`$
In the last equality we used the fact that $`v_j=A_XY_j`$ is a Killing vector field along the fibre $`\pi ^1(\pi (p))`$ (see \[Bis\] or \[Bes\]). Thus
$$A_{A_{L_\alpha }v_i}A_{L_\beta }v_j=g(L_\alpha ,L_\beta )\widehat{}_{v_i}v_j.$$
###### Lemma 3.6.
The following assertions are true$`:`$
* $`r2`$.
* If $`r=1`$, then $`A_{A_{L_\alpha }v_1}A_{L_\beta }v_1=0`$.
* If $`r=3`$ and if we set $`v_{3p}=(\widehat{}_{v_1}v_2)(p)`$, then $`v_3=\widehat{}_{v_1}v_2`$ and
$$g(\widehat{}_{v_i}v_j,v_k)=\{\begin{array}{cc}0& \text{if two of }i\text{}j\text{}k\text{ are equal,}\hfill \\ \epsilon \left(\genfrac{}{}{0pt}{}{\mathrm{1\; 2\; 3}}{ijk}\right)g(v_3,v_3)& \text{if }\{i,j,k\}=\{1,2,3\},\hfill \end{array}$$
where $`\epsilon \left(\genfrac{}{}{0pt}{}{\mathrm{1\; 2\; 3}}{ijk}\right)`$ is the signature of the permutation $`\left(\genfrac{}{}{0pt}{}{\mathrm{1\; 2\; 3}}{ijk}\right)`$.
###### Proof of Lemma 3.6.
Since $`v_1,\mathrm{},v_r`$ are Killing vector fields along $`\pi ^1(\pi (p))`$ and $`g(v_i,v_i)\{1,1\}`$ for every $`i`$, we get
$$g(\widehat{}_{v_i}v_j,v_i)=g(\widehat{}_{v_i}v_i,v_j)=g(\widehat{}_{v_j}v_i,v_i)=0$$
for every $`i,j\{1,\mathrm{},r\}`$.
(a) The case $`r=2`$ is not possible. Indeed, if $`r=2`$, then the relation $`g(_{v_1}v_2,v_1)=g(_{v_1}v_2,v_2)=0`$ implies $`_{v_1}v_2=0`$. On the other hand,
$$g(_{v_1}v_2,_{v_1}v_2)=g(\widehat{}_{v_1}\widehat{}_{v_2}v_2,v_1)+\widehat{R}(v_1,v_2,v_1,v_2)=g(v_1,v_1)g(v_2,v_2)\{1,1\},$$
since $`\widehat{}_{v_2}v_2=g(X,X)^1A_{A_Xv_2}A_Xv_2=0`$ and each fibre has constant curvature $`1`$. So we get a contradiction.
(b) If $`r=1`$, then $`A_{A_{L_\alpha }v_1}A_{L_\beta }v_1=0`$ for every $`\alpha `$ and $`\beta `$, because $`0=A_{A_Xv_1}A_Xv_1=g(X,X)_{v_1}v_1`$ implies $`_{v_1}v_1=0`$.
(c) In the case $`r=3`$ we shall prove $`g(\widehat{}_{v_1}v_2,v_3)`$ is constant along the fibre $`\pi ^1(\pi (p))`$. Since O’Neill’s integrability tensor $`A`$ is skew-symmetric, it follows that $`\widehat{}_{v_i}v_j=\widehat{}_{v_j}v_i`$. Then $`\widehat{}_{v_i}v_j=(1/2)[v_i,v_j]`$ is a Killing vector field along $`\pi ^1(\pi (p))`$. We then obtain
$`v_1g(\widehat{}_{v_1}v_2,v_3)`$ $`=`$ $`g(\widehat{}_{v_1}\widehat{}_{v_1}v_2,v_3)+g(\widehat{}_{v_1}v_2,\widehat{}_{v_1}v_3)`$
$`=`$ $`g(\widehat{}_{v_3}\widehat{}_{v_1}v_2,v_1)+g(\widehat{}_{v_1}v_2,\widehat{}_{v_1}v_3)`$
$`=`$ $`v_3g(\widehat{}_{v_1}v_2,v_1)+g(\widehat{}_{v_1}v_2,\widehat{}_{v_1}v_3+\widehat{}_{v_3}v_1)=0.`$
Analogously, we get $`v_2g(\widehat{}_{v_1}v_2,v_3)=v_2g(\widehat{}_{v_2}v_1,v_3)=0`$. We also obtain
$$v_3g(\widehat{}_{v_1}v_2,v_3)=g(\widehat{}_{v_3}\widehat{}_{v_1}v_2,v_3)+g(\widehat{}_{v_1}v_2,\widehat{}_{v_3}v_3)=0,$$
since $`\widehat{}_{v_3}v_3=0`$ and $`\widehat{}_{v_1}v_2`$ is a Killing vector field along $`\pi ^1(\pi (p))`$. It is easy to see that
$`g(\widehat{}_{v_1}v_2,v_3)`$ $`=`$ $`g(\widehat{}_{v_2}v_1,v_3)=g(\widehat{}_{v_2}v_3,v_1)`$
$`=`$ $`g(\widehat{}_{v_3}v_2,v_1)=g(\widehat{}_{v_3}v_1,v_2)=g(\widehat{}_{v_1}v_3,v_2).`$
Thus $`g(\widehat{}_{v_i}v_j,v_l)`$ is constant along the fibre $`\pi ^1(\pi (p))`$ for each $`i,j,l\{1,2,3\}`$. Therefore
$$g(A_XA_{A_Xv_i}A_Xv_j,A_Xv_l)=g(X,X)g(A_{A_Xv_i}A_Xv_j,v_l)=g(X,X)^2g(\widehat{}_{v_i}v_j,v_l)$$
is constant along $`\pi ^1(\pi (p))`$. Also, we compute for $`\alpha 1`$
$$g(A_XA_{A_Xv_i}A_Xv_j,A_{L_\alpha }v_l)=g(A_{A_Xv_i}A_Xv_j,A_XA_{L_\alpha }v_l)=0,$$
$$g(A_XA_{A_Xv_i}A_Xv_j,L_\alpha )=g(A_{A_Xv_i}A_Xv_j,A_XL_\alpha )=0.$$
Hence $`A_XA_{A_Xv_i}A_Xv_j=g(X,X)A_X\widehat{}_{v_i}v_j`$ is a basic vector field for each $`i,j\{1,\mathrm{},r\}`$.
We choose $`v_{3p}=(\widehat{}_{v_1}v_2)(p)`$. Since $`A_X\widehat{}_{v_1}v_2`$ is a basic vector field along $`\pi ^1(\pi (p))`$, we get the horizontal lifting along $`\pi ^1(\pi (p))`$ of $`\pi _{}(A_X\widehat{}_{v_1}v_2(p))=\pi _{}A_Xv_{3p}`$ is $`g(X,X)^1A_X\widehat{}_{v_1}v_2`$. On the other hand, $`Y_3`$ is, by definition, the horizontal lifting of $`g(X,X)^1\pi _{}A_Xv_{3p}`$ along $`\pi ^1(\pi (p))`$. It follows that $`Y_3=g(X,X)^1A_X\widehat{}_{v_1}v_2`$ along $`\pi ^1(\pi (p))`$. Thus
$$v_3=A_XY_3=g(X,X)^1A_XA_X\widehat{}_{v_1}v_2=\widehat{}_{v_1}v_2$$
along the fibre $`\pi ^1(\pi (p))`$. ∎
For $`r=3`$, we choose $`v_{3p}^{}=(\widehat{}_{v_1^{}}v_2^{})(p)`$. If we repeat the argument above for the basis $`\{v_1^{},v_2^{},v_3^{}\}`$, by Lemma 3.6, we get $`v_3^{}=\widehat{}_{v_1^{}}v_2^{}`$ along the fibre $`\pi ^1(\pi (p))`$. It follows that $`g(\widehat{}_{v_i}v_j,v_l)=g(\widehat{}_{v_i^{}}v_j^{},v_l^{})`$ for each $`i,j,l\{1,2,3\}`$.
Returning to the computation of $`g(A_{A_{L_\alpha }v_i}A_{L_\beta }v_j,v_l)`$, in both cases $`r=1`$ and $`r=3`$, we get for every $`\alpha ,\beta 0`$ and $`i,j,k\{1,\mathrm{},r\}`$
$`g(A_{A_{L_\alpha }v_i}A_{L_\beta }v_j,v_l)`$ $`=`$ $`g(L_\alpha ,L_\beta )g(\widehat{}_{v_i}v_j,v_l)`$
$`=`$ $`g(L_\alpha ^{},L_\beta ^{})g(\widehat{}_{v_i^{}}v_j^{},v_l^{})=g(A_{A_{L_\alpha ^{}}v_i^{}}A_{L_\beta ^{}}v_j^{},v_l^{}).`$
Hence $`\varphi (A_{A_{L_\alpha }v_i}A_{L_\beta }v_j)=A_{\varphi (A_{L_\alpha }v_i)}\varphi (A_{L_\beta }v_j)`$ and $`\varphi (A_{A_{L_\alpha }v_i}v_j)=A_{\varphi (A_{L_\alpha }v_i)}\varphi (v_j).`$
By Corollary 2.3.14 in \[Wol\] we see that $`\varphi :T_pMT_pM`$ extends to an isometry on $`M`$, denoted by $`f:MM`$, such that $`f(p)=p`$ and $`f_p=\varphi `$. Hence $`f_pX=Y`$ and $`f_{}(_p)=_p`$. Since $`f_{}A_EF=A_{f_{}E}f_{}F`$ for every $`E,FT_pM`$, we see, by Theorem 2.4, that there is an isometry $`\stackrel{~}{f}:BB`$ such that $`\stackrel{~}{f}\pi =\pi f`$. Thus $`\stackrel{~}{f}_{}X^{}=\stackrel{~}{f}_{}\pi _{}X=\pi _{}f_{}X=\pi _{}Y=Y^{}`$ and $`\stackrel{~}{f}(x)=\stackrel{~}{f}(\pi (p))=\pi (f(p))=\pi (p)=x`$.
Therefore $`B`$ is an isotropic semi-Riemannian manifold. This completes the proof of Theorem 3.3. ∎
If the metric on the base space is negative definite, the following lemma follows from Theorem 3.3.
###### Lemma 3.7.
If $`\pi :MB`$ is a semi-Riemannian submersion with connected totally geodesic fibres from an $`(n+r)`$-dimensional semi-Riemannian manifold $`M`$ of index $`r^{}+n`$ and of constant negative curvature onto an $`n`$-dimensional semi-Riemannian manifold $`B`$ of index $`n`$, then $`r^{}=r`$.
###### Proof.
By Theorem 3.3, we have $`n=q_1(r^{}+1)+q_2(rr^{})=(q_1+q_2)(r+1)`$ for some nonnegative integers $`q_1`$ and $`q_2`$. Hence $`0=q_1(rr^{})+q_2(r^{}+1)`$. Since the right hand side is the sum of two non-negative numbers, it follows that $`q_1(rr^{})=0`$ and $`q_2(r^{}+1)=0`$. Therefore $`q_2=0`$. This implies $`r^{}=r`$. ∎
###### Remark.
Changing simultaneously the signs of metrics on the total space and on the base space, any semi-Riemannian submersion, under the assumptions of Lemma 3.7, becomes a Riemannian submersion with totally geodesic fibres from a sphere onto a Riemannian manifold. This case was completely classified by Escobales (see \[Esc1\]) and Ranjan (see \[Ran1\]).
###### Proposition 3.8.
Let $`\pi :MB`$ be a semi-Riemannian submersion with connected totally geodesic fibres from a complete simply connected semi-Riemannian manifold $`M`$ onto a semi-Riemannian manifold $`B`$. Then $`B`$ is simply connected and complete.
###### Proof.
If $`M`$ is geodesically complete, then so is $`B`$ (see \[Bes\] or \[Ba-Ia\]). Since $`M`$ is a complete semi-Riemannian manifold and the fibres are totally geodesic, any fibre is also geodesically complete. By a theorem in \[Rec\], it follows that the horizontal distribution $``$ is an Ehresmann connection. Therefore, by \[Ehr\], we see that $`\pi `$ is a fibre bundle. So we obtain an exact homotopy sequence:
$$\mathrm{}\pi _2(M)\pi _2(B)\pi _1(fibre)\pi _1(M)\pi _1(B)0.$$
Thus $`\pi _1(B)=0`$. ∎
By Theorem 12.3.2 in \[Wol\], we know that any connected, simply connected isotropic semi-Riemannian manifold is isometric to one of the following semi-Riemannian manifolds:
* $`_t^m`$ or the universal semi-Riemannian covering of the pseudo-hyperbolic space $`H_t^m(c)`$ with constant sectional curvature $`c<0`$, or of the pseudo-sphere $`S_t^m(c)`$ with constant sectional curvature $`c>0`$.
* The complex pseudo-hyperbolic space $`H_t^m(c)`$ with constant holomorphic sectional curvature $`c<0`$, or the complex pseudo-projective space $`P_t^m(c)`$ with constant holomorphic sectional curvature $`c>0`$.
* The quaternionic pseudo-hyperbolic space $`H_t^m(c)`$ with constant quaternionic sectional curvature $`c<0`$, or the quaternionic pseudo-projective space $`P_t^m(c)`$ with constant quaternionic sectional curvature $`c>0`$.
* The Cayley pseudo-hyperbolic plane $`aH_t^2(c)`$ with Cayley sectional curvature $`c<0`$, or the Cayley pseudo-projective plane $`aP_t^2(c)`$ with Cayley sectional curvature $`c>0`$.
###### Lemma 3.9.
(a) If $`B`$ is a semi-Riemannian manifold isometric to one of the semi-Riemannian manifolds $`P_t^m(c)`$, $`P_t^m(c)`$, $`aP_t^2(c)`$ $`(c>0)`$, then the curvature tensor satisfies the inequality
(3.2)
$$R^{}(X^{},Y^{},X^{},Y^{})\frac{c}{4}(g^{}(X^{},X^{})g^{}(Y^{},Y^{})g^{}(X^{},Y^{})^2)$$
for each tangent vectors $`X^{}`$, $`Y^{}`$ of $`B`$.
(b) If $`B`$ is a semi-Riemannian manifold isometric to one of the semi-Riemannian manifolds $`H_t^m(c)`$, $`H_t^m(c)`$, $`aH_t^2(c)`$ $`(c<0)`$, then the curvature tensor satisfies the inequality
(3.3)
$$R^{}(X^{},Y^{},X^{},Y^{})\frac{c}{4}(g^{}(X^{},X^{})g^{}(Y^{},Y^{})g^{}(X^{},Y^{})^2)$$
for each tangent vectors $`X^{}`$, $`Y^{}`$ of $`B`$.
###### Proof.
For each tangent vectors $`X^{}`$, $`Y^{}`$ of $`B`$, we have the following formulas for the curvature tensors:
* If $`B\{P_t^m(c),H_t^m(c)\}`$ and $`I_0`$ is the natural complex structure on $`B`$, then
(3.4)
$$R^{}(X^{},Y^{},X^{},Y^{})=\frac{c}{4}(g^{}(X^{},X^{})g^{}(Y^{},Y^{})g^{}(X^{},Y^{})^2+3g^{}(X^{},I_0Y^{})^2).$$
* If $`B\{P_t^m(c),H_t^m(c)\}`$ and $`I_0,J_0,K_0`$ are local almost complex structures which give rise to the quaternionic structure on $`B`$, then
$`R^{}(X^{},Y^{},X^{},Y^{})`$ $`=`$ $`(c/4)(g^{}(X^{},X^{})g^{}(Y^{},Y^{})g^{}(X^{},Y^{})^2`$
$`+3g^{}(X^{},I_0Y^{})^2+3g^{}(X^{},J_0Y^{})^2+3g^{}(X^{},K_0Y^{})^2).`$
By these explicit formulas for curvature tensors, in all cases we obtain the inequalities (3.2) and (3.3). ∎
First, we shall discuss the case of a base space with nonconstant curvature.
###### Lemma 3.10.
If $`\pi :H_{s+r^{}}^{n+r}B_s^n`$ is a semi-Riemannian submersion with connected totally geodesic fibres from an $`(n+r)`$-dimensional pseudo-hyperbolic space $`H_{s+r^{}}^{n+r}`$ of index $`s+r^{}>1`$ onto an $`n`$-dimensional isotropic semi-Riemannian manifold $`B_s^n`$ of index $`s`$ with nonconstant curvature, then the induced metrics on the fibres are negative definite and $`B`$ is isometric to one of the following semi-Riemannian manifolds$`:`$
* $`H_t^m`$, $`m>1`$,
* $`H_t^m`$, $`m>1`$,
* $`aH_t^2`$.
###### Proof.
Since $`dim=k(dim𝒱+1)`$ for some positive integer $`k`$, we get $`dimdim𝒱+1`$. Let $`X`$ be a horizontal vector field along a fibre $`\pi ^1(\pi (p))`$ such that $`g(X,X)0`$ and $`X`$ is the horizontal lifting of some tangent vector of $`B`$.
First, we shall prove that
$$dim>dim𝒱+1.$$
Suppose that $`dim=dim𝒱+1`$. Then $`A_X:𝒱X^{}=\{Y|g(X,Y)=0\}`$ is bijective. For every $`YX^{}`$ we get $`Y=A_XV`$ for some vertical vector $`V`$. It follows that
$`g(A_XY,A_XY)`$ $`=`$ $`g(A_XA_XV,A_XA_XV)=g(X,X)^2g(V,V),`$
$`g(Y,Y)`$ $`=`$ $`g(A_XV,A_XV)=g(X,X)g(V,V).`$
Thus $`g(A_XY,A_XY)=g(X,X)g(Y,Y)`$ for every $`YX^{}`$. By O’Neill’s equations, we have
$`R^{}(\pi _{}X,\pi _{}Y,\pi _{}X,\pi _{}Y)`$ $`=`$ $`g(X,X)g(Y,Y)+g(X,Y)^2+3g(A_XY,A_XY)`$
$`=`$ $`4(g(X,X)g(Y,Y)g(X,Y)^2)`$
for every horizontal vector field $`Y`$ along $`\pi ^1(\pi (p))`$. Hence $`B`$ has constant curvature is a contradiction.
We established that $`dim>dim𝒱+1.`$ So we can find a horizontal vector field $`Z`$ along the fibre $`\pi ^1(\pi (p))`$ such that $`Z\mathrm{ker}A_X^{}`$, $`g(X,Z)=0`$, $`g(Z,Z)0`$ and $`Z`$ is the horizontal lifting of some $`Z^{}T_{\pi (p)}B`$. We then have
$`R^{}(\pi _{}X,\pi _{}Z,\pi _{}X,\pi _{}Z)`$ $`=`$ $`g(X,X)g(Z,Z)+g(X,Z)^2+3g(A_XZ,A_XZ)`$
$`=`$ $`g(X,X)g(Z,Z).`$
Since $`B`$ is a simply connected isotropic semi-Riemannian manifold with nonconstant curvature, we see that $`B`$ is isometric to one of the following semi-Riemannian manifolds:
* $`P_t^m(c)`$, $`P_t^m(c)`$, $`aP_t^2(c)`$, or
* $`H_t^m(c)`$, $`H_t^m(c)`$, $`aH_t^2(c)`$.
We shall prove that only the case (b) is possible.
First, we suppose that $`B`$ is isometric to one of the following semi-Riemannian manifolds:
$`P_t^m(c)`$, $`P_t^m(c)`$, $`aP_t^2(c)`$ $`(c>0)`$.
By the inequality (3.2), we get
$`R^{}(\pi _{}X,\pi _{}A_XV,\pi _{}X,\pi _{}A_XV)`$ $`=`$ $`4g(X,X)g(A_XV,A_XV)`$
$`=`$ $`4g(X,X)^2g(V,V)(c/4)g(X,X)^2g(V,V).`$
Therefore
(3.6)
$$g(V,V)0$$
for every vertical vector $`V`$. Since $`X`$ and $`Z`$ are basic vector fields along $`\pi ^1(\pi (p))`$ with $`g(X,Z)=0`$ and $`A_XZ=0`$ along $`\pi ^1(\pi (p))`$, it follows from the relation (3.1) that $`A_ZV\mathrm{ker}A_X^{}`$. On the other hand, by the inequality (3.2), we get
$`R^{}(\pi _{}X,\pi _{}Z,\pi _{}X,\pi _{}Z)`$ $`=`$ $`g(X,X)g(Z,Z)(c/4)g(X,X)g(Z,Z),`$
$`R^{}(\pi _{}X,\pi _{}A_ZV,\pi _{}X,\pi _{}A_ZV)`$ $`=`$ $`g(X,X)g(A_ZV,A_ZV)(c/4)g(X,X)g(A_ZV,A_ZV).`$
Hence $`g(X,X)g(Z,Z)0`$ and $`g(X,X)g(A_ZV,A_ZV)0`$. Thus
$$0g(Z,Z)g(A_ZV,A_ZV)=g(Z,Z)^2g(V,V).$$
So for any vertical vector $`V`$ we get
(3.7)
$$g(V,V)0.$$
Since the induced metrics on fibres are nondegenerate, it is not possible to have both (3.6) and (3.7). So we obtain the required contradiction. It follows that $`B`$ is isometric to one of the following semi-Riemannian manifolds:
$`H_t^m(c)`$, $`H_t^m(c)`$, $`aH_t^2(c)`$ $`(c<0)`$.
We shall now prove that $`c=4`$. Suppose $`(c/4)+10`$. By the inequality (3.3), we get
(3.8)
$$R^{}(\pi _{}X,\pi _{}Z,\pi _{}X,\pi _{}Z)=g(X,X)g(Z,Z)(c/4)g(X,X)g(Z,Z),$$
$$R^{}(\pi _{}X,\pi _{}A_ZV,\pi _{}X,\pi _{}A_ZV)=g(X,X)g(A_ZV,A_ZV)(c/4)g(X,X)g(A_ZV,A_ZV).$$
Hence
(3.9)
$$((c/4)+1)^2g(X,X)^2g(Z,Z)g(A_ZV,A_ZV)0,$$
from which follows that $`0g(Z,Z)g(A_ZV,A_ZV)=g(Z,Z)^2g(V,V)`$. Therefore $`g(V,V)0`$ for every vertical vector field $`V`$. In particular, we have $`g(A_XY,A_XY)0`$, which implies
(3.10)
$$R^{}(\pi _{}X,\pi _{}Y,\pi _{}X,\pi _{}Y)g(X,X)g(Y,Y)g(X,Y)^2$$
for every horizontal vectors $`X`$ and $`Y`$. We have the following cases:
Case (a) $`0<\mathrm{index}B<dimB`$. We can choose vector fields $`X^{}`$, $`Y^{}`$ on $`B`$ such that $`g^{}(X^{},X^{})g^{}(Y^{},Y^{})<0`$ and that one of the following conditions is satisfied:
* $`Y^{}\{X^{},I_0X^{}\}^{}`$ if $`B=H_s^m(c)`$, where $`I_0`$ is the natural complex structure on $`H_s^m(c)`$,
* $`Y^{}\{X^{},I_0X^{},J_0X^{},K_0X^{}\}^{}`$ if $`B=H_s^m(c)`$, where $`\{I_0,J_0,K_0\}`$ are local almost complex structures which give rise to the quaternionic structure on $`H_s^m(c)`$, or
Let $`X`$, $`Y`$ be the horizontal liftings of $`X^{}`$, $`Y^{}`$. The inequality (3.10) then implies
$$\frac{c}{4}g(X,X)g(Y,Y)g(X,X)g(Y,Y).$$
Hence $`((c/4)+1)g(X,X)g(Y,Y)0`$. Therefore $`(c/4)+1>0`$. On the other hand, we can choose horizontal vector fields $`X`$, $`Z`$ such that $`g(X,Z)=0`$, $`Z\mathrm{ker}A_X^{}`$ and $`g(X,X)g(Z,Z)<0`$, because $`0<\mathrm{index}B<dimB`$. Then the inequality (3.8) becomes $`(c/4)+1<0`$. So we get a contradiction.
Case (b) $`\mathrm{index}B\{0,dimB\}`$. Similarly, we can choose vector fields $`X^{}`$, $`Y^{}`$ on $`B`$ such that $`g^{}(X^{},Y^{})=0`$ and $`R^{}(X^{},Y^{},X^{},Y^{})=(c/4)g^{}(X^{},X^{})g^{}(Y^{},Y^{})`$. The inequality (3.10) then implies $`((c/4)+1)g^{}(X^{},X^{})g^{}(Y^{},Y^{})0`$. By the hypothesis of Case (b), we get $`(c/4)+10`$. On the other hand, the inequality (3.8) becomes $`(c/4)+1>0`$. So we get a contradiction.
We have proved $`c=4`$. The inequality (3.3) then becomes
(3.11)
$$R^{}(X^{},Y^{},X^{},Y^{})g^{}(X^{},X^{})g^{}(Y^{},Y^{})+g^{}(X^{},Y^{})^2$$
for tangent vector fields $`X^{}`$, $`Y^{}`$ on $`B`$. Then we have
$$R^{}(\pi _{}X,\pi _{}A_XV,\pi _{}X,\pi _{}A_XV)=4g(X,X)g(A_XV,A_XV)g(X,X)g(A_XV,A_XV)$$
for a vertical vector field $`V`$ and for a horizontal vector field $`X`$. Hence
$$0g(X,X)g(A_XV,A_XV)=g(X,X)^2g(V,V).$$
Therefore the induced metrics on fibres are negative definite. ∎
By Lemma 3.10, we deduce the following proposition.
###### Proposition 3.11.
If $`\pi :H_{s+r^{}}^{n+r}B_s^n`$ is a semi-Riemannian submersion with connected totally geodesic fibres from an $`(n+r)`$-dimensional pseudo-hyperbolic space $`H_{s+r^{}}^{n+r}`$ of index $`s+r^{}`$ onto an $`n`$-dimensional isotropic semi-Riemannian manifold $`B_s^n`$ of index $`s`$ with nonconstant curvature, and if the fibres are negatively definite then one of the following holds$`:`$
* $`n=2m>2`$, $`s=2t`$, $`r=r^{}=1`$ for some non-negative integers $`m`$, $`t`$, and $`B_s^n`$ is isometric to $`H_t^m`$.
* $`n=4m>4`$, $`s=4t`$, $`r=r^{}=3`$ for some non-negative integers $`m`$, $`t`$, and $`B_s^n`$ is isometric to $`H_t^m`$.
* $`n=16`$, $`s\{0,8,16\}`$, $`r=r^{}=7`$, and $`B_s^n`$ is isometric to $`aH_{s/8}^2`$.
###### Proof.
First, we shall discuss the case $`s+r^{}>1`$. By Lemma 3.10, $`B`$ is isometric to one of the semi-Riemannian manifolds $`H_t^m`$, $`H_t^m`$, $`aH_t^2`$ for some $`m>1`$.
Let $`xB`$ and let $`X^{}T_xB`$ such that $`g^{}(X^{},X^{})0`$, and let $`_X^{}`$ be the subspace in $`T_xB`$ given by
$$_X^{}=\{Y^{}T_xB|R^{}(X^{},Y^{})X^{}=g^{}(X^{},Y^{})X^{}+g^{}(X^{},X^{})Y^{}\}.$$
Let $`p\pi ^1(x)`$ and let $`X`$ be the horizontal lifting vector at $`p`$ of $`X^{}`$. By O’Neill’s equations, we have $`R^{}(\pi _{}X,\pi _{}Y,\pi _{}X,\pi _{}Z)=R(X,Y,X,Z)+3g(A_X^{}Y,A_X^{}Z)`$ for horizontal vectors $`Y`$, $`Z`$. Since $`A_X^{}:_p𝒱_p`$ is surjective and since the induced metrics on fibres are nondegenerate, we get $`Y\mathrm{ker}A_X^{}`$ if and only if $`\pi _{}Y_X^{}`$. Thus
$$dim\mathrm{ker}A_X^{}=dimdim𝒱=dim_X^{}.$$
We have the following possibilities:
* $`B_s^n`$ is isometric to $`H_t^m`$. So $`n=2m`$, $`s=2t`$. From the geometry of the complex pseudo-hyperbolic space (see relation (3.4)), we get $`dim_X^{}=dim1`$. It follows that $`r=r^{}=dim𝒱=1`$.
* $`B_s^n`$ is isometric to $`H_t^m`$. So $`n=4m`$, $`s=4t`$. From the geometry of the quaternionic pseudo-hyperbolic space (see relation ((ii))), we get $`dim_X^{}=dim3`$. It follows that $`r=r^{}=dim𝒱=3`$.
* $`B_s^n`$ is isometric to the Cayley pseudo-hyperbolic plane $`aH_t^2`$. So $`n=16`$, $`s\{0,8,16\}`$. From the geometry of the Cayley pseudo-hyperbolic plane, we obtain $`dim_X^{}=dim7`$. Hence $`r=r^{}=dim𝒱=7`$.
Now, we discuss the remaining case $`s+r^{}=1`$. From $`s+r^{}=1`$, we have either
* $`s=0`$, $`r^{}=1`$, or
* $`s=1`$, $`r^{}=0`$.
If $`s=0`$, $`r^{}=1`$, then $`\pi :H_1^{n+r}B^n`$ is a semi-Riemannian submersion with totally geodesic fibres from an anti-de Sitter space onto a Riemannian manifold. In this case, investigated by Magid in \[Mag\], it follows that $`B`$ is isometric to the complex hyperbolic space $`H^m`$ and $`r=r^{}=1`$.
For $`s=1`$, $`r^{}=0`$, we get, by Theorem 3.3, $`1=q_1+q_2rq_1+q_2`$ with $`q_1+q_2=k=n/(r+1)`$. Thus $`q_1+q_2=1`$. It follows that $`n=r+1`$. Hence $`A_X:𝒱X^{}`$ is bijective. Since $`R^{}(\pi _{}X,\pi _{}A_XV,\pi _{}X,\pi _{}A_XV)=4g(X,X)g(A_XV,A_XV)`$, we see that $`B`$ has constant curvature $`4`$, which contradicts our assumption of nonconstant curvature of the base space. ∎
We shall now discuss the case where the base space is of constant curvature.
###### Proposition 3.12.
If $`\pi :H_{s+r^{}}^{n+r}B_s^n`$ is a semi-Riemannian submersion with connected totally geodesic fibres from an $`(n+r)`$-dimensional pseudo-hyperbolic space of index $`s+r^{}`$ onto an $`n`$-dimensional semi-Riemannian manifold of index $`s`$ with constant curvature, and if the fibres are negatively definite, then one of the following holds$`:`$
* $`n=s=2^t`$, $`r=r^{}=n1`$, $`B`$ is isometric to $`H_{2^t}^{2^t}(4)`$ and $`t\{1,2,3\}`$.
* $`n=2^t`$, $`s=0`$, $`r=r^{}=n1`$, $`B`$ is isometric to $`H^{2^t}(4)`$ and $`t\{1,2,3\}`$.
###### Proof.
Since $`B`$ has constant curvature, the curvature of $`B`$ is $`4`$ and $`n=r+1`$. By Theorem 3.3, $`s=q_1(r^{}+1)+q_2(rr^{})=q_1(r+1)`$ and $`q_1+q_2=n/(r+1)=1`$. Then either $`q_1=0`$ or $`q_1=1`$. If $`q_1=0`$, then $`s=0`$. If $`q_1=1`$ then $`s=r+1=n`$. Summarizing, we have $`\mathrm{index}(B)\{0,dimB\}`$.
If $`\mathrm{index}(B)=dimB`$, then, by Lemma 3.7, we obtain $`r=r^{}`$. Hence, by \[Ran1\], we have (1).
If $`\mathrm{index}(B)=0`$, then, by \[Ba-Ia\], we have (2).
The idea of the proof in \[Ran1\] and \[Ba-Ia\] is to see that the tangent bundle of any fibre is trivial and that fibres are diffeomorphic to spheres, and then to apply a well-known result of Adams which claims that the spheres of dimensions 1, 3 and 7 are the only spheres with trivial tangent bundle. ∎
The next theorems solve the equivalence problem of semi-Riemannian submersions from real and complex pseudo-hyperbolic spaces.
###### Theorem 3.13.
If $`\pi _1,\pi _2:H_{s+r^{}}^{n+r}B_s^n`$ are two semi-Riemannian submersions with connected totally geodesic fibres from a pseudo-hyperbolic space of index $`s+r^{}>1`$, if the fibres are negative definite, and if the dimension of the fibres is $`r\{1,3\}`$, then $`\pi _1`$ and $`\pi _2`$ are equivalent.
###### Proof.
Let $`p,qH_{s+r^{}}^{n+r}`$. Let
$$=\{L_0,A_{1L_0}v_1,\mathrm{},A_{1L_0}v_r,\mathrm{},L_{k1},A_{1L_{k1}}v_1,\mathrm{},A_{1L_{k1}}v_r\},$$
$$^{}=\{L_0^{},A_{2L_0^{}}v_1^{},\mathrm{},A_{2L_0^{}}v_r^{},\mathrm{},L_{k1}^{},A_{2L_{k1}^{}}v_1^{},\mathrm{},A_{2L_{k1}^{}}v_r^{}\}$$
be two orthonormal bases of $`_1`$ along $`\pi _1^1(\pi _1(p))`$ and of $`_2`$ along $`\pi _2^1(\pi _2(q))`$ constructed as in the proof of Theorem 3.3 such that $`g_p(L_\alpha ,L_\beta )=g_q(L_\alpha ^{},L_\beta ^{})=\epsilon _\alpha \delta _{\alpha \beta }`$ for $`\alpha ,\beta \{0,\mathrm{},k1\}`$, $`g_p(v_i,v_j)=g_q(v_i^{},v_j^{})=\epsilon _i\delta _{ij}`$ for $`i,j\{1,\mathrm{},r\}`$ and for $`r=3`$, $`v_{3p}=(\widehat{}_{v_1}v_2)(p)`$ and $`v_{3q}^{}=(\widehat{}_{v_1^{}}v_2^{})(q)`$.
Let $`\varphi :T_pH_{s+r^{}}^{n+r}T_qH_{s+r^{}}^{n+r}`$ be the linear map given by $`\varphi (L_\alpha )=L_\alpha ^{}`$, $`\varphi (A_{1L_\alpha }v_i)=A_{2L_\alpha ^{}}v_i^{}`$, $`\varphi (v_i)=v_i^{}`$ for every $`\alpha `$ and $`i`$. In a manner similar to the proof of Theorem 3.3, we obtain $`\varphi (A_{1E}F)=A_{2\varphi (E)}\varphi (F)`$ for every $`E`$, $`FT_pH_{s+r^{}}^{n+r}`$. By Corollary 2.3.14 in \[Wol\], $`\varphi `$ extends to an isometry on $`H_{s+r^{}}^{n+r}`$, denoted by $`f:H_{s+r^{}}^{n+r}H_{s+r^{}}^{n+r}`$, satisfying $`f(p)=q`$ and $`f_p=\varphi `$. From Theorem 2.4 it follows that $`f`$ induces an isometry $`\stackrel{~}{f}`$ on $`B`$, such that $`\stackrel{~}{f}\pi =\pi f`$. Hence $`\pi _1`$ and $`\pi _2`$ are equivalent. ∎
###### Theorem 3.14.
If $`\pi _1,\pi _2:H_{2s+1}^{2n+1}H_s^n`$ are two semi-Riemannian submersions with connected complex totally geodesic fibres from a complex pseudo-hyperbolic space, and if the fibres are negative definite, then $`\pi _1`$ and $`\pi _2`$ are equivalent.
###### Proof.
Let $`\theta :H_{4s+3}^{4n+3}H_{2s+1}^{2n+1}`$ be the canonical semi-Riemannian submersion. By Theorem 2.5 in \[Esc2\], we see that $`\stackrel{~}{\pi }_1=\pi _1\theta :H_{4s+3}^{4n+3}H_s^n`$ and $`\stackrel{~}{\pi }_2=\pi _2\theta :H_{4s+3}^{4n+3}H_s^n`$ are semi-Riemannian submersions with totally geodesic fibres. We denote by $`\stackrel{~}{A}_1`$, $`\stackrel{~}{A}_2`$, $`A_1`$, $`A_2`$, $`A`$ O’Neill’s integrability tensors of $`\stackrel{~}{\pi }_1`$, $`\stackrel{~}{\pi }_2`$, $`\pi _1`$, $`\pi _2`$, $`\theta `$, respectively. In order to reduce the proof of the equivalence theorem of semi-Riemannian submersions from a complex pseudo-hyperbolic space to that from a pseudo-hyperbolic space, we need to establish relations among the integrability tensors $`\stackrel{~}{A}_1`$, $`A_1`$, $`A`$.
First, we prove that $`\theta _{}\stackrel{~}{A}_{1X}Y=A_{1\theta _{}X}\theta _{}Y`$ for $`\stackrel{~}{\pi }_1`$-basic vector fields $`X`$ and $`Y`$. Let $`pH_{4s+3}^{4n+3}`$. Let $`w_1^{}`$, $`w_2^{}`$ be two orthonormal $`\pi _1`$-vertical vectors in $`T_{\theta (p)}H_{2s+1}^{2n+1}`$ and let $`w_1`$, $`w_2`$ be the $`\theta `$-horizontal liftings at $`p`$ of $`w_1^{}`$, $`w_2^{}`$, respectively. Let $`w_3`$ be a unit $`\theta `$-vertical vector in $`T_pH_{4s+3}^{4n+3}`$. Then $`\{w_1,w_2,w_3\}`$ gives an orthonormal basis of $`\stackrel{~}{𝒱}_{1p}`$. Since the induced metrics on the fibres of $`\stackrel{~}{\pi }_1`$ are negative definite, we have
$$\stackrel{~}{A}_{1X}Y=g(_XY,w_1)w_1g(_XY,w_2)w_2g(_XY,w_3)w_3.$$
Thus
$$\theta _{}\stackrel{~}{A}_{1X}Y=g^{}(_{\theta _{}X}^{}\theta _{}Y,\theta _{}w_1)\theta _{}w_1g^{}(_{\theta _{}X}^{}\theta _{}Y,\theta _{}w_2)\theta _{}w_2=A_{1\theta _{}X}\theta _{}Y$$
for $`\stackrel{~}{\pi }_1`$-basic vector fields $`X`$ and $`Y`$, where $`g^{}`$ denotes the metric on $`H_{2s+1}^{2n+1}`$ and $`^{}`$ is the Levi-Civita connection of $`g^{}`$.
Let $`X`$ be the $`\stackrel{~}{\pi }_1`$-horizontal lifting along the fibre $`\stackrel{~}{\pi }_1^1(\stackrel{~}{\pi }_1(p))`$ of some unit vector in $`T_{\stackrel{~}{\pi }_1(p)}H_s^n`$. Let $`Y_1`$, $`Y_2`$, $`Y_3`$ be the $`\stackrel{~}{\pi }_1`$-horizontal liftings along the fibre $`\stackrel{~}{\pi }_1^1(\stackrel{~}{\pi }_1(p))`$ of $`\stackrel{~}{\pi }_1\stackrel{~}{A}_{1X}w_1`$, $`\stackrel{~}{\pi }_1\stackrel{~}{A}_{1X}w_2`$, $`\stackrel{~}{\pi }_1\stackrel{~}{A}_{1X}w_3`$, respectively. Let $`v_i=\stackrel{~}{A}_{1X}Y_i`$ for $`i\{1,2,3\}`$. As in Theorem 3.3, we choose $`w_3=g(X,X)^1\left(_{v_1}v_2\right)(p)`$, which implies that $`v_3=_{v_1}v_2`$ (see Lemma 3.6).
We remark that $`v_3=\stackrel{~}{A}_{1X}Y_3`$ is a $`\theta `$-vertical vector field along the fibre $`\theta ^1(\theta (p))`$. Indeed, we have
$`\theta _{}\left(\stackrel{~}{A}_{1X}Y_3(p^{})\right)`$ $`=`$ $`\left(A_{1\theta _{}X}\theta _{}Y_3\right)(\theta (p^{}))=\left(A_{1\theta _{}X}\theta _{}Y_3\right)(\theta (p))=\theta _{}\left(\stackrel{~}{A}_{1X}Y_3(p)\right)`$
$`=`$ $`\theta _{}(A_{1X}A_{1X}w_3)=g(X,X)\theta _{}w_3=0`$
for any $`p^{}\theta ^1(\theta (p))`$.
Since $`v_1`$, $`v_2`$ are orthogonal to the vertical vector field $`v_3`$ along $`\theta ^1(\theta (p))`$, we see that $`v_1`$, $`v_2`$ are $`\theta `$-horizontal. Since $`\theta _{}\left(\stackrel{~}{A}_{1X}Y_i(p^{})\right)=\left(A_{1\theta _{}X}\theta _{}Y_i\right)(\theta (p^{}))`$ for $`p^{}\theta ^1(\theta (p))`$ and for $`i\{1,2\}`$, we obtain that $`v_1`$, $`v_2`$ are $`\theta `$-basic vector fields along $`\theta ^1(\theta (p))`$. Thus $`h_{v_3}v_1=A_{v_1}v_3`$ along $`\theta ^1(\theta (p))`$. Here $`h`$ and $`v`$ denote the $`\theta `$-horizontal and $`\theta `$-vertical projections, respectively. We also obtain that $`v_{v_3}v_1=g(_{v_3}v_1,v_3)v_3=0`$. Therefore, $`A_{v_1}v_3=_{v_3}v_1=v_2`$ along $`\theta ^1(\theta (p))`$.
We shall prove that $`\stackrel{~}{A}_{1X}v_3=A_Xv_3`$ along $`\theta ^1(\theta (p))`$ for every $`\stackrel{~}{\pi }_1`$-basic vector field $`X`$ along $`\stackrel{~}{\pi }_1^1(\stackrel{~}{\pi }_1(p))`$. We first obtain along $`\theta ^1(\theta (p))`$ that
$$\stackrel{~}{A}_{1X}v_3=_Xv_3+g(_Xv_3,v_1)v_1+g(_Xv_3,v_2)v_2+g(_Xv_3,v_3)v_3,$$
$`g(_Xv_3,v_1)`$ $`=`$ $`g(A_Xv_3,v_1)=g(v_3,A_Xv_1)=g(v_3,A_{v_1}X)=g(A_{v_1}v_3,X)`$
$`=`$ $`g(v_2,X)=0`$
for a $`\stackrel{~}{\pi }_1`$-basic vector field $`X`$ along $`\stackrel{~}{\pi }_1^1(\stackrel{~}{\pi }_1(p))`$. Analogously, we get $`g(_Xv_3,v_2)=0`$. Thus
$$\stackrel{~}{A}_{1X}v_3=_Xv_3+g(_Xv_3,v_3)v_3=A_Xv_3$$
along $`\theta ^1(\theta (p))`$ for every $`\stackrel{~}{\pi }_1`$-basic vector field $`X`$ along $`\stackrel{~}{\pi }_1^1(\stackrel{~}{\pi }_1(p))`$.
Let $`\stackrel{~}{}=\{L_0=X,\stackrel{~}{A}_{1L_0}v_1,\stackrel{~}{A}_{1L_0}v_2,\stackrel{~}{A}_{1L_0}v_3,\mathrm{},L_{n1},\stackrel{~}{A}_{1L_{n1}}v_1,\stackrel{~}{A}_{1L_{n1}}v_2,\stackrel{~}{A}_{1L_{n1}}v_3\}`$ be an orthonormal basis of $`\stackrel{~}{}_1`$ along the fibre $`\stackrel{~}{\pi }_1^1(\stackrel{~}{\pi }_1(p))`$ constructed as in Theorem 3.3, for the semi-Riemannian submersion $`\stackrel{~}{\pi }_1`$. From the proof of Theorem 3.3, we have
$$g(\stackrel{~}{A}_{1\stackrel{~}{A}_{1L_j}v_1}v_3,\stackrel{~}{A}_{1L_l}v_2)=0$$
for $`jl`$, and
$$g(\stackrel{~}{A}_{1\stackrel{~}{A}_{1L_j}v_1}v_3,L_t)=0$$
for $`0j,tn1`$. We then obtain along $`\stackrel{~}{\pi }_1^1(\stackrel{~}{\pi }_1(p))`$ that
$`g(\stackrel{~}{A}_{1\stackrel{~}{A}_{1L_j}v_1}v_3,\stackrel{~}{A}_{1L_j}v_2)`$ $`=`$ $`g(v_3,\stackrel{~}{A}_{1\stackrel{~}{A}_{1L_j}v_1}\stackrel{~}{A}_{1L_j}v_2)`$
$`=`$ $`g(v_3,_{v_1}v_2)g(L_j,L_j)`$
$`=`$ $`g(v_3,v_3)g(L_j,L_j)=g(v_2,v_2)g(L_j,L_j)`$
$`=`$ $`g(\stackrel{~}{A}_{1L_j}v_2,\stackrel{~}{A}_{1L_j}v_2),`$
from which follows $`\stackrel{~}{A}_{1L_j}v_2=\stackrel{~}{A}_{1\stackrel{~}{A}_{1L_j}v_1}v_3`$. Hence $`\stackrel{~}{A}_{1L_j}v_2=A_{\stackrel{~}{A}_{1L_j}v_1}v_3`$, because $`\stackrel{~}{A}_{1L_j}v_1`$ is $`\stackrel{~}{\pi }_1`$-basic. We also have $`\stackrel{~}{A}_{1L_j}v_3=A_{L_j}v_3`$.
Let $`=\stackrel{~}{}\{v_1,v_2\}`$. Summarizing all the above, we obtain that
$$=\{L_0,A_{L_0}v_3,\stackrel{~}{A}_{1L_0}v_1,A_{\stackrel{~}{A}_{1L_0}v_1}v_3,\mathrm{},L_{n1},A_{L_{n1}}v_3,\stackrel{~}{A}_{L_{n1}}v_1,A_{\stackrel{~}{A}_{1L_{n1}}v_1}v_3,v_1,A_{v_1}v_3\}$$
is an orthonormal basis of the $`\theta `$-horizontal space $``$ along the fibre $`\theta ^1(\theta (p))`$ and $``$ satisfies all conditions imposed in the construction of the basis $``$ in the proof of Theorem 3.3. We notice that $`v_3=A_XY_3`$ along $`\theta ^1(\theta (p))`$, and that along $`\theta ^1(\theta (p))`$, $`Y_3`$ is equal to the $`\theta `$-horizontal lifting of $`\theta _{}A_Xw_3`$.
Let $`qH_{4s+3}^{4n+3}`$. Let
$$\stackrel{~}{}^{}=\{L_0^{},\stackrel{~}{A}_{2L_0^{}}v_1^{},\stackrel{~}{A}_{2L_0^{}}v_2^{},\stackrel{~}{A}_{2L_0^{}}v_3^{},\mathrm{},L_{n1}^{},\stackrel{~}{A}_{2L_{n1}^{}}v_1^{},\stackrel{~}{A}_{2L_{n1}^{}}v_2^{},\stackrel{~}{A}_{2L_{n1}^{}}v_3^{}\}$$
be an orthonormal basis of $`\stackrel{~}{}_2`$ along $`\stackrel{~}{\pi }_2^1(\stackrel{~}{\pi }_2(q))`$ constructed in the same way as $`\stackrel{~}{}`$, but for the semi-Riemannian submersion $`\stackrel{~}{\pi }_2`$ (see the proof of Theorem 3.3), in such a way that $`g_p(L_\alpha ,L_\beta )=g_q(L_\alpha ^{},L_\beta ^{})`$ for $`0\alpha ,\beta n1`$, $`g_p(v_i,v_j)=g_q(v_i^{},v_j^{})`$ for $`1i,j3`$, and $`v_3^{}(q)=\left(_{v_1^{}}v_2^{}\right)(q)`$. Let $`\varphi :T_pH_{4s+3}^{4n+3}T_qH_{4s+3}^{4n+3}`$ be the linear map given by $`\varphi (v_i)=v_i^{}`$, $`\varphi (\stackrel{~}{A}_{1L_\alpha }v_i)=\stackrel{~}{A}_{2L_\alpha ^{}}v_i^{}`$ for $`0\alpha n1`$ and for $`1i3`$.
By Corollary 2.3.14 in \[Wol\], $`\varphi `$ extends to an isometry $`f:H_{4s+3}^{4n+3}H_{4s+3}^{4n+3}`$ such that $`f(p)=q`$ and $`f_p=\varphi `$. By the proof of Theorem 3.3, we have $`f_{}\stackrel{~}{A}_{1E}F=\stackrel{~}{A}_{2f_{}E}f_{}F`$ for every $`E`$, $`FT_pH_{4s+3}^{4n+3}`$. By the proof of Theorem 3.13 and by Theorem 2.4, $`f`$ induces an isometry on $`H_{2s+1}^{2n+1}`$, denoted by $`\stackrel{~}{f}:H_{2s+1}^{2n+1}H_{2s+1}^{2n+1}`$, such that $`\theta f=\stackrel{~}{f}\theta `$. Since the $`\pi _1`$-vertical space at $`\theta (p)`$ is spanned by $`\{\theta _{}v_1,\theta _{}v_2\}`$, since the $`\pi _2`$-vertical space at $`\theta (q)`$ is spanned by $`\{\theta _{}v_1^{},\theta _{}v_2^{}\}`$, and since $`\stackrel{~}{f}_{}(\theta _{}v_i)=\theta _{}v_i^{}`$, for $`i\{1,2\}`$, we see that $`\stackrel{~}{f}_{}`$ maps the $`\pi _1`$-vertical space at $`\theta (p)`$ into the $`\pi _2`$-vertical space at $`\theta (q)`$. For $`\stackrel{~}{\pi }_1`$-horizontal vectors $`X`$ and $`Y`$ we obtain
$`\stackrel{~}{f}_{}A_{1\theta _{}X}\theta _{}Y`$ $`=`$ $`\stackrel{~}{f}_{}\theta _{}\stackrel{~}{A}_{1X}Y=\theta _{}f_{}\stackrel{~}{A}_{1X}Y`$
$`=`$ $`\theta _{}\stackrel{~}{A}_{2f_{}X}f_{}Y=A_{2\theta _{}f_{}X}\theta _{}f_{}Y`$
$`=`$ $`A_{2\stackrel{~}{f}_{}(\theta _{}X)}\stackrel{~}{f}_{}(\theta _{}Y).`$
Therefore, by Theorem 2.4, we see that $`\pi _1`$ and $`\pi _2`$ are equivalent. ∎
###### Remark.
We notice that our equivalence theorems can be applied, in particular, to Riemannian submersions from a sphere with totally geodesic fibres of dimension less than or equal to $`3`$, and for Riemannian submersions with complex totally geodesic fibres from a complex projective space. Unlike those in \[Esc1\], \[Esc2\], \[Ran1\], our proofs of the equivalence theorems are intrinsic, we do not need to assume the existence of any specific structure on the base space, such as complex or quaternionic one. In Theorem 3.14, we need to assume only that the fibres are 2-dimensional and that the induced metrics on fibres are negative definite.
Summarizing all results above, we now prove the main theorems.
###### Proof of Theorem 1.1.
If $`s+r^{}>1`$, then $`H_{s+r^{}}^{n+r}`$ is simply connected and hence, by Theorem 3.3, $`B`$ is an isotropic semi-Riemannian manifold and $`r\{1,3\}`$. By Propositions 3.11 and 3.12, we see that the base space of the semi-Riemannian submersion is isometric to a complex pseudo-hyperbolic space if the dimension of fibres is one, or to a quaternionic pseudo-hyperbolic space if the dimension of fibres is $`3`$. In Theorem 3.13 we solved the equivalence problem. The existence problem is solved by the explicit construction given in the preliminaries (see Examples 1 and 2).
If $`s+r^{}=1`$, then either (i) $`s=1`$, $`r^{}=0`$, or (ii) $`s=0`$, $`r^{}=1`$. Since the fibres are assumed to be negative definite, (i) cannot occur.
(ii) If $`s=0`$, $`r^{}=1`$, then $`\pi `$ is a semi-Riemannian submersion from an anti-de Sitter space onto a Riemannian manifold. By \[Mag\], $`\pi `$ is equivalent to the canonical submersion $`\pi :H_1^{2m+1}H^m`$. This falls in the case (a). ∎
###### Proof of Theorem 1.2.
If the dimension of the fibres is less than or equal to $`3`$, then, by Theorem 1.1, $`\pi `$ is equivalent to the canonical semi-Riemannian submersions:
* $`H_{2t+1}^{2m+1}H_t^m`$, $`0tm,`$ or
* $`H_{4t+3}^{4m+3}H_t^m`$, $`0tm.`$
Now we assume that the dimension of the fibres is greater than or equal to 4.
(A) If we assume that the dimension of the fibres is greater than or equal to $`4`$ and $`B`$ is an isotropic semi-Riemannian manifold with non-constant curvature, then, by Proposition 3.11, $`B`$ is isometric to $`aH_t^2`$, $`t\{0,1,2\}`$, and the dimension of the fibres is $`r=r^{}=7`$. By Proposition 2.7, there are no such semi-Riemannian submersions with base space $`aH_t^2`$. Therefore, the assumptions (A) and $`r4`$ imply that $`B`$ has constant curvature, and hence, by Proposition 3.12, we obtain $`s=\mathrm{index}(B)\{0,dim(B)\}`$.
(B) If $`\mathrm{index}(B)=0`$ and $`r4`$, then, by \[Ba-Ia\], the semi-Riemannian submersion $`\pi `$ is equivalent to the canonical semi-Riemannian submersion $`H_7^{15}H^8(4)`$. If $`\mathrm{index}(B)=dim(B)`$, then, by Lemma 3.7, we get $`r^{}=r`$. By changing the signs of the metrics on the base and on the total space, $`\pi `$ becomes a Riemannian submersion with connected totally geodesic fibres from a sphere onto a Riemannian manifold. So, by \[Esc1\] and \[Ran1\], one obtains the conclusion. ∎
###### Proof of Theorem 1.3.
Let $`\theta :H_{2s+1}^{2n+1}H_s^n`$ be the canonical semi-Riemannian submersion. By Theorem 2.5 in \[Esc2\], one obtains that $`\pi \theta :H_{2s+1}^{2n+1}B`$ is a semi-Riemannian submersion with connected totally geodesic fibres.
(A) If the dimension of the fibres of $`\pi `$ is $`r`$ and $`1r2`$, then the dimension of the fibres of the semi-Riemannian submersion $`\pi \theta `$ is less than or equal to $`3`$ and greater than or equal to $`2`$. By Theorem 1.1, $`B`$ is isometric to $`H_t^m`$ and $`2n+1=4m+3`$, $`2s+1=4t+3`$. Then $`n=2m+1`$, $`s=2t+1`$. By Theorem 3.14, we see that $`\pi :H_{2t+1}^{2m+1}H_t^m`$ is equivalent to the canonical semi-Riemannian submersion.
(B) and (C) If $`B`$ is an isotropic semi-Riemannian manifold or if $`\mathrm{index}(B)\{0,dimB\}`$, then, by Theorem 1.2, $`\pi \theta `$ is equivalent to one of the following canonical semi-Riemannian submersions:
* $`H_{2t+1}^{2m+1}H_t^m`$, $`0tm`$;
* $`H_{4t+3}^{4m+3}H_t^m`$, $`0tm`$;
* $`H_{7+8t}^{15}H_{8t}^8(4)`$, $`t\{0,1\}`$.
If the dimension of the fibres of $`\pi `$ is greater than or equal to $`3`$, then the dimension of the fibres of $`\pi \theta `$ is greater than or equal to $`4`$. Hence, in this case, $`\pi \theta `$ is equivalent to $`H_{7+8t}^{15}H_{8t}^8(4)`$, $`t\{0,1\}`$. For $`t=1`$, the semi-Riemannian submersion $`\pi `$ is, after a change of signs of the metrics on the total space and on the base space, of type $`\pi :P^7S^8(4)`$. For $`t=0`$, $`\pi `$ is of type $`\pi :H_3^7H^8(4)`$. In \[Ran1\] (for case t=1) and \[Ba-Ia\] (for case t=0), it is proved that there are no such semi-Riemannian submersions with totally geodesic fibres. We proved that the dimension of fibres of $`\pi `$ is less than or equal to $`2`$. ∎
###### Proof of Theorem 1.4.
We suppose that there are such semi-Riemannian submersions. It is well-known that any quaternionic submanifold in $`H_s^n`$ is totally geodesic. Let $`\eta :H_{4s+3}^{4n+3}H_s^n`$, $`\xi :H_{2s+1}^{2n+1}H_s^n`$, be the canonical semi-Riemannian submersions. By Theorem 2.5 in \[Esc2\], we see that $`\pi \eta :H_{4s+3}^{4n+3}B`$ is a semi-Riemannian submersion with connected totally geodesic fibres. We remark that the dimension of the fibres of $`\pi \eta `$ is greater than or equal to $`4`$. Thus, by Theorem 1.2, we see that $`\pi \eta `$ is equivalent to the canonical semi-Riemannian submersion
$$H_7^{15}H^8(4),\mathrm{or}H_{15}^{15}H_8^8(4).$$
It follows that $`\pi `$ is one of the following types:
* $`\pi :H_1^3H^8(4)`$, or
* $`\pi :H_3^3H_8^8(4)`$.
In \[Ucc\], Ucci proved that there are no Riemannian submersions with fibres $`P^1`$ from $`P^3`$ onto $`S^8(4)`$. Therefore, Case (ii) is not possible.
The fibres of semi-Riemannian submersion $`\pi \xi :H_3^7H^8(4)`$ are totally geodesic by Theorem 2.5 in \[Esc2\], and complex submanifolds, since the horizontal lifting of the tangent space of the quaternionic line $`\pi ^1(\pi (p))`$ is invariant under the canonical complex structure on $`H_3^7`$. By \[Ba-Ia\], there are no semi-Riemannian submersions with complex totally geodesic fibres from $`H_3^7`$ onto $`H^8(4)`$. Thus Case (i) is impossible. ∎
|
warning/0005/cond-mat0005376.html
|
ar5iv
|
text
|
# Asymptotic step profiles from a nonlinear growth equation for vicinal surfaces
## I Introduction
Ten years ago, Bales and Zangwill predicted that a growing vicinal surface should undergo a step meandering instability when kinetic step edge barriers suppress the attachment of atoms to descending steps . The instability has meanwhile been observed in experiments and Monte Carlo simulations , and a number of theoretical studies have been devoted to the nonlinear evolution of the surface both in the presence and absence of desorption .
Since linear stability analysis shows the in-phase mode of the collective step meander to be the most unstable , the two-dimensional surface morphology can be represented by a one-dimensional function $`\zeta (x,t)`$ describing the displacement of the common step profile from the flat straight reference configuration $`\zeta =0`$, with the $`x`$-axis oriented along the step . For the case of infinite step edge barriers, attachment-detachment kinetics and no desorption, the nonlinear evolution equation
$$\zeta _t=\left\{\frac{\alpha \zeta _x}{1+\zeta _x^2}+\frac{\beta }{1+\zeta _x^2}\left[\frac{\zeta _{xx}}{(1+\zeta _x^2)^{3/2}}\right]_x\right\}_x$$
(1)
was proposed in Ref. (subscripts denote derivatives). It can be derived from the Burton-Cabrera-Frank (BCF) theory of growth on vicinal surfaces using a singular multiscale expansion in $`ϵ^{1/2}`$, where $`ϵ=\mathrm{\Omega }F\mathrm{}^2/D`$ is the Péclet number. Here $`F`$ is the deposition flux, $`D`$ the in-plane surface diffusion coefficient, $`\mathrm{}`$ the nominal step spacing, and $`\mathrm{\Omega }`$ the atomic area. The coefficients in Eq.(1) are given by $`\alpha =\mathrm{\Omega }F\mathrm{}^2/2`$ and $`\beta =\mathrm{\Omega }^2D\mathrm{}\gamma c_{\mathrm{eq}}/k_BT`$, with $`\gamma `$ and $`c_{\mathrm{eq}}`$ referring to the step stiffness and the equilibrium adatom density, respectively.
According to (1), the straight step is linearly unstable against perturbations with wavelengths larger than $`\lambda _c=2\pi \sqrt{\beta /\alpha }`$, with a fastest growing wavelength $`\lambda _u=\sqrt{2}\lambda _c`$. To explore the nonlinear regime, in Ref. a numerical integration of Eq.(1) was carried out which showed an increase of the meander amplitude as $`\sqrt{t}`$ at fixed wavelength $`\lambda _u`$, as well as the formation of spike singularities at maxima and minima of $`\zeta `$. The latter is surprising because the second term on the right hand side of (1) would be expected to suppress such rapid variations of the step curvature.
Here we revisit the problem using a more accurate numerical algorithm . We demonstrate that the step profile remains smooth near maxima and minima, where it approaches asymptotically a stationary (time-independent) solution of (1), while the sides of the profile follow a separable solution with an amplitude of order $`\sqrt{t}`$. The matching of the two solutions occurs near the point of maximum slope. We further show heuristically how the effect of step edge diffusion can be included in the theory, and introduce a generalized evolution equation which contains edge diffusion and attachment-detachment kinetics as special cases. Finally, we address the question to what extent an initially imposed meander wavelength different from $`\lambda _u`$ is preserved under the time evolution. This is relevant in view of the recent experiments of Maroutian et al. .
## II Shape selection
Before presenting the numerical results we recapitulate the two classes of analytic solutions to (1) which were found in . Stationary solutions are obtained by setting the mass current along the step (the quantity inside to curly brackets on the right hand side of (1)) to zero. In terms of $`m(x)=\zeta _x/\sqrt{1+\zeta _x^2}`$ the stationarity condition reduces to Newton’s equation $`\beta d^2m/dx^2=dU/dm`$ for a classical particle of mass $`\beta `$ moving in the potential $`U(m)=\alpha \sqrt{1m^2}`$, which can be solved by quadratures. One thus obtains a one-parameter family of periodic profiles $`\zeta _S(x)`$ which are most conveniently parameterized by the maximum slope $`S\mathrm{max}_x\zeta _x`$, and which have been described previously in the context of a different surface evolution equation . The amplitude $`A(S)`$
is an increasing function of $`S`$, while the wavelength $`\mathrm{\Lambda }(S)`$ decreases with increasing $`S`$, starting out at $`\mathrm{\Lambda }(0)=\lambda _c`$. For $`S\mathrm{}`$ finite limiting values $`A(\mathrm{})=\sqrt{8\beta /\alpha }`$, $`\mathrm{\Lambda }(\mathrm{})=\sqrt{2\pi \beta /\alpha }\mathrm{\Gamma }(3/4)/\mathrm{\Gamma }(5/4)0.5393527..\lambda _c`$ are approached.
The separable solution of interest reads
$$\zeta (x,t)=2\sqrt{\alpha t}\mathrm{erf}^1\left(14|x|/\lambda _s\right),$$
(2)
$`\lambda _s/2<x<\lambda _s/2`$, where $`\mathrm{erf}(z)=(2/\sqrt{\pi })_0^z𝑑ye^{y^2}`$, and the wavelength $`\lambda _s`$ is arbitrary. Equation (2) solves (1) exactly in the limit $`t\mathrm{}`$, when the second term on the right hand side becomes negligible compared to the first, and the evolution equation reduces to $`\zeta _t=(\alpha /\zeta _x)_x`$. The solution (2) is singular near the maxima and minima, where it diverges as $`\zeta \pm \sqrt{\mathrm{ln}(1/|xx_0|)}`$, $`x_0=0,\pm \lambda _s/2`$.
In Figure 1 we show results of a numerical solution of (1), starting from a small amplitude random initial condition. To secure good numerical stability we used a fully implicit, backwards Euler algorithm for integration. The algorithm was implemented on an adaptive grid in order to obtain sufficient lateral resolution at the singular points.
A regular meander pattern of wavelength $`\lambda _u`$ develops, with an amplitude growing indefinitely as $`\sqrt{t}`$. Closer inspection reveals that the sides of the profile follow the separable solution (Figure 2), while near the maxima and minima smooth caps appear which approach pieces of the
stationary solutions (Figure 3). This can be understood by noting that the mass current along the sides of the profile vanishes as $`1/\sqrt{t}`$ according to (2), and therefore the stationarity condition is asymptotically satisfied; we have checked that the deviation from the stationary profile which is discernible in Figure 3 vanishes as $`1/\sqrt{t}`$. Since the slope of (2) increases monotonically upon approaching an extremum while it decreases for the stationary profiles, the matching of the two solutions occurs near the point of maximum slope. For $`t\mathrm{}`$ the slope of the separable solution diverges, hence the cap profile approaches the limiting stationary solution $`\zeta _{\mathrm{}}(x)`$, and the length of the cap becomes $`\mathrm{\Lambda }(\mathrm{})/2`$. The rescaled step profile $`\zeta (x,t)/\sqrt{t}`$ approaches an invariant shape in which the cap appears as a flat facet. The wavelength $`\lambda _s`$ of the separable solution depends on the cap length and on the total meander wavelength $`\lambda `$, and is fixed by mass balance requirements ; for large total wavelength $`\lambda _s\lambda `$ (see Figure 2).
## III Step edge diffusion and a generalized evolution equation
On many fcc metal surfaces, diffusion along step edges is the fastest kinetic process which therefore provides the dominant step smoothening mechanism . To see how Eq.(1) has to be modified to take this effect into account, we note first that the second, relaxational term on the right hand side can be rewritten in a geometrically covariant form as $`(\sigma \mu _s)_x`$, where $`\mu =\mathrm{\Omega }\gamma \kappa `$ is the step chemical potential , $`\kappa =(1+\zeta _x^2)^{3/2}\zeta _{xx}`$ is the step curvature, $`s=𝑑x\sqrt{1+\zeta _x^2}`$ is the arclength along the step, and
$$\sigma =\frac{D\mathrm{\Omega }c_{\mathrm{eq}}}{k_BT}\frac{\mathrm{}}{\sqrt{1+\zeta _x^2}}$$
(3)
is a mobility. The $`\mathrm{}`$-dependence in (3) reflects the assumed relaxation kinetics , in which mass exchange between different parts of the step occurs through detachment followed by diffusion over the terrace, reflection at the descending step, and re-attachment (case E of ). The factor $`1/\sqrt{1+\zeta _x^2}`$ has a simple geometric interpretation : For a deformed in-phase step train the distance to the nearest step, measured along the step normal, is $`\mathrm{}/\sqrt{1+\zeta _x^2}`$ rather than $`\mathrm{}`$.
For relaxation through step edge diffusion the mobility is clearly independent of the step distance, and is given by $`\stackrel{~}{\sigma }=D_\mathrm{e}\mathrm{\Omega }c_\mathrm{e}/k_BT`$, where $`D_\mathrm{e}`$ and $`c_\mathrm{e}`$ denote the edge diffusion coefficient and the equilibrium concentration of edge atoms, respectively. When edge diffusion dominates (i.e. $`\stackrel{~}{\sigma }\sigma `$), the appropriate nonlinear growth equation should thus be given by (1) with the second term replaced by $`(\stackrel{~}{\sigma }\mu _s)_x=([1+\zeta _x^2]^{1/2}\stackrel{~}{\sigma }\mu _x)_x`$. This is confirmed by the explicit derivation of Gillet et al. , who also studied the crossover between attachment-detachment kinetics and edge diffusion. Here our primary goal is to gain further insight into the shape selection mechanism. This has lead us to consider the generalized class of equations
$$\zeta _t=\left\{\frac{\alpha \zeta _x}{1+\zeta _x^2}+\frac{\beta }{(1+\zeta _x^2)^n}\left[\frac{\zeta _{xx}}{(1+\zeta _x^2)^{3/2}}\right]_x\right\}_x,$$
(4)
which reduces to (1) for $`n=1`$ and describes relaxation through step edge diffusion when $`n=1/2`$ and $`\beta =\mathrm{\Omega }^2D_\mathrm{e}c_\mathrm{e}\gamma /k_BT`$. Below we discuss the properties of (4) for general $`n`$, keeping in mind that the cases $`n=1/2`$ and $`n=1`$ are of immediate physical relevance.
The separable solution (2) becomes exact in a limit where the relaxation term in (1) can be neglected, hence it remains a valid asymptotic solution also of (4) for $`n>1/2`$; for $`n1/2`$ the relaxation term can never be ignored. The stationary solutions of (4) can be analyzed in terms of the same mechanical analogy described above, the particle potential being given by $`U(m)=\alpha (1m^2)^{3/2n}/(32n)`$. For $`1/2<n<3/2`$ the behavior is analogous to that for $`n=1`$: The wavelength $`\mathrm{\Lambda }(S)`$ is a decreasing function of the maximal slope $`S`$, and wavelength and amplitude reach finite values $`A(\mathrm{})=\sqrt{8\beta /\alpha }\sqrt{32n}/(2n1)`$ and
$$\mathrm{\Lambda }(\mathrm{})=\sqrt{2\pi (32n)(\beta /\alpha )}\frac{\mathrm{\Gamma }[(2n+1)/4]}{\mathrm{\Gamma }[(2n+3)/4]}$$
(5)
for $`S\mathrm{}`$. Thus the asymptotic step profiles look similar to those generated by Eq.(1), with the length of the cap decreasing with increasing $`n`$. As $`n3/2`$ the cap length, given by $`\mathrm{\Lambda }(\mathrm{})/2`$, vanishes. For $`n3/2`$ we therefore expect true spike singularities to develop at the maxima and minima of the profile. Using that the slope imposed by the separable solution (2) grows as $`S\sqrt{t}`$, we predict that the curvature at the extrema diverges as $`t^{(2n3)/4}`$.
A numerical solution for the case of edge diffusion ($`n=1/2`$) is shown in the lower panel of Figure 1. For $`n=1/2`$ the potential $`U(m)`$ is harmonic and hence the wavelength $`\mathrm{\Lambda }(S)=\lambda _c`$ independent of $`S`$. The amplitude of the stationary profiles diverges as $`A(S)\mathrm{ln}S`$, leading to a corresponding increase of the cap height as $`\mathrm{ln}t`$. Since this is still small compared to the overall profile amplitude, the caps nevertheless appear as flat in the rescaled shape $`\zeta /\sqrt{t}`$ (Figure 2; a detailed view of the cap is shown in Figure 4). This remains true in the entire interval $`1/2<n1/2`$, where the cap length (5) remains finite and the cap height grows as $`t^{(12n)/4}`$. However a qualitative change in the profile evolution occurs at the value $`n_c0.2283`$ where the asymptotic stationary wavelength (5) becomes equal to the most unstable wavelength $`\lambda _u`$, which sets the lateral length scale in the
early stages of growth. For $`n<n_c`$ we expect to see an intermediate coarsening regime in which the lateral length scale increases from $`\lambda _u`$ to $`\mathrm{\Lambda }(\mathrm{})`$. Coarsening to arbitrarily large length scales, similar to what is observed in related evolution equations for one-dimensional unstable growth , sets in at $`n=1/2`$, where (5) diverges. Throughout the regime $`n<n_c`$ the evolving profile is describable in terms of stationary solutions, and the separable solution (2) no longer plays any role (Figure 4).
## IV Persistence of the initial wavelength
Finally, we address the recent experiments on surfaces vicinal to Cu(100), in which the meander wavelength was measured as a function of temperature, and it was concluded that the observed behavior is inconsistent with the theoretical prediction for the linearly most unstable wavelength $`\lambda _u`$. Maroutian et al. therefore proposed that the meander wavelength is set by the nucleation length describing the distance between the one-dimensional nuclei appearing on a flat step in the early stage of growth, which can be considerably larger than $`\lambda _u`$.
A necessary consistency requirement for this scenario is that an initially imposed meander wavelength $`\lambda _i>\lambda _u`$ persists during the nonlinear evolution. We have therefore numerically integrated Eqs.(1, 4) starting from a sinusoidal initial condition with varying wavelength $`\lambda _i`$. We do find that a range of wavelengths can be preserved during growth. This is reasonable in view of the analysis presented above, which shows that asymptotic profiles, composed of the separable solution (2) and a stationary cap, can in principle be constructed for arbitrary wavelength (see e.g. Figure 2). However, when $`\lambda _i`$ exceeds $`\lambda _c`$ by more than a factor of 3, so that an additional meander fits between the maxima and minima of the profile, the wavelength spontaneously decreases to a value near $`\lambda _u`$ (Figure 5). This result contradicts the assumption of that initial wavelengths much larger than $`\lambda _u`$ persist, but it should not be overemphasized: Clearly processes which involve a change in the collective meander wavelength may not be accurately described in a model which assumes in-phase meandering from the outset.
## V Outlook
In conclusion, we have described an unusual shape selection scenario for a class of physically motivated growth equations. A number of issues remain to be clarified. Mathematically, the behavior in the region where separable and stationary solutions match needs further investigation; our numerical work indicates the appearance of singularities in higher derivatives of $`\zeta `$. Also the dynamics in the singular regime $`n3/2`$ and in the coarsening regime $`n<n_c`$ of Eq.(4) deserves attention. Physically, it is imperative that the predictions of the one-dimensional equations for the in-phase step meander be confirmed by more complete descriptions of the growing surface, as provided by two-dimensional continuum equations and Monte Carlo models , in order to assess their ultimate relevance for the experimentally observed morphologies.
## Acknowledgements
We are much indebted to Jens Eggers for help with the numerical algorithm and useful discussions. Olivier Pierre-Louis and Chaouqi Misbah kindly supplied us with a copy of prior to publication. Support by DFG/SFB 237 is gratefully acknowledged.
|
warning/0005/cond-mat0005420.html
|
ar5iv
|
text
|
# Andreev Peaks and Massive Magnons in Cuprate SNS junctions
\[
## Abstract
The projected SO(5) theory (pSO(5)) is used to resolve the puzzle of two distinct energy gaps in high T<sub>c</sub> Superconductor-Normal-Superconductor junctions. Counter to conventional theory of multiple Andreev reflections (MAR), the differential resistance peaks are associated with the antiferromagnetic resonance observed in neutron scattering, and not with Cooper pair breaking. The pSO(5) and MAR theories differ by the expected tunneling charges at the peaks. We propose that shot noise experiments could discriminate against the conventional interpretation.
PACS numbers: 74.20.-z,74.65.+n
\]
In current transport through high $`T_c`$ superconductor junctions, there seem to be two energy scales. The upper energy is seen in tunneling conductance, and is identified with the “pseudogap” $`\mathrm{\Delta }_p`$ which appears in magnetic resonance and photoemmission. A lower gap, which scales differently with hole doping, manifests as peaks in the differential resistance of Superconducting-Normal-Superconducting (SNS) Josephson junctions. These peaks have been interpreted using the conventional theory of multiple Andreev reflections, following Klapwijk, Blonder, and Tinkham (KBT).
KBT theory treats two conventional superconductors with a single $`s`$-wave BCS quasiparticle gap $`\mathrm{\Delta }`$, separated by a free electron metal. Electrons traversing the metal are Andreev reflected back as holes, gaining energy increments $`eV`$ at each traversal (as depicted in Fig. 1). Peaks in the differential resistance appear at voltages $`2\mathrm{\Delta }/ne`$, and are due to the $`(E\mathrm{\Delta })^{1/2}`$ singularity in the quasiparticles’ density of states.
However, in cuprate SNS junctions, such as YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6.6</sub> \- YBa<sub>2</sub>Cu<sub>2.55</sub>Fe<sub>0.45</sub>O<sub>y</sub> -YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6.6</sub> examined by Nesher and Koren , application of KBT theory is problematic. A naive fit to KBT expression faces the two gaps puzzle, i.e. an “Andreev gap” is of order $`\mathrm{\Delta }16`$meV, while the tunneling gap is about three times larger, and scales differently with $`T_c`$. Without perfect alignment of the interfaces, it is hard to understand the observed sharpness of peaks since the $`d`$-wave gap is modulated for different directions. Moreover, the barrier is by no means a “normal” metal devoid of interactions: it is an underdoped cuprate with antiferromagnetic correlations and strong pairing interactions as evidenced by a large proximity effect.
The purpose of this Letter is to provide an alternative explanation for the differential resistance peaks series, which takes into account the strong correlations in the pseudogap regime. Our analysis resolves the two energy scales puzzle.
We employ the projected SO(5) (pSO(5)) model, which is a strong coupling effective Hamiltonian. It describes the dynamics and interactions of four primary bosonic modes of cuprates: preformed hole pairs and massive spin one magnons.
A differential resistance peaks series is found at bias voltages $`V_n=\mathrm{\Delta }_s/(en)`$, $`n=1,2,\mathrm{}`$ where $`\mathrm{\Delta }_s`$ is the antiferromagnetic resonance energy. This resonance has been directly measured by inelastic neutron scattering. The peaks are thus associated with emmission of magnon pairs at the resonance threshold, and not with pair breaking, as in KBT theory. We note that other predictions to observe magnons (also called $`\pi `$-modes) in various cuprate junctions were made, but await experimental confirmation. We propose that measurement of the excess shot noise below the peaks, could discriminate against the latter interpretation. pSO(5) theory predicts tunneling charge $`2ne`$ below the $`n`$-th peak, while KBT theory expects charge $`ne`$.
Degrees of freedom: At energies below the pseudogap $`\mathrm{\Delta }_p`$, preformed hole pairs (with internal $`d`$-wave symmetry), describe the primary charge degrees of freedom in the underdoped regime. The hole pairs are bosons, and their phase fluctuations are controlled by the two dimensional superconducting stiffness $`\rho _c`$, as measured by the London penetration depth. At $`T_c`$, the pairs Bose condense and long range phase coherence is established. This scenario can explain the empirical relations $`T_c\rho _c`$, which have been observed in cuprates at low doping concentrations. The other low energy charge excitations are fermionic quasiparticles near the $`d`$-wave nodes. These have a smooth density of states which decreases below $`\mathrm{\Delta }_p`$.
Additional bosonic excitations below the pseudogap energy scale, are antiferromagnetic spin fluctuations i.e. magnons. Massive spin one magnons have been observed in inelastic neutron scattering in YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6+δ</sub>. They manifest as a sharp resonance in the spin correlation function $`S_{\alpha \alpha ^{}}`$, which near the antiferromagnetic wavevector $`𝐪\stackrel{}{\pi }`$ has the form
$$S_{\alpha \alpha ^{}}(\omega ,𝐪)s_0\frac{\delta _{\alpha \alpha ^{}}}{\omega ^2c^2(𝐪\stackrel{}{\pi })^2\mathrm{\Delta }_s^2}$$
(1)
Here $`c`$ is the spin wave velocity, and $`s_0`$ is a normalization factor. The doping dependent resonance energy $`\mathrm{\Delta }_s(\delta )`$ increases between $`\mathrm{\Delta }_s(0.5)=25`$meV, (with $`T_c=52^{}K`$), and $`\mathrm{\Delta }_s(1)=40`$meV, (at $`T_c=92^{}K`$). The projected SO(5) theory. The large onsite Hubbard repulsion between electrons is imposed by an apriori projection of doubly occupied states from the Hilbert space.
The undoped vacuum $`|0`$ is a half filled Mott insulator in a quantum spin liquid state. The pSO(5) vacuum possesses short range antiferromagnetic correlations. A translationally invariant realization of $`|0`$ on the microscopic square lattice, is the short range resonating valence bonds state.
Out of this undoped vacuum, $`b_h^{}`$ create charge $`2e`$ bosons (hole pairs) with internal $`d`$-wave symmetry under rotations, and $`b_\alpha ^{}`$, $`\alpha =x,y,z`$ create a triplet of antiferromagnetic, spin one magnons.
The lattice pSO(5) Hamiltonian is
$`^{pSO(5)}`$ $`=`$ $`^{charge}+^{spin}+^{int}+^{Coul}+^{ferm}`$ (2)
$`^{charge}`$ $`=`$ $`(ϵ_c2\mu ){\displaystyle \underset{i}{}}b_{hi}^{}b_{hi}^{}{\displaystyle \frac{J_c}{2}}{\displaystyle \underset{ij}{}}\left(b_{hi}^{}b_{hi}^{}+\text{h.c.}\right)`$ (3)
$`^{spin}`$ $`=`$ $`ϵ_s{\displaystyle \underset{i\alpha }{}}b_{\alpha i}^{}b_{\alpha i}^{}J_s{\displaystyle \underset{\alpha ij}{}}n_i^\alpha n_j^\alpha `$ (4)
$`^{int}`$ $`=`$ $`W{\displaystyle \underset{i}{}}:\left(b_{hi}^{}b_{hi}^{}+{\displaystyle \underset{\alpha }{}}b_{\alpha i}^{}b_{\alpha i}^{}\right)^2:,`$ (5)
where $`:():`$ denotes normal ordering, and $`n_i^\alpha =(b_{i\alpha }^{}+b_{i\alpha }^{})/\sqrt{2}`$ is the Néel spin field. $`^{int}`$ describes short range interactions between bosons, and $`^{Coul}`$ describes the long range Coulomb interactions. $`H^{ferm}`$ describes coupling to the nodal (fermionic) quasiparticles, which contribute to a large, but smooth, conductance background. Here we will concentrate on the conductance singularities, and will not compute the fermionic background.
The mean field approximation to Eq. (LABEL:pSO5) is straightforward. It amounts to replacing $`b_{\gamma i}^{}b_{\gamma i}^{}`$, $`\gamma =h,\alpha `$, and minimizing $`^{charge}+^{spin}+^{int}`$ with respect to $`b_{\gamma i}^{}`$. There is a first order transition between the two primary mean field phases on the square lattice at $`\mu =\mu _c`$, where
$$\mu _c=\frac{1}{2}(ϵ_cϵ_s)(J_c2J_s),$$
(7)
$`\mu _c`$ is of the order of the Hubbard interaction scale.
At $`\mu <\mu _c`$ we have an undoped Mott insulator with no hole pair bosons, and where the magnons Bose-condense. The condensate supports a finite staggered magnetization
$$|n^\alpha |^2=(2J_s\frac{1}{2}ϵ_s)/W\mu <\mu _c$$
(8)
There are two linear spin wave modes $`\omega =c|𝐪|`$, with where $`c=\sqrt{2}J_s/\mathrm{}`$ is the semiclassical spinwave velocity of the Heisenberg antiferromagnet.
At $`\mu >\mu _c`$ the ground state becomes doped with hole pairs which Bose-condense into a superconducting phase with an order parameter
$$|b_i^{}|^2=(\mu \mu _c+2J_sϵ_s/2)/W\mu >\mu _c$$
(9)
Long range interactions in $`^{Coul}`$, frustrate the first order transition and create intermediate (possibly incommensurate) phases, which we shall not discuss here.
The mean field phase stiffness is given by $`\rho _c=J_cb_i^{}^2`$, and therefore Eq. (9) explains why $`\rho _c`$ increases with chemical potential (and doping) in the underdoped superconducting regime, as observed experimentally.
Analysis of the linear quantum fluctuations about mean field theory determines the magnon dispersion i.e. the poles of Eq. (1). The mean field magnon gap is found to be
$$\mathrm{\Delta }_s=2\sqrt{(\mu \mu _c)(\mu \mu _c+4J_s)}$$
(10)
which by Eq. (9) implies that $`\mathrm{\Delta }_s^2\rho _c,T_c`$. Thus the pSO(5) mean field theory can explain the systematic increase of $`\mathrm{\Delta }_s`$ with $`T_c`$ which is observed by Fong et. al..
The cuprate SNS junction. We consider a junction, where the barrier (N) has no superconducting or magnetic order $`b_h^{}=0,n^\alpha =0`$. We derive on general grounds the form of the effective tunneling Hamiltonian between superconductors as follows. An integration of the barrier’s charged bosons $`b_h`$ out of the path integral results in an effective action $`𝒮^{tun}`$ which couples the charges of the two superconductors . $`𝒮^{tun}[b_{h_L},b_{h_R},b_\alpha ]`$ explicitly depends on the hole pairs bosons on the left and right interfaces, and on the magnons in the barrier. By charge conservation, an expansion of $`𝒮^{tun}`$ as a power series leaves only terms with equal number of $`b_h`$’s and $`b_h^{}`$’s. By spin conservation, the magnon terms are singlets and hence at least bilinear in $`n^\alpha `$.
This expansion leads to a series of tunneling terms. For the Andreev peaks we retain only the leading order terms (in $`b^{},b`$) which are
$`^{tunmag}`$ $`=`$ $`{\displaystyle \underset{n}{}}(𝒜_n^{}+𝒜_n^{})`$ (11)
$`𝒜_n`$ $`=`$ $`{\displaystyle \underset{y_1\mathrm{}y_{2n},𝐱,𝐱^{}}{}}T_nb_{h_L,1}^{}\mathrm{}b_{h_L,n}^{}b_{h_Rn+1}^{}\mathrm{}b_{h_R,2n}^{}`$ (13)
$`\times \left({\displaystyle \underset{\alpha }{}}n^\alpha (𝐱)n^\alpha (𝐱^{})\right)`$
$`𝒜_n^{}`$ describes a simultaneous tunneling of $`n`$ hole pairs from the left to the right superconductor, coupled to a magnon pair excitation. $`T_n`$ is the tunneling vertex function, which depends on the bosons positions.
The energy transfer mechanism is depicted diagrammatically in Fig.2. We do not compute $`T_n`$’s which depend on the details of the barrier and the interfaces. A “good” N barrier is defined to have sizeable $`T_n`$, if multiple pair tunneling terms are to be observed. This requires a thin barrier with slowly decaying spin and charge correlations. It is important to note that multiple pair tunneling, i.e. the differential resistance peaks at $`n>1`$, depends on strong anharmonic interactions between the hole pairs and magnons. These interactions are an essential part of the pSO(5) theory as modelled by $`^{int}`$ in Eq. (LABEL:pSO5).
The junction’s conductance is calculated in the standard fashion: the bias voltage $`V`$ transforms the left bosons $`b_{h_L}e^{i2eVt}b_{h_L}`$, which yields time dependent operators $`𝒜_n(t)`$. The current is calculated by second order perturbation theory in $`^{tunmag}`$ yielding
$`I`$ $`=`$ $`{\displaystyle \underset{n}{}}2neX_n^{ret}(2eV)`$ (14)
$`X_n^{ret}(\omega )`$ $`=`$ $`i{\displaystyle _0^{\mathrm{}}}𝑑te^{i\omega t}[A_n^{}(t),A_n^{}]`$ (15)
For singular contributions $`I^{sing}`$, we ignore superconducting condensate fluctuations $`b_h^{}b_h^{}`$, which have a smooth spectrum. Similarly, we ignore the frequency dependence of $`T_n(\omega )`$. Setting $`b_R^{}b_h^{}`$ and $`b_L^{}e^{i2eVt}b_h^{}`$ leads to
$`I^{sing}`$ $`=`$ $`{\displaystyle \underset{n}{}}2ne{\displaystyle \underset{|q_x|\pi /d,|q_y|\pi /W}{}}b_h^{}^{4n}|T_n[𝐪]|^2`$ (17)
$`\times \mathrm{}{\displaystyle \underset{\omega }{}}S(𝐪,i\omega +2neV+i0^+)S(𝐪,i\omega )`$
where the barrier dimensions are $`d\times W`$ (see Fig.2), and $`_\omega `$ is a Matsubara sum.
For a nearly antiferromagnetic “N” barrier, $`T_n(𝐱𝐱^{})`$ in (13) decays slowly with the distance between magnons. Thus for a narrow barrier $`d<<W`$, the magnons are excited at $`q_y0`$, and the momentum sum reduces to a one dimensional sum over $`q_x`$. At zero temperature we obtain
$`I^{sing}`$ $`=`$ $`{\displaystyle \underset{n}{}}2neb_h^{}^{4n}|T_n[0]|^2`$ (19)
$`\times s_0^2{\displaystyle \frac{dq_x}{2\pi }\frac{\delta (2neV2\sqrt{c^2q_x^2+\mathrm{\Delta }_s^2})}{2(\mathrm{\Delta }_s^2+c^2q_x^2)}}`$
$``$ $`{\displaystyle \underset{n}{}}t_n{\displaystyle \frac{\theta (neV\mathrm{\Delta }_s)}{\mathrm{\Delta }_s^{3/2}\sqrt{neV\mathrm{\Delta }_s}}}`$ (20)
The last expression emphasizes the singular form of $`I^{sing}(V,\mathrm{\Delta }_s)`$ at the peaks. For a large background conductance $`dI/dV>>dI^{sing}/dV`$, the inverse square root singularities in $`I^{sing}`$ create peaks in the differential resistance $`dV/dI`$ at voltages
$$V_n=\mathrm{\Delta }_s/(ne),n=1,2,\mathrm{},Q_n=2ne$$
(21)
where $`Q_n`$ is the excess tunneling charge below the $`n`$-th peak. Note that $`Q_n`$ changes in increments of $`2e`$. The differential resistance peak series is depicted in Fig. 3, for weak broadening of the singularities and an arbitrary set of coefficients $`t_n`$.
Discussion. We have seen that magnon pair creation induces peaks in the differential resistance which are similar in appearance to the Andreev peaks of the KBT mechanism. The crucial difference is that here the singular dissipative process does not involve Cooper pair breaking, but low energy antiferromagnetic excitations.
In KBT theory for two identical superconductors, the peaks appear at voltages $`V_n^{KBT}=2\mathrm{\Delta }/(ne),n=1,2,\mathrm{}`$ which are the upper threshold for tunneling of charges $`Q_n=ne`$. Thus, KBT allows both even and odd number of electron charges to participate in the multiple Andreev reflection process, as depicted in Fig. 1. These charges change in increments of $`e`$ at each peak. Therefore a decisive discrimination between the processes of Fig.1 and Fig.2 would be measurements of the excess tunneling charge increments at the peaks. A feasible method would perhaps be low temperature shot noise $`S`$ which measures the tunneling charges via the relation $`S=2Q_nI(V_n)`$. We eagerly look forward to the results of such experiments.
Acknowledgements. We thank G. Deutscher, G. Koren, A. Mizel, O. Nesher and E. Polturak for useful discussions. Support from the Israel Science Foundation and the Fund for Promotion of Research at Technion is acknowledged.
|
warning/0005/quant-ph0005050.html
|
ar5iv
|
text
|
# Transient interference of transmission and incidence
## Abstract
Due to a transient quantum interference during a wavepacket collision with a potential barrier, a particular momentum, that depends on the potential parameters but is close to the initial average momentum, becomes suppressed. The hole left pushes the momentum distribution outwards leading to a significant constructive enhancement of lower and higher momenta. This is explained in the momentum complex-plane language in terms of a saddle point and two contiguous “structural” poles, which are not associated with resonances but with incident and transmitted components of the wavefunction.
The traditional formulation of quantum scattering theory in terms of an “$`S`$-matrix” assumes that only the results of the collision can be observed at “asymptotic” distances and times, but that the collision itself cannot be observed. This perspective is justified to analyze the products of standard collisions of atomic or molecular beams. But the $`S`$-matrix approach is not enough to describe modern experiments where the collision complex can be observed by means of femtosecond laser pulses or “spectroscopy of the transition state” . Also, in quantum kinetic theory of gases accurate treatments must abandon the “completed collision” approximation and use a full description, e.g. in terms of Möller wave operators as in the Waldmann-Snider equation and its generalizations for moderately dense gases . In any case, it is important to understand the collision process itself to control or modify the products. This has motivated a recent trend of theoretical and experimental work to investigate the details of the collision itself, and not only its asymptotics. In particular, a quantum effect has been recently described by Brouard and Muga in which the probability to find the particle with a momentum above a given value is larger, in the midst of the collision, than the quantity allowed classically by energy conservation. The effect belongs to a group where the conservation of classical energy seems to be violated. Well known examples are the tunnel effect, or in general the non vanishing probability to find the particle in evanescent regions beyond the classical turning points.
This transient enhancement of the momentum tail may in principle be observed by collisions of ultracold atoms with a laser field that can be turned off suddenly in the time scale of the atomic motion , and implies as a macroscopic consequence deviations from the Maxwellian velocity distribution . We initiated the research of the present work looking for conditions that increase the effect and favor its observability. In so doing we have found an unexpected regime where the effect is much higher than in previously studied cases. In this letter we shall describe such a regime and analyze its physical origin, namely a transient interference between transmission and incidence components of the wavepacket. Let us first review briefly the main aspects of the classically forbiden increase of high-momenta. Brouard and Muga have studied several examples where the quantity
$$G^\mathrm{q}(p,t)_p^{\mathrm{}}\left\{|\psi (p^{},t)|^2|\psi (p^{},0)|^2\right\}𝑑p^{}$$
takes on positive values for positive potentials (the corresponding classical quantity is negative or zero due to energy conservation) . An important aspect of this effect is its transient character, $`G^\mathrm{q}0`$ before and after the collision. The effect is also generic , because the stationary components of the wavepacket have, in momentum representation, a tail due to the resolvent which is always present in the Lippmann-Schwinger equation. This tail goes beyond the maximum value allowed by the conservation of energy.
For a Gaussian wavepacket colliding with an infinite wall, maximum values of $`G_{max}^\mathrm{q}(p,t)0.05`$ have been reported . Also a “delta” potential was used to analyze the influence of the opacity of the barrier. For the cases examined, an increase of $`G^\mathrm{q}`$ with the opacity was observed up to a saturation level where the infinite wall results were recovered . This suggested that the observability of the effect would improve with strongly opaque conditions. In a complementary study we have systematically varied the spatial variance of the wavepacket, $`\delta _x`$, and the height of a square barrier, $`V_0`$, for a fixed average initial momentum $`p_c`$. We have found, contrary to previous expectations, that the maximum effect corresponds to energies well above the barrier and to large values of $`\delta _x`$. In this regime the barrier is not at all opaque and essentially the full wave is finally transmitted. Moreover, $`G_{max}^\mathrm{q}`$ is as large as $`0.27`$.
The numerical effort to perform these calculations by ordinary propagation methods (such as the split operator method) is rather heavy, since large values of $`\delta _x`$ and the need to discern fine details of the momentum distribution require an extense and dense grid. In fact for very large values of $`\delta _x`$ this numerical route has to be eventually abandoned. But even if one gets numerical results with enough computer power, they will not provide any explanation of the unexpectedly high $`G^\mathrm{q}`$ values. Fortunately these two difficulties can be overcome with an approximate analytical solution. Here we shall sketch its obtention, a more detailed account will be given elsewhere. First the momentum representation of the wavefunction is expressed using the basis of stationary eigenstates of $`H`$, $`|p^{}_{}{}^{}+`$, corresponding to incident momentum $`p^{}`$, and energy $`E^{}=p^2/(2m)`$,
$$\psi (p,t)=_{\mathrm{}}^+\mathrm{}p|p^{}_{}{}^{}+e^{iE^{}t/\mathrm{}}p^{}_{}{}^{}+|\psi (t=0)𝑑p^{}.$$
(1)
If the initial state at time $`t=0`$ does not overlap with the potential and has negligible negative momentum components we can write
$$\psi (p,t)=_0^+\mathrm{}p|p^{}_{}{}^{}+e^{iE^{}t/\mathrm{}}p^{}|\psi (t=0)𝑑p^{}.$$
(2)
To facilitate the treatment of the integral in the $`p^{}`$-complex plane we may now extend the lower limit to $`\mathrm{}`$ using the analytical continuation of $`p|p^{}_{}{}^{}+`$, $`p^{}>0`$, over $`p^{}<0`$ (and later over the whole complex plane).
For a square barrier of height $`V_0`$ and width $`d`$, centered at the coordinate origin, the delta-normalized stationary wavefunction with incident momentum $`p^{}`$ is
$$x|p_{}^{}{}_{}{}^{+}=\frac{1}{h^{1/2}}\{\begin{array}{cc}Ie^{ik^{}x}+Re^{ik^{}x},\hfill & x<d/2\hfill \\ Ce^{ik^{\prime \prime }x}+De^{ik^{\prime \prime }x},\hfill & d/2<x<d/2\hfill \\ Te^{ik^{}x},\hfill & x>d/2\text{,}\hfill \end{array}$$
(3)
where $`I=1`$, $`k^{}=p^{}/\mathrm{}`$ and $`k^{\prime \prime }=\sqrt{p^{}_{}{}^{}22mV_0}/\mathrm{}`$. The coefficientes $`R,C,D`$ and $`T`$ are determined by continuity of the wavefunction and its derivative. The momentum representation $`p|p^{}{}_{}{}^{+}`$ will correspondingly have five terms. The terms with $`I`$, $`R`$ and $`T`$ have poles in the $`p^{}`$-complex momentum plane at
$`p_I^{}=p+i0`$ (4)
$`p_R^{}=pi0`$ (5)
$`p_T^{}=pi0,`$ (6)
while the terms with $`C`$ and $`D`$ do not have these structural poles , which are not related to resonances or to the potential profile. The four functions $`R,C,D`$ and $`T`$ present an infinite series of resonance and anti-resonance poles in the third and fourth quadrants due to the zeros of a common denominator
$$\mathrm{\Omega }(p^{})=\mathrm{cos}\left(k^{\prime \prime }d\right)\frac{i}{2}\left(\frac{k^{\prime \prime }}{k^{}}+\frac{k^{}}{k^{\prime \prime }}\right)\mathrm{sin}\left(k^{\prime \prime }d\right).$$
(7)
(In particular $`T(p^{})=\mathrm{exp}(ik^{}d)/\mathrm{\Omega }(p^{})`$.) The conditions examined in this work correspond however to “direct scattering”, where these resonance singularities do not play any significant role.
The initial state is taken as a minimum-uncertainty-product Gaussian centered at the position $`\alpha \delta _x`$, $`\alpha >0`$, with average momentum $`p_c`$,
$$p^{}|\psi (t=0)=\left(\frac{2\delta _x}{\pi \mathrm{}^2}\right)^{1/4}e^{\frac{\delta _x(p^{}p_c)^2}{\mathrm{}^2}+\frac{ip^{}\alpha \delta _x}{\mathrm{}}}.$$
(8)
This expression and the momentum representation of (3) are inserted in (2) to obtain five integrals. The full treatment of the resulting integrals follows closely ref. . The integrals with $`C`$ and $`D`$ may be evaluated with the steepest descent method for large values of $`\delta _x`$. The steepest descent path (SDP) is a straight line with slope $`t\mathrm{}/(2m\delta _x)`$, with a saddle point close to $`p_c`$ in the midst of the collision. We shall always assume that the slope is small so that when the integration contour is deformed along this path it “cuts” the resonance poles of the fourth quadrant far from the real axis, and their residues can be neglected (“direct scattering” conditions).
Because of the interference between the saddle and the structural poles, the simple steepest descent treatment valid for $`C`$ and $`D`$ is not valid for the other terms. A uniform expression for a smooth treatment of the pole crossing of the SDP is provided by the $`w`$-function, $`w(z)=e^{z^2}erfc(iz)`$, which may also be defined by its integral expression
$$w(z)=\frac{1}{i\pi }_\mathrm{\Gamma }_{}𝑑u\frac{e^{u^2}}{uz},$$
(9)
where $`\mathrm{\Gamma }_{}`$ goes from $`\mathrm{}`$ to $`\mathrm{}`$ passing below the pole. Since we are interested in wavepackets with energy well above the barrier maximum the “reflection term” with $`R`$ may be neglected. The remaining contribution is
$$\psi _{IT}=ih^{1/2}\mathrm{}\tau _{\mathrm{}}^{\mathrm{}}\left[g_I(p^{})+g_T(p^{})\right]e^{\varphi (p^{})}𝑑p^{},$$
(10)
where
$`\tau `$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\pi \mathrm{}}}}\left({\displaystyle \frac{2\delta _x}{\pi \mathrm{}^2}}\right)^{1/4}`$ (11)
$`g_I(p^{})`$ $`=`$ $`{\displaystyle \frac{e^{ipd/2\mathrm{}}}{pp^{}+i0}}`$ (12)
$`g_T(p^{})`$ $`=`$ $`{\displaystyle \frac{T(p^{})\mathrm{exp}\left[i\left(2p^{}p\right)d/2\mathrm{}\right]}{(pp^{}i0)}},`$ (13)
and
$$\varphi \left(p^{}\right)=\frac{ip_{}^{}{}_{}{}^{2}t}{2m\mathrm{}}\frac{\delta _x\left(p^{}p_c\right)^2}{\mathrm{}^2}+\frac{ip^{}\left(\alpha \delta _xd/2\right)}{\mathrm{}}.$$
(14)
The functions $`g_I(p^{})`$ and $`g_T(p^{})`$ present structural poles at $`p_I^{}`$ and $`p_T^{}`$; in addition $`g_T(p^{})`$ has resonance and anti-resonance poles.
The SDP is a straight line with the same negative slope as before, passing through the saddle point,
$`s`$ $`=`$ $`{\displaystyle \frac{m}{4m^2\delta _x^2+t^2\mathrm{}^2}}\{4mp_c\delta _x^2+(\alpha \delta _xd/2)\mathrm{}^2t`$ (15)
$`+`$ $`i2\mathrm{}[m\delta _x(\alpha \delta _xd/2)p_c\delta _xt]\}.`$ (16)
To integrate (10), the contour is deformed to the SDP passing over the saddle. The same reasons to neglect the residues from the resonance poles in the $`C`$ and $`D`$ terms are now applicable. To introduce the $`w`$-functions, the integrand must be put in the form (9). We complete the square in (14) and use the change of variable
$$u=\frac{p^{}s}{f},f=\left(\frac{\delta _x}{\mathrm{}^2}+i\frac{t}{2m\mathrm{}}\right)^{1/2}$$
(17)
to obtain
$`p|\psi (t)`$ $``$ $`if\tau h^{1/2}\mathrm{}e^{\left(\delta _xp_c^2/\mathrm{}^2\right)+\eta ^2}`$ (18)
$`\times `$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}\left[g_I\left(u\right)+g_T\left(u\right)\right]e^{u^2}𝑑u,`$ (19)
where $`g(u)g(p^{}(u))`$, and
$$\eta =\left(\frac{2p_c\delta _x}{\mathrm{}^2}+i\frac{(\alpha \delta _xd/2)}{\mathrm{}}\right)\left[4\left(\frac{\delta _x}{\mathrm{}^2}+i\frac{t}{2m\mathrm{}}\right)\right]^{1/2}.$$
We may retain the main contribution from $`g_I`$ from its behaviour near the pole by approximating $`g_I(u)_I/(uu_I)`$, where $`_I`$ is residue of $`g_I(u)`$ at the point $`u=u_I=(p_I^{}l)/f`$, and similarly for $`g_F`$. For an approximate expression of $`p|\psi `$, and considering that the wave is much more extended in space than the barrier we may neglect the contribution from $`C`$ and $`D`$ and retain only the incidence and transmission terms,
$`p|\psi (t)`$ $``$ $`h^{1/2}\pi \tau \mathrm{}e^{\left(\delta _xp_c^2/\mathrm{}^2\right)+\eta ^2}e^{ipd/2\mathrm{}}`$ (20)
$`\times `$ $`\left[w(u_I)+T(p)w(u_T)\right]\psi _{IT}^0(p,t).`$ (21)
A more precise expression including a term $`\psi _{RCD}`$ and corrections to the zeroth order $`\psi _{IT}^0`$ is given elsewhere and allows to obtain the wave function and $`G_{max}^q`$ accurately for large values of $`\delta _x`$ with small computational effort. However (21) captures the essential, and provides a simple, explanatory picture of the phenomenon we want to discuss.
Fig. 1 shows the distribution of momenta $`|p|\psi (t)|^2`$ for different instants of time, from the initial one to a time after the collision has been completed, passing through the instant where $`G^\mathrm{q}=0.27`$ is maximum. In all figures the numerical values correspond to collisions of ultracold Rubidium atoms with an effective laser potential. The observed behaviour does not have a classical explanation. Note that the wavepacket is considerably broader than the barrier. A classical ensemble of particles with the same Gaussian phase-space (Wigner) distribution as (8) would only be slightly deformed due to the small fraction of particles located on the barrier at a given time, and would keep the maximum at the average momentum $`p_c`$. Moreover, there could not be any spectacular acceleration or deceleration as the one seen in the two peaks of the quantum distribution. We shall see that the zero of the quantum momentum distribution, which forbids in this case the initially dominant momentum $`p_c`$, is due to a destructive interference whereas the two new peaks correspond to momentum regions of constructive interference.
In Fig. 2 the Argand diagrams of the two terms are represented, namely the imaginary versus the real parts obtained by varying $`p`$ at equal intervals. Each lobule corresponds to one of the terms. The “motion” as $`p`$ increases begins close to the origin, downwards in both diagrams. The left peak of the momentum distribution corresponds to the zone where the two moduli increase together and are approximately in phase. After the descending motion there is a fast, aproximately circular motion where the phases become opposed (destructive interference). Finally, the two curves meet again in phase in the upper part of the lobules, this momentum interval corresponds to the right peak of the momentum distribution. The described behaviour is essentially due to the two $`w`$-functions $`w(u_I)`$ y $`w(u_T)`$, as shown in Fig. 3, where the two Argand diagrams of the two $`w`$-functions and of the factors that multiply them are represented between the momenta of the two maxima. Clearly the efect of the factors, whose phases remain essentially constant around $`\pi `$, is simply to invert the two lobules of the $`w`$s. The fast motion of the $`w`$-functions is due to the pass of the two contiguous structural poles $`u_I=(p+i0)/f`$ and $`u_T=(pi0)/f`$ near the saddle point at $`u=0`$. A sweep from smaller to larger values of $`p`$ moves the couple of poles along the real $`p^{}`$ axis from left to right, while, for fixed $`t`$, the saddle point and the steepest descent path do not depend of $`p`$. Since $`u_Iu_T`$ we can write, using the relation between $`w`$-functions of argument of opposite sign, see (9),
$$w(u_I)=e^{u_I^2}w(u_T).$$
(22)
During the collision, the saddle point is very close to the real axis of the $`p^{}`$-plane, only slightly below in Fig. 4, and the slope of the SDP is very small. This means that the difference between the two $`w`$-contributions is essentially a real exponential, which implies a “simultaneous motion”, with equal maginary parts, along the two lobules of the Argand diagram. The rapid variation of the phases of the $`w`$s when passing close to the saddle point follows from its integral expression (9). When $`u_I`$ passes close to $`u=0`$ and close to the real axis of the $`u`$-plane, the denominator is essentially real and changes its sign quickly, so there is a rapid change by $`\pi `$ in the phases of $`w(u_I)`$ and $`w(u_T)`$. The phase oposition alone does not explain however why the interference is totally destructive. It is also necessary to have equal moduli of the two incidence and transmission terms of (21) for an exact cancellation. Actually the equality is obtained only transitorily, since before and after the collision only one lobule remains, the one for incidence before the collision, and the one for transmission after the collision. Along the collision the incident lobule decreases and the transmission one grows, until they equilibrate and give a perfect cancellation and the two constructive interference zones of Fig. 1.
By changing the barrier height the fases of the factors that multiply the $`w`$s change, the lobules rotate with respect to each other, and one of the two in-phase regions grows while the other diminishes, so that the two peaks of the momentum distribution become asymetric, see Figures 5 and 6, where the momentum distributions and the corresponding lobules of the Argand diagrams are shown, compare also with Fig. 2. Note that these factors do not depend on time and therefore the angle between the lobules remains constant throughout the collision. This means that the positions of the maxima and minima do not change significantly for a given collision.
An important point is that the interference effect described does not depend critically on the square barrier potential, and we have observed it in particular for a Gaussian barrier. Note that the arguments leading to Eq. (21) are in fact of general validity and independent of the potential shape, with $`d/2`$ and $`d/2`$ being points where the potential may be assumed to be essentially zero, and $`T`$ being the corresponding transmission amplitude. The possibility to observe this effect with ultracold atoms rests on the ability to prepare appropriate initial states. Turning off the laser potential during the collision will leave a two peaked momentum distribution that implies at later times a visible spatial separation between two wave components, one faster than the other.
We thank A. Steinberg for many useful discussions, and acknowledge support by Ministerio de Educación y Ciencia (PB97-1482).
Figure captions
FIG. 1. $`\left|p|\psi (t)\right|^2`$ for differents values of $`t`$: $`t=0`$ (dotted-dashed line); $`t=2.333t_u`$ (solid line); $`t=2.731t_u`$ (dashed line); and $`t=3.233t_u`$ (dotted line). $`m=1.558023m_u`$, $`V_0=102.5e_u`$, $`d=2.5l_u`$, $`\alpha \delta _x=50l_u`$, $`\delta _x=107.99l_u^2`$, with an average momentum $`p_c=28.48p_u`$ well above the classical threshold $`(2mV_0)^{1/2}=17.87p_u`$. The units are scaled for numerical convenience in the computations as $`e_u=10^{13}`$ a.u. of energy, $`p_u=10^4`$ a.u. of momentum, $`l_u=2\times 10^6`$ a.u. of lenght, $`m_u=10^5`$ a.u. of mass, and $`t_u=2\times 10^{15}`$ a.u. of time.
FIG. 2. Imaginary versus real parts of the incident contribution to $`\psi _{IT}^0(p,t)`$ (empty circles), and of the transmision contribution (filled circles), for $`t=2.731t_u`$ and different values of $`p`$ equally spaced between $`p=28p_u`$ and $`p=29p_u`$. Other parameters as in Fig. 1.
FIG. 3. Imaginary versus real parts of $`\omega (u_I)`$ (empty circles) and $`\omega (u_T)`$ (filled circles) for $`t=2.731t_u`$ and different values of $`p`$ equally spaced between the two peaks of Fig. 1, see the text. The prefactors corresponding to $`\omega (u_I)`$ and $`\omega (u_T)`$ in $`\psi _{IT}^0(p,t)`$ are also shown for the same momentum interval, solid and dashed lines respectively. Other parameters as in Fig. 1.
FIG. 4. Integration contour in the complex $`p^{}`$-plane when the SDP crosses the pair of structural poles $`p_I^{}`$ and $`p_T^{}`$. Also shown the structural pole $`p_R^{}`$, the saddle point of equation (15), and the incident average momentum $`p_c`$.
FIG. 5. $`\left|p|\psi (t)\right|^2`$ as a function of $`p`$, for two different values of $`V_0`$: $`102.5e_u`$ (solid line), and $`105e_u`$ (dashed line). The value of $`t`$ is selected to get the maximum effect, $`G^\mathrm{q}0.24`$; $`t=2.731t_u`$. Other parameters as in Fig. 1.
FIG. 6. Imaginary versus real parts of the incident contribution to $`\psi _{IT}^0(p,t)`$ (empty circles), and of the transmision contribution (filled circles), for $`V_0=105e_u`$, the value of $`t`$ for which the effect is maximum ($`t=2.731t_u`$), and different values of $`p`$. Other parameters as in Fig. 1.
|
warning/0005/cond-mat0005146.html
|
ar5iv
|
text
|
# Franck-Condon-Broadened Angle-Resolved Photoemission Spectra Predicted in LaMnO3
## Abstract
The sudden photohole of least energy created in the photoemission process is a vibrationally excited state of a small polaron. Therefore the photoemission spectrum in LaMnO<sub>3</sub> is predicted to have multiple Franck-Condon vibrational sidebands. This generates an intrinsic line broadening $``$ 0.5 eV. The photoemission spectral function has two peaks whose central energies disperse with band width $``$ 1.2 eV. Signatures of these phenomena are predicted to appear in angle-resolved photoemission spectra.
The colossal magnetoresistance (CMR) effect in doped manganese oxides has attracted a lot of attention. The interplay of charge, orbital and magnetic order results in a very rich phase diagram . The parent compound LaMnO<sub>3</sub> has orthorhombic symmetry at low temperature. The Mn<sup>+3</sup> ion has $`d^4`$ ($`t_{2g}^3`$,$`e_g^1`$) configuration with an “inert” $`t_{2g}`$ core (spin 3/2) and a half-filled doubly degenerate $`e_g`$-type $`d`$ orbital which is Jahn-Teller (JT) unstable. Ignoring rotation of the MnO<sub>6</sub> octahedra, which occurs below 1010 K, the JT symmetry breaking is cubic to tetragonal at $`\mathrm{T}_{\mathrm{JT}}=`$750 K. The corresponding orbital order has $`x`$\- and $`y`$-oriented E<sub>g</sub> orbitals alternating in the $`xy`$ plane with wave vector $`\stackrel{}{Q}=(\pi ,\pi ,0)`$. This in turn causes layered antiferromagnetic (AFA) order to set in at T$`{}_{\mathrm{N}}{}^{}=`$140 K.
The electronic structure of LaMnO<sub>3</sub> has been studied, for example, by photoemission and by first principles calculations . Still there is controversy about the nature of the low energy excitations, arising from the interplay between strong on-site Coulomb repulsion (which leads to magnetic order) and strong electron-phonon (e-p) interactions (which lead to orbital order).
When an electron is removed from the JT-ordered ground state, e-p coupling causes the hole to self-localize in an “anti-JT” small polaron state. In a previous paper we have described the localized polaron in adiabatic approximation. Residual non-adiabatic coupling allows the hole to disperse with band width narrowed by Huang-Rhys factor $`e^{3\mathrm{\Delta }/4\mathrm{}\omega }10^4`$. The photoemission process is sudden. The emitted electron with wavevector $`\stackrel{}{k}`$ leaves a hole in a lattice “frozen” in the unrelaxed JT state. Ignoring lattice relaxation, this hole would disperse with band width $`2t`$1 eV ($`t`$ is the hopping parameter), as shown on Fig. 1. However, this is not a stationary state and must be regarded as a superposition of exponentially narrowed small polaron bands. Such bands have anti-JT oxygen distortions at each site, but a sufficient number of vibrational quanta are also excited such that the anti-JT distortion at time zero is “undone”. This is “Franck-Condon principle”.
The measured spectrum at wavevector $`\stackrel{}{k}`$ will consist of a central $`\delta `$-function at the energy of the frozen lattice (dispersive) band $`\epsilon _{1,2}(\stackrel{}{k})`$, plus multiple vibrational side-bands at energy $`\epsilon _{1,2}(\stackrel{}{k})\pm n\mathrm{}\omega `$, with an overall Gaussian envelope whose width is approximately the polaron binding energy. Franck-Condon broadening has been seen in photoemission spectra of solid nitrogen and oxygen .
We find at each wavevector $`\stackrel{}{k}`$ the photoemission spectrum has an intrinsic Franck-Condon broadening indicated by error bars in Fig. 1. The position of the maximum disperses with $`\stackrel{}{k}`$-vector close to the “frozen” lattice spectrum. A qualitative picture of this process has been given by Sawatzky in the context of high temperature superconductors and by Dessau and Shen for the manganites. The present paper gives an exact algebraic prediction for the Angle-Resolved-Photoemission-Spectra (ARPES) of a model Hamiltonian for LaMnO<sub>3</sub>.
Our model Hamiltonian , first introduced by Millis , has hopping $`_{\mathrm{el}}`$, electron-phonon $`_{\mathrm{ep}}`$, and lattice $`_\mathrm{L}`$ energies:
$`_{\mathrm{el}}`$ $`=`$ $`t{\displaystyle \underset{\mathrm{},\pm }{}}\left\{[c_x^{}(\mathrm{})c_x(\mathrm{}\pm \widehat{x})]+[xy]+[yz]\right\}`$ (1)
$`_{\mathrm{ep}}`$ $`=`$ $`g{\displaystyle \underset{\mathrm{},\alpha }{}}\widehat{n}_{\mathrm{},\alpha }(u_{\mathrm{},\alpha }u_{\mathrm{},\alpha })`$ (2)
$`_\mathrm{L}`$ $`=`$ $`{\displaystyle \underset{\mathrm{},\alpha }{}}(P_{\mathrm{},\alpha }^2/2M+Ku_{\mathrm{},\alpha }^2/2).`$ (3)
In these formulas $`c_\alpha ^{}(\mathrm{})`$ creates a state with orbital $`\psi _\alpha =|3\alpha ^2r^2>`$, where $`\alpha =x,y,z`$. These three orbitals span the two dimensional $`e_g`$ subspace and can be expressed in terms of the conventional orthogonal basis $`\mathrm{\Psi }_2=d_{x^2y^2}`$, $`\mathrm{\Psi }_3=d_{3z^2r^2}=\mathrm{\Psi }_z`$; specifically $`\mathrm{\Psi }_{x,y}=\pm \sqrt{3}/2\mathrm{\Psi }_2\mathrm{\Psi }_3/2`$. The resulting $`_{\mathrm{el}}`$ coincides with the nearest-neighbor two-center Slater-Koster hopping Hamiltonian with overlap integral $`t=(dd\sigma )`$ and $`(dd\delta )=0`$. The hopping parameter $`t=0.5`$ eV is chosen to agree with an ab initio $`e_g`$ band width of 1 eV . The e-p interaction $`_{\mathrm{ep}}`$ is modeled by a linear energy reduction of an occupied $`\psi _x`$ orbital ($`\widehat{n}_{\mathrm{},x}=c_x^{}(\mathrm{})c_x(\mathrm{})`$) if the corresponding two oxygens in the $`\pm \widehat{x}`$ direction expand outwards, and similarly for $`\widehat{y}`$ and $`\widehat{z}`$ oxygens if $`\psi _y`$ or $`\psi _z`$ orbitals are occupied. The strength of the e-p coupling $`g`$ determines the JT splitting $`2\mathrm{\Delta }=1.9`$ eV, which is fitted to agree with the lowest optical conductivity peak . Static oxygen distortions $`2u=\sqrt{2\mathrm{\Delta }/M\omega ^2}`$=0.296 Å given by our model agree well with neutron diffraction data 0.271 Å . For the lattice term $`_\mathrm{L}`$ we use a simplified model where oxygen vibrations along Mn-O-Mn bonds are local Einstein oscillators. The displacement $`u_{\mathrm{},\alpha }`$ is measured from cubic perovskite position of the nearest oxygen in the $`\widehat{\alpha }`$-direction to the Mn atom at $`\mathrm{}`$. The oxygen vibrational energy $`\mathrm{}\omega =\mathrm{}\sqrt{K/M}=0.075`$ eV is taken from Raman data . In addition there is a large on-site Coulomb repulsion U and a large Hund energy. These terms inhibit hopping except to empty sites where the t<sub>2g</sub> core spins are aligned correctly.
In adiabatic approximation one can solve this problem for U=0 or U=$`\mathrm{}`$. Both cases give a good description of the observed cooperative JT order. When U=0, the ground state wavefunction is:
$$|\mathrm{GS},0>=\underset{\stackrel{}{k}}{}c_{\stackrel{}{k}1}^{}c_{\stackrel{}{k}2}^{}|\mathrm{vac}>.$$
(4)
A JT gap $`2\mathrm{\Delta }`$ opens and the lower two bands of energy $`\lambda _{\stackrel{}{k}1}`$, $`\lambda _{\stackrel{}{k}2}`$ are filled. The photohole as initially created has energy:
$`\lambda _{1,2}^2=\mathrm{\Delta }^2+t^2(2C_x^2+C_xC_y+2C_y^2)\pm t|C_x+C_y|`$ (5)
$`\sqrt{\mathrm{\Delta }^2+4t^2(C_x^2C_xC_y+C_y^2)},`$ (6)
where $`C_{x,y}=\mathrm{cos}k_{x,y}`$ and $`C_z`$ not entering at T=0 K. These bands are shown in Fig. 1(a) as dashed lines.
At this point, at least in principle, could proceed numerically to find the polaronic energy lowering. However, it is both easier and more realistic to switch to U=$`\mathrm{}`$. The ground state wavefunction:
$$|\mathrm{GS},\mathrm{}>=\underset{\mathrm{}}{\overset{A}{}}c_X^{}(\mathrm{})\underset{\mathrm{}^{}}{\overset{B}{}}c_Y^{}(\mathrm{}^{})|\mathrm{vac}>,$$
(7)
has orbitals $`\mathrm{\Psi }_{X,Y}=(\mathrm{\Psi }_3\mathrm{\Psi }_2)/\sqrt{2}`$ occupied singly on interpenetrating $`A`$ and $`B`$ sublattices. This is a fully correlated state with zero double occupancy, while Eq. (4) is a band wavefunction in which two electrons are found on the same Mn atom with non-zero probability. For U$``$ 6t $`3`$ eV, the state (7) has lower energy than state (4) . Neglecting creation of orbital defects with energy $`2\mathrm{\Delta }`$ the “frozen” lattice approximation predicts photohole energies $`\mathrm{\Delta }+\epsilon _{1,2}(\stackrel{}{k})`$, where $`\epsilon _{1,2}(\stackrel{}{k})=\pm (t/2)(C_x+C_y)`$. We need to add a non-adiabatic treatment of the e-p coupling. The effective Hamiltonian $`_{\mathrm{eff}}=_{\mathrm{el}}+_{\mathrm{ep}}+_\mathrm{L}`$ for the single hole is:
$`_{\mathrm{el}}^\mathrm{A}`$ $`=`$ $`{\displaystyle \underset{\mathrm{}A}{}}{\displaystyle \frac{t}{4}}\left(d_Y^{}(\mathrm{}\pm x)d_X(\mathrm{})+d_Y^{}(\mathrm{}\pm y)d_X(\mathrm{})\right)`$ (8)
$`_{\mathrm{ep}}^\mathrm{A}+_\mathrm{L}^\mathrm{A}`$ $`=`$ $`{\displaystyle \underset{\mathrm{}A}{}}d_X^{}(\mathrm{})d_X(\mathrm{})[\mathrm{\Delta }+{\displaystyle \underset{\alpha }{}}\kappa _\alpha (a_\alpha (\mathrm{})+`$ (9)
$`a_\alpha ^{}(\mathrm{})b_\alpha (`$ $`\mathrm{}`$ $`\alpha )b_\alpha ^{}(\mathrm{}\alpha ))]+{\displaystyle }_\alpha a_\alpha ^{}(\mathrm{})a_\alpha (\mathrm{}).`$ (10)
Here the operator $`d_X^{}(\mathrm{})=c_X(\mathrm{})`$ creates a hole in the JT ground state by destroying an electron on orbital $`X`$ at site $`\mathrm{}`$ (if $`\mathrm{}A`$ sublattice), and the operator $`d_Y^{}(\mathrm{})=c_Y(\mathrm{})`$ creates a hole on $`B`$ sublattice (if $`\mathrm{}B`$). The phonon operators $`a_\alpha ^{}(\mathrm{})`$ or $`b_\alpha ^{}(\mathrm{})`$ create vibrational quanta on the $`\mathrm{}+\widehat{\alpha }/2`$ oxygen, if $`\mathrm{}A`$ or $`\mathrm{}B`$ respectively. The e-p coupling constants are $`\kappa _{x,y,z}=\sqrt{\mathrm{\Delta }/12}(1+\sqrt{3}/2;1\sqrt{3}/2;1)`$. The Hamiltonian and all other energy parameters $`\mathrm{\Delta }`$ and $`t`$ in Eq. 10 are in units of $`\mathrm{}\omega `$. The total Hamiltonian has an additional term $`_{\mathrm{el}}^\mathrm{B}+_{\mathrm{ep}}^\mathrm{B}+_\mathrm{L}^\mathrm{B}`$ which is obtained from Eq. (10) by interchanging operators:
$$d_Yd_X,a_xb_y,a_yb_x,a_zb_z$$
(11)
and summing over the $`B`$ sublattice.
Following Cho and Toyozawa we are able to diagonalize Hamiltonian (10) in a very large truncated basis of functions with a hole present on site $`\mathrm{}`$ and an arbitrary number of vibrational quanta $`\mathrm{p}_{\pm x},\mathrm{p}_{\pm y},\mathrm{p}_{\pm z}`$ on the six displaced neighboring oxygens:
$`|\mathrm{\Psi }^\mathrm{A}(\mathrm{},\{p\})>=d_X^{}(\mathrm{}){\displaystyle \underset{\alpha }{}}\mathrm{U}_{\mathrm{}}^{a_\alpha }(\kappa _\alpha ){\displaystyle \frac{(a_\alpha ^{}(\mathrm{}))^{p_{+\alpha }}}{\sqrt{p_{+\alpha }!}}}`$ (12)
$`\mathrm{U}_\mathrm{}\alpha ^{b_\alpha }(\kappa _\alpha ){\displaystyle \frac{(b_\alpha ^{}(\mathrm{}\alpha ))^{p_\alpha }}{\sqrt{p_\alpha !}}}|\mathrm{GS},\mathrm{}>.`$ (13)
The displacement operator $`\mathrm{U}_{\mathrm{}}^a(\kappa )=\mathrm{exp}[\kappa (a_{\mathrm{}}a_{\mathrm{}}^{})]`$ makes the $`_{\mathrm{ep}}+_\mathrm{L}`$ part of the Hamiltonian diagonal. To get basis functions $`|\mathrm{\Psi }^\mathrm{B}(\mathrm{},\{p\})>`$ for holes on the $`B`$ sublattice, the operators in Eq. (13) should be interchanged according to Eq. (11). The next step is to build Bloch wavefunctions by Fourier transformation of the basis functions Eq. (13). Then the Hamiltonian (10) will be diagonal with respect to $`\stackrel{}{k}`$-vector. The hopping term of the Hamiltonian $`_{\mathrm{el}}`$ couples the $`|\mathrm{\Psi }^\mathrm{A}(\mathrm{},\{p\})>`$ and $`|\mathrm{\Psi }^\mathrm{B}(\mathrm{}^{},\{p\})>`$-wavefunctions on the neighboring sites. The vibrational wavefunctions give a product of Huang-Rhys factors, but a shared oxygen contributes a non-factorizable overlap integral. However if one treats this shared oxygen as two independent atoms, one coupled to each site, then the Hamiltonian has a simple form:
$`_{pp^{}}^{\mathrm{AA}}(\stackrel{}{k})`$ $`=`$ $`_{pp^{}}^{\mathrm{BB}}(\stackrel{}{k})=\delta _{\{p\}\{p^{}\}}\left[{\displaystyle \frac{\mathrm{\Delta }}{4}}+{\displaystyle \underset{\alpha }{}}(p_{+\alpha }+p_\alpha )\right]`$ (14)
$`_{pp^{}}^{\mathrm{AB}}(\stackrel{}{k})`$ $`=`$ $`_{pp^{}}^{\mathrm{BA}}(\stackrel{}{k})=\epsilon (\stackrel{}{k}){\displaystyle \underset{\alpha }{}}(1)^{p_\alpha +p_\alpha ^{}}`$ (15)
$`[P(p_{+\alpha }`$ $`,\kappa _\alpha `$ $`)P(p_\alpha ,\kappa _\alpha )P(p_{+\alpha }^{},\kappa _\alpha )P(p_\alpha ^{},\kappa _\alpha )]^{1/2},`$ (16)
where $`P(p,\kappa )=\mathrm{exp}(\kappa ^2)\kappa ^{2p}/p!`$ is a Poisson distribution. Since off-diagonal terms factorize, the analytical solution is available in this approximation:
$`\mathrm{\Psi }_\lambda ^{1,2}(\stackrel{}{k})={\displaystyle \underset{\{p\}=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{\alpha }{}}(1)^{p_\alpha }\left[P(p_{+\alpha },\kappa _\alpha )P(p_\alpha ,\kappa _\alpha )\right]^{1/2}`$ (17)
$`{\displaystyle \frac{\mathrm{\Psi }_\stackrel{}{k}^\mathrm{A}(\{p\})\pm \mathrm{\Psi }_\stackrel{}{k}^\mathrm{B}(\{p\})}{\sqrt{2G^{}(x_\lambda )}\left(_\alpha ^{}(p_{+\alpha ^{}}+p_\alpha ^{})x_\lambda \right)}}.`$ (18)
The corresponding eigenvalues are:
$`E_\lambda ^{1,2}(\stackrel{}{k})={\displaystyle \frac{\mathrm{\Delta }}{4}}+x_\lambda ^{1,2}(\stackrel{}{k}),1+\epsilon _{1,2}(\stackrel{}{k})G(x_\lambda ^{1,2})=0`$ (19)
$`G(x_\lambda )=e^{3\mathrm{\Delta }/4}{\displaystyle \underset{p=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(3\mathrm{\Delta }/4)^p}{p!}}{\displaystyle \frac{1}{px_\lambda }}.`$ (20)
The $`G^{}(x_\lambda )`$ function in Eq. (18) is a first derivative of $`G(x_\lambda )`$ and makes wavefunctions normalized. The correct solution needs a numerical diagonalization of the Hamiltonian which explicitly includes vibrational states of the four shared oxygens. As can be seen on Fig. 2 the difference is negligible between a typical spectrum obtained using approximation (20) and correct numerical treatment. The ground state of the Hamiltonian (16), with energy (from Eq. 20) $`E_0(\stackrel{}{k})=\mathrm{\Delta }/4+x_0(\stackrel{}{k})`$, corresponds to the anti-JT polaron. Its effective mass, deduced from $`d^2x_0(\stackrel{}{k})/d\stackrel{}{k}^2`$, provides a realistic alternative (exact for $`\mathrm{\Delta }0`$ or $`\mathrm{}`$) to the available variational approaches or exact quantum Monte Carlo simulations .
An ARPES experiment measures the spectral function $`A(\stackrel{}{k},\omega )=\frac{1}{\pi }G(\stackrel{}{k},\omega )`$ with momentum $`\stackrel{}{k}`$ fully resolved, provided there is no dispersion in the direction perpendicular to the surface. Although LaMnO<sub>3</sub> is cubic, because of the layered AFA magnetic structure, at low temperatures Mn $`e_g`$ electrons are two-dimensional and the spectrum can be measured:
$$A(\stackrel{}{k},\omega )=\underset{f}{}|<f|d^{}(\stackrel{}{k})|\mathrm{GS},\mathrm{}>|^2\delta (EE_f).$$
(21)
The operator $`d^{}(\stackrel{}{k})`$ excites a hole from the JT ground state. Summation over final eigenstates $`|f>`$ includes summation over branch index $`i=1,2`$ and number of phonons $`\lambda =0,1\mathrm{}\mathrm{}`$. The sequence of delta functions in Eq. (21) should be replaced by convolved local densities of phonon states, which we approximate by a Gaussian, $`\delta (E)\mathrm{exp}(E^2/2\gamma ^2)/\sqrt{2\pi }\gamma `$. Substituting solution (18), (20) into equation (21), we obtain the spectral function:
$$A(\stackrel{}{k},\omega )=\underset{\lambda ,i}{}\frac{G^2(x_\lambda ^i)}{G^{}(x_\lambda )}\delta (EE_\lambda ^i).$$
(22)
Equation (22) along with (20) gives the ARPES spectrum normalized to $`𝑑\omega A(\stackrel{}{k},\omega )=1`$. The first energy moment of the spectrum coincides with the free hole energy calculated in the “frozen” lattice approximation $`\mathrm{\Delta }+\epsilon _{1,2}(\stackrel{}{k})`$ shown on Fig. 2. The edge of the spectrum corresponds to polaron creation at energy $`\mathrm{\Delta }/4`$. This transition is weaker by 3 orders of magnitude than the peak at $`\mathrm{\Delta }+\epsilon _{1,2}(\stackrel{}{k})`$.
At room temperature magnetic order is lost. The paramagnetic state is modeled by a mean field approximation, namely scaling the effective hopping integral by 2/3 and allowing hopping in $`\pm \widehat{z}`$ direction. This modifies the single particle energy band entering Eq. (16) to $`\epsilon _{1,2}(\stackrel{}{k})=t/3(2C_z\pm (C_x+C_y))`$. But the JT orbital order is not destroyed at T=300 K and Franck-Condon broadening is still expected. When spins are disordered, $`\stackrel{}{k}`$ is not a good quantum number and additional broadening is expected. Only phonon broadening of the ARPES along with peak positions are shown on Fig. 1(b).
The angle-integrated spectrum, shown on Fig. 3 for low and high temperatures, has a width of about 1.2 eV and is almost temperature independent. The uncorrelated (U=0) band structure, shown for comparison, is sensitive to magnetic order and therefore temperature dependent.
The existing photoemission data are consistent with our predictions. Higher resolution experiments are needed to test the theory and to unravel the nature of the lowest energy excitations in the LaMnO<sub>3</sub>. To make such an experiment possible, a single domain sample (having 2D dispersion at T=0K) is needed, with good control of oxygen concentration .
###### Acknowledgements.
We thank P. D. Johnson for helpful conversations. This work was supported in part by NSF Grant No. DMR-9725037.
|
warning/0005/astro-ph0005488.html
|
ar5iv
|
text
|
# LOPSIDEDNESS IN DWARF IRREGULAR GALAXIES
## 1 Introduction
One of the least understood aspects of galaxy evolution is the onset of lopsidedness in the gaseous and stellar distributions of disk galaxies. Recent models of disk galaxies suggest that the presence of a dominant dark halo can both produce and help sustain asymmetries in the gaseous and stellar components. For instance, a lopsided gravitational potential of a dark matter halo can produce an asymmetric galaxy as the gas surface density responds to the overall asymmetry (Jog 1997). Alternatively, a symmetric dark matter halo can produce an asymmetric galaxy if the disk orbits off–center of the overall potential (Levine & Sparke 1998). In addition to these “intrinsic” models, the environment and merger history of a galaxy can affect its present appearance. For example, recent dynamical simulations of the effect of an infalling satellite indicate that tidal interactions are yet another mechanism by which asymmetric galaxies can be formed (Walker, Mihos, & Hernquist 1996; Zaritsky & Rix 1997).
However, all the above models were designed to account for asymmetries observed in massive spiral galaxies. The question of asymmetries in dwarf irregular galaxies (dIs) may demand a different approach, as dIs lack spiral density waves and tidal shear forces that contribute to induce gas instabilities. In dIs, the gas appears to be close to stability throughout the disk, even though star formation is occurring (Hunter et al. 1998). Moreover, there is growing evidence (Mihos, McGaugh & de Blok 1997) that low-surface-brightness (LSB) disks are reasonably stable and remain structurally intact during tidal encounters. In addition, various tests of Virgo dIs favor internal over external mechanisms of star formation (Heller et al. 1998), with the implicit conclusion that asymmetries may also form through internal mechanisms.
We conclude that, while theoretical arguments, such as the presence of a dominant dark matter halo potential in dIs, may contribute to the long-persistence of the asymmetries, other processes of star formation, such as the Stochastic Self-Propagating Star Formation (SSPSF, Gerola, Seiden & Schulman 1980), or alternatively random gas compression from turbulence, or random collisions of ISM clouds (Larson 1986; Elmegreen 1998) may be the dominant regulators of the star formation in dIs. This conclusion is supported by recent star formation histories (SFH) of Sextans A and GR 8 derived from HST observations. A series of chronological frames showing the spatial distribution of blue HeB stars indicate that, chronologically, the star formation activity is propagating around in these galaxies with typical sizes of $``$100 pc and lifetimes of order 100 Myr (Dohm-Palmer et al. 1997, 1998a, b). The question remains as to whether random mechanisms may introduce temporary asymmetries in the stellar and gaseous components of low-mass systems. We address this issue by (a) analyzing the light distribution in deep narrow-band H$`\alpha `$ and continuum images of a large sample of star forming dIs, (b) developing a new impartial algorithm to compute the lopsidedness of star-forming regions and simultaneously compare it to the distribution of the stellar component, (c) constructing 1000 model galaxies and showing that it is the “discrete” behavior of random star forming regions that produce the asymmetric structure observed in most of the dIs.
The plan of the paper is as follows: we first describe the sample of galaxies, which is a collection of objects with previously published observations. The analysis method is described next, then the results are presented. Finally, we describe the simulations performed to understand the observational results and their implications.
## 2 The sample
The galaxy sample studied here consists of 78 dIs observed with the Wise Observatory (WO) 1.0 m telescope or with the Kitt Peak National Observatory (KPNO) 0.9 m telescope. The galaxies are classified in the original publications as dIs, with absolute blue magnitudes $`<`$ -18, and are smaller than 2 arcmin. The only restrictions to the inclusion in the sample are the availability of CCD H$`\alpha `$ images with detected HII regions, and $`v_{}3,000`$ km sec<sup>-1</sup> (except for UM408, which has $`v_{}=3,492`$ km sec<sup>-1</sup>). All the selected galaxies appear to be isolated, with the exception of objects marked with an asterisk in Table 1. Representative H$`\alpha `$ images, catalog references, and extensive additional details, can be found in van Zee et al. (1997a, b, c), Almoznino & Brosch (1998), Heller et al. (1999, 2000), and Norton & Salzer (2000).
In order to test for dependence of the SFR on the lopsidedness we divided the sample in two sub-groups. The first, called here BCD, is represented by 33 blue compact dwarf galaxies (classified morphologically as BCD or anything+BCD: references 5, 6, and 7 in Table 1). These are galaxies whose optical light output is often dominated by the strong starburst component. The second group, called LSB, is represented by 45 low surface brightness dwarf galaxies (references 1, 2, 3, and 4). This group includes dIs, primarily from standard catalogues, which are gas-rich and, in general, have central surface brightness fainter than $`23.0`$ mag arcsec<sup>-2</sup>. Some of the more luminous LSB galaxies show evidence of spiral features and may belong to the “dwarf spiral” class (UGC numbers 191, 634, 3050, 3174, 4660, 5716, 7178, 9762, 10281, and 11820). The typical SFR for the LSB group is $`7\times 10^3`$ M yr<sup>-1</sup>; this is, on average, one order of magnitude weaker than for BCD objects, although there is overlap in SFR between the brighter LSBs and the fainter BCDs.
## 3 Analysis and results
### 3.1 Method
In general, the lopsidedness of a galaxy is measured on a broad-band image (usually in the red). Some authors (Zaritsky & Rix 1997; Rudnick & Rix 1998) use the ratio of the m=1 to m=0 Fourier amplitudes of the image as a quantitative measure of lopsidedness in early-type disk galaxies. Others take a more direct approach of comparing the integrated light within specified regions of the galaxy. For instance, Kornreich et al. (1998) compare the relative fluxes within trapezoidal sectors arranged symmetrically about the galaxy’s center of light. Similarly, Abraham et al. (1996) define the rotational asymmetry parameter as half the ratio of the absolute value of the difference between the original galaxy image and the image rotated by a half-turn about its center, to the original image. The rotational asymmetry parameter, together with the central concentration of the emitted flux, has proven to be an important tool to extend the morphological classification of the galaxies from the nearby Universe to high redshifts. Likewise, color-asymmetry diagrams, when combined with information about the axial ratio, can be used to disentangle interacting galaxies from non-interacting, face-on systems at high redshift (Conselice 1997; Conselice & Bershady 1999).
We have developed a new method to evaluate the variation of asymmetry with azimuthal position angle, and also the concentration of the star forming regions and of the general stellar distribution. In our method, the H$`\alpha `$ line emission represents the distribution of the recently formed massive stars that are younger than a few tens of Myrs, while the red continuum emission represents the distribution of the integrated stellar populations. This latter component may be contaminated by nebular continuum emission from HII regions, but nebular emission amounts to only 30$`\%`$ of the light during the first few Myrs of a starburst, and becomes negligible as soon as the first red supergiants appear, at about 10<sup>6.9</sup> yr (Leitherer & Heckman 1995). We investigate both the asymmetry and concentration properties of these two components, as well as look for correlations between them.
For most galaxies in the sample, two red, narrow-band images were used: one centered on the rest-frame H$`\alpha `$ line ($`H\alpha _{on}`$) and the other sampling the continuum ($`Cont`$) region near H$`\alpha `$. For some of the LSBs, the narrow-band continuum images were no longer available at the time of the present analysis. For these systems (marked “+” in Table 1), sky-subtracted B broad-band images were only used to trace the ellipse contour of the galaxies (see below) but are not included in the statistical results for $`Cont`$. In all cases, the sky background was subtracted in each band and the net-H$`\alpha `$ ($`H\alpha `$) images were derived by subtracting $`Cont`$ from $`H\alpha _{on}`$ images with proper scaling. Details are given in Heller et al. (1999).
Due to the lack of obvious central concentration and the irregular shape of the galaxies, we performed different fits of elliptical isophotes, allowing the ellipticity and the position angle to vary, and fitting out to 25 mag arcsec<sup>-2</sup> on the continuum image. The convergence criterion for the final parameters was set when the outer isophote retained the position angle and ellipticity of the ellipse traced at half the major axis. From then on, the position angle, the ellipticity, and the extent of the galaxy were held fixed, and the outer ellipse contour was transposed to the H$`\alpha `$ images. For those objects without calibrated images, (refs. 2, 5, and 6 in Table 1), the outermost isophote was adopted at the level where the mean intensity reached the sky fluctuations. This choice presumably depends on the depth of the exposure and on systematic errors in the subtraction of the sky background. The reduction was done with IRAF<sup>1</sup><sup>1</sup>1IRAF is distributed by the National Optical Astronomy Observatories. and the fitted ellipse parameters are listed in Table 1.
We integrated the fluxes in the two halves of the galaxy separated by a bisector line, represented by the major-axis of the outer ellipse, and computed the ratio of the lower flux to the higher flux from the two galaxy halves. The resulting ratio defines one asymmetry index (AI<sub>i</sub>). After this, using the maximum allowed 90 vertexes of the ellipse contour produced from the ELLIPSE task of IRAF in the plane of the galaxy, we rotated the bisector anti-clockwise around the center of the ellipse to the line defined by the next two opposite vertexes, and thus obtained 90 asymmetry indices, one for each position angle $`\mathrm{\Phi }_i`$ of the vertex. For the maximum ellipticity (e=0.75) measured in the sample, we reach an upper spatial resolution of $`\mathrm{\Phi }_i`$= 0.4 degrees, and a lower resolution of $`\mathrm{\Phi }_i`$=4.4 degrees. For example, in a perfect circle (ellipticity e=0) the resolution is $`\mathrm{\Phi }_i`$ = 4 degrees. This method is more useful than the usual one constructed from two asymmetry indices because it covers the full range of possibilities in azimuthal angle. Moreover, the presence of faint HII regions is emphasized by the irregularities in the luminosity profiles. The variation of AI with position angle (the ”lopsidedness distribution”, LD) is plotted for each galaxy in the sample in Figures 1.1 to 2.3.
A representative lopsidedness index (A) for each galaxy was computed by normalizing the total lopsidedness range (the difference between maximum and minimum AI) to the maximum asymmetry index: $`A=\frac{AI_{max}AI_{min}}{AI_{max}}`$. The mean asymmetry index $`<`$AI$`>`$, the lopsidedness index (A), and the asymmetry amplitude $`ampl=\frac{AI_{max}AI_{min}}{2}`$ are listed in Tables 2.1, 2.2, and 2.3. A symmetric distribution of the light is represented by A=0 and AI<sub>max</sub>=1, while an extremely asymmetric distribution will have A=1 and AI<sub>min</sub>=0.
However, galaxies in general may have bulge and disk components and, therefore, a large range in scale lengths. In order to enhance the light distribution analysis we utilize a second parameter: the concentration index (CI). We calculated the CI index as the ratio of the flux from the inner part of the galaxy to $`1/3`$ of the flux from its outer annulus. The one-third factor brings the comparison to an equal-area basis, and makes it independent of the distance of the galaxy. The outer aperture was defined as the ellipse fitted to calculate the LD. The inner aperture was chosen as a smaller ellipse, half the size of the outer one. The annulus is the space between the inner and the outer apertures. As defined, CI can range between zero and infinity.
The two structural indices are similar to those used in Brosch et al. (1998), with the differences being: (a) the use of the $`H\alpha `$ flux instead of the number counts of HII regions, (b) the application of an objective automatic algorithm instead of eyeball recognition, (c) the derivation of the full LD instead of only indices, and (d) the use of a normalization factor of $`1/3`$ instead of $`1/4`$ for CI. Another difference is that here we calculate the CI and AI indices for the continuum, as well as for the net line emission.
### 3.2 Results
In Fig. 3 we plot the structural indices vs. the number of HII regions and ellipticity of the ellipse contour. The number of HII regions in BCDs ranges from one to three, and for LSBs from one to twelve. These numbers represent the number of resolved peaks detected in the LDs by an automatic algorithm that searches for slope changes in the LDs. The main limitation of the algorithm is the lack of resolution in special cases of multiple HII regions perfectly aligned in the radial direction. Since the ‘clumpiness’ of a galaxy depends on the seeing, the resolution at which the image is sampled, and on the resolution of the LDs, we cannot derive the number of HII regions in real galaxies as an absolute parameter; for example, nearby systems will appear clumpier than more distant ones. We show below that a change in the resolution, or a difference in the number of HII regions between BCDs and LSBs, cannot explain the differences in concentration indices between the types. Due to the intrinsic irregular shape of these galaxies, the ellipticity (e=1 - b/a) is also uncertain, but it does provide some measure of the inclination. We can see in Fig. 3b and Fig. 3d that the derived quantities are not simply the result of projection effects or affected by extinction through the disk.
The profile of the lopsidedness distribution appears to be related to the central surface brightness of a galaxy. A characteristic feature of the low surface brightness (LSB) sub-group is the multi-component structure of the LD, with sharp features shown in the $`H\alpha `$ profiles, while the continuum LD is smoother and with shallower features (Figs. 2.1 through 2.3), but not fully symmetric. The multiplicity of the AI profiles indicates that a number of individual HII regions with different luminosities and sizes are distributed over the galaxy; some may even not be resolved or recognized in our images but their contribution to the local H$`\alpha `$ flux is counted by the algorithm.
The interpretation of the LD profile widths depends not only on the sizes of the HII regions but also on their radial location; single HII regions closer to the center produce a wide peak, while those further out show narrow peaks. At the same radial distance, the bigger the size of the HII region, the wider the LD profile will appear. A nuclear HII region of some extent will present a flat profile. In the BCD sub-sample, many of the LDs show profiles that are mostly smooth, free of multi component structure, generally symmetric and wide (Figs. 1.1 through 1.3), as expected for single HII regions located near the centers. BCDs tend to be more concentrated than the LSBs; the median CI is 8.56 for BCDs and 2.25 for the LSBs. We found a strong correlation between log(CI<sub>Cont</sub>) and log(CI) (Fig. 4a) with the correlation coefficient cc=0.61 (F=36)<sup>2</sup><sup>2</sup>2F is the ratio between the mean square deviation due to the regression and the mean square deviation due to the residual variation. For a linear regression, which is the present situation, F=t<sup>2</sup> and this is the equivalent of a t-test. For more details see Draper & Smith (1981).. Linear regression tests between other data sets are listed in Table 3. At the same CI, both sub-samples reach similar degrees of asymmetry. Note that both galaxy types tend to clump at log(CI<sub>Cont</sub>)=0.5$`\pm `$0.2 and A<sub>Cont</sub>=0.2$`\pm `$0.1 (Fig. 4c). The entire sample has a median A of 0.69; the median for BCDs is 0.71 and for LSBs is 0.69.
We find an apparent upper limit for the asymmetry of the continuum light; 97$`\%`$ of the galaxies have A$`{}_{Cont}{}^{}`$ 0.5 (Fig. 4c). We also find an apparent lower limit of the emission line asymmetry; 97$`\%`$ of the galaxies have A$`{}_{H\alpha }{}^{}`$ 0.3 (Fig. 4d). A perusal of Tables 2.1, 2.2, and 2.3 shows that the $`Cont`$ asymmetry is always smaller than the H$`\alpha `$ asymmetry. A $`\chi ^2`$-test of the cumulative histograms of A and A<sub>Cont</sub> indicates that the two data sets originate from different distributions ($`\chi ^2`$=244, with 18 degrees of freedom). The median $`Cont`$ asymmetry is 0.25 for BCDs, 0.21 for LSBs, and 0.23 for all the objects with narrow-band images for the $`Cont`$. Note that objects with blue images for the continuum were not included in this analysis. Median values of the structure parameters are listed in Table 4.
A fundamental issue is whether the continuum and line-emission LDs are correlated in angular phase. We should expect a correlation if the locus of recent star formation, as witnessed by the H$`\alpha `$ emission, responds with a delay to some disturbance of the stellar distribution (the continuum light), producing a lag in the angular distribution of the azimuthal indices. This can be understood in a scenario of rotating disk-like systems. In fact, HI synthesis maps of a number these galaxies (Skillman et al. 1987; van Zee et al. 1997c, 1998a, b) show rotation-dominated systems with maximal rotation velocities of 40 – 100 km sec<sup>-1</sup> and with slowly rising rotation curves, typical of very late-type spirals (some appear to be undergoing small differential rotation), or systems with velocity still increasing beyond the optical disk, characteristic of the solid-body rotation found in many low-mass systems. In those cases, the angular phase correlation may depend on many factors, such as differences in the angular momentum of the stellar and gas masses, rotational speed, disk shear, and external SF triggers. Note that GR 8, Leo A and DDO 210 do not have well defined rotation curves (Carignan et al. 1990; Young & Lo 1996; Young, van Zee, & Lo 2000). The mismatch of the velocity gradient and the HI major axis in Leo A hints that some very low mass systems may be tumbling rather than spinning. In such cases, a delay between the past and the recent onset of star formation should also be expected.
To measure the correlation in the phase space of the line azimuthal asymmetry distribution with the distribution of the continuum light we used a similar analysis to that applied to the study of AGNs variability in the time-frequency domain (Netzer et al. 1996; Kaspi et al. 2000). The technique is the derivation of the cross-correlation function (CCF), which is a set of correlation coefficients, giving a measure of the correlation between two data sets. We used here two methods: the first one is the discrete correlation function (DCF; Edelson & Krolik 1988), which we applied after interpolating 45 continous data points every $`4^{}`$. The second method is the Z-transformed discrete correlation function (ZDCF) of Alexander (1997), which is an improvement of the DCF. For unevenly-sampled sets of data the ZDCF has the advantage that it avoids interpolation and reduces the resulting uncertaintly in the position of the peak. This is a consequence of the Fisher Z-transformation to the correlation coefficients and of the binning by equal population, rather than by equal separation.
The typical errors in the lags rage between $`0.52\mathrm{°}`$, with the exception of $`40\mathrm{°}`$ in IIZW40. The two methods (DCF and ZDCF) gave consistent results for our data, and we will refer to the ZDCF results in the following analyses. The uncertainties in the cross-correlation lags were conservatively over-estimated by the Monte Carlo-averaged ZDCF with simulated random errors and were provided by the ZDCF procedure of Alexander (1997). The results of the cross-correlation analysis are presented in Tables 2.1 and 2.2. The columns labeled “r<sub>zdcf</sub>” show the peaks of the CCFs, defined as the point of maximum correlation; a high value of “r<sub>zdcf</sub>” implies a good correlation between the two azimuthal indices at the listed “lag”. For the definition of “r<sub>zdcf</sub>” see Alexander (1997). The sign of the lag is defined as AI<sub>Cont</sub> \- AI, that is AI<sub>Cont</sub> lags after AI.
The cross-correlation (CC) analysis of the H$`\alpha `$ vs. continuum distribution of AIs indicates a very high CC for a broad range of angular phase lags. The CC is higher for BCDs than for LSBs, and there is a trend for smaller angular phase lags in BCDs than in LSBs (Figs. 5a and 5b). In fact, $`62\%`$ of the BCDs having peak CC coefficient above 0.8 show lags smaller than $`|\mathrm{\Delta }\mathrm{\Phi }|<30\mathrm{°}`$ compared to $`33\%`$ of LSBs. The higher CC is explained by increased CIs (Fig. 5c), however the distribution of lags seems to be independent of the concentration parameter (Fig. 5d).
## 4 A random distribution of star formation regions?
In this section we explore the possibility that the properties found for star-forming regions in dIs can be produced by random processes that engulfs the full scale of a galaxy. We tested a model of random star formation by constructing 1000 images of galaxies, which simulate the observed net H$`\alpha `$-flux and off-band red emission of dIs as found above, without distinguishing between LSB and BCD types. The model was created with the ARTDATA package in IRAF and included atmospheric seeing effects and detector readout noise.
A galaxy was modeled as a disk centered on a 256$`\times `$256 pixel image with zero background. The intensity profile was that of an exponential disk $`I=I_{}exp(1.6783R/R_{})`$, with the scale radius R containing half the total flux. The apparent integrated magnitudes, scale-lengths, position angles of the major axis, and ellipticities were allowed to change randomly. The total magnitudes followed a Schechter (1976) luminosity function with $`\alpha `$=1.6 and $`M_{}=21.41`$ in the red continuum, covering the apparent magnitude range from 17 to 19, similar to that of the objects in our sample. The maximum semi-major axis at half-flux was set to 30 pixels. The ellipticity was allowed to vary between 0.05 and 1.00. Random noise was added to the image by using Poisson statistics; a similar process was followed for the net and continuum images described below. At this stage, the output parameters of the disk were recorded as an ellipse contour with a semi-major axis twice the derived scale length. That is, CI$`{}_{Cont}{}^{}=3`$ and A$`{}_{Cont}{}^{}=0`$, by the definition of the underlying exponential disk.
The $`H\alpha `$ emission image was created by random generation of coordinates of up to 15 objects within the ellipse derived from the disk on a mean zero background. This range (1-15) covers the number of resolved peaks detected in the LDs with the algorithm, as explained before. We will show that changing the total resolution, or the ratio between maximum and minimum resolution, cannot explain the differences in concentration indices between BCDs and LSBs. The objects simulate HII regions, whose apparent magnitudes were allowed to change randomly between 18 and 23 following a shallow power law with index 0.1. This range of magnitudes reproduces the $`H\alpha `$ flux densities observed for the HII regions of our sample of galaxies and yields total line fluxes in the range 10<sup>-15</sup> \- 10<sup>-13</sup> erg cm<sup>-2</sup> s<sup>-1</sup>. We assumed a simplified profile for an individual HII cloud as a spherical distribution with a star-like Moffat profile ($`\beta `$ =2.5). A Moffat profile (Moffat 1969) appears more natural than a Gaussian, because it produces a sharper boundary to an HII region, as expected for a Strömgren sphere. However, this choice does not appear to affect the results described below.
The red continuum image was simulated by adding to the smooth underlying exponential disk the same list of objects coordinates and flux densities that represented the HII regions, but this time simulating red star clusters (or super star clusters) distributed on the disk. Keeping the same distribution (with zero angular phase lag) implies an a-priori correlation of HII regions with the red star clusters restricted to zero angular phase lag. Keeping the same flux densities implies a uniform equivalent width (EW) for all individual HII regions, limited to the FWHM of the narrow filters used for the observed galaxies, that is 50 - 89Å . This assumption is justified by our finding for dIs in the Virgo Cluster where we showed that individuals HII regions are restricted to EW=10 to 100Å (Heller et al. 1998). We found that, in this way, the images and the LD profiles of the simulated net and continuum images reproduced the patterns observed in the real images. The simulation results in high degrees of star formation lopsidedness with a median A$`{}_{H\alpha }{}^{}=0.77`$, for CI ranging from 0.01 to 30 (median CI$`{}_{H\alpha }{}^{}=1.05`$). A comparison set of net H$`\alpha `$ and continuum images, as well as plots of the azimuthal asymmetry of a real galaxy and a simulated one is shown in Fig. 6.
We plot in Fig. 7 the dependence of A and CI of the simulated galaxies on the number of HII regions (N) and on the integrated H$`\alpha `$ fluxes. The results show no dependence on the total flux. Changing N does not affect the lopsidedness range of possibilities, but there is a clear trend to CI=1 as the number of HII regions increases. This effect is reflected in Figure 8 where we plot A vs. log(CI) for all simulated galaxies (filled circles). Note that these galaxies are distributed around CI=1. The actual galaxies (represented by triangles and squares) show, in general, higher values of CI than the simulated galaxies.
A closer look at the plots of A vs. log(CI) for different number of HII regions (Fig. 9) helps us interpret this effect. We see that the “phase space” accessible to simulated galaxies in the AI-CI plane becomes more restricted, the more HII regions a galaxy has. While for N=1 objects almost one half of the plane is populated, at N=12 the distribution is concentrated mostly at CI=1 for a large spread of AI’s.
A trend of reduced concentration with increasing number of HII regions is visible for the right side of the distribution in Fig. 9. This is explained as the result of the fact that the more HII regions a (simulated) galaxy has, the more “balanced” is the distribution of these HII regions. Another trend is visible for the left side of the distribution in Fig. 9. The fewer HII regions a galaxy has, the better the chance to find these “unbalanced”, more to one side of a galaxy than the other. This means that galaxies with few HII regions will be more asymmetric than galaxies with many HII regions. A test for 15-20 HII regions did not change the distribution for N=12, but it is obvious that by increasing N it will finally converge to the point (A=0, CI=1).
In order to test if the number of HII regions of the actual galaxies was exaggerated by the number of irregularities detected in the LDs we plotted in Fig. 10 only simulated galaxies with N=1, 2 and 3. We can see that reducing the number of HII regions shifts the simulated galaxies to a higher mean CI=1.3. This fits better the LSB sub-group, however, it is not enough in order to explain the general shift of the BCD galaxies. We discuss this in the next section.
Summarizing, we have shown how the degree of asymmetry and concentration index of star forming regions in simulated galaxies change with the total number of HII regions and their luminosity distribution. A similar asymmetry behaviour occurs for the continuum, but to a lesser degree, due to the relatively smaller contribution of the young stellar clusters over the disk brightness.
## 5 Discussion
Our analysis indicates that most dIs show a lopsided morphology in their recent star formation and in the distribution of red light. Since the analysis was performed on the basis of structural indices that are independent of distance, angular size, and/or inclination of the galaxies, we believe that this is a intrinsic property of dwarf-irregular galaxies. The entire sample has a median lopsidedness index of 0.69 in their star forming distribution; similar results are obtained for LSBs and BCDs. For the same concentration index, LSB and BCD galaxies reach similar degrees of lopsidedness. The correlation detected between the continuum and the line emission concentration is supported by a strong correlation between on-going star formation regions and the red stellar population in the angular-phase domain. The correlation is stronger in BCDs, with a trend for smaller angular phase lags than LSBs. The results are consistent with the correlation found between line and continuum fluxes of individual HII regions in dIs in the Virgo Cluster, which is much stronger for BCDs than for LSBs (Heller et al. 1999).
We mentioned already an important difference between the BCD and LSB galaxies: BCDs exhibit stronger H$`\alpha `$ concentration in their nuclear regions than do LSBs. This is emphasized by the profiles of the lopsidedness distribution and is best seen in the distribution of CI (see e.g., Fig. 4a), where the squares tend more to the right than the filled triangles. This tendency is not shown as strongly in the distribution of CI<sub>Cont</sub>; the degree of concentration of the red continuum light in BCDs is rather similar to that in LSB galaxies (Fig. 4c).
Interpreting the continuum light in the simulated galaxies as showing the distribution of previous stellar generations implies that in this aspect LSBs and BCDs are similar, but it is clear that shorter scale-length exponential disks are needed in compact galaxies. The implication is the presence of at least one past major star formation event in the central regions. This result is consistent with detailed surface brightness fitting of BCDs (e.g., Salzer & Norton 1998, Norton & Salzer 2000).
The differences become more evident as one compares the distribution of newly-formed stars, as measured by the H$`\alpha `$ emission in the real galaxies, with that in the simulated galaxies. BCDs have more concentrated emission and do not fit the median CI of the simulated galaxies. The different concentration indices measured for the two types are not merely a result of fewer HII regions per BCD than per LSB; reducing to 1-3 the number of HII regions in simulated galaxies increased the mean of the distribution by 30%, but the lower limit in CI for BCDs stayed were the model predicts the mean. We conclude that randomly generated star formation may be proceeding through the disk in LSB dwarf-irregular galaxies, but probably not in BCDs.
The higher correlation of the line and continuum LDs in BCDs is reminiscent of the similarity of blue and near-IR images of galaxies in the HDF (Richard Ellis, private communication). In the rest frame of these galaxies, imaged with the WFPC-2 and NICMOS, the optical colors correspond to rest-frame UV (i.e., the young stars) and the near-IR correspond to the optical continuum light used here to estimate the distribution of the older stellar populations. Their correlation shows that at z $``$2 the young stars form where there are more of the old stars, as we found here for dwarf irregulars.
From the diagram in Fig. 10 we learn that there is a limit to the degree of concentration a real LSB galaxy can have; higher CI indices would imply non-realistic, extremely extended galaxies with a dominant optical core, such as the very rare Malin I-types. In our LSB sample, the upper concentration limit is set by the extended-HI galaxy sub-group of van Zee et al. (primary). The behavior should be similar at the lower limit, but this limit is not well-defined.
## 6 Conclusions
We analyzed images of 78 dIs and measured the concentration and asymmetry of the H$`\alpha `$ line and red continuum emission by applying an objective automatic algorithm and by tracing the asymmetry along the azimuthal direction. Our findings show a high degree of asymmetry of the H$`\alpha `$ emission, which follows a milder asymmetry in the distribution of the red light. Both concentration indices (line and continuum) are highly correlated. The continuum and line emission lopsidedness distribution are correlated in angular phase, and there is a trend for higher correlation and smaller angular phase lags in BCDs than in LSBs. We found considerable differences between these two types of dwarf galaxies in terms of lopsidedness distribution profiles and concentration index.
We showed that a random distribution of HII regions can produce the observed lopsidedness of low surface brightness disk-like systems. The key parameters that most affect the model are the scale lengths, the number of HII regions per galaxy, and the luminosity distribution of the HII regions. The model fits well most of the observables: the frequency, strength, and profiles of the lopsidedness are recovered, both in the line emission and in the continuum. The model matches the distribution observed in normal LSB dwarf galaxies in the lopsidedness-concentration plane, but short scale-length exponential disks and some central diffuse light components are called for in LSB galaxies with extended HI envelopes. We showed that reducing the number of HII regions cannot explain the higher concentrations observed in BCDs. It seems that a random distribution of star formation may explain the patterns observed in LSB dIs, but not in BCDs.
## Acknowledgments
AH and NB acknowledge support from the US-Israel Binational Science Foundation. EA is supported by a special grant from the Ministry of Science and the Arts to develop TAUVEX, a UV space imaging experiment. NB is grateful for the continued support of the Austrian Friends of Tel Aviv University. Astronomical research at Tel Aviv University is partly supported by a Center of Excellence Award from the Israel Academy of Sciences.
## References
Abraham, R. G., van den Bergh, S., Glazebrook, K., Ellis, R. S., Santiago, B. X., Surma, P., & Griffiths, R. E., 1996, ApJS, 107, 1.
Alexander, T., in “Astronomical Time Series”, edited by D. Maoz, A. Sternberg & E. M Leibowitz, 1997, Astrophysics and Space Science Library, 218, Kluwer: Dordrecht.
Almoznino, E. & Brosch, N., 1998, MNRAS, 298, 920.
Brosch, N., Heller, A.B. & Almoznino, E., 1998, MNRAS, 300,1091.
Carignan, C., Beaulieu, S., & Freeman, K. C., 1990, AJ, 99, 178.
Conselice, C. J., 1997, PASP, 109, 1251.
Conselice, C. J. & Bershady, M.A., in “After the Dark Ages: When Galaxies were Young (the Universe at $`2<z<5`$).”, 9th Annual October Astrophysics Conference in Maryland, held 12-14 October,1998. , edited by S. Holt & E. Smith., American Institute of Physics Press, 1999, 225.
Dohm-Palmer, R. C., Skillman, E. D., Saha, A., Tolstoy, E., Mateo, M., Gallagher, J., Hoessel., J., Chiosi, C., Dufour, R. J., 1997, AJ, 114, 2527.
Dohm-Palmer, R. C., Skillman, E. D., Saha, A., Tolstoy, E., Mateo, M., Gallagher, J., Hoessel., J., Chiosi, C., Dufour, R. J., 1998a, AJ, 115, 152.
Dohm-Palmer, R. C., Skillman, Gallagher, J., Tolstoy, E., Mateo, M., Dufour, R. J., Saha, A., Hoessel., J., Chiosi, C., 1998b, AJ, 116, 1227.
Draper, N.R. & Smith, H., 1981, Applied Regression Analysis, New York: Wiley, 32.
Edelson, R., & Krolik, J. H., 1988, ApJ, 333, 646.
Elmegreen, B. G., 1998, in ”Origins of Galaxies, Star, Planets and Life” (C. E. Woodward, H. A. Thronson & M. Shull, eds.), ASP series.
Gerola, H., Seiden, P. E., Schulman, L., S., 1980, ApJ, 242, 517.
Heller, A. B., Almoznino, E. & Brosch, N. 1998, IAU Colloquium No. 171: “The Low Surface Brightness Universe”, ed. J. I. Davies, C. Impey, & S. Phillips, ASP Conf. Series, 170, 282.
Heller, A. B., Almoznino, E. & Brosch, N., 1999, MNRAS, 304, 8.
Heller, A. B., Almoznino, E. & Brosch, N., 2000, in preparation.
Hunter, D. A., Elmegreen, B. G., & Baker, A. L., 1998, ApJ, 493, 595.
Jog, C. J., 1997, ApJ, 488, 642.
Kaspi, S., Smith, P. S., Netzer, H., Maoz, D., Jannuzi, B. T., & Giveon, U., 2000, ApJ, in press.
Kornreich, D. A., Haynes, M. P., & Lovelace, R. V. E., 1998, AJ, 116, 2154.
Larson, R. B., 1986, MNRAS, 218, 409.
Leitherer, C. & Heckman T.M., 1995, ApJS. 96, 9.
Levine, S. E. & Sparke, L.S., 1998, ApJ, 496L, 13L.
Netzer, A., Heller, A., Loinger, F., Alexander, T., Baldwin, J. A., Wills, B. J., Han, M., Frueh, M., & Higdon, J. L., 1996, MNRAS, 279, 429.
Mihos, J. C., McGaugh, S. S., & de Blok, W. J. G., 1997, AA, 477, L79.
Moffat, A. F. J., 1969, A& A, 3, 455.
Norton, S. A. & Salzer, J. J., 2000, in preparation.
Rudnick, G. & Rix. H.W., 1998, AJ, 116, 1163.
Salzer, J. J. & Norton, S.A., 1998, IAU Colloquium No. 171: “The Low Surface Brightness Universe”, ed. J. I. Davies, C. Impey, & S. Phillips, ASP Conf. Series, 170, 253.
Schechter, 1976, ApJ, 203, 297.
Skillman, E. D., Bothun, G. D., Murray, M. A., & Warmels, R.H., 1987, A&A , 185, 61.
van Zee, L., Haynes, M. P., & Salzer, J. J., 1997a, AJ, 114, 2479.
van Zee, L., Haynes, M. P., & Salzer, J. J., 1997b, AJ, 114, 2497.
van Zee, L., Haynes, M. P., Salzer, J. J., & Broeils, A. H., 1997c, AJ, 113, 1618.
van Zee, L., Westpfahl, D., Haynes, M. P., Salzer, J. J., 1998a, AJ, 115,1000.
van Zee, L., Skillman, E. D., & Salzer, J. J., 1998b, AJ, 116,1186.
Walker, I. R., Mihos, J. C. & Hernquist, L. , 1996, ApJ, 460, 121.
Young, L. M., Lo, K. Y., 1996, ApJ, 462, 203.
Young, L. M., van Zee, L., & Lo, K. Y., 2000, in preparation.
Zaritsky, D. & Rix. H. W., 1997, ApJ, 477, 118.
|
warning/0005/hep-th0005230.html
|
ar5iv
|
text
|
# 1 Introduction
## 1 Introduction
After the Randall-Sundurm (RS) suggestion of a new compactification to confine gravity to four-dimensions, there is an explosive activity of studying various generalizations. One possible application of the RS scenario is to solve the cosmological constant problem which is the key obstacle to make the models of particle physics that can be derived from string theory more realistic . Recently a simple self-tuning mechanism has been suggested in to at least improve on the cosmological constant problem (see also for an earlier mechanism involving extra dimensions) where all order standard model loop contributions are off-set by the parameters appearing in the solution to the five dimensional gravity-scalar system. The authors of showed that one can find static solutions to the classical equations of motion for the coupled gravity-scalar system. However, all the solutions found in (see also ) were obtained either by having a constant potential for the scalar or by making a simple ansatz. In some integrable bulk potentials are obtained and a general and exact solution was obtained for the exponential potential by using a first order formalism . Most of these solutions have naked singularities.
Leaving aside the physical interpretation of the naked singularity one would like to understand better these solutions. However the “explicit” solution given in is not quite explicit and it is given only as an implicit function of the coordinate $`r`$ (see below). Even if this is quite enough for numerical analysis they are not quite suitable for ananlytic calculations. Some solutions are quite explicit and reasonably simple. But these solutions are just invented purposely (to be solvable and simple) and the origin for the scalar potential is not clear. (The potentials in gauged supergravity are also quite complicated and the brane world is not yet realized as a BPS or non-BPS configuration of supersymmetric theory (however see for a cosine potential).) In this paper we will fill this gap and show that there does exist a simple and exact solution with the simple exponential scalar potential. Amazingly, the exponent for which we can find such a solution is exactly the value which comes from string theory as noted in .
The organization of this paper is as follows: in Section 2 we present full details of our exact solution. Here the solution is not obtained by using the first order formalism as was done in . Instead we cast the coupled system of equations into a simple form of (effectively) two first order differential equations, given as eqs. (11) and (12). From here we found the special value which made the equations solvable by elementary method and function. This value $`c=1`$ ($`a=\sqrt{\frac{2}{3}}`$) corresponds exactly to the potential which was reduced from the compactification of string theory with a non-vanishing cosmological constant ( $`a=\frac{4}{3}`$ in their normalization, see eq. (1) below). We also show that this solution is self-tuning when a 3-brane is included. In Section 3 we turn to a more general problem of finding solutions with horizons. Here the equations can be simplified as in Section 2 and we also found a first integration, eq. (49). The equations obtained have the same structure as the equations without horizons. By using the same method we found a similar exact solution with fine-tuned exponent coefficient with an integration constant (eq. (50)). Failing to find a solution with horizon we prove the non-existence of horizons. The naked singularities cannot be resolved by horizons. We comment and conclude in Section 4.
## 2 The self-tuning exact solution
Our starting point is the following action for five-dimensional gravity coupled to a single real scalar with an exponential potential :
$$S=d^Dx\sqrt{|G|}\left(R\frac{1}{2}(\varphi )^2\mathrm{\Lambda }e^{a\varphi }\right).$$
(1)
Brane sources are not included at this moment but they can easily be added (see below). From this action we have the following equations of motion:
$`R_{MN}{\displaystyle \frac{1}{2}}_M\varphi _N\varphi ={\displaystyle \frac{\mathrm{\Lambda }}{D2}}e^{a\varphi }G_{MN},`$ (2)
$`^2\varphi a\mathrm{\Lambda }e^{a\varphi }=0.`$ (3)
In this paper we will first study solutions with Poincaré-invariant $`(D1)`$-dimensional spacetime and we have the following ansatz for the metric:
$$ds^2=e^{2A(r)}\left(\eta _{\mu \nu }dx^\mu dx^\nu \right)+dr^2,$$
(4)
where the function $`A`$ depends only on $`r`$<sup>3</sup><sup>3</sup>3We use the mostly positive convention for the metric..
Substituting this ansatz into eqs. (2)-(3), these equations are given as follows (setting $`D=5`$):
$`A^{\prime \prime }+4(A^{})^2+{\displaystyle \frac{\mathrm{\Lambda }}{3}}e^{a\varphi }=0,`$ (5)
$`4A^{\prime \prime }+4(A^{})^2+{\displaystyle \frac{1}{2}}(\varphi ^{})^2+{\displaystyle \frac{\mathrm{\Lambda }}{3}}e^{a\varphi }=0,`$ (6)
$`\varphi ^{\prime \prime }+4A^{}\varphi ^{}a\mathrm{\Lambda }e^{a\varphi }=0.`$ (7)
where prime denotes differentiation with respect to $`r`$, and we have assumed that the scalar field also depends only on $`r`$.
By simple algebra the above three equations can be recasted into the following form
$`A^{\prime \prime }+{\displaystyle \frac{1}{6}}(\varphi ^{})^2=0,`$ (8)
$`(A^{\prime \prime }+{\displaystyle \frac{1}{3a}}\varphi ^{\prime \prime })+4(A^{}+{\displaystyle \frac{1}{3a}}\varphi ^{})A^{}=0,`$ (9)
$`A^{\prime \prime }+4(A^{})^2+{\displaystyle \frac{\mathrm{\Lambda }}{3}}e^{a\varphi }=0,`$ (10)
By doing some simple rescaling and setting $`\varphi =\sqrt{6c}\mathrm{\Phi }`$ and $`a=\sqrt{2c/3}`$, the above equations are simplified to the following form:
$`A^{\prime \prime }+c(\mathrm{\Phi }^{})^2=0,`$ (11)
$`(A^{\prime \prime }+\mathrm{\Phi }^{\prime \prime })+4(A^{}+\mathrm{\Phi }^{})A^{}=0,`$ (12)
$`A^{\prime \prime }+4(A^{})^2+{\displaystyle \frac{\mathrm{\Lambda }}{3}}e^{2c\mathrm{\Phi }}=0,`$ (13)
Adding together (11) and (12) gives the following equation
$$(2A^{\prime \prime }+\mathrm{\Phi }^{\prime \prime })+(2A^{}+\mathrm{\Phi }^{})^2+(c1)(\mathrm{\Phi }^{})^2=0.$$
(14)
Now we observe that for $`c=1`$ the above equation simplifies and we can solve it to obtain
$$2A^{}+\mathrm{\Phi }^{}=\frac{1}{rr_0},$$
(15)
where $`r_0`$ is an arbitrary (integration) constant.
Now we use the above relation in eq. (11) to eliminate $`A`$ to get the following equation for $`\mathrm{\Phi }`$ (remember $`c=1`$):
$$\mathrm{\Phi }^{\prime \prime }2(\mathrm{\Phi }^{})^2+\frac{1}{(rr_0)^2}=0.$$
(16)
To solve this equation we first find a special solution by trying the ansatz $`\mathrm{\Phi }^{}=a_1/(rr_0)`$ to get $`a_1=1`$. Setting
$$\mathrm{\Phi }^{}=\stackrel{~}{\mathrm{\Phi }}^{}\frac{1}{rr_0},$$
(17)
the equation for $`\stackrel{~}{\mathrm{\Phi }}`$ is given as follows
$$\stackrel{~}{\mathrm{\Phi }}^{\prime \prime }+\frac{4}{rr_0}\stackrel{~}{\mathrm{\Phi }}^{}2(\stackrel{~}{\mathrm{\Phi }}^{})^2=0,$$
(18)
which can easily be solved, first by dividing both sides by $`\mathrm{\Phi }^{}`$ and then combining the first two terms into a total derivative to yield
$$\frac{1}{\stackrel{~}{\mathrm{\Phi }}^{}(rr_0)^4}=\frac{2/3}{(rr_0)^3}+\text{const.}.$$
(19)
Substituting the above result back to eq. (17) we finally get
$`\mathrm{\Phi }(r)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\mathrm{ln}|rr_0|\mathrm{ln}|(r_1r_0)^3(rr_0)^3|\right)+\mathrm{\Phi }_0,`$ (20)
$`A(r)`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left(\mathrm{ln}|rr_0|+\mathrm{ln}|(r_1r_0)^3(rr_0)^3|\right)+A_0,`$ (21)
where $`r_0`$, $`r_1`$, $`\mathrm{\Phi }_0`$ and $`A_0`$ are integration constants. Substituting these results back into eq. (13) fixes the integration constant $`\mathrm{\Phi }_0`$:
$$\mathrm{\Phi }_0=\frac{1}{2}\mathrm{ln}\frac{9}{|\mathrm{\Lambda }|}=\mathrm{ln}3\frac{1}{2}\mathrm{ln}|\mathrm{\Lambda }|.$$
(22)
For $`\mathrm{\Lambda }>0`$ the consistency of the equation ( $`A^{\prime \prime }+4(A^{})^20`$) requires $`0<rr_0<r_1r_0`$ for $`r_1>r_0`$ and $`r_1r_0<rr_0<0`$ for $`r_1<r_0`$. It is always possible to set $`r_0=0`$ by using the translation invariance along the fifth dimension. We also assume $`r_1>0`$ as the other case can be obtained from this one by a reflection operation $`rr`$. In this case our solution is interpolating between two naked singularities at $`r=0`$ and $`r=r_1`$.
On the other hand for $`\mathrm{\Lambda }<0`$ the consistency of the equation requires the following inequality (after setting $`r_0=0`$):
$$\frac{r}{r^3r_1^3}>0.$$
(23)
For $`r_1>0`$ this gives two regions: $`r>r_1`$ and $`r<0`$. For $`r_1<0`$ the two regions are $`r>0`$ and $`r<r_1`$. In both cases, the five dimensional spacetime consists of two disconnected pieces. Each piece has a naked singularity at $`r=\pm r_0`$ and extends to $`\pm \mathrm{}`$.
To get some idea of how the solution looks like, in Fig. 1 we show the profile of the volume of the 4d spacetime along the fifth dimension $`r`$. The height is the relative volume of the 4d space-time which is given by $`e^{4A(r)}=|r(1r^3)|`$. The middle part is the region $`\mathrm{\Lambda }>0`$. The disconnected left and right parts are the two regions for $`\mathrm{\Lambda }<0`$.
As one can see from Fig. 1, this solution is not symmetric between the two naked singularities. The singularity ar $`r=0`$ is at weak coupling $`g=e^\varphi =0`$, and the singularity at $`r=R`$ is at strong coupling $`g=+\mathrm{}`$. We will discuss the possible physical interpretation of this asymmetry in Section 4.
With this exact (bulk) solution at hand one can also study the self-tuning mechanism of refs. by putting a 3-brane along the 5th dimension, say at $`r=0`$. As the self-tuning mechanism is not evident here, we will analyse it in some details. For definiteness we will consider the case of $`\mathrm{\Lambda }<0`$. The solution on the two sides of the 3-brane is given as follows<sup>4</sup><sup>4</sup>4This solution doesn’t confine gravity because $`A(r)+\mathrm{}`$ as $`r\pm \mathrm{}`$. We use it here just to show the self-tuning mechanism.:
$`r`$ $`>`$ $`0,`$ (25)
$`A(r)={\displaystyle \frac{1}{4}}\left(\mathrm{ln}(r+r_1)+\mathrm{ln}((r+r_1)^3R_1^3)\right)+A_1,`$
$`\mathrm{\Phi }(r)={\displaystyle \frac{1}{2}}\left(\mathrm{ln}(r+r_1)\mathrm{ln}((r+r_1)^3R_1^3)\right)+\mathrm{\Phi }_0,`$
$`r`$ $`<`$ $`0,`$ (27)
$`A(r)={\displaystyle \frac{1}{4}}\left(\mathrm{ln}(r_2r)+\mathrm{ln}((r_2r)^3R_2^3)\right)+A_2,`$
$`\mathrm{\Phi }(r)={\displaystyle \frac{1}{2}}\left(\mathrm{ln}(r_2r)\mathrm{ln}((r_2r)^3R_2^3)\right)+\mathrm{\Phi }_0,`$
where $`r_1`$, $`r_2`$, $`R_1`$, $`R_2`$, $`A_1`$ and $`A_2`$ are constants and $`\mathrm{\Phi }_0`$ is given by eq. (22). We take $`r_1`$ and $`r_2`$ to be positive and they satisfy the following inequalities: $`r_1>R_1`$ and $`r_2>R_2`$. ($`R_1`$ and $`R_2`$ are not necessarily positive and can take any (real) values.) One can check that these are solutions to the bulk equations in their respective valid regions.
The matching conditions are as follows <sup>5</sup><sup>5</sup>5Please do the rescaling $`\varphi \frac{3}{2}\mathrm{\Phi }`$ and set $`a=\frac{4}{3}`$.
$`\mathrm{\Phi }^{}(r)|_{r=0+}\mathrm{\Phi }^{}(r)|_{r=0}={\displaystyle \frac{bV}{4}}e^{\frac{3b}{2}\mathrm{\Phi }(0)},`$ (28)
$`A^{}(r)|_{r=0+}A^{}(r)|_{r=0}={\displaystyle \frac{V}{6}}e^{\frac{3b}{2}\mathrm{\Phi }(0)},`$ (29)
$`\mathrm{\Phi }(r)|_{r=0+}=\mathrm{\Phi }(r)|_{r=0}\mathrm{\Phi }(0),`$ (30)
$`A(r)|_{r=0+}=A^{}(r)|_{r=0}.`$ (31)
One can always use the last equation to solve $`A_1`$ and the other constant $`A_2`$ is an overall constant. Noticing the particular combinations $`2A^{}\pm \mathrm{\Phi }^{}`$ we obtain the following relations from the first three equations in the above:
$`{\displaystyle \frac{1}{r_1}}+{\displaystyle \frac{1}{r_2}}={\displaystyle \frac{V}{4}}({\displaystyle \frac{4}{3}}b)e^{\frac{3b}{2}\mathrm{\Phi }(0)},`$ (32)
$`{\displaystyle \frac{3r_1^2}{r_1^3R_1^3}}+{\displaystyle \frac{3r_2^3}{r_2^3R_2^3}}={\displaystyle \frac{V}{4}}({\displaystyle \frac{4}{3}}+b)e^{\frac{3b}{2}\mathrm{\Phi }(0)},`$ (33)
$`\mathrm{\Phi }(0)\mathrm{\Phi }_0={\displaystyle \frac{1}{2}}(\mathrm{ln}r_1\mathrm{ln}(r_1^3R_1^3))={\displaystyle \frac{1}{2}}(\mathrm{ln}r_2\mathrm{ln}(r_2^3R_2^3)).`$ (34)
Because of $`r_{1,2}>0`$, we must have $`V<0`$ and $`\frac{4}{3}<b<\frac{4}{3}`$ in order to have solutions. From eq. (34) one can solve $`R_2`$ in terms of $`r_{1,2}`$ and $`R_1`$ uniquely. By using eq. (34), Eq. (33) can be simplified as follows:
$$\frac{3r_1}{r_1^3R_1^3}(r_1+r_2)=\frac{V}{4}(\frac{4}{3}+b)e^{\frac{3b}{2}\mathrm{\Phi }(0)},$$
(35)
which can be combined with eq. (32) to yield
$$r_2=\frac{(4+3b)}{3(43b)}\frac{r_1^3R_1^3}{r_1^2}>0.$$
(36)
Substituting this result into eq. (32) and using eq. (34) we have
$$\frac{3}{r_1^{(1+\frac{3b}{4})}(r_1^3R_1^3)^{(1\frac{3b}{4})}}\left[\left(1+\frac{(4+3b)}{3(43b)}\right)r_1^3R_1^3\right]=\frac{V}{4}(\frac{4}{3}b)e^{\frac{3b}{2}\mathrm{\Phi }_0}.$$
(37)
For $`\frac{4}{3}<b<\frac{2}{3}`$ and taking any value of $`R_1>0`$, as $`r_1`$ varies from $`R_1`$ to $`+\mathrm{}`$, the left-hand side of the above equation varies from $`+\mathrm{}`$ to $`0`$. So for any value ($`>0`$) of the right-hand side, we can find a value of $`r_1>0`$ which solves the above equation. This will give a unique solution with $`r_2`$ and $`R_2`$ which depend on an arbitrary $`R_1`$ and the various parameters appearing in the action: $`V`$, $`b`$ and $`\mathrm{\Lambda }`$. So we have a flat 4 dimensional spacetime which is self-tuning. Presumably other regions of the parameters will be covered by different choices of pasting the bulk solutions .
As a final note, we point out that we can also get the solution (II) obtained in by taking the limit $`\mathrm{\Lambda }0`$. Here we must adjust the integration constant $`rR`$ appropriately with $`\mathrm{\Lambda }`$ so that one can take the limit smoothly. An expansion of the solution
$`\mathrm{\Phi }(r)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\mathrm{ln}|r|\mathrm{ln}|R^3r^3|\right)+\mathrm{\Phi }_0,`$ (38)
$`A(r)`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left(\mathrm{ln}|r|+\mathrm{ln}|R^3r^3|\right)+A_0,`$ (39)
around $`r=0`$ gives the solution (2.37) and (2.38) with the positive sign in , and an expansion around $`r=R`$ gives the solution with the minus sign there.
## 3 The non-existence of solutions with horizon
It is important to understand better the physical significance of the singularities in our solution. Normally these naked singularities would be discarded as unphysical, but in some instances there are reasons to believe that considering these singularities may be meaningful .
One possible resolution of these naked singularities is to cover them with horizons which have been studied in the appendix of . Here we will prove the non-existence of such solutions with (generic) exponential potential. We will comment on other interpretations of these singularities in Section 4.
As in the Schwarzchild solution in general relativity we make the following ansatz for the 5 dimensional metric :
$`ds^2`$ $`=`$ $`e^{2A_0(r)}(dt)^2+e^{2A_1(r)}((dx^1)^2+(dx^2)^2+(dx^3)^2)+e^{2(A_1(r)A_0(r)}(dr)^2`$ (40)
$`=`$ $`e^{2A_1(r)}\left(h(r)(dt)^2+(dx^1)^2+(dx^2)^2+(dx^3)^2\right)+{\displaystyle \frac{(dr)^2}{h(r)}},`$
where $`h(r)=e^{2A_0(r)2A_1(r)}`$. Horizon appears at the point where the function $`h(r)`$ has a simple zero.
By using this ansatz we obtain the following equations of motion from eqs. (2)-(3):
$`A_0^{\prime \prime }+2A_0^{}(A_0^{}+A_1^{})+{\displaystyle \frac{\mathrm{\Lambda }}{3}}e^{a\varphi +2(A_1A_0)}=0,`$ (41)
$`A_1^{\prime \prime }+2A_1^{}(A_0^{}+A_1^{})+{\displaystyle \frac{\mathrm{\Lambda }}{3}}e^{a\varphi +2(A_1A_0)}=0,`$ (42)
$`A_0^{\prime \prime }+3A_1^{\prime \prime }+2A_0^{}(A_0^{}+A_1^{})+{\displaystyle \frac{1}{2}}(\varphi ^{})^2+{\displaystyle \frac{\mathrm{\Lambda }}{3}}e^{a\varphi +2(A_1A_0)}=0,`$ (43)
$`\varphi ^{\prime \prime }+2\varphi ^{}(A_0^{}+A_1^{})a\mathrm{\Lambda }e^{a\varphi +2(A_1A_0)}=0,`$ (44)
A similar rescaling of the field $`\varphi `$ and some simple algebras can bring the above equations to the following equivalent form:
$`A_1^{\prime \prime }+c(\mathrm{\Phi }^{})^2=0,`$ (45)
$`(A_1^{\prime \prime }+\mathrm{\Phi }^{\prime \prime })+2(A_1^{}+\mathrm{\Phi }^{})(A_0^{}+A_1^{})=0,`$ (46)
$`(A_0^{\prime \prime }A_1^{\prime \prime })+2(A_0^{}A_1^{})(A_0^{}+A_1^{})=0,`$ (47)
$`A_1^{\prime \prime }+2A_1^{}(A_0^{}+A_1^{})+{\displaystyle \frac{\mathrm{\Lambda }}{3}}e^{2c\mathrm{\Phi }+2(A_1A_0)}=0,`$ (48)
Now we can use eq. (46) and eq. (47) to get
$$A_0^{}A_1^{}=d(A_1^{}+\mathrm{\Phi }^{}),$$
(49)
where $`d`$ is an integration constant. By using this relation we can eliminate $`A_0^{}`$ from eqs. (45)-(46) to arrive at a similar system of equations as discussed in the previous section. We will not go through the various steps as we did in the previous section, as we have intentionally presented the full details arriving at the special value and solving the resulting equations there. The final result is that if we choose $`d`$ carefully (a fine-tuning) we can actually use the method in Section 2 to obtain an exact solution. We have
$`d`$ $`=`$ $`{\displaystyle \frac{2(c1)}{2c}},`$ (50)
$`\mathrm{\Phi }(r)`$ $`=`$ $`{\displaystyle \frac{(2c)}{2c}}\mathrm{ln}|r|{\displaystyle \frac{1}{2}}\mathrm{ln}\left|r_1^{(\frac{4}{c}1)}r^{(\frac{4}{c}1)}\right|+\mathrm{\Phi }_0,`$ (51)
$`A_0(r)`$ $`=`$ $`{\displaystyle \frac{(c^26c+4)}{4c}}\mathrm{ln}|r|`$ (52)
$`+{\displaystyle \frac{(2c)}{4}}\mathrm{ln}\left|r_1^{(\frac{4}{c}1)}r^{(\frac{4}{c}1)}\right|+a_0,`$
$`A_1(r)`$ $`=`$ $`{\displaystyle \frac{(2c)^2}{4c}}\mathrm{ln}|r|+{\displaystyle \frac{c}{4}}\mathrm{ln}\left|r_1^{(\frac{4}{c}1)}r^{(\frac{4}{c}1)}\right|+a_1.`$ (53)
Note that the above solution is not valid for $`c=2`$, for this will give $`d=\mathrm{}`$ which is not a well defined mathematical operation. In this special case we have $`A_1^{}=\mathrm{\Phi }^{}`$. This kind of ansatz was used in and has been well studied there.
With the above special solution at hand one can easily see that there is no horizon. The function $`h(r)`$ is given as
$$h(r)e^{2(A_0(r)A_1(r))}=\text{(const.)}\times r^{\frac{(c^25c+4)}{2c}}(r_1^{(\frac{4}{c}1)}r^{(\frac{4}{c}1)})^{1c},$$
(54)
which can only have possible zeroes at $`r=0`$ and/or $`r=r_1`$. These two points are also the singularities of the metric (curvature). So it is not a horizon.
Actually one can prove that there is no exact solutions with horizons. This proof doesn’t depend on the above special solution. The crucial observation is eq. (49). Assuming we do have a solution with a horizon at $`r=r_h`$:
$$h(r)=e^{2(A_0(r)A_1(r))}=a_3((rr_h)+O((rr_h)^2)).$$
(55)
Then we have
$$A_0^{}A_1^{}=\frac{1/2}{rr_h}+\mathrm{},$$
(56)
around $`r=r_h`$. If we require that $`r_h`$ is different from any singularities of $`A_1(r)`$, $`A_1(r)`$ is well behaved around $`r=r_h`$ and so is $`A_1^{}`$:
$$A_1^{}(r)=A_1^{}(r_h)+A_1^{\prime \prime }(r_h)(rr_h)+\mathrm{}.$$
(57)
By using eq. (49) we have the following expansion for $`\mathrm{\Phi }^{}`$:
$$\mathrm{\Phi }^{}(r)=\frac{1}{2d}\frac{1}{rr_h}+\mathrm{},$$
(58)
but this expansion is in sharp contradiction with eq. (45) ($`c0`$). So we proved that there is no solution with horizon for the five dimensional gravity coupled with a scalar with exponential potential. For a discussion with a generic potential with more scalar fields, please see .
## 4 Comment and conclusion
Five dimensional gravity (with or without matters) is an interesting subject which was resurrected many times. The original Kaluza-Klein idea involves only 5 dimensions. But non-abelian gauge field requires more dimensions. With the advent of superstring and M theory, one can view the extra dimensions as a dynamical thing: the radii of the extra dimensions could depend on coupling constant. Recent development with AdS/CFT correspondence gives a new look at the fifth dimension: it is identified with the energy scale of the 4d field theory and the Hamiltonian-Jocobi equation was used to derive the renormalization group flow equations (see for related works and, for example, for later developments). In RS scenarios the compactification and the confining of gravity to 4 dimensions was achieved by an exponential warped factor and one has the freedom to put 3-branes at some points along the fifth dimension by fine-tuning various parameters . According to this fine-tuning is not required because the parameters in the solution will adjust themselves (a self-tuning mechanism). But here we have a plethora of different solutions and some solutions have naked singularities. The puzzle is: what principle should we use to select the right solution?
As we said in the introduction, solutions with singularities cannot be simply discarded as unphysical. Some solutions with singularities have physical interpretations in type IIB superstring theory (see , for example). In fact the solution we found in eqs. (20)-(21) was argued as some solution-independent behaviour and corresponds to the deformation of non-conformal field theory . We may associate the two naked singularities to two different field theories (or putting it differently: the two naked singularities are resolved by two well-defined field theories). The problem with this interpretation is that it is difficult to think that these two non-conformal field theories can be related by a renormalization group flow. The different behaviour of the coupling constants between the two singularities offer a possible way out: either the 5d gravity description breaks down at one of the singularities or these two non-conformal field theories are in fact strong-weak dual pairs. It is interesting to study the higher embedding of these naked singularities.
For $`\mathrm{\Lambda }>0`$ our exact solution has an interesting behaviour: as $`r`$ varies from $`0`$ to $`r_1`$, $`A(r)`$ (so is $`e^{4A(r)}`$) first increases, reaches a maximum, and then decreases, all within a finite interval of $`r`$. In the original RS scenario this behaviour was introduced by hand by introducing 3-brane. Here we have no 3-branes. We can only attribute this behaviour to the inclusion of scalar field. In the AdS/CFT correspondence , scalar fields are interpreted as deformations of four dimensional field theory. The profile of these scalar fields are just running coupling constants. This correspondence is valid in the strong coupling. Recent suggestions to extend this correspondence to weak coupling seem to be in contradiction with our exact solution. A general programme of Hamiltonian-Jacobi equation/Renormalization group flow equation correspondence seems only an approximation.
## Acknowledgments
I would like to thank Roberto Iengo, Shamit Kachru and K. S. Narain for helpful discussions and comments. This work is supported in part by funds from National Natural Science Foundation of China and Pandeng Project. The author would also like to thank Prof. S. Randjbar-Daemi and the hospitality at Abdus Salam International Centre for Theoretical Physics, Trieste, Italy.
|
warning/0005/cond-mat0005324.html
|
ar5iv
|
text
|
# Stability of a Thin Solid Film with Interactions
## Abstract
We investigate the question of stability of a solid thin film which experiences external interactions such as van der Waals forces from a contacting surface or forces from an external electric field. Both perfectly elastic and viscoelastic material behaviours are considered in linear stability analysis performed here. These analyses indicate that for sufficiently soft (shear modulus between 1 and 10 MPa) and nearly incompressible films (Poisson’s ratio close to 0.5), bifurcations are possible, i. e., the surface of the film becomes non-planar. The modes of bifurcation and rates of growth of perturbations are determined as a function of material parameters. The results of this study are of significance in understanding the adhesive properties between a soft material (such as rubber) and a comparatively rigid solid (such as steel), and the behaviour of soft solid films in an electric field.
Instabilities and pattern formation in thin solid and liquid films are of interest both from a scientific and technological view point. Morphological instabilities in thin liquid films occur due to causes such as competition between capillary forces and van der Waals interactions Safran1993 or an external electric field Schaffer2000 and can often lead to dewetting leading to interesting patterns. Morphological instabilities are also common in solid films; for example, in stressed solid films the strain energy drives the surface roughening in competition with the surface energy with surface mass diffusion being the dissipative mechanismRuckenstein1978 ; Freund1994 .
Analysis of interacting thin films has hitherto been restricted to fluid films. Here, we pose the question of stability of a thin solid film bonded to a rigid substrate whose free surface experiences an effective force. This force may arise from any of the various causes such as a van der Waals interaction with another contacting surface nearby and/or with the substrate, an external electric field, etc. The theoretical analysis presented in this paper indicates that for a soft and nearly incompressible solid thin film, instabilities are possible and that the film “buckles”. Physically, this instability occurs because it is possible, for sufficiently large interaction forces, to reduce the net potential energy of the system (the elastic strain energy and the surface energy of the film + potential energy of interaction of the surface) by a periodic non homogeneous deformation in the film. We believe that these results could be useful in understanding phenomena of adhesion between materials (such as rubber and steel), behaviour of thin films in an electric field etc.
The system considered here is shown in fig. 1 – a film of height $`h`$ bonded to a rigid substrate described by coordinates $`(x_1,x_2)`$ such that surface of the film $`S`$ interacting with external agency has $`x_2=0`$ and that bonded with the rigid substrate has $`x_2=h`$. We restrict attention to plane strain deformations of the film for the sake of mathematical simplicity and to understand the essential physics. The total potential energy of this system system is
$`{\displaystyle _V}W\left(\mathit{ϵ}\right)\text{d}V+{\displaystyle _S}\left(\gamma \sqrt{1+\left(u_{2,1}\right)^2}U\left(𝒖𝒏\right)\right)\text{d}S`$ (1)
where $`\mathit{ϵ}`$ is the strain tensor, $`W(\mathit{ϵ})`$ is the elastic strain energy density, $`\gamma `$ is the surface energy, $`U(𝒖𝒏)`$ is the interaction potential between the surface of the film and the external agency such as a contactor or an electric field, $`𝒖`$ is the displacement vector and $`𝒏`$ is the outward normal to the surface. Linearised analysis is performed by expanding the interaction term $`U(𝒖𝒏)`$ in a power series about $`𝒖=\mathrm{𝟎}`$ and retaining all terms up to quadratic order in $`𝒖`$. The resulting approximate energy functional is
$`{\displaystyle _V}W\left(\mathit{ϵ}\right)\text{d}V+{\displaystyle _S}\gamma \sqrt{1+\left(u_{2,1}\right)^2}\text{d}S`$
$`{\displaystyle _S}\left(U_0+F_{}𝒖𝒏+{\displaystyle \frac{1}{2}}Y\left(𝒖𝒏\right)^2\right)\text{d}S`$ (2)
where
$`U_0=U\left(0\right),F_{}=U^{}\left(0\right)\text{and}Y=U^{\prime \prime }\left(0\right).`$ (3)
The equilibrium stress field $`𝝈`$ in the film (which minimises the potential energy (2) over an appropriate length of the film) satisfies the equilibrium equation $`𝝈=\mathrm{𝟎}`$ in $`V`$ and the boundary condition
$`𝝈𝒏=\gamma u_{2,11}𝒏+F_{}𝒏+Y\left(𝒖𝒏\right)𝒏`$ (4)
on $`S`$. Taking the film to be an isotropic linear elastic solid with shear modulus $`\mu `$ and Poisson’s ratio $`\nu `$, gives a standard expression for the strain energy density Landau1989 with a resulting expression for the stress tensor expressed in terms of the gradient of displacement. Thus the problem can be cast into a boundary value problem for the unknown displacement field with the boundary condition of vanishing displacements at $`x_2=h`$ at the the film substrate interface in addition to (4).
The Homogeneous Solution: A solution to the above boundary value problem exists such that the stresses in the film are equal everywhere. This homogeneous solution $`(𝒖^h)`$ is $`u_1^h=0`$ everywhere, and $`u_2^h`$ has a linear variation with $`x_2`$ starting from 0 at $`x_2=h`$, i.e.,
$`u_2^h(x_1,0)=u_{}={\displaystyle \frac{F_o}{\left(\frac{2\left(1\nu \right)\mu }{\left(12\nu \right)h}Y\right)}}.`$ (5)
For the case when $`\nu =0.5`$, i. e., the incompressible limit, the homogeneous solution is such that the displacement vanishes everywhere in the film, and a pressure field $`p`$ develops such that $`p(x_1,x_2)=F_o`$. So long as
$`Y<Y_m,{\displaystyle \frac{hY_m}{\mu }}={\displaystyle \frac{2\left(1\nu \right)}{\left(12\nu \right)}}`$ (6)
the homogeneous solution is meaningful in that $`u_{}`$ has the same sign as $`F_{}`$. This conditions on $`Y`$ is most easily met when $`\nu `$ is close to $`0.5`$ (the r.h.s. of (6) tends to $`\mathrm{}`$ as $`\nu `$ tends to 0.5), i. e., when the material in nearly incompressible. It is this class of materials that the focus of this paper. Nevertheless, results are presented for all values of $`\nu `$ for the sake of completeness.
Bifurcations: What are the conditions (on $`Y,\mu ,\nu ,h,\gamma `$) for another solution (inhomogeneous state) to exist? If such a solution exists, it can be taken to be of the form $`𝒖^h+𝒖`$, where the symbol $`𝒖`$ now stands for a “bifurcation” displacement field. This bifurcation field must satisfy the equilibrium equations in the bulk and the rigid boundary condition at the film substrate interface, just as the homogeneous solution. On the surface of the film at $`x_2=0`$, the bifurcation field satisfies (here $`𝝈`$ is the additional stresses due to $`𝒖`$),
$`𝝈𝒏`$ $`=`$ $`\gamma u_{2,11}𝒏+Y\left(𝒖𝒏\right)𝒏,`$ (7)
instead of (4). To investigate the existence of a nontrivial solution to the problem defined above, the bifurcation fields are assumed to have the form
$`u_j(x_1,x_2)=e^{ikx_1}u_j\left(x_2\right)`$ (8)
where $`k`$ is a real positive wavenumber. The problem of finding nontrivial bifurcation fields can be cast into the problem of finding those values of $`k`$ such that the functions $`u_j(x_2)`$ are nontrivial. It can be shown that (a detailed account will be published elsewhere) nontrivial bifurcation fields of the form (8) exist for those values of $`k`$ that satisfy the equation
$`(k[4e^{2hk}hk^2(h\mu (1\nu )\gamma )+(e^{4hk}1)k\gamma (37\nu +4\nu ^2)`$
$`+\mu ((34\mu )(1+e^{4hk})2e^{2hk}(512\nu +8\nu ^2))])/`$
$`\left(\left(1\nu \right)\left[\left(34\nu \right)\left(e^{4hk}1\right)4hke^{2hk}\right]\right)=Y`$ (9)
This relation is valid for the incompressible case as well (i. e., when $`\nu =0.5`$). Real roots of (9) are sought when $`Y<Y_m`$ which is the range of $`Y`$ for which the homogeneous solution is valid.
We first focus attention on the case when $`\gamma `$ vanishes. Fig. 2a depicts graphically the solution to (9), i. e., for a given value of $`\nu `$, the values of $`k`$ that solve (9) are plotted as a function of $`Y`$ ($`hY/\mu `$ in non-dimensional terms). The important results may be noted: (i) There are no bifurcation modes for any value of $`\nu `$ when $`hY/\mu <2`$. (ii) For all values of $`\nu `$, $`k=0`$ is a bifurcation mode when $`Y=Y_m`$. (iii)When $`\nu 0.25`$, there are no bifurcation modes for $`Y<Y_m`$. (iv) When $`\nu >0.25`$, there are two modes starting from a critical value $`Y_c`$ (such as the point $`C`$ shown in fig. 2a) that depends on the value of $`\nu `$ until $`Y`$ reaches $`Y_m`$. When the film is incompressible $`hY_c/\mu =6.22`$ and the corresponding bifurcation mode has $`hk_c=2.12`$. For this case bifurcations are possible for all values of $`Y`$ greater than $`6.22\mu /h`$, with two possible values of $`k`$ as shown in the fig. 2a.
Next, we consider the case when $`\gamma 0`$. Fig. 2b shows a plot of the possible wavenumbers of bifurcation modes for various values of $`\gamma `$ with $`\nu =0.4`$. The key effect of the surface energy on the bifurcation modes are noted as follows: (i) Surface energy inhibits bifurcation, in that a larger value of $`Y_c`$ is effected with a non zero value of $`\gamma `$. The critical mode $`k_c`$ decreases with increasing $`\gamma `$. Both of these results are as expected since a larger value of $`k`$ implies a larger energy penalty in terms of surface energy. (ii) As $`\gamma `$ gets larger $`Y_c`$ approaches $`Y_m`$. In fact, it can be shown that $`Y_c`$ equals $`Y_m`$ when $`\gamma =\gamma _m`$ where
$`{\displaystyle \frac{\gamma _m}{\mu h}}={\displaystyle \frac{2\nu \left(4\nu 1\right)}{3\left(12\nu \right)^2}},`$ (10)
a result which is pertinent when $`\nu >0.25`$. The curve for $`\gamma /\mu h=4.0`$ for the case of $`\nu =0.4`$ shown in fig. 2, graphically illustrates this point. If $`\gamma >\gamma _m`$, then there are no bifurcations in the physically meaningful range $`Y<Y_m`$.
A more detailed analysis gives the following formulae for $`Y_c`$ and $`k_c`$ as a function of $`\gamma `$ and $`\nu `$ when $`\gamma /\mu h1`$ and $`\nu 0.5`$:
$`{\displaystyle \frac{h}{\mu }}Y_c(\nu ,\gamma /\mu h)`$ $`=`$ $`6.2210.46\left(12\nu \right)+4.49{\displaystyle \frac{\gamma }{\mu h}},`$
$`hk_c(\nu ,\gamma /\mu h)`$ $`=`$ $`2.122.86\left(12\nu \right)2.42{\displaystyle \frac{\gamma }{\mu h}}.`$ (11)
It is also interesting to consider the time evolution of deformation in the film so as to obtain the dominant or the fastest growing mode. To this end, the film is considered to be viscoelastic with a constitutive relation of the form
$`𝝈`$ $`=`$ $`2\mu \left({\displaystyle \frac{1}{2}}\left(𝒖+𝒖^T\right)+{\displaystyle \frac{\nu }{12\nu }}𝒖𝑰\right)`$ (12)
$`+2\eta \left({\displaystyle \frac{1}{2}}\left(\dot{𝒖}+\dot{𝒖}^T\right){\displaystyle \frac{1}{3}}\dot{𝒖}𝑰\right),`$
where $`(\dot{})`$ stands for the time derivative, $`\eta `$ is a viscosity parameter and $`𝑰`$ is the second order identity tensor. In the consideration of the time evolution of the system, inertial effects are neglected since the time scale of interest is much larger than the time scale of the propagation of an elastic wave through the thickness of the film.
The Homogeneous Viscoelastic Solution: The homogeneous solution of the field equations with the viscoelastic constitutive relation (12) is
$`u_1^h=0,u_2^h(x_1,x_2,t)=u_{}\left(1+{\displaystyle \frac{x_2}{h}}\right)\left(1e^{\omega ^ht}\right)`$ (13)
where $`\omega ^h`$ is given by
$`\omega ^h={\displaystyle \frac{3}{4\eta }}\left({\displaystyle \frac{2\left(1\nu \right)\mu }{\left(12\nu \right)h}}Y\right)={\displaystyle \frac{3}{4\eta }}\left(Y_mY\right).`$ (14)
From (14) it is evident that the time dependent homogeneous solution tends to the elastic homogeneous solution as $`t\mathrm{}`$ when $`Y<Y_m`$. If $`Y>Y_m`$, the present analysis indicates that the homogeneous solution blows up as $`t\mathrm{}`$.
Growth of Perturbations: Just as in the case of the elastic film, it is of interest to investigate the growth of perturbations of the homogeneous solution. The perturbations $`𝒖`$ are assumed to be of the form
$`u_j(x_1,x_2,t)=e^{ikx_1}u_j\left(x_2\right)e^{\omega t}.`$ (15)
For a given $`k`$, the rate of growth $`\omega `$ is determined by insisting that the the perturbation satisfies equilibrium equations and boundary conditions and that they be nontrivial. The relation between $`\omega `$ and $`k`$ can be obtained by replacing $`\mu `$ and $`\nu `$ in (9) respectively by $`\mu ^{}`$ and $`\nu ^{}`$ where,
$`\mu ^{}=\mu +\eta \omega ,\nu ^{}={\displaystyle \frac{3\nu \mu \left(12\nu \right)\eta \omega }{3\mu +\left(12\nu \right)\eta \omega }}.`$ (16)
This procedure results in a cubic equation for $`\omega `$. The solution of this equation is obtained by numerical means.
The solution for $`\omega `$ indicates that for $`Y_c<Y<Y_m`$, all perturbation modes with wavenumbers between the two bifurcation modes given by the elastic analysis are unstable i. e., $`\omega `$ for these modes are positive. Indeed, there is a mode with wavenumber ($`k_m`$) between wavenumbers of the two elastic bifurcation modes such that the rate of growth ($`\omega `$) is a maximum. Fig. 3(a) shows a plot of $`k_m`$ as a function of $`Y`$ ($`Y_cYY_m)`$ for various values of $`\nu `$ (with $`\gamma /\mu h=0)`$. When $`\nu <0.5`$, the value of $`k_m`$ starts at $`k_c`$ when $`Y=Y_c`$ and monotonically falls with increasing $`Y`$. For the case of $`\nu =0.5`$, $`k_m=k_c`$ for all values of $`Y`$. When $`\gamma 0`$, $`k_m`$ is smaller as is evident from fig. 3(b); the effect of surface energy on the fastest growing mode becomes increasingly less significant for large values of $`Y`$. Just as in (11), an analytic result can be derived for $`k_m`$ for small values of $`\gamma /\mu h`$, $`\nu 0.5`$ and $`h(YY_c)/\mu 1`$:
$`hk_m(\nu ,{\displaystyle \frac{\gamma }{\mu h}})`$ $`=`$ $`hk_c(\nu ,{\displaystyle \frac{\gamma }{\mu h}})`$ (17)
$`+\left(0.39{\displaystyle \frac{\gamma }{\mu h}}0.46\left(12\nu \right)\right){\displaystyle \frac{h}{\mu }}\left(YY_c\right)`$
Instability in a thin film whose surface experiences forces depends on three key sets of non-dimensional parameters namely the Poisson’s ratio $`\nu `$, the normalised second derivative of the interaction potential $`hY/\mu `$ and the normalised surface energy $`\gamma /\mu h`$. The whole picture of stability and bifurcation in this system and its dependence on the nondimensional parameters can be depicted pictorially as shown in fig. 4. Region $`I`$ in fig. 4 is where the homogeneous solution is unique and stable while region marked $`III`$ in the figure corresponds to the case when the homogeneous solution is “unphysical”, i.e., this analysis is not adequate. Region $`II`$ is the most interesting – this corresponds to nearly incompressible material behaviour. In this region the homogeneous solution is unstable, with two possible elastic bifurcation modes; a viscoelastic analysis predicts a fastest growing mode with a wave vector that lies between the two elastic bifurcation modes.
We now turn to specific cases of the type of system considered in this paper. First, we consider a rigid contactor interacting with the film via van der Waals forces. Assuming that the contactor is at a distance $`d`$ above the undeformed surface of the film, the interaction potential $`U`$ can can be taken to be $`U(𝒖𝒏)=\frac{A}{12\pi (𝒖𝒏d)^2}`$ with $`F_{}=\frac{A}{6\pi d^3}`$, $`Y=\frac{A}{2\pi d^4}`$. Taking the film to be made of rubber ($`\mu =1`$ MPa, $`\nu =0.5`$, $`\gamma =0.1`$J/m<sup>2</sup>) and $`h=1`$micron with $`A1`$eV. When $`d=10`$nanometers we get $`hY/\mu =1.6`$ and for $`d=5`$nanometers, $`hY/\mu =25.6`$. Since the latter value is greater than $`hY_c/\mu `$ which is 6.63 when $`\gamma /\mu h=0.1`$ (which is the present case), it is clear that the condition for bifurcation will be achieved as $`d`$ is reduced from $`10`$nm to $`5`$nm. Thus as the contactor approaches the film, the film would buckle. This implies that the contact that forms between the contactor surface and the film will not be planar. We are not aware of any experimental work that can corroborate our results. We do, however, hope that the contents of this paper will be useful in designing experiments to verify our conclusions.
The second case considered is that of a film interacting with an external electric field. The system consists of two plates separated by a distance $`d`$; the bottom plate is coated with a nearly incompressible polymeric film of height $`h`$. A potential difference of $`V`$ is applied between the two plates. The quantity of interest is the value of the gap thickness $`dh`$ at which instability occurs in the film. The potential of interaction for this case is given by $`U(𝒖𝒏)=\frac{\epsilon _0\epsilon _pV^2}{2(\epsilon _pd(\epsilon _p1)(h+𝒖𝒏))}`$ where $`\epsilon _0`$ is the permittivity of free space, $`\epsilon _p`$ is the dielectric constant of the polymer. Taking the mechanical properties of the polymer to be same as in the previous case, and taking $`\epsilon _p=3`$, we get that the critical gap thickness $`dh`$ of $`0.05`$micron for a film of height $`0.1`$micron with the applied voltage of 100V. A gap thickness smaller than 0.05micron will cause the film to buckle. It is evident that large electric fields are required to cause the instability.
We do wish to point out that this analysis is based on a linearised model, and will only provide the modes of instability, i. e., the wavelength of surface undulation and not the magnitude. A nonlinear analysis is required to obtain such a quantity and will be pursued in subsequent papers.
|
warning/0005/math0005207.html
|
ar5iv
|
text
|
# Divisorial contractions in dimension 3 which contract divisors to smooth points
## 0 Introduction
Divisorial contractions play a major role in the minimal model program (\[KMM87\]). Now that we know this program works in dimension $`3`$ (\[M88\]), it is desirable to describe them explicitly in dimension $`3`$. Moreover also in view of the Sarkisov program (\[Co95\]) and its applications (for example \[CPR99\]), we can recognize the importance of such description since Sarkisov links of types I and II in this program start from the converse of divisorial contractions.
Now we concentrate on divisorial contractions in dimension $`3`$. Let $`f:(YE)(XP)`$ be such a contraction. There are two ways to deal with $`f`$, that is to say, one starting from $`Y`$, and the other from $`X`$. From the former standpoint, S. Mori classified them in the case when $`Y`$ is smooth (\[M82\]), and S. Cutkosky extended this result to the case when $`Y`$ has only terminal Gorenstein singularities (\[Cu88\]). On the other hand, from the latter standpoint, Y. Kawamata showed that $`f`$ must be a certain weighted blow-up when $`P`$ is a terminal quotient singularity (\[K96\]), and A. Corti showed that $`f`$ must be the blow-up when $`P`$ is an ordinary double point (\[Co99, Theorem 3.10\]).
While it seems that singularities on $`Y`$ make it hard to tackle the problem in the former case, the singularity of $`P`$ may be useful in the latter case because it gives a special filtration in the tangent space at $`P`$. In this paper we treat the case when $`P`$ is a smooth point and prove the following theorem:
###### Theorem 1.2.
Let $`Y`$ be a $`3`$-dimensional $``$-factorial normal variety with only terminal singularities, and let $`f:(YE)(XP)`$ be an algebraic germ of a divisorial contraction which contracts its exceptional divisor $`E`$ to a smooth point $`P`$. Then we can take local parameters $`x,y,z`$ at $`P`$ and coprime positive integers $`a`$ and $`b`$, such that $`f`$ is the weighted blow-up of $`X`$ with its weights $`(x,y,z)=(1,a,b)`$.
Now we explain our approach to the problem. Y. Kawamata adopted the method of comparing discrepancies of exceptional divisors, and A. Corti applied Shokurov’s connectedness lemma (\[$`\mathrm{K}^+`$92, Theorem 17.4\]). But in the case when $`P`$ is a smooth point, these methods do not work well if the center of $`E`$ on $`\mathrm{Bl}_P(X)`$ is a point. Our main tools are the singular Riemann-Roch formula (\[R87, Theorem 10.2\]) on $`Y`$ and a relative vanishing theorem (\[KMM87, Theorem 1-2-5\]) with respect to $`f`$. First with them we derive a rather simple formula for $`dim_k𝒪_X/f_{}𝒪_Y(iE)`$’s and an upper-bound of the number of fictitious non-Gorenstein points of $`Y`$ (Proposition 2.7). Next using this upper-bound, we show that the coefficient of $`E`$ in the pull-back of a general prime divisor through $`P`$ is $`1`$ (Subsection 2.3). And finally investigating the values of $`dim_k𝒪_X/f_{}𝒪_Y(iE)`$’s more carefully, we prove the theorem (Subsection 2.4).
I wish to express my gratitude to Professor Yujiro Kawamata for his valuable comments and warm encouragement. He also recommended me to read the papers \[CPR99\] and \[Co99\]. In fact I found the problem treated here as \[Co99, Conjecture 3.11\].
## 1 Statement of the theorem
We work over an algebraically closed field $`k`$ of characteristic zero. A variety means an integral separated scheme of finite type over $`\mathrm{Spec}k`$. We use basic terminologies in \[$`\mathrm{K}^+`$92, Chapters 1, 2\].
Before we state the theorem, we have to define a divisorial contraction. In this paper it means a morphism which may emerge in the minimal model program (see \[KMM87\]).
###### Definition 1.1.
Let $`f:YX`$ be a morphism with connected fibers between normal varieties. We call $`f`$ a divisorial contraction if it satisfies the following conditions:
1. Y is $``$-factorial with only terminal singularities.
2. The exceptional locus of $`f`$ is a prime divisor.
3. $`K_Y`$ is $`f`$-ample.
4. The relative Picard number of $`f`$ is $`1`$.
Now it is the time when we state the theorem precisely.
###### Theorem 1.2.
Let $`Y`$ be a $`3`$-dimensional $``$-factorial normal variety with only terminal singularities, and let $`f:(YE)(XP)`$ be an algebraic germ of a divisorial contraction which contracts its exceptional divisor $`E`$ to a smooth point $`P`$. Then we can take local parameters $`x,y,z`$ at $`P`$ and coprime positive integers $`a`$ and $`b`$, such that $`f`$ is the weighted blow-up of $`X`$ with its weights $`(x,y,z)=(1,a,b)`$.
## 2 Proof of the theorem
### 2.1 Strategy for its proof
We may assume that $`X`$ is projective and smooth, and consider its algebraic germ if necessary. First we construct a series of birational morphisms.
###### Construction 2.1.
We construct birational morphisms $`g_i:X_iX_{i1}`$ between smooth varieties, integral closed subschemes $`Z_iX_i`$, and prime divisors $`F_i`$ on $`X_i`$ inductively, and define positive integers $`n,m`$, with the following procedure:
1. Define $`X_0`$ as $`X`$ and $`Z_0`$ as $`P`$.
2. Let $`b_i:\mathrm{Bl}_{Z_{i1}}(X_{i1})X_{i1}`$ be the blow-up of $`X_{i1}`$ along $`Z_{i1}`$, and let $`b_{}^{}{}_{i}{}^{}:X_i\mathrm{Bl}_{Z_{i1}}(X_{i1})`$ be a resolution of $`\mathrm{Bl}_{Z_{i1}}(X_{i1})`$, that is, a proper birational morphism from a smooth variety $`X_i`$ which is isomorphic over the smooth locus of $`\mathrm{Bl}_{Z_{i1}}(X_{i1})`$. We note that $`b_{}^{}{}_{i}{}^{}`$ is isomorphic at the generic point of the center of $`E`$ on $`\mathrm{Bl}_{Z_{i1}}(X_{i1})`$. We define $`g_i=b_ib_{}^{}{}_{i}{}^{}:X_iX_{i1}`$.
3. Define $`Z_i`$ as the center of $`E`$ on $`X_i`$ with the reduced induced closed subscheme structure, and $`F_i`$ as the only $`g_i`$-exceptional prime divisor on $`X_i`$ which contains $`Z_i`$.
4. We stop this process when $`Z_n=F_n`$. This process must terminate after finite steps (see Remark 2.1.2) and thus we get the sequence $`X_n\mathrm{}X_0`$.
5. We define $`mn`$ as the largest integer such that $`Z_{m1}`$ is a point.
6. We define $`g_{ji}(ji)`$ as the morphism from $`X_i`$ to $`X_j`$.
###### Remark 2.1.1.
We remark that $`f_{}𝒪_Y(iE)=g_{0n}𝒪_{X_n}(iF_n)`$ for any $`i`$ because $`E`$ and $`F_n`$ are the same as valuations.
###### Remark 2.1.2.
We prove the termination of the process. Assume that we have the sequence $`X_l\mathrm{}X_0`$ and $`Z_lF_l`$. We take common resolutions of $`X_l`$ and $`Y`$ over $`X`$, that is, birational morphisms $`h:WX_l`$ and $`h^{}:WY`$ from a smooth variety $`W`$ such that $`g_{0l}h=fh^{}`$. We put
$`K_Y`$ $`=f^{}K_X+aE,`$
$`K_{X_l}`$ $`=g_{0l}^{}K_X+sF_l+(\mathrm{others}),`$
$`K_W`$ $`=h^{}K_{X_l}+c(h_{}^{}{}_{}{}^{1})_{}E+(\mathrm{others}),`$
$`h^{}F_l`$ $`=(h^1)_{}F_l+t(h_{}^{}{}_{}{}^{1})_{}E+(\mathrm{others}).`$
We note that $`a,s,c`$ and $`t`$ are positive integers. Then
$`K_W`$ $`=h_{}^{}{}_{}{}^{}(f^{}K_X+aE)+(\mathrm{others})`$
$`=h^{}(g_{0l}^{}K_X+sF_l+(\mathrm{others}))+c(h_{}^{}{}_{}{}^{1})_{}E+(\mathrm{others})`$
$`=h^{}g_{0l}^{}K_X+s(h^1)_{}F_l+(st+c)(h_{}^{}{}_{}{}^{1})_{}E+(\mathrm{others}).`$
Comparing the coefficients of $`(h_{}^{}{}_{}{}^{1})_{}E`$, we have $`a=st+c`$ and especially $`a>s`$. On the other hand because we know $`sl+1`$ by the construction of $`F_l`$, we get $`a>l+1`$. It shows that the above process terminates with $`na1`$. ∎
We state an easy lemma.
###### Lemma 2.2.
Let $`f_i:(Y_iE_i)(Xf_i(E_i))`$ with $`i=1,2`$ be algebraic germs of divisorial contractions. Assume that $`E_1`$ and $`E_2`$ are the same as valuations. Then $`f_1`$ and $`f_2`$ are isomorphic as morphisms over $`X`$.
###### Proof.
Let $`g_i:ZY_i`$ with $`i=1,2`$ be common resolutions and $`h=f_ig_i`$. We choose $`g_i`$-exceptional effective $``$-divisors $`F_i(i=1,2)`$ and a $``$-divisor $`G`$ on $`Z`$ such that $`G=g_1^{}E_1+F_1=g_2^{}E_2+F_2`$. Then,
$`Y_i=\mathrm{Proj}_X_{j0}f_i𝒪_{Y_i}(jE_i)=\mathrm{Proj}_X_{j0}h_{}𝒪_Z(jG).`$
For weighted blow-ups in dimension $`3`$, we have a criterion on terminal singularities.
###### Theorem 2.3.
Let $`XP`$ be an algebraic germ of a smooth $`3`$-dimensional variety with local parameters $`x,y,z`$ at $`P`$, let $`r,a,b`$ be positive integers with $`rab`$, and let $`YX`$ be the weighted blow-up of $`X`$ with its weights $`(x,y,z)=(r,a,b)`$. Then $`Y`$ has only terminal singularities if and only if $`r=1`$ and $`a,b`$ are coprime.
By the above lemma and theorem, the problem is reduced to proving that $`F_n`$ equals, as valuations, an exceptional divisor obtained by a weighted blow-up of $`X`$. We restate this in terms of ideal sheaves of $`𝒪_X`$.
###### Proposition 2.4.
(Notation as above). $`F_n`$ equals, as valuations, an exceptional divisor obtained by a weighted blow-up of $`X`$ with its weights $`(x,y,z)=(1,m,n)`$ for suitable local parameters $`x,y,z`$ at $`P`$, if and only if the following conditions hold:
1. $`f_{}𝒪_Y(2E)𝔪_P`$, that is, $`g_{0n}𝒪_{X_n}(2F_n)𝔪_P`$.
2. $`f_{}𝒪_Y(nE)𝔪_P^2`$, that is, $`g_{0n}𝒪_{X_n}(nF_n)𝔪_P^2`$.
Here $`𝔪_P𝒪_X`$ is the ideal sheaf of $`P`$.
###### Proof.
The “only if” part is obvious taking it into account that for any $`i`$ $`g_{0n}𝒪_{X_n}(iF_n)=(x^sy^tz^u|s+mt+nui)`$. Actually $`xg_{0n}𝒪_{X_n}(2F_n)`$ and $`zg_{0n}𝒪_{X_n}(nF_n)`$.
Now we prove the “if” part. The condition 1 means that the coefficient of $`F_n`$ in $`g_{1n}^{}F_1`$ is $`1`$. This says that for any $`i1`$, $`F_i`$ is the only $`g_{0i}`$-exceptional prime divisor on $`X_i`$ containing $`Z_i`$ and the coefficient of $`F_n`$ in $`g_{in}^{}F_i`$ is $`1`$.
We consider a prime divisor $`DP`$ on $`X`$ which is smooth at $`P`$ and define $`1ln`$ as the largest integer such that $`Z_{l1}(g_{0,l1}^1)_{}D`$. Then $`(g_{0i}^1)_{}D`$ is smooth at the generic point of $`Z_i`$ for any $`i<l`$, and so we get $`g_{0l}^{}D=(g_{0l}^1)_{}D+_{i=1}^li(g_{il}^1)_{}F_i+(\mathrm{others})`$. Therefore the coefficient of $`F_n`$ in $`g_{0n}^{}D`$ is $`l`$. By the condition 2, we can choose $`z𝔪_P𝔪_P^2`$ such that $`g_{0n}^{}\mathrm{div}(z)nF_n`$, that is, $`Z_{n1}(g_{0,n1}^1)_{}\mathrm{div}(z)`$ because of the above argument. Adding $`x,y𝔪_P𝔪_P^2`$ such that $`Z_{m1}(g_{0,m1}^1)_{}\mathrm{div}(y)`$, we can take local parameters $`x,y,z`$ at $`P`$. Then $`F_i(1in)`$ equals, as valuations, the exceptional divisor obtained by the weighted blow-up of $`X`$ with its weights $`(x,y,z)=(1,\mathrm{min}\{i,m\},i)`$, and especially $`F_n`$ is obtained by the weighted blow-up of $`X`$ with its weights $`(x,y,z)=(1,m,n)`$. ∎
So we prove the above two conditions.
### 2.2 Preliminaries
Let $`K_Y=f^{}K_X+aE`$, and let $`r`$ be the global Gorenstein index of $`Y`$, that is, the smallest positive integer such that $`rK_Y`$ is Cartier. Since $`a`$ equals the discrepancy of $`F_n`$ with respect to $`K_X`$, $`a_2`$.
###### Lemma 2.5.
(Notation as above). $`a`$ and $`r`$ are coprime.
###### Proof.
Let $`s`$ be the greatest common divisor of $`a`$ and $`r`$, and let $`a=sa^{},r=sr^{}.`$ Since $`r^{}aE=a^{}rE`$ is Cartier by \[K88, Corollary 5.2\], so is $`r^{}K_Y`$. Hence $`r^{}=r`$ and $`s=1`$. ∎
We recall the singular Riemann-Roch formula (\[R87, Theorem 10.2\]).
###### Theorem 2.6.
Let $`X`$ be a projective $`3`$-dimensional variety with only canonical singularities, and let $`D`$ be a Weil divisor on $`X`$ such that for any $`PX`$ there exists an integer $`i_P`$ satisfying $`(𝒪_X(D))_P(𝒪_X(i_PK_X))_P`$. Then there is a formula of the form
$`\chi (𝒪_X(D))`$ $`=\chi (𝒪_X)+{\displaystyle \frac{1}{12}}D(DK_X)(2DK_X)`$
$`+{\displaystyle \frac{1}{12}}Dc_2(X)+{\displaystyle \underset{P}{}}c_P(D),`$
where the summation takes place over singular points of $`X`$, and $`c_P(D)`$ is a contribution depending only on the local analytic type of $`PX`$ and $`𝒪_X(D)`$.
If $`P`$ is a terminal quotient singularity of type $`\frac{1}{r_P}(1,1,b_P)`$, then
$`c_P(D)=\overline{i_P}{\displaystyle \frac{r_P^21}{12r_P}}+{\displaystyle \underset{j=1}{\overset{\overline{i_P}1}{}}}{\displaystyle \frac{\overline{jb_P}(r_P\overline{jb_P})}{2r_P}},`$
where $`\overline{}`$ denotes the smallest residue modulo $`r_P`$, that is, $`\overline{j}=j\frac{j}{r_P}r_P`$ in terms of the round down $``$. The definition of the round down $``$ is $`j=\mathrm{max}\{k|kj\}`$.
And for any terminal singularity $`P`$,
$`c_P(D)={\displaystyle \underset{\alpha }{}}c_{P_\alpha }(D_\alpha ),`$
where $`\{(P_\alpha ,D_\alpha )\}_\alpha `$ is a flat deformation of $`(P,D)`$ to terminal quotient singularities.
###### Remark 2.6.1.
If $`X`$ has only terminal singularities, then we can write the contribution term $`_Pc_P(D)`$ as $`_Qc_Q(D)`$, where
$`c_Q(D)=\overline{i_Q}{\displaystyle \frac{r_Q^21}{12r_Q}}+{\displaystyle \underset{j=1}{\overset{\overline{i_Q}1}{}}}{\displaystyle \frac{\overline{jb_Q}(r_Q\overline{jb_Q})}{2r_Q}}.`$
For its summation takes place over points which need not lie on $`X`$ but may lie on deformed varieties of $`X`$, $`Q`$’s are called “fictitious” points in the sense of M. Reid. This description holds even though $`X`$ has canonical singularities, but in this case $`Q`$’s may lie on deformed varieties of crepant blown-up varieties of $`X`$ (see \[R87\] for details).
By Lemma 2.5, we can take an integer $`e`$ such that $`ae1`$ modulo $`r`$. Then $`(𝒪_Y(E))_Q(𝒪_Y(eK_Y))_Q`$ for any $`QE`$. Using the singular Riemann-Roch formula, we get
(2.1) $`\chi (𝒪_Y(iE))`$ $`=\chi (𝒪_Y)+{\displaystyle \frac{1}{12}}i(ia)(2ia)E^3`$
$`+{\displaystyle \frac{1}{12}}iEc_2(Y)+A_i,`$
where $`A_i`$ is the contribution term and has the below description:
$$A_i=\underset{QI}{}c_Q(iE),$$
$$c_Q(iE)=\overline{ie}\frac{r_Q^21}{12r_Q}+\underset{j=1}{\overset{\overline{ie}1}{}}\frac{\overline{jb_Q}(r_Q\overline{jb_Q})}{2r_Q}.$$
Here $`QI`$ are fictitious singularities. The type of $`Q`$ is $`\frac{1}{r_Q}(1,1,b_Q)`$, $`(𝒪_{Y_Q}(E_Q))_Q(𝒪_{Y_Q}(eK_{Y_Q}))_Q`$ where $`(Y_Q,E_Q)`$ is the fictitious pair for $`Q`$, and $`\overline{}`$ denotes the smallest residue modulo $`r_Q`$. We note that $`b_Q`$ is coprime to $`r_Q`$ and also $`e`$ is coprime to $`r_Q`$ because $`r|(ae1)`$. So $`v_Q=\overline{eb_Q}`$ is coprime to $`r_Q`$. With this description, $`r=1`$ if $`I`$ is empty, and otherwise $`r`$ is the lowest common multiple of $`\{r_Q\}_{QI}`$. We note that $`c_Q(iE)`$ depends only on $`i`$ mod $`r_Q`$ and equals $`0`$ if $`r_Q|i`$. Especially $`A_i`$ depends only on $`i`$ mod $`r`$ and equals $`0`$ if $`r|i`$.
We put $`B_i=(A_i+A_i)`$. Because
$`c_Q(iE)+c_Q(iE)`$ $`=\left(\overline{ie}{\displaystyle \frac{r_Q^21}{12r_Q}}+{\displaystyle \underset{j=1}{\overset{\overline{ie}1}{}}}{\displaystyle \frac{\overline{jb_Q}(r_Q\overline{jb_Q})}{2r_Q}}\right)`$
$`+\left(\overline{ie}{\displaystyle \frac{r_Q^21}{12r_Q}}+{\displaystyle \underset{j=1}{\overset{\overline{ie}1}{}}}{\displaystyle \frac{\overline{jb_Q}(r_Q\overline{jb_Q})}{2r_Q}}\right)`$
$`={\displaystyle \frac{r_Q^21}{12}}+\left({\displaystyle \underset{j=1}{\overset{r_Q}{}}}{\displaystyle \frac{\overline{jb_Q}(r_Q\overline{jb_Q})}{2r_Q}}\right){\displaystyle \frac{\overline{ieb_Q}(r_Q\overline{ieb_Q})}{2r_Q}}`$
$`={\displaystyle \frac{r_Q^21}{12}}+\left({\displaystyle \underset{j=1}{\overset{r_Q}{}}}{\displaystyle \frac{j(r_Qj)}{2r_Q}}\right){\displaystyle \frac{\overline{iv_Q}(r_Q\overline{iv_Q})}{2r_Q}}`$
$`={\displaystyle \frac{\overline{iv_Q}(r_Q\overline{iv_Q})}{2r_Q}}`$
where the third equality comes from the property that $`b_Q`$ and $`r_Q`$ are coprime, we have
(2.2) $`B_i={\displaystyle \underset{QI}{}}(c_Q(iE)+c_Q(iE))={\displaystyle \underset{QI}{}}{\displaystyle \frac{\overline{iv_Q}(r_Q\overline{iv_Q})}{2r_Q}}.`$
###### Proposition 2.7.
(Notation as above).
(A) $`rE^3_{>0}.`$
(B) $`1={\displaystyle \frac{1}{2}}aE^3+{\displaystyle \underset{QI}{}}{\displaystyle \frac{v_Q(r_Qv_Q)}{2r_Q}}.`$
(C) $`dim_k𝒪_X/f_{}𝒪_Y(iE)=`$
$`i^2{\displaystyle \frac{1}{2}}{\displaystyle \underset{QI}{}}\underset{0j<i}{\mathrm{min}}\{(1+j)jr_Q+i(i12j)v_Q\}(1ia).`$
(D) $`{\displaystyle \underset{QI}{}}\mathrm{min}\{v_Q,r_Qv_Q\}=dim_kf_{}𝒪_Y(2E)/𝔪_P^2.`$
###### Remark 2.7.1.
In particular (A), (C) and (D) are essential. We use (A) to bound the value of $`a`$ from above and use (C) to control the values of $`r_Q`$’s. (D) shows that the number of fictitious non-Gorenstein points of $`Y`$ is at most $`3`$. We prove the conditions 1 and 2 in Proposition 2.4 according to the value of $`dim_kf_{}𝒪_Y(2E)/𝔪_P^2`$.
###### Remark 2.7.2.
In fact, because of (2.2) and (2.9) the right hand side of (C) is the same if we replace $`v_Q`$ by $`r_Qv_Q`$.
###### Proof.
We consider the exact sequence:
(2.3) $`0𝒪_Y((i1)E)𝒪_Y(iE)𝒬_i0.`$
By (2.1), we get
(2.4) $`\chi (𝒬_i)`$ $`=\chi (𝒪_Y(iE))\chi (𝒪_Y((i1)E))`$
$`={\displaystyle \frac{1}{12}}\{2(3i^23i+1)3(2i1)a+a^2\}E^3`$
$`+{\displaystyle \frac{1}{12}}Ec_2(Y)+A_iA_{i1}.`$
Since $`\chi (𝒬_i)\chi (𝒬_{r+i})=\frac{r}{2}(a+1r2i)E^3`$ is an integer for any $`i`$ and $`E^3`$ is positive, we have (A).
By (2.4),
(2.5) $`\chi (𝒬_i)\chi (𝒬_{i+1})=(i+{\displaystyle \frac{1}{2}})aE^3+B_{i+1}B_i.`$
Let $`d(i)=dim_kf_{}𝒪_Y(iE)/f_{}𝒪_Y((i1)E)`$. We note that $`d(i)=0`$ if $`i1`$, and $`d(0)=1`$. Because $`(Y,\epsilon E)`$ is weak KLT and $`iE(K_Y+\epsilon E)`$ is $`f`$-ample for a sufficiently small positive rational number $`\epsilon `$ and an integer $`ia`$, using \[KMM87, Theorem 1-2-5\], we have $`R^jf_{}𝒪_Y(iE)=0`$ for $`ia,j1`$. So by (2.3), for any $`ia`$,
* $`H^0(Y,𝒬_i)=f_{}𝒬_i=f_{}𝒪_Y(iE)/f_{}𝒪_Y((i1)E),`$
* $`H^j(Y,𝒬_i)=R^jf_{}𝒬_i=0`$ for $`j1`$,
and therefore $`d(i)=\chi (𝒬_i)`$.
Putting $`i=0`$ in (2.5), we get
(2.6) $`1={\displaystyle \frac{1}{2}}aE^3+B_1.`$
Combining this and (2.2) with $`i=1`$, we get (B).
With (2.5), we obtain for $`1ia`$,
(2.7) $`{\displaystyle \underset{1j<i}{}}d(j)`$ $`={\displaystyle \underset{1j<i}{}}\{\chi (𝒬_j)\chi (𝒬_{j+1})\}`$
$`={\displaystyle \underset{1j<i}{}}\{(j+{\displaystyle \frac{1}{2}})aE^3+B_{j+1}B_j\}`$
$`={\displaystyle \frac{1}{2}}(i^21)aE^3+B_iB_1.`$
Eliminating $`\frac{1}{2}aE^3`$ with (2.6), we obtain
(2.8) $`{\displaystyle \underset{1j<i}{}}d(j)=(i^21)+B_ii^2B_1(1ia).`$
Since for $`i1`$,
$`{\displaystyle \frac{\overline{iv_Q}(r_Q\overline{iv_Q})}{2r_Q}}i^2{\displaystyle \frac{v_Q(r_Qv_Q)}{2r_Q}}`$
$`={\displaystyle \frac{1}{2}}\{r_Q({\displaystyle \frac{iv_Q\overline{iv_Q}}{r_Q}}{\displaystyle \frac{iv_Q}{r_Q}}+{\displaystyle \frac{1}{2}})^2+i^2{\displaystyle \frac{v_Q(r_Qv_Q)}{r_Q}}{\displaystyle \frac{r_Q}{4}}\}`$
$`={\displaystyle \frac{1}{2}}\{r_Q({\displaystyle \frac{iv_Q}{r_Q}}{\displaystyle \frac{iv_Q}{r_Q}}+{\displaystyle \frac{1}{2}})^2+i^2{\displaystyle \frac{v_Q(r_Qv_Q)}{r_Q}}{\displaystyle \frac{r_Q}{4}}\}`$
$`={\displaystyle \frac{1}{2}}\underset{0j<i}{\mathrm{min}}\{r_Q(j{\displaystyle \frac{iv_Q}{r_Q}}+{\displaystyle \frac{1}{2}})^2+i^2{\displaystyle \frac{v_Q(r_Qv_Q)}{r_Q}}{\displaystyle \frac{r_Q}{4}}\}`$
$`={\displaystyle \frac{1}{2}}\underset{0j<i}{\mathrm{min}}\{(1+j)jr_Q+i(i12j)v_Q\},`$
with (2.2) we have
(2.9) $`B_ii^2B_1={\displaystyle \frac{1}{2}}{\displaystyle \underset{QI}{}}\underset{0j<i}{\mathrm{min}}\{(1+j)jr_Q+i(i12j)v_Q\}(i1).`$
Of course because $`dim_k𝒪_X/f_{}𝒪_Y(iE)=1+_{j=1}^{i1}d(j)`$, combining this with (2.8) and (2.9), we obtain (C).
Putting $`i=2`$ in (2.8) and (2.9), we have
$`d(1)=3{\displaystyle \underset{QI}{}}\mathrm{min}\{v_Q,r_Qv_Q\}.`$
Since $`dim_kf_{}𝒪_Y(2E)/𝔪_P^2=3d(1)`$, we get (D). ∎
### 2.3 Proof of $`𝐟_{}𝒪_𝐘(\mathrm{𝟐}𝐄)𝔪_𝐏`$
Assuming that $`f_{}𝒪_Y(2E)=𝔪_P`$, we will derive a contradiction. The assumption means that the coefficient of $`F_n`$ in $`g_{1n}^{}F_1`$ is bigger than $`1`$, so there exists a $`Z_i`$ which is contained in at least two $`g_{0i}`$-exceptional prime divisors on $`X_i`$. The minimum value of $`a`$ in this case occurs when $`Z_1`$ is a curve, $`Z_2=(g_2^1)_{}F_1F_2`$, and $`n=3`$, and the minimum value is $`6`$. So we get $`a6`$. By the assumption and (D), we obtain $`_{QI}\mathrm{min}\{v_Q,r_Qv_Q\}=3`$. Thus we have only to consider the three cases:
| Case 1. $`\{(r_Q,\overline{v_Q})\}_{QI}=\{(r,\overline{\pm 3})\},r7.`$ |
| --- |
| Case 2. $`\{(r_Q,\overline{v_Q})\}_{QI}=\{(r_1,\overline{\pm 1}),(r_2,\overline{\pm 2})\},r_12,r_25.`$ |
| Case 3. $`\{(r_Q,\overline{v_Q})\}_{QI}=\{(r_1,\overline{\pm 1}),(r_2,\overline{\pm 1}),(r_3,\overline{\pm 1})\},2r_1r_2r_3.`$ |
Here $`\pm `$ means that one of these occurs for each $`\overline{v_Q}`$. We remark that $`v_Q`$ is coprime to $`r_Q`$.
Since $`{\displaystyle \underset{QI}{}}{\displaystyle \frac{v_Q(r_Qv_Q)}{2r_Q}}<1`$ from (B), we have the below inequalities:
| Case 1. $`3/29/2r<1`$. |
| --- |
| Case 2. $`3/2(1/2r_1+2/r_2)<1`$. |
| Case 3. $`3/2(1/2r_1+1/2r_2+1/2r_3)<1`$. |
Using this evaluation, we can restrict possible values of $`r_Q`$’s. Below we show all the possible values and the corresponding values of $`aE^3`$:
| Case 1. | $`r`$ | : | $`7`$ | $`8`$ |
| --- | --- | --- | --- | --- |
| | $`aE^3`$ | : | $`2/7`$ | $`1/8`$ |
| Case 2. | $`(r_1,r_2)`$ | : | $`(2,5)`$ | $`(3,5)`$ | $`(4,5)`$ | $`(2,7)`$ |
| --- | --- | --- | --- | --- | --- | --- |
| | $`aE^3`$ | : | $`3/10`$ | $`2/15`$ | $`1/20`$ | $`1/14`$ |
| Case 3. | $`(r_1,r_2,r_3)`$ | : | $`(2,2,r_3)`$ | $`(2,3,3)`$ | $`(2,3,4)`$ | $`(2,3,5)`$ |
| --- | --- | --- | --- | --- | --- | --- |
| | $`aE^3`$ | : | $`2/2r_3`$ | $`1/6`$ | $`1/12`$ | $`1/30`$ |
Recalling that $`r`$ is the lowest common multiple of $`\{r_Q\}_{QI}`$, with (A) we have $`a3`$ for all the above cases. This contradicts $`a6`$. ∎
### 2.4 Proof of $`𝐟_{}𝒪_𝐘(\mathrm{𝐧𝐄})𝔪_𝐏^\mathrm{𝟐}`$
Because $`f_{}𝒪_Y(2E)𝔪_P`$, we have $`g_{1n}^{}F_1=_{i=1}^n(g_{in}^1)_{}F_i+(\mathrm{others})`$ and,
1. $`F_i(1im)`$ is obtained as a valuation by the weighted blow-up of $`X`$ with its weights $`(x,y,z)=(1,i,i)`$ for local parameters $`x,y,z`$ at $`P`$ such that $`Z_{m1}(g_{0,m1}^1)_{}\mathrm{div}(y)(g_{0,m1}^1)_{}\mathrm{div}(z)`$.
We divide the proof according to the value of $`dim_kf_{}𝒪_Y(2E)/𝔪_P^22`$.
| Case 1. | $`dim_kf_{}𝒪_Y(2E)/𝔪_P^2=0`$. |
| --- | --- |
| | This is the case when $`Z_1F_1`$ is neither a line nor a point. |
| Case 2. | $`dim_kf_{}𝒪_Y(2E)/𝔪_P^2=1`$. |
| | This is the case when $`Z_1F_1`$ is a line. |
| Case 3. | $`dim_kf_{}𝒪_Y(2E)/𝔪_P^2=2`$. |
| | This is the case when $`Z_1F_1`$ is a point. |
Since
$`dim_kf_{}𝒪_Y(2E)/𝔪_P^2`$ $`=dim_k\mathrm{Im}[(v𝔪_P|Z_1(g_1^1)_{}\mathrm{div}(v))𝔪_P/𝔪_P^2]`$
$`=dim_k\{v\mathrm{\Gamma }(F_1,𝒪_{F_1}(1))|v=0\mathrm{or}Z_1\mathrm{div}(v)\},`$
the value of $`dim_kf_{}𝒪_Y(2E)/𝔪_P^2`$ decides the type of $`Z_1F_1_k^2`$ as above.
In Case 1, $`_{QI}\mathrm{min}\{v_Q,r_Qv_Q\}=0`$ by (D). Therefore $`I`$ is empty and thus $`Y`$ is Gorenstein. By \[Cu88, Theorem 5\], $`f`$ must be the blow-up of $`X`$ along $`P`$, that is, $`f=g_1`$, and so we have nothing to do. Thus we have only to consider Cases 2 and 3. In these cases we investigate the values of $`dim_k𝒪_X/f_{}𝒪_Y(iE)`$’s carefully.
###### Proposition 2.8.
(Notation as above). Let $`2ln`$ be an integer such that $`g_{0i}(iF_i)𝔪_P^2`$ for any $`i<l`$.
(1) If $`g_{0l}(lF_l)𝔪_P^2`$, then
$`dim_k𝒪_X/f_{}𝒪_Y(lE)l{\displaystyle \frac{1}{2}}\underset{0j<l}{\mathrm{min}}\{((1+j)m2l)j\}.`$
(2) If $`g_{0l}(lF_l)𝔪_P^2`$ (in this case we have $`l>m`$ by $`()`$), then
$`dim_k𝒪_X/f_{}𝒪_Y(lE)>l{\displaystyle \frac{1}{2}}\underset{0j<l}{\mathrm{min}}\{((1+j)m2l)j\}.`$
###### Remark 2.8.1.
In the case when $`m=1`$ because
$`\underset{0j<l}{\mathrm{min}}\{((1+j)m2l)j\}=\underset{0j<l}{\mathrm{min}}\{(j(2l1))j\}=l(l1),`$
we can simplify the above inequalities:
(1) $`dim_k𝒪_X/f_{}𝒪_Y(lE){\displaystyle \frac{1}{2}}l(l+1).`$
(2) $`dim_k𝒪_X/f_{}𝒪_Y(lE)>{\displaystyle \frac{1}{2}}l(l+1).`$
###### Proof.
(1) By the assumption and $`f_{}𝒪_Y(2E)𝔪_P`$, the proof of Proposition 2.4 says that we can take local parameters $`x,y,z`$ at $`P`$ such that $`Z_{\mathrm{min}\{l,m\}1}(g_{0,\mathrm{min}\{l,m\}1}^1)_{}\mathrm{div}(y)`$ and $`Z_{l1}(g_{0,l1}^1)_{}\mathrm{div}(z)`$. Then for $`1il`$, $`F_i`$ equals, as valuations, the exceptional divisor obtained by the weighted blow-up of $`X`$ with its weights $`(x,y,z)=(1,\mathrm{min}\{i,m\},i)`$.
Hence
$`f_{}𝒪_Y(lE)=g_{0n}𝒪_{X_n}(lF_n)`$
$`g_{0l}𝒪_{X_l}(lF_l)=(x^sy^tz^u|s+\mathrm{min}\{l,m\}t+lul),`$
and so
$`dim_k𝒪_X/f_{}𝒪_Y(lE)`$ $`dim_k𝒪_X/(x^sy^tz^u|s+\mathrm{min}\{l,m\}t+lul)`$
$`=l{\displaystyle \frac{1}{2}}\underset{0j<l}{\mathrm{min}}\{((1+j)m2l)j\}.`$
Here we used Lemma 2.9 proved later.
(2) As in the proof of (1), we can take local parameters $`x,y,z`$ at $`P`$ such that $`Z_{m1}(g_{0,m1}^1)_{}\mathrm{div}(y)`$ and $`Z_{l2}(g_{0,l2}^1)_{}\mathrm{div}(z)`$. Then for $`1i<l`$, $`F_i`$ equals, as valuations, the exceptional divisor obtained by the weighted blow-up of $`X`$ with its weights $`(x,y,z)=(1,\mathrm{min}\{i,m\},i)`$.
We have
$`f_{}𝒪_Y(lE)=g_{0n}𝒪_{X_n}(lF_n)`$
$`g_{0,l1}𝒪_{X_{l1}}(lF_{l1})+(v𝔪_P|Z_{l1}(g_{0,l1}^1)_{}\mathrm{div}(v)).`$
But since
$`(v𝔪_P|Z_{l1}(g_{0,l1}^1)_{}\mathrm{div}(v))g_{0l}𝒪_{X_l}(lF_l)𝔪_P^2,`$
for any $`v𝔪_P`$ such that $`Z_{l1}(g_{0,l1}^1)_{}\mathrm{div}(v)`$ we have
$`g_{0,l1}^{}\mathrm{div}(v)g_{1,l1}^{}(2F_1+(g_1^1)_{}\mathrm{div}(v))(2+(l2))F_{l1}=lF_{l1}.`$
Thus
$`(v𝔪_P|Z_{l1}(g_{0,l1}^1)_{}\mathrm{div}(v))g_{0,l1}𝒪_{X_{l1}}(lF_{l1}),`$
and hence
$`f_{}𝒪_Y(lE)g_{0,l1}𝒪_{X_{l1}}(lF_{l1})=(x^sy^tz^u|s+mt+(l1)ul).`$
Therefore with Lemma 2.9,
$`dim_k𝒪_X/f_{}𝒪_Y(lE)`$ $`dim_k𝒪_X/(x^sy^tz^u|s+mt+(l1)ul)`$
$`>dim_k𝒪_X/(x^sy^tz^u|s+mt+lul)`$
$`=l{\displaystyle \frac{1}{2}}\underset{0j<l}{\mathrm{min}}\{((1+j)m2l)j\}.`$
We used the following lemma in the above proof.
###### Lemma 2.9.
Let $`XP`$ be an algebraic germ of a smooth $`3`$-dimensional variety with local parameters $`x,y,z`$ at $`P`$, and let $`l,m`$ be positive integers. Then
$`dim_k𝒪_X/(x^sy^tz^u|s+\mathrm{min}\{l,m\}t+lul)=l{\displaystyle \frac{1}{2}}\underset{0j<l}{\mathrm{min}}\{((1+j)m2l)j\}.`$
###### Proof.
$`dim_k𝒪_X/(x^sy^tz^u|s+\mathrm{min}\{l,m\}t+lul)`$
$`=dim_k\mathrm{Span}_kx^sy^t|s+\mathrm{min}\{l,m\}t<l`$
$`={\displaystyle \underset{0t\frac{l}{\mathrm{min}\{l,m\}}}{}}(l\mathrm{min}\{l,m\}t)`$
$`={\displaystyle \underset{0t\frac{l}{m}}{}}(lmt)`$
$`=l{\displaystyle \frac{m}{2}}\{({\displaystyle \frac{l}{m}}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{l}{m}})^2({\displaystyle \frac{1}{2}}{\displaystyle \frac{l}{m}})^2\}`$
$`=l{\displaystyle \frac{m}{2}}\underset{0j<l}{\mathrm{min}}\{(j+{\displaystyle \frac{1}{2}}{\displaystyle \frac{l}{m}})^2({\displaystyle \frac{1}{2}}{\displaystyle \frac{l}{m}})^2\}`$
$`=l{\displaystyle \frac{1}{2}}\underset{0j<l}{\mathrm{min}}\{((1+j)m2l)j\}.`$
Now we prove $`f_{}𝒪_Y(nE)𝔪_P^2`$ in Cases 2 and 3.
###### Proof in Case 2.
For $`Z_1F_1`$ is a line in this case and $`a`$ is the discrepancy of $`F_n`$ with respect to $`K_X`$, we get $`m=1`$ and
(2.10) $`a=n+1(n2).`$
By (D), $`_{QI}\mathrm{min}\{v_Q,r_Qv_Q\}=1`$ and thus $`\{(r_Q,\overline{v_Q})\}_{QI}=\{(r,\overline{\pm 1})\}`$. From (B), we obtain $`aE^3=(r+1)/r`$. By (A),
(2.11) $`ar+1.`$
From (C), Remark 2.7.2, and (2.10), for $`1in+1`$ we have
$`dim_k𝒪_X/f_{}𝒪_Y(iE)`$ $`=i^2{\displaystyle \frac{1}{2}}\underset{0j<i}{\mathrm{min}}\{(1+j)jr+i(i12j)\}`$
$`={\displaystyle \frac{1}{2}}i(i+1){\displaystyle \frac{1}{2}}\underset{0j<i}{\mathrm{min}}\{((1+j)r2i)j\}.`$
Hence for $`1in+1`$,
(2.12) $`dim_k𝒪_X/f_{}𝒪_Y(iE){\displaystyle \frac{1}{2}}i(i+1),`$
where the equality holds if and only if $`ir`$.
If there exists a positive integer $`2ln`$ such that $`g_{0l}𝒪_{X_l}(lF_l)𝔪_P^2`$ and $`g_{0i}𝒪_{X_i}(iF_i)𝔪_P^2`$ for any $`i<l`$, then by Proposition 2.8, Remark 2.8.1, and the condition of the equality in (2.12), we obtain $`l=r+1`$. Thus with (2.10), we have $`r+1=ln=a1`$, that is, $`ar+2`$. This contradicts (2.11) and hence we get $`g_{0n}𝒪_{X_n}(nF_n)𝔪_P^2`$. ∎
###### Proof in Case 3.
In this case we use essentially the same idea as in Case 2, but it is a little more complicated. By (D), $`_{QI}\mathrm{min}\{v_Q,r_Qv_Q\}=2`$. Thus we have only to consider the two subcases:
* Subcase 1. $`\{(r_Q,\overline{v_Q})\}_{QI}=\{(r,\overline{\pm 2})\},r5.`$
* Subcase 2. $`\{(r_Q,\overline{v_Q})\}_{QI}=\{(r_1,\overline{\pm 1}),(r_2,\overline{\pm 1})\},2r_1r_2.`$
In Subcase 1, we have $`aE^3=4/r`$ by (B) and thus $`a4`$ from (A). But since $`Z_1F_1`$ is a point, we get $`n=2`$ and $`a=4`$. Then choosing local parameters $`x,y,z`$ at $`P`$ such that $`Z_1(g_1^1)_{}\mathrm{div}(y)(g_1^1)_{}\mathrm{div}(z)`$, $`F_2`$ equals, as valuations, the exceptional divisor obtained by the weighted blow-up of $`X`$ with its weights $`(x,y,z)=(1,2,2)`$. So we have only to investigate Subcase 2.
Recalling that $`a`$ is the discrepancy of $`F_n`$ with respect to $`K_X`$, we have
(2.13) $`a=m+n(2mn).`$
Calculating with (B) we obtain $`aE^3=(r_1+r_2)/r_1r_2`$, and thus by (A),
(2.14) $`ar_1+r_2.`$
From (C), Remark 2.7.2, and (2.13), for $`1im+n`$ we have
$`dim_k𝒪_X/f_{}𝒪_Y(iE)`$
$`=i^2{\displaystyle \frac{1}{2}}\underset{0j<i}{\mathrm{min}}\{(1+j)jr_1+i(i12j)\}`$
$`{\displaystyle \frac{1}{2}}\underset{0j<i}{\mathrm{min}}\{(1+j)jr_2+i(i12j)\}`$
$`=i{\displaystyle \frac{1}{2}}\left(\underset{0j<i}{\mathrm{min}}\{((1+j)r_12i)j\}+\underset{0j<i}{\mathrm{min}}\{((1+j)r_22i)j\}\right).`$
Hence for $`1im+n`$,
(2.15) $`dim_k𝒪_X/f_{}𝒪_Y(iE)`$ $`i{\displaystyle \frac{1}{2}}\underset{0j<i}{\mathrm{min}}\{((1+j)r_12i)j\}`$
$`i,`$
where the equality of the first inequality holds if and only if $`ir_2`$, and the second holds if and only if $`ir_1`$.
###### Claim 2.10.
$`r_1=m`$.
###### Proof of the claim.
Utilizing Proposition 2.8 (1) with $`l=m`$, we have
(2.16) $`dim_k𝒪_X/f_{}𝒪_Y(mE)m{\displaystyle \frac{1}{2}}\underset{0j<m}{\mathrm{min}}\{(j1)jm\}=m.`$
We take local parameters $`x,y,z`$ at $`P`$ as in $`()`$, satisfying $`Z_m(g_{0m}^1)_{}\mathrm{div}(z)`$ if $`Z_mF_m_k^2`$ is a line. We have
$`f_{}𝒪_Y((m+1)E)=g_{0n}𝒪_{X_n}((m+1)F_n)`$
$`g_{0m}𝒪_{X_m}((m+1)F_m)+(v𝔪_P|Z_m(g_{0m}^1)_{}\mathrm{div}(v)).`$
But since
$`(v𝔪_P|Z_m(g_{0m}^1)_{}\mathrm{div}(v))g_{0m}𝒪_{X_m}(mF_m)=(x^m,y,z),`$
we get
$`(v𝔪_P|Z_m(g_{0m}^1)_{}\mathrm{div}(v))(z)+g_{0m}𝒪_{X_m}((m+1)F_m),`$
and thus
$`f_{}𝒪_Y((m+1)E)`$ $`(z)+g_{0m}𝒪_{X_m}((m+1)F_m)`$
$`=(z)+(x^sy^tz^u|s+mt+mum+1).`$
Hence,
(2.17) $`dim_k𝒪_X/f_{}𝒪_Y((m+1)E)`$
$`dim_k𝒪_X/((z)+(x^sy^tz^u|s+mt+mum+1))`$
$`=dim_k\mathrm{Span}_kx^s,y|sm`$
$`=m+2.`$
From (2.16), (2.17), and the condition of the second equality in (2.15), we have $`r_1=m`$. ∎
If there exists a positive integer $`ln`$ such that $`g_{0l}𝒪_{X_l}(lF_l)𝔪_P^2`$ and $`g_{0i}𝒪_{X_i}(iF_i)𝔪_P^2`$ for any $`i<l`$, then by Proposition 2.8, Claim 2.10, and the condition of the first equality in (2.15), we obtain $`l=r_2+1`$. Thus with (2.13) and Claim 2.10, we have $`r_1+r_2+1=m+lm+n=a`$. This contradicts (2.14) and hence we get $`g_{0n}𝒪_{X_n}(nF_n)𝔪_P^2`$. ∎
Department of Mathematical Sciences, University of Tokyo, Komaba, Meguro, Tokyo 153-8914, Japan
kawakita@ms.u-tokyo.ac.jp
|
warning/0005/astro-ph0005572.html
|
ar5iv
|
text
|
# 1 Introduction
## 1 Introduction
Though the number of secure associations between NSs and SNRs continues to grow (Caraveo 1993, 1995; Allakhverdiyev et al. 1995, Kaspi 1996, 1998, 2000; Frail 1998; Helfand 1998, Mereghetti 1998, 1999; Marsden et al. 1999), it is considered that many of claimed associations are spurious and are merely the results of superposition (e.g. Gaensler & Johnston 1995a,b; cf. Lorimer, Lyne & Camilo 1998). It were proposed five criteria for the evaluation of possible NS/SNR associations which come to the following questions (Kaspi 1996):
– do independent distance estimates agree?
– do independent age estimates agree?
– is the implied transverse velocity reasonable?
– is there evidense for any interaction between the NS and SNR?
– does the proper motion vector of the NS point away from the SNR centre?
The last question is considered the most important one since ”a proper motion measurement has the potential to disprove an association regardless of the answers to the other questions” (Kaspi 1996).
In this Letter we point out that a cavity SN explosion of a moving massive star could result in a significant offset of the NS birth-place from the geometrical centre of the SNR (Sect. 2). Three important consequences can be drawn from this: 1) the implied transverse velocity of the NS (i.e. the velocity derived from the displacement of the NS from the geometrical centre of the SNR) could be significantly overestimated; 2) the proper motion vector of the NS should not necessarily point away from the geometrical centre of the associated SNR (it even could be directed to the centre of the SNR!); 3) the circle of possible NS/SNR associations could be enlarged (Sect. 3). These facts are quite obvious, but have been largely overlooked in studies of NS/SNR associations. It is also suggested that the birth-place of a NS could be marked by a nebula of thermal X-ray emission (Sect. 3). The possible detection of such nebulae will allow to get more reliable estimates of transverse velocities of NSs associated with middle-aged SNRs.
## 2 Off-centred cavity SN explosion
Massive stars (the progenitors of most of SN stars; e.g. van den Berg & Tammann 1991; Tammann, Löffler & Schröder 1994) strongly modify, during their evolution from the main-sequence (MS) to the SN explosion, the ambient medium by ionizing emission and winds, that results in the origin of a system of cavities and shells (Avedisova 1972; Dyson & de Vries 1972; Dyson 1975; Castor, McCray & Weaver 1975; Steigmann, Strittmatter & Williams 1975; Weaver et al. 1977; McCray 1983; McKee, Van Buren & Lazareff 1984; D’Ercole 1992). The subsequent interaction of the SN blast wave with the reprocessed circumstellar and interstellar medium results in the origin of a SNR (e.g. Fabian, Brinkmann & Stewart 1983; Shull et al. 1985; McKee 1988; Ciotti & D’Ercole 1989; Chevalier & Liang 1989; Chevalier & Emmering 1989; Franco et al. 1991; McCray 1993; Brighenti & D’Ercole 1994). The structure of a young ($`10^3`$ years) SNR is mostly determined by the interaction of the SN blast wave with circumstellar structures created during the late evolutionary stages of the SN progenitor star (e.g. McCray 1993; Garsia-Segura, Langer & Mac Low 1996; Borkowski et al. 1996), while the appearance of a middle-aged SNR could be affected by the interaction of the SN blast wave with large-scale structures created in the interstellar medium by the stellar ionizing emission (the shell of neutral gas around a Strömgren sphere; e.g. Shull et al. 1985) and/or (fast) stellar wind (the density jump at the edge of a stalled wind-driven bubble or the dense large-scale shell swept-up from the interstellar medium by an expanding bubble; e.g. Ciotti & D’Ercole 1989; Franco et al. 1991; D’Ercole 1992; Gvaramadze 1999a,b).
It is clear that the SN explodes in the centre of the system of cavities and shells if the space velocity of the SN progenitor star is equal to zero. In this case the birth-place of the SN stellar remnant (e.g. a NS) coincides with the geometrical centre of the future SNR, and therefore the implied transverse velocity of the stellar remnant is equal to the true one.
But the SN explosion site could be significantly offset from the geometrical centre of large-scale structures created in the reprocessed ambient medium if the massive star moves relative to the interstellar medium (it is known that most of massive stars have a space velocity of a few $`\mathrm{km}\mathrm{s}^1`$; e.g. Vanbeveren, De Loore & Van Rensbergen 1998). E.g. a 25 $`M_{}`$ star moving with the velocity of 2 $`\mathrm{km}\mathrm{s}^1`$ travels from the centre of the MS bubble for about 18 parsecs. The stellar motion does not affect the spherical shape of the bubble (here and below we assume that there is no large-scale density inhomogeneities in the ambient interstellar medium) since the sound speed in the hot interior of the bubble is about two orders of magnitude larger than the velocity of the star (e.g. Weaver et al. 1977). Correspondingly, the (middle-aged) SNR also acquires the spherical shape even if the SN exploded far from the centre of the wind-driven bubble (e.g. Różychka et al. 1993). The off-centred cavity SN explosion results, however, in the inhomogeneous distribution of the surface brightness over the SNR’s shell, that could explain the arc-like appearance of some of middle-aged SNRs (see Różychka et al. 1993, Brighenti & D’Ercole 1994). Note that the standard explanation of the origin of incomplete shells implies the interaction of the (Sedov-Taylor) blast wave with the inhomogeneous (e.g. cloudy) interstellar medium. Therefore, the non-detection of a dense large-scale cloud nearby to the arc-like SNR could serve as an indirect evidence that this SNR is generated by a moving massive star (cf. Brighenti & D’Ercole 1994). The anonymous referee mentioned that ”the direction of motion of the NS ought to depend on its location relative to the bright part of the shell”. This is correct, however, only in the absence of other factors affecting the brightness distribution over the SNR’s shell. The existence of large-scale density gradients in the interstellar medium and/or magnetized wind-driven shells leads to the more complex appearance of SNRs (see e.g. Gvaramadze 1999b), and therefore the enhanced brightness of a part of the shell not necessarily ought to be due to the proximity of the SN explosion site. The detailed study of this important issue is beyond the scope of this Letter and will be carried out elsewhere.
It should be mentioned that only a fraction of middle-aged and old SNRs is the result of cavity SN explosions. Indeed, one can show (e.g. Brighenti & D’Ercole 1994) that only slowly moving ($`12\mathrm{km}\mathrm{s}^1`$) and/or very massive stars explode inside the wind-driven bubbles created during the MS stage of their evolution. But even if a massive star is fast enough to cross the stalled MS bubble it could again find itself in the bubble interior if it ends its evolution as a red supergiant (RSG) star (i.e. if the zero age main sequence mass of the star is $`1520\mathrm{M}_{}`$; e.g. Vanbeveren et al. 1998). During the RSG stage the stalled bubble can re-expand (D’Ercole 1992) and catch up the moving star, provided that there are no external sources of ionizing emission. The high-velocity massive stars also could explode inside the large-scale bubbles, but this happens only for stars whose mass $`1520\mathrm{M}_{}`$. In this case, a massive star before it exploded as a SN becomes a Wolf-Rayet (WR) star (e.g. Vanbeveren et al. 1998), whose fast wind blows up a new large-scale bubble surrounded by a dense shell. And again, the proper motion of the WR star results in the significant offset of the SN explosion site from the centre of the wind-driven bubble. E.g. for the duration of the WR stage of $`23\times 10^5`$ yr and the stellar velocity of $`30\mathrm{km}\mathrm{s}^1`$ the SN explosion site will be displaced from the centre of the bubble for about $`69`$ pc (see Arnal 1992 for examples of non-central location of WR stars in bubbles created by their winds). On the other hand, most of massive stars are less massive than $`15M_{}`$, and therefore do not evolve through the WR stage before the SN explosion. Thus, if a massive star of mass $`<15M_{}`$ is fast enough to cross the MS bubble before it exploded as a SN, the SN blast wave mainly interacts with the unperturbed interstellar medium and the SN explosion site coincides with the geometrical centre of the resulting (middle-aged) SNR. In this case the shell of the SNR could appear as an incomplete circle, that is due to the presence of a low-density tunnel created by the stellar wind behind the moving star (see Brighenti & D’Ercole 1994).
It follows from the above discussion that though many of SNRs are produced by SN explosions outside of large-scale wind-driven bubbles (and therefore their structure could be discribed in the framework of the standard model based on the Sedov-Taylor solution), there are could exist SNRs whose origin is connected with off-centred cavity SN explosions<sup>2</sup><sup>2</sup>2Though the knowledge of the fraction of these SNRs is very important for statistical studies of NS/SNR associations, we are now not in a position to quantify it. Two main difficulties on the way to do this are the absence of self-consisting evolutionary models for rotating stars (the evolutionary paths of rotating stars differ from those of non-rotating ones) and wind-driven bubbles (taking into account the heat conduction and magnetic effects). These difficulties along with uncertainties in the velocity distribution and initial mass function of massive stars do not allow to solve the problem properly.. Therefore the high transverse velocities inferred for a number of NSs through their association with SNRs could be reduced (see next Sect.). The velocity reduction could be large enough for NSs associated with middle-aged and old SNRs, i.e. the SNRs whose origin could be connected with the interaction of SN blast waves with large-scale structures created by the fast stellar wind, and where the SN explosion sites could be significantly offset from the centres of the wind-driven bubbles. But the velocity reduction should be less considerable in the case of young ($`<10^3`$ years) SNRs, whose appearance is mostly determined by the interaction of SN blast waves with circumstellar (i.e. small-scale) structures created during the late (RSG and WR) evolutionary stages of SN progenitor stars (see Sect. 3). These stages are significantly shorter than the MS stage and only sufficiently fast stars have time to became significantly displaced from the centres of associated circumstellar structures. Therefore, the geometrical centres of young SNRs better correspond to the SN explosion sites, that explains the quite symmetric appearance of these SNRs (a marked exception is the Kepler’s SNR, whose asymmetric shell is due to the very fast motion of the SN progenitor star (Bandiera 1987)). But even the slow motion of the SN progenitor star results in an appreciable asymmetry of circumstellar structures (e.g. the mass distribution over the nearly spherical circumstellar shell becomes inhomogeneous), that, for instance, results in the asymmetric expansion of young SNRs. A good example of such a young SNR is the Cas A. A compact X-ray source was recently discovered near the geometrical centre of this SNR (Tananbaum 1999). The implied transverse velocity of the compact source derived through the various determinations of the expansion centre of Cas A (e.g. van den Berg & Kamper 1983, Reed et al. 1995) ranges from $`50`$ to $`1000\mathrm{km}\mathrm{s}^1`$ (Pavlov et al. 2000). We suggest that the separation of the SN explosion site (and the compact source) from the geometrical centre of Cas A could be caused to a large extent by the proper motion of the SN progenitor star. This motion (with the velocity less than the velocity of the RSG wind) will result in the deviation from spherical (or axial) symmetry of the circumstellar matter, that in its turn could be responsible for the observed (see Vink et al. 1998, and references therein) expansion asymmetry of Cas A.
## 3 Discussion
In Sect. 2 we showed that the large displacement of a NS from the geometrical centre of the associated SNR does not inevitably mean that this NS is moving with high transverse velocity. This could have an important impact on the understanding of the origin of the phenomenon of anomalous X-ray pulsars and soft gamma-ray repeaters (SGRs) since the high implied velocities of some of these objects were interpreted as a sign that they represent a high-velocity ($`1000\mathrm{km}\mathrm{s}^1`$) population of NSs (e.g. Thompson & Duncan 1995; Marsden et al. 1999). It seems that the recent association of two of known SGRs with clusters of massive stars (Fuchs et al. 1999, Vrba et al. 2000) should reduce the acuteness of the problem of high implied velocities of these objects, but the large angular offset of the SGR 0525-66 from the centre of the nearly spherical (in X-rays) SNR N 49 in LMC (e.g. Rothschild, Kulkarni & Lingenfelter 1994) still continues to raise doubts in the association of these two objects (e.g. Kaspi 2000, Kaplan et al. 2001). Note also that the high transverse velocities derived by Frail, Goss & Whiteoak (1994) for pulsars associated with SNRs were used to put forward a number of quite strong suggestions, e.g. that ”SNRs are produced preferentially by the (SN) explosions that yield fast kicks” (Cordes & Chernoff 1998).
The high implied transverse velocities of NSs are sometimes used to discard the possible NS/SNR associations. E.g. Stappers, Gaensler & Johnston (1999) suggested that the lack of a pulsar wind radio nebula around the PSR B 1610-50 means that the maximum space velocity $`v_\mathrm{p}`$ of this pulsar is $`450(d/5\mathrm{kpc})\mathrm{km}\mathrm{s}^1`$, where $`d`$ is the distance to the pulsar, and therefore it could not be associated with the nearby SNR Kes 32 since this association implies the transverse velocity of the pulsar of $`2000\mathrm{km}\mathrm{s}^1`$ (Caraveo 1993). The implied transverse velocity, however, could be reduced two times simply due to the possible off-centred SN explosion, and once again two or even more times if the braking index of the pulsar is similar respectively to that of the PSR B 0540-69 (Boyd et al. 1995) or the Vela pulsar (Lyne et al. 1996).
The association of PSR B 1610-50 with SNR Kes 32 was also recently questioned by Pivovaroff, Kaspi & Gotthelf (2000). They used the non-detection of an X-ray nebula around the PSR B 1610-50 to estimate $`v_\mathrm{p}`$ to be less than $`170(d/7.3\mathrm{kpc})^2(n/1\mathrm{cm}^3)^{1/2}\mathrm{km}\mathrm{s}^1`$. This estimate was derived under the assumption that the wind of the PSR B 1610-50 has the same characteristics as that of the Crab pulsar, and for the number density of the ambient medium $`n=1\mathrm{cm}^3`$. One can, however, show that reasonable variations of the assumed parameters allow to increase the estimated velocity of the pulsar. E.g. for $`n10^2\mathrm{cm}^3`$ and $`d=5`$ kpc (Stappers et al. 1999), one has $`v_\mathrm{p}780\mathrm{km}\mathrm{s}^1`$.
The high transverse velocity was also inferred for the pulsar PSR B 1757-24, which lies well outside the shell of the SNR G 5.4-1.2 (e.g. Caswell et al. 1987). The physical association of these two objects was firmly established after the discovery (Frail & Kulkarni 1991; see also Manchester et al. 1991) of a tail of radio emission connecting the pulsar with the SNR. The pulsar PSR B 1757-24 is, however, more interesting in that that its proper motion vector does not point away from the geometrical centre of the nearly circular shell of the remnant (the radius of which is about 16 arcmin), but misses it by nearly 5 arcmin (Frail, Kassim & Weiler 1994). To explain this inconsistency, Frail et al. (1994) suggested that the SN exploded in an exponentially stratified medium and used the Kompaneets (1960) solution to fit the shape of the remnant. This allowed them to put the possible SN explosion site closer to the present position of the pulsar, that reduces the implied velocity of the pulsar to the value between 1300 and 1700 $`\mathrm{km}\mathrm{s}^1`$. We propose an alternative explanation and suggest that the SNR G 5.4-1.2 is the result of the off-centred SN explosion in the pre-existing wind-driven bubble surrounded by a massive shell (Gvaramadze, in preparation)<sup>3</sup><sup>3</sup>3The mass of the shell is a very important parameter since it determines the evolution of the SNR. If the mass of the shell is larger than about 50 times the mass of the SN ejecta (e.g. Franco et al. 1991), the SN blast wave merges with the shell and evolves into a momentum-conserving stage (i.e. it skips the Sedov-Taylor stage). In this case, even a young NS moving with a moderate velocity ($`200\mathrm{km}\mathrm{s}^1`$) is able to overrun the SNR’s shell (cf. Gaensler & Johnston 1995a), provided that it was born not far from the edge of the wind-driven bubble.. This suggestion allows to reduce considerably the transverse velocity of the pulsar and naturally explains why the tail behind the pulsar does not point back to the centre of the remnant. An indirect support of our suggestion comes from the recent observations of the radio nebula surrounding the pulsar (Gaensler & Frail 2000). These observations put an upper limit on the pulsar proper motion, which turns out to be much smaller than that expected if the pulsar velocity is indeed in the range derived by Frail et al. (1994). Assuming that the pulsar was born in the geometrical centre of the associated SNR, Gaensler & Frail (2000) argued that the true age of the pulsar should be more than 10 times larger than the characteristic age (cf. Istomin 1994). The off-centred cavity SN explosion provides another possible explanation for the low value of the pulsar proper motion.
Let us discuss the third criterion for the evaluation of NS/SNR associations suggested by Kaspi (1996). There are three factors which could affect the estimates of the characteristic age of a NS (e.g. Camilo 1996). First, the braking index could be different from that follows from the simplest spindown models (see e.g. Lyne et al. 1993, Kaspi et al. 1997, Marshall et al. 1998). Second, the NS could be born with the large initial spin period (e.g. Spruit & Phinney 1998), and therefore the true age could be much smaller than the characteristic one. Third, the spindown torque (as well as the braking index) could be a function of time. E.g. if the spindown of a NS is due to the magnetic dipole radiation, than the secular increase of the magnetic moment of the NS results in the increase of the braking torque (e.g. Blandford & Romani 1988). The spindown rate of a NS could be also enhanced due to the interaction of its magnetosphere with the dense ambient medium (Istomin 1994, Yusifov et al. 1995, Gvaramadze 1999c, 2001, Menou, Perna & Hernquist 2001). In both cases the true age of the NS could be much larger than the characteristic age derived from the present value of the spin period derivative. These arguments were used to reconcile the ages of the pulsar PSR B 1509-58 and the associated SNR MSH 15-52 (Blandford & Romani 1988, Gvaramadze 1999c, 2001), or to show that the implied transverse velocity of the pulsar PSR B 1757-24 could be reduced (Istomin 1994).
In Gvaramadze (1999c, 2001) we suggested that the high spin-down rate of the pulsar PSR B 1509-58 is inherent only for a relatively short period of its present spin history, and that the enhanced braking torque is caused by the interaction of the pulsar’s magnetosphere with the material of a dense circumstellar clump created during the late stages of evolution of the SN progenitor star. The origin of dense circumstellar clumps could be explained in the framework of the three-wind model (e.g. Garsia-Segura et al. 1996). The fast (WR) wind sweeps the slow (RSG) wind and creates a low-density cavity surrounded by a shell of swept-up circumstellar matter. This shell expands with the nearly constant velocity $`v_{\mathrm{sh}}(\dot{M}_{\mathrm{WR}}v_{\mathrm{WR}}^2v_{\mathrm{RSG}}/3\dot{M}_{\mathrm{RSG}})^{1/3}`$, where $`\dot{M}_{\mathrm{WR}},\dot{M}_{\mathrm{RSG}}`$ and $`v_{\mathrm{WR}},v_{\mathrm{RSG}}`$ are, correspondingly, the mass-loss rates and wind velocities during the WR and RSG stages (e.g. Dyson 1981), until it catches up the shell separating the RSG wind from the MS bubble (the characteristic radius of this shell is a few pc; the high-pressure gas of the MS bubble interior hinders the free expansion of the RSG wind (e.g. Chevalier & Emmering 1989, D’Ercole 1992)). The interaction of two circumstellar shells results in the Rayleigh-Taylor and other dynamical instabilities, whose development is accompanied by the formation of dense clumps moving with radial velocities of $`v_{\mathrm{cl}}v_{\mathrm{sh}}`$ (Garsia-Segura et al. 1996). For parameters typical for RSG and WR winds, one has $`v_{\mathrm{cl}}100200\mathrm{km}\mathrm{s}^1`$. The radial velocity of quasi-stationary flocculy in Cas A (whose origin could be attributed to the processes discussed above; e.g. Garsia-Segura et al. 1996) ranges from $`80`$ to $`400\mathrm{km}\mathrm{s}^1`$. The dense clumps could originate much closer to the SN progenitor star due to the stellar wind acceleration during the transition from the RSG to the WR stage (Brighenti & D’Ercole 1997). After the SN exploded, the SN blast wave propagates through the tenuous interclump medium, leaving behind the dense clumps embedded in the hot shocked interclump gas (the filling factor of the clumps is small (e.g. Gvaramadze 2001) and therefore they do not affect considerably the dynamics of the SN blast wave). The gradual evaporation of the dense material of radially moving clumps results in the origin of an expanding nebula of thermal X-ray emission, which marks the SN explosion site.
It is clear from the above discussion that the nebulae of thermal X-ray emission should exist only in those SNRs, whose origin is connected with explosions of massive stars with zero age main sequence mass $`1520M_{}`$ (only in these cases one can expect that the (clumpy) circumstellar material will survive the passage of the SN blast wave). It is clear also that the motion of the SN progenitor star could result in a significant offset of the compact region of dense circumstellar matter from the centre of the MS bubble (the RSG and WR stages are about 10-20 times shorter than the MS stage), and correspondingly in the offset of the nebula of thermal X-ray emission from the geometrical centre of the associated middle-aged SNR. The possible detection of such nebulae will provide the direct observational test for our proposal, and could be used for the re-estimation of transverse velocities of the already known NSs, or for the search of new stellar remnants possibly associated with these SNRs.
To find a crude order of magnitude estimate of the luminocity of the nebula of thermal X-ray emission, we assume that all the X-ray emitting interclump material (including the gas already evaporated from the clumps) is at the same temperature between $`10^7`$ and $`10^8`$ K, and that this material is uniformly dispersed over a sphere of radius $`R=R_0+v_{\mathrm{cl}}t_{\mathrm{SNR}}`$, where $`R_0(12`$ pc) is the radius of the region occupied by dence clumps at the moment of SN explosion, $`t_{\mathrm{SNR}}`$ is the age of the SNR, then one has $`L_\mathrm{x}1.2\times 10^{33}n^2R_{\mathrm{pc}}^3\mathrm{ergs}\mathrm{s}^1`$, where $`n`$ is the number density of the emitting gas, $`R_{\mathrm{pc}}=R/1\mathrm{pc}`$ (Gorenstein & Tucker 1976). For $`R_0=2`$ pc, $`v_{\mathrm{cl}}=100\mathrm{km}\mathrm{s}^1`$ and $`t_{\mathrm{SNR}}=10^4`$ years, and assuming that the mass of the emitting gas is $`5M_{}`$ (i.e. about a half of the mass lost by a $`15M_{}`$ star during the RSG stage), one has $`R3`$ pc and $`L_\mathrm{x}4\times 10^{34}\mathrm{ergs}\mathrm{s}^1`$. The similar estimates were used by Gvaramadze (1999a) to show that the nebula of hard X-ray emission found by Willmore et al. (1992) around the Vela pulsar could be the dense material lost by the SN progenitor star in the form of the RSG wind and heated to the observed temperature after the SN exploded. The more detailed analysis of this problem constitutes a part of a project underway to study the origin of mixed-morphology SNRs (Rho & Petre 1998) and will be published elsewhere (for a different point of view see e.g. White & Long 1991 and Petruk 2000).
In conclusion we note that in Gvaramadze (1999c) we interpreted a bright X-ray spot (which nearly coincides with the error box for the SGR 0525-66; Rothschild et al. 1994) on the periphery of the SNR N 49 as an X-ray nebula marking the SN explosion site and suggested that the large implied transverse velocity of the NS associated with the SGR could be reduced about ten times. In our analysis we assumed that the spot is a thermal feature (cf. Dickel et al. 1995) and that the radius of the spot is about $`5^{\prime \prime }10^{\prime \prime }`$ (i.e. $`12`$ pc; see Rothschild et al. 1994 and Dickel et al. 1995). We found that to explain the X-ray luminocity of the spot of $`10^{36}\mathrm{ergs}\mathrm{s}^1`$ (Rothschild et al. 1994), the mass of the X-ray emitting gas should be in a range $`410M_{}`$ (i.e. a reasonable value, provided that the zero age main sequence mass of the SN progenitor star was $`15M_{}`$). However, recent high-resolution Chandra X-ray Observatory observations of the X-ray spot in N 49 (Kaplan et al. 2001 and references therein) showed that this source is pointlike and could be considered as the X-ray counterpart of SGR 0525-66. Though this result discards our interpretation of the spot as an X-ray nebula, we believe that the large angular displacement of SGR 0525-66 from the centre of N 49 is due to the effect discussed in this Letter.
## 4 Summary
A cavity SN explosion of a moving massive star could result in a significant offset of the NS birth-place from the centre of the nearly spherical middle-aged SNR. Therefore: a) the high transverse velocities inferred for a number of NSs (e.g. PSR B 1610-50, PSR B 1757-24, SGR 0525-66) through their association with SNRs could be reduced; b) the proper motion vector of the NS should not necessarily point away from the geometrical centre of the associated SNR. These two facts allow to enlarge the circle of possible NS/SNR associations and should be taken into account in evaluating of their reliability. The birth-place of the NS could be marked by a (compact) nebula of thermal X-ray emission. The discovery of such nebulae in middle-aged SNRs could be used for the re-estimation of transverse velocities of the already known NSs, or for the search of new stellar remnants possibly associated with these SNRs.
Acknowledgements. I am grateful to N.D’Amico and A.D’Ercole for their interest to this work and to the anonymous referee for her/his questions and comments allowing me to clarify some points discussed in the Letter.
|
warning/0005/hep-th0005138.html
|
ar5iv
|
text
|
# 1 Introduction
## 1 Introduction
Noncommutative geometry deals with functions on deformation of ordinary space, such that coordinates on it do not commute<sup>2</sup><sup>2</sup>2In the sequel we use the same notation $`[,]`$ both for ordinary and for star-commutator. To avoid confusions, we supply all noncommutative quantities with the hats.:
$$\begin{array}{c}[\widehat{x}_\mu ,\widehat{x}_\nu ]=2\pi i\theta _{\mu \nu },\mu ,\nu =1,\mathrm{}d\end{array}$$
(1.1)
The antisymmetrical tensor $`\theta _{\mu \nu }`$ is called noncommutativity parameter. Such deformed flat ($`\theta _{\mu \nu }=const`$) and compact space is called noncommutative (quantum) torus $`𝐓_\theta ^d`$. In the last few years noncommutative geometry, and especially the noncommutative torus has been realized to play an important role in compactifications of M-theory and in string theory (see and references therein). It also turned to be very useful in compactification of instanton’s moduli spaces . The way to deal with the curved quantum spaces is provided by the Kontsevich’s deformation quantization.
A very intriguing subject from noncommutative geometry is so-called Morita equivalence . Roughly speaking, it states that certain bundles on different noncommutative tori are dual to each other. From the physical point of view it results in equivalence between certain noncommutative and ordinary gauge theories. In what follows we try to clarify this statement using a set of simple examples.
## 2 Notations
The algebra $`𝒜_\theta `$ of smooth functions on the noncommutative torus is defined using the Moyal star product:
$$\begin{array}{c}fg(\widehat{𝐱})=e^{i\pi \theta _{\mu \nu }\frac{}{\xi _\mu }\frac{}{\eta _\nu }}f(\xi )g(\eta )|_{\xi =\eta =\widehat{𝐱}}\end{array}$$
(2.1)
The main property of this product is associativity. In applications it is useful to decompose functions on noncommutative torus into the Fourier components<sup>3</sup><sup>3</sup>3Without loss of generality we can consider a torus of size $`2\pi `$. :
$$\begin{array}{c}f(\widehat{𝐱})=_{𝐤\text{Z}\text{Z}^d}f_𝐤e^{i𝐤\widehat{𝐱}}\end{array}$$
(2.2)
This corresponds to the Weil or symmetric ordering of coordinates. Exponents $`\widehat{U}_𝐤=e^{i𝐤\widehat{𝐱}}`$ may serve as a basis elements for the algebra $`𝒜_\theta `$.
A very intriguing thing happens when components of the $`\theta `$tensor becomes rational. Let us first consider two-torus $`𝐓^2`$:
$$\begin{array}{c}[\widehat{x}_\mu ,\widehat{x}_\nu ]=2\pi i\theta ϵ_{\mu \nu },\mu ,\nu =1,2\end{array}$$
(2.3)
with the rational noncommutativity parameter $`\theta =\frac{M}{N}`$, where $`M`$ and $`N`$ are relatively prime integers. Then
$$\begin{array}{c}[\widehat{U}_𝐧,\widehat{U}_𝐧^{}]=2i\mathrm{sin}\left(\pi M\frac{n_2n_1^{}n_1n_2^{}}{N}\right)\widehat{U}_{𝐧+𝐧^{}}=2i\mathrm{sin}(𝐧\times 𝐧^{})\widehat{U}_{𝐧+𝐧^{}}\end{array}$$
(2.4)
where by definition, $`𝐧\times 𝐧^{}\pi \theta _{\mu \nu }n_\mu n_\nu ^{}`$. Note that elements $`\widehat{U}_{N𝐤}`$ generate a center of the $`𝒜_\theta `$, that is for any $`f(\widehat{𝐱})`$:
$$\begin{array}{c}[e^{iN𝐤\widehat{𝐱}},f(\widehat{𝐱})]=0\end{array}$$
(2.5)
This means that one can treat exponents $`\{\widehat{U}_𝐤,𝐤=0|_{modN}\}`$ in the decomposition (2.2) as if they are ordinary exponents defined on ordinary (commutative) space. Other $`N^21`$ exponents, obtained from the set $`\{\widehat{U}_𝐤,𝐤0|_{modN}\}`$ after factorization over commutative part, generates closed algebra under star-commutator. This algebra is isomorphic to the algebra of $`SU(N)`$, as we will see in a moment. Therefore, at the rational value of the noncommutativity parameter one can identify algebra of functions on the noncommutative torus with the algebra of matrix-valued functions on commutative torus.
We conclude this section by giving an explicit matrix representation for the noncommutative exponents algebra (see also ). Such a representation has been indeed well-known for many years . Let us introduce the following clock and shift generators
$$\begin{array}{c}Q=\left(\begin{array}{ccccc}1\hfill & & & & \\ & \omega \hfill & & & \\ & & \omega ^2\hfill & & \\ & & & \mathrm{}\hfill & \\ & & & & \omega ^{N1}\hfill \end{array}\right)P=\left(\begin{array}{cccccc}0\hfill & 1\hfill & & & & 0\hfill \\ & 0\hfill & 1\hfill & & & \\ & & \mathrm{}\hfill & \mathrm{}\hfill & & \\ & & & \mathrm{}\hfill & & 1\hfill \\ 1\hfill & & & & & 0\hfill \end{array}\right)\end{array}$$
(2.6)
where $`\omega =e^{2\pi i\theta }`$. Matrices $`P`$ and $`Q`$ are unitary, traceless and satisfy:
$$\begin{array}{c}P^N=Q^N=\mathrm{𝟏},PQ=\omega QP\end{array}$$
(2.7)
Moreover,
$$\begin{array}{c}Tr(P^nQ^m)=\{\begin{array}{c}N,ifn=0|_{modN}andm=0|_{modN}\hfill \\ 0,ifn0|_{modN}orm0|_{modN}\hfill \end{array}\end{array}$$
(2.8)
It is straightforward to check that the generators, defined as
$$\begin{array}{c}J_𝐧=\omega ^{\frac{n_1n_2}{2}}Q^{n_1}P^{n_2},𝐧=(n_1,n_2)\end{array}$$
(2.9)
satisfy commutation relations (2.4):
$$\begin{array}{c}[J_𝐧,J_𝐧^{}]=2i\mathrm{sin}\left(𝐧\times 𝐧^{}\right)J_{𝐧+𝐧^{}}\end{array}$$
(2.10)
This identity can be tautologically rewritten in the form of the Lie algebra commutation relations:
$$\begin{array}{c}[J_𝐧,J_𝐦]=f_{\mathrm{𝐧𝐦}}^𝐤J_𝐤,\end{array}$$
(2.11)
where the structure constants $`f_{\mathrm{𝐧𝐦}}^𝐤`$ are
$$\begin{array}{c}f_{\mathrm{𝐧𝐦}}^𝐤=2i\delta _{𝐧+𝐦,𝐤}\mathrm{sin}(𝐧\times 𝐦)\end{array}$$
(2.12)
The set of unitary unimodular $`N\times N`$ matrices (2.9) suffices to span the algebra of $`SU(N)`$.
## 3 Morita Equivalence
### 3.1 Two-torus. $`U(1)|_{\theta =\frac{M}{N}}U(N)`$
To define Morita map we make an additional decomposition of the function (2.2) on the noncommutative two-torus:
$$\begin{array}{c}\widehat{f}=_{k\text{Z}\text{Z}^2}e^{iN𝐤\widehat{𝐱}}_{n_1,n_2=0}^{N1}f_{𝐤,𝐧}e^{in_1\widehat{x}_1+in_2\widehat{x}_2}\end{array}$$
(3.1)
Then we define corresponding $`U(N)`$-valued function on the ordinary two-torus as follows:
$$\begin{array}{c}f=_{k\text{Z}\text{Z}^2}e^{iN\mathrm{𝐤𝐱}}_{n_1,n_2=0}^{N1}f_{𝐤,𝐧}e^{i\mathrm{𝐧𝐱}}J_𝐧\end{array}$$
(3.2)
Because of the relation
$$\begin{array}{c}J_𝐧J_𝐧^{}=e^{i𝐧\times 𝐧^{}}J_{𝐧+𝐧^{}}\end{array}$$
(3.3)
Morita map (3.1, 3.2) takes star-product to the matrix product. Obviously, $`U(N)`$-valued function of general type can not be represented in the form (3.2). It turns out that this particular form corresponds to the functions with nontrivial boundary conditions. Namely, under shifts these functions transforms as
$$\begin{array}{c}f(x_1+2\pi \frac{M}{N},x_2)=\mathrm{\Omega }_1f(x_1,x_2)\mathrm{\Omega }_1^{},f(x_1,x_2+2\pi \frac{M}{N})=\mathrm{\Omega }_2f(x_1,x_2)\mathrm{\Omega }_2^{}\end{array}$$
(3.4)
where
$$\begin{array}{c}\mathrm{\Omega }_1=(P)^M,\mathrm{\Omega }_2=(Q^{})^M\end{array}$$
(3.5)
This can be treated as a constant gauge transformation. The size $`2\pi \frac{M}{N}`$ of the dual torus can be fixed by the requirement for the Morita map to be single-valued<sup>4</sup><sup>4</sup>4I am indebted to K. Selivanov for this comment.. To illustrate this, let us consider a torus of the size $`2\pi \frac{M}{N}n`$ (where $`n𝐍`$; there are no other possibilities if we want functions of the type (3.2) to be gauge transformed by the constant matrix when $`x`$ is shifted by a period of the torus.) Obviously, in this case there are functions which cannot be represented in the form (3.2). Such functions do not conjugates when translated along the vectors $`(2\pi \frac{M}{N},0)`$ and $`(0,2\pi \frac{M}{N})`$.
Therefore, having a set of Fourier coefficients $`f_{𝐤,𝐧}`$, we can construct both a function on the noncommutative torus of size $`l`$ and a matrix-valued function with twisted boundary conditions (3.4) on the commutative torus of size $`\frac{M}{N}l`$ by the following rule:
$$\{\begin{array}{c}e^{i𝐧\widehat{𝐱}}e^{i\mathrm{𝐧𝐱}}J_𝐧,n_1,n_2<N\hfill \\ \\ e^{iN𝐤\widehat{𝐱}}e^{iN\mathrm{𝐤𝐱}}\mathbf{\hspace{0.33em}1}\hfill \end{array}$$
(3.6)
### 3.2 $`𝐓^d`$. $`U(1)|_\theta U(N_1)\times \mathrm{}\times U(N_r)`$.
Generalization to the $`d`$-dimensional case goes by simple modifications in formulas from the previous subsection. It is always possible to rotate $`\theta _{\mu \nu }`$ into a canonical skew-diagonal form:
$$\begin{array}{c}\theta _{\mu \nu }=\left(\begin{array}{ccccccc}0\hfill & \theta _1\hfill & & & & & \\ \theta _1\hfill & 0\hfill & & & & & \\ & & & \mathrm{}\hfill & & & \\ & & & & 0\hfill & \theta _r\hfill & \\ & & & & \theta _r\hfill & 0\hfill & \\ & & & & & & \mathrm{𝟎}_{d2r}\hfill \end{array}\right)\end{array}$$
(3.7)
where $`r`$ is the rank of $`\theta _{\mu \nu }`$. Thus, algebra of higher dimensional noncommutative torus becomes embedded into a $`d`$-fold tensor product of $`r`$ noncommutative two-tori algebras and ordinary $`(d2r)`$-torus commutative algebra. This immediately leads to other examples of Morita equivalence, when some of these noncommutative two-tori are mapped to the commutative ones using relations (3.6). If $`\theta _i=\frac{M_i}{N_i}`$, after Morita map we obtain an ordinary YM theory with the gauge group $`U(N_1)\times \mathrm{}\times U(N_r)`$.
### 3.3 $`𝐓^d`$. $`U(1)|_\theta U(N)`$.
Algebra of noncommutative exponents can also be realized using a set of $`SU(N)`$valued matrices $`\mathrm{\Omega }_\mu ,\mu =1,\mathrm{},d`$ obeying the following relations:
$$\begin{array}{c}\mathrm{\Omega }_\mu \mathrm{\Omega }_\nu =e^{2\pi i\theta _{\mu \nu }}\mathrm{\Omega }_\nu \mathrm{\Omega }_\mu \end{array}$$
(3.8)
Explicit construction of such matrices can be found in . Define generators $`J_𝐧`$ as follows:
$$\begin{array}{c}J_𝐧=\mathrm{exp}\left(\underset{\nu <\mu }{}\theta _{\nu \mu }n_\nu n_\mu \right)\mathrm{\Omega }_1^{n_1}\mathrm{}\mathrm{\Omega }_d^{n_d}\end{array}$$
(3.9)
Then
$$\begin{array}{c}[J_𝐧,J_𝐦]=2i\mathrm{sin}\left(𝐧\times 𝐦\right)J_{𝐧+𝐦}\end{array}$$
(3.10)
which coincides with the algebra of noncommutative exponents. Therefore, in this case Morita map takes the form:
$$\begin{array}{c}\widehat{f}=_{k\text{Z}\text{Z}^d}e^{iN𝐤\widehat{𝐱}}_{𝐧<N^d}f_{𝐤,𝐧}e^{i𝐧\widehat{𝐱}}f=_{k\text{Z}\text{Z}^d}e^{iN\mathrm{𝐤𝐱}}_{𝐧<N^d}f_{𝐤,𝐧}e^{i\mathrm{𝐧𝐱}}J_𝐧\end{array}$$
(3.11)
## 4 Noncommutative YM vs Ordinary YM
Let us now turn to the physical applications of the Morita map. One can define noncommutative version of the Yang-Mills theory with the action
$$\begin{array}{c}S_{YM}=\frac{1}{4\pi g_{YM}^2}𝑑𝐱Tr(F_{\mu \nu }F^{\mu \nu })\end{array}$$
(4.1)
just by replacing in all formulas matrix product by the Moyal star-product and supplementing all quantities with the hats. Therefore, noncommutative $`U(1)`$ Yang-Mills action is
$$\begin{array}{c}\widehat{S}=\frac{1}{4\pi g_{NCYM}^2}𝑑\widehat{𝐱}\widehat{F}_{\mu \nu }\widehat{F}^{\mu \nu }\end{array}$$
(4.2)
where $`\widehat{F}_{\mu \nu }=_\mu \widehat{A}_\nu _\nu \widehat{A}_\mu i[\widehat{A}_\mu ,\widehat{A}_\nu ]_{}`$. For simplicity in this section we consider only two-torus. Generalization to the higher-dimensional case is straightforward.
Morita map takes NC $`U(1)`$ gauge fields to the $`U(N)`$ gauge fields with nontrivial boundary conditions. Generally, functions on torus can became gauge transformed when shifted by a period of the torus:
$$\begin{array}{c}A_\lambda (𝐱+𝐥_\mu )=\mathrm{\Omega }_\mu (𝐱)A_\lambda (𝐱)\mathrm{\Omega }_\mu ^1(𝐱)+i\mathrm{\Omega }_\mu (𝐱)_\lambda \mathrm{\Omega }_\mu ^1(𝐱)\end{array}$$
(4.3)
where $`\mathrm{\Omega }_\mu (𝐱)`$ are elements of the $`U(N)`$ group, known as twist matrices. They should satisfy a consistency conditions:
$$\begin{array}{c}\mathrm{\Omega }_\mu (𝐱+𝐥_\nu )\mathrm{\Omega }_\nu (𝐱)=e^{2\pi i\frac{M}{N}ϵ_{\mu \nu }}\mathrm{\Omega }_\nu (𝐱+𝐥_\mu )\mathrm{\Omega }_\mu (𝐱)\end{array}$$
(4.4)
An integer $`M`$ in this formula is so-called ’t Hooft’s flux. It is known only three types of possible boundary conditions (solutions of the eqs (4.4) ):
1. twist eaters: $`\mathrm{\Omega }_\mu =const`$
2. abelian twists
3. nonabelian twists
For more details see the recent review .
The map (3.6) corresponds exactly to the first case. It is not well understood how to realize Morita map corresponding to the other boundary conditions. Roughly speaking, when working in the Fourier basis (2.2), after shifts one can only multiply functions on numbers and cannot add something like $`\mathrm{\Omega }_\mu (𝐱)_\lambda \mathrm{\Omega }_\mu ^1(𝐱)`$. To do this, one needs another basis for the functions on noncommutative torus (creation/annihilation operators, noncommutative theta-functions?).
Under Morita map, defined in the previous section, actions go to the actions, equations of motions go to the equations of motions, and solutions (e.g. instantons) also go to the solutions, even at the quantum level. These properties of the Morita map can be encoded in the following identity:
$$\begin{array}{c}𝑑\widehat{𝐱}\widehat{A}_\mu \widehat{A}_\nu \mathrm{}\widehat{A}_\lambda =\frac{1}{N}𝑑𝐱Tr(A_\mu A_\nu \mathrm{}A_\lambda )\end{array}$$
(4.5)
which is straightforward to prove using the definition
$$\begin{array}{c}𝑑\widehat{𝐱}e^{i𝐤\widehat{𝐱}}=\delta _{𝐤,\mathrm{𝟎}}\end{array}$$
(4.6)
and the property (2.8) of the clock and shift generators. In fact, one can insert arbitrary number of derivatives into the integrals in (4.5) and thus obtain equivalent gauge invariant quantities in noncommutative and ordinary gauge theories. Due to the identity (4.5) we can establish the following correspondence between correlators:
$$\begin{array}{c}𝒟A_{𝐤,𝐧}^\mu e^{\widehat{S}[\theta =\frac{M}{N}]}\widehat{𝒪}_1\mathrm{}\widehat{𝒪}_l=𝒟A_{𝐤,𝐧}^\mu e^{S_{YM}}|_{fxdbndryconds,flux=M}𝒪_1\mathrm{}𝒪_l\end{array}$$
(4.7)
where $`g_{NCYM}^2=Ng_{YM}^2`$, and
$$\begin{array}{c}\widehat{𝒪}=𝑑\widehat{𝐱}(\widehat{F}_{\mu \nu })^n,𝒪=\frac{1}{N}𝑑𝐱Tr(F_{\mu \nu })^n\end{array}$$
(4.8)
Other important gauge invariant quantities of the YM theory are the Wilson loops:
$$\begin{array}{c}W[C]=TrP\mathrm{exp}\left(i_CA_\mu (𝐱)𝑑x_\mu \right)\end{array}$$
(4.9)
which corresponds to the closed path C. On torus there are paths from the different homotopy classes, which can be classified by winding numbers $`w_\mu `$ around the $`\mu `$-th direction. The corresponding Wilson loops are called Polyakov loops. The simplest Polyakov loop corresponds to the straight line along the $`\mu `$-th direction:
$$\begin{array}{c}W_P[𝐱,\mu ]=Tr\left[P\mathrm{exp}\left(i\underset{𝐱}{\overset{𝐱+𝐥_\mu }{}}A_\mu (𝐱)𝑑x_\mu \right)\mathrm{\Omega }_\mu e^{ix_\mu }\right]\end{array}$$
(4.10)
where insertion of the twist matrix (3.5) is necessary to guarantee gauge invariance.
Wilson lines in noncommutative Yang-Mills theory were constructed by Ishibashi, Iso, Kawai and Kutazawa (see also ). This construction goes as follows. First, introduce an oriented curve $`C`$ in auxiliary commutative two-dimensional space parametrized by the functions $`\xi (\sigma )`$ with $`0\sigma 1`$. Fix the starting point $`\xi _\mu (0)=0`$ end the endpoint $`\xi _\mu (1)=v_\mu `$. Then, assign to this curve a noncommutative analog of the parallel transport operator:
$$\begin{array}{c}𝒰[\widehat{𝐱},C]=1+_{n=1}^{\mathrm{}}i^n\underset{0}{\overset{1}{}}𝑑\sigma _1\underset{\sigma _1}{\overset{1}{}}𝑑\sigma _2\mathrm{}\underset{\sigma _{n1}}{\overset{1}{}}𝑑\sigma _n\frac{d\xi _{\mu _1}(\sigma _1)}{d\sigma _1}\mathrm{}\frac{d\xi _{\mu _n}(\sigma _n)}{d\sigma _n}\\ \times A_{\mu _1}(\widehat{𝐱}+\xi (\sigma _\mathrm{𝟏}))\mathrm{}A_{\mu _n}(\widehat{𝐱}+\xi (\sigma _𝐧))\end{array}$$
(4.11)
The series in (4.11) is noncommutative analog of the $`P`$-exponent. The star-gauge invariant quantity is then
$$\begin{array}{c}\widehat{𝒪}[C]=𝑑\widehat{𝐱}𝒰[\widehat{𝐱},C]S[\widehat{𝐱},C]\end{array}$$
(4.12)
where $`S[\widehat{𝐱},C]=1`$ if the path $`C`$ is closed and
$$\begin{array}{c}S[\widehat{𝐱},C]=e^{i(\theta ^1)_{\mu \nu }v_\nu \widehat{x}_\mu }\end{array}$$
(4.13)
if the path is open. Gauge invariance requires that the coordinates of the endpoint must be equal to $`v_\mu =2\pi r_\mu \frac{M}{N},r_\mu =0,\mathrm{},N1`$. In the simplest case, when $`C_\mu `$ is the straight line along the $`\mu `$-th direction and $`v_\mu =2\pi \frac{M}{N}`$, the function $`S[\widehat{𝐱},C_\mu ]`$ under Morita map (3.6) go to the twist function $`\mathrm{\Omega }_\mu e^{ix_\mu }`$. Therefore, using identity (4.5) we obtain the following relation between the Polyakov loops in the ordinary YM theory and open noncommutative Wilson loops:
$$\begin{array}{c}\frac{1}{N}𝑑𝐱W_P[𝐱,\mu ]=\widehat{𝒪}[C_\mu ]\end{array}$$
(4.14)
## 5 Conclusions
In this paper we have made some comments on the Morita equivalence between noncommutative and ordinary gauge theories. We present a simple prescription how to identify gauge fields and correlators of the gauge invariant observables in the $`U(1)`$ NC YM theory on torus at the rational value of the $`\theta `$-parameter with those ones in the ordinary $`U(N)`$ or $`U(N_1)\times \mathrm{}\times U(N_r)`$ YM theory with nontrivial boundary conditioxns on the dual torus. The size of the dual torus is determined by the requirement for the Morita map to be single-valued. We also show that under Morita map Polyakov loops in the ordinary YM theory go to the open noncommutative Wilson loops<sup>5</sup><sup>5</sup>5This fact firstly was mentioned in .
An open question is to generalize Morita equivalence to the case of the non-twist-eater’s type boundary conditions. Another interesting direction is to link three different descriptions of the Morita equivalence: field theory approach using the Fourier components, string theory approach using T-duality and brane language , and mathematical approach via the twisted bundles over the noncommutative torus .
Acknowledgements
I am grateful to D. Belov and N. Nekrasov for important comments reviving the interest in the subject. I thank S. Gorsky for numerous discussions and correspondence, and A. Morozov for his interest in this work and encouragement. I am also grateful to I. Polyubin, A. Rosly and especially to K. Selivanov for useful comments on the manuscript. I acknowledge Y. Makeenko for helpful discussion. I would like to thank Ira Vashkevich for technical support. The work was partly supported by the Russian President’s grant 00-15-99296 and RFBR grant 98-02-16575.
|
warning/0005/hep-ph0005251.html
|
ar5iv
|
text
|
# Effect of light 𝜎-meson Production in 𝑝𝑝̄→3𝜋⁰ at rest
## 1 Introduction
The light iso-singlet scalar $`\sigma `$ meson plays an important role in the mechanism of spontaneous breaking of chiral symmetry, and to confirm its real existence is one of the most important topics in hadron physics. Recently in various $`\pi \pi `$-production experiments a broad peak is observed in mass spectra below 1 GeV. Conventionally this peak was regarded as a mere non-resonant background, basing on the “universality argument,” since no $`\sigma `$ was seen in $`\pi \pi `$ scatering at that time. However, at present the $`\pi \pi `$-scattering phase shift $`\delta _S^{I=0}`$ is reanalyzed by many authors and the existence of $`\sigma `$ meson is strongly suggested. A reason of missing $`\sigma `$ in the conventional analysis is pointed out to be due to overlooking the cancellation mechanism, which is guaranteed by chiral symmetry, between the effects of $`\sigma `$ and those of repulsive $`\pi \pi `$-interaction. Moreover, the conventional treatment, based on the universality argument, of the low mass broad peak was shown not to be correct, and a new effective method, variant mass and width(VMW), is proposed to analyze resonance productions. In this method the production amplitude $``$ is directly represented by the sum of Breit-Wigner amplitudes with production couplings and phase factors (, including initial strong phases) of relevant resonances. The consistency of this method with the unitarity is seen from the following field theoretical viewpoint.
Presently, after knowing the quark physics, the strong interaction $`_{\mathrm{str}}`$ among hadrons (in our example, mesons) is regarded as a residual interaction of QCD among color singlet $`q\overline{q}`$-bound states, the “bare states,” denoted as $`\pi =\overline{\pi },\overline{\sigma },\overline{f}_0`$ and $`\overline{f}_2`$. In switching off $`_{\mathrm{str}}`$, the bare states appear as stable particles with zero widths. In switching on $`_{\mathrm{str}}`$, they change into the physical states with finite widths. The unitarity of $`S`$-matrix is guaranteed automatically by the hermiticity of $`_{\mathrm{str}}`$. The VMW method is obtained directly as the representation of production amplitude by physical state bases, with a diagonal mass and width.
The VMW method is already applied to the analyses of $`pp`$-central collision $`pppp\pi ^0\pi ^0`$ and $`J/\psi \omega \pi \pi `$ decay, leading to a strong evidence of existence of the light $`\sigma `$ meson. In this paper we apply this method to analysis of the high statistics data on the process $`p\overline{p}3\pi ^0`$ at rest obtained in the Crystal Barrel experiment.
## 2 Amplitude describing $`p\overline{p}3\pi ^0`$ at rest by VMW method
### 2.1 $`_{\mathrm{str}}`$ for $`p\overline{p}3\pi ^0`$ at rest
We apply the iso-bar model, describing the process in following two steps: In the first step the $`p\overline{p}`$ annihilates into the resonance $`f(f_0\mathrm{and}f_2)`$ and $`\pi ^0`$, and the $`f`$ decays into $`2\pi ^0`$ in the second step. Since both $`p`$ and $`\overline{p}`$ are at rest, the relative momentum $`p_\mu =p_{p\mu }p_{\overline{p}\mu }`$ (and the relative angular momentum $`L_{p\overline{p}}`$) between $`p`$ and $`\overline{p}`$ is 0. Charge conjugation parity $`P_C`$ of $`f\pi ^0`$-system is +1. Thus the three types of $`(\overline{p},p)`$ bi-linear forms with $`P_C`$=+1 are possible: $`\overline{p}i\gamma _5p,\overline{p}i\gamma _5\gamma _\mu p`$ and $`\overline{p}p`$. The second type reduces to the first type by using the equation $`_\mu (\overline{p}i\gamma _5p)=2m_p\overline{p}i\gamma _5\gamma _\mu p+\overline{p}i\gamma _5i\sigma _{\mu \nu }\stackrel{}{_\nu }p=2m_p\overline{p}i\gamma _5\gamma _\mu p`$ (, derived from Dirac equation and the above “rest condition”), and the third type is forbidden by parity. Thus only the first type remains.<sup>\*)</sup><sup>\*)</sup>\*)This implies that the initial $`p\overline{p}`$ is in the $`{}_{}{}^{1}S_{0}^{}`$ state. However, in the original analysis, the phenomenological parameters related with the $`{}_{}{}^{3}P_{1}^{}`$ and $`{}_{}{}^{3}P_{2}^{}`$ states, which are considered not to contribute since of the above “rest condition,” are included. The most simple form of $`_{\mathrm{str}}^{1,2}`$ describing the 1st step $`p\overline{p}f\pi ^0`$ and the 2nd step $`f2\pi ^0`$ is given, respectively, by
$`_{\mathrm{str}}^1`$ $`=`$ $`{\displaystyle \underset{\overline{f}_0,\overline{f}_2}{}}(\overline{\xi }_{\overline{f}_0}\overline{p}i\gamma _5p\overline{f}_0\pi ^0+\overline{\xi }_{\overline{f}_2}\overline{p}i\gamma _5p\overline{f}_{2\mu \nu }_\mu _\nu \pi ^0),`$
$`_{\mathrm{str}}^2`$ $`=`$ $`{\displaystyle \underset{\overline{f}_0,\overline{f}_2}{}}(\overline{g}_{\overline{f}_0}\overline{f}_0\pi ^2+\overline{g}_{\overline{f}_2}\overline{f}_{2\mu \nu }(\pi \stackrel{}{_\mu }\stackrel{}{_\nu }\pi )).`$ (2.1)
### 2.2 Amplitude by VMW method
First denoting the three $`\pi ^0`$ as $`\pi _1,\pi _2`$ and $`\pi _3`$ with momenta $`p_1,p_2`$ and $`p_3`$, respectively, we consider the $`\pi _1`$ and $`\pi _2`$ forming the resonance $`f`$ with squared mass $`s_{12}`$ (where $`s_{ij}(p_i+p_j)^2`$) and with momentum $`|𝐩|`$ in z-direction. The $`\pi _1`$ has momentum $`|𝐪|`$ and polar angle $`\theta `$ in the $`f`$ rest frame. In the lowest order in bare state representation the amplitude is <sup>\**)</sup><sup>\**)</sup>\**) $`N(s_{12},cos\theta _{12})`$ is obtained by the calculation of $`N=p_{3\mu }p_{3\nu }𝒫_{\mu \nu ;\lambda \kappa }(p_1p_2)_\lambda (p_1p_2)_\kappa `$, where we use the tensor projection operator $`𝒫_{\mu \nu ;\lambda \kappa }`$ with mass squared $`s_{12}`$ instead of $`m_{f_2}^2`$: $`𝒫_{\mu \nu ;\lambda \kappa }=\frac{1}{2}(\stackrel{~}{\delta _{\mu \lambda }}\stackrel{~}{\delta _{\nu \kappa }}+\stackrel{~}{\delta _{\mu \kappa }}\stackrel{~}{\delta _{\nu \lambda }})\frac{1}{3}\stackrel{~}{\delta _{\mu \nu }}\stackrel{~}{\delta _{\lambda \kappa }}`$, where $`\stackrel{~}{\delta _{\mu \nu }}=\delta _{\mu \nu }+\frac{P_\mu P_\nu }{s_{12}};P_\mu =p_{1\mu }+p_{2\mu }`$.
$`2im_pf^{s_ps_{\overline{p}}}({\displaystyle \underset{\overline{f}_0}{}}{\displaystyle \frac{\overline{\xi }_{\overline{f}_0}\overline{g}_{\overline{f}_0}}{\overline{m}_{\overline{f}_0}^2s_{12}}}+{\displaystyle \underset{\overline{f}_2}{}}{\displaystyle \frac{\overline{\xi }_{\overline{f}_2}\overline{g}_{\overline{f}_2}N(s_{12},\mathrm{cos}\theta _{12})}{\overline{m}_{\overline{f}_2}^2s_{12}}}),`$
$`\mathrm{where}`$ (2.2)
$`2im_pf^{s_ps_{\overline{p}}}\overline{p}(\mathrm{𝟎},s_{\overline{p}})i\gamma _5p(\mathrm{𝟎},s_p)(s_{\overline{p}}\mathrm{and}s_p\mathrm{being}\mathrm{spin}\mathrm{of}\overline{p}\mathrm{and}p)`$
$`f^{++}=f^{}=0,f^+=f^+=1.`$
$`N(s_{12},\mathrm{cos}\theta _{12})={\displaystyle \frac{(s_{23}s_{31})^2}{4}}+{\displaystyle \frac{16m_p^2𝒑^2𝒒^2}{3s_{12}}},`$
$`|𝒑|={\displaystyle \frac{\sqrt{(4m_p^2s_{12}m_\pi ^2)^24s_{12}m_\pi ^2}}{4m_p}},|𝒒|=\sqrt{{\displaystyle \frac{s_{12}}{4}}m_\pi ^2}`$
Owing to the effect of final (and initial) state interaction, the “full order” of the amplitude in physical state representation is given by
$`A_{s_ps_{\overline{p}}}(s_{12},\mathrm{cos}\theta _{12})`$ $`=`$ $`2im_pf^{s_ps_{\overline{p}}}({\displaystyle \underset{f_0}{}}{\displaystyle \frac{r_{f_0}e^{i\theta _{f_0}}}{m_{f_0}^2s_{12}i\sqrt{s_{12}}\mathrm{\Gamma }_{f_0}(s_{12})}}`$ (2.3)
$`+{\displaystyle \underset{f_2}{}}{\displaystyle \frac{r_{f_2}e^{i\theta _{f_2}}N(s_{12},\mathrm{cos}\theta _{12})}{m_{f_2}^2s_{12}i\sqrt{s_{12}}\mathrm{\Gamma }_{f_2}}}).`$
The symmetric amplitude, satisfying the statistics property of $`3\pi ^0`$ system, $`_{s_ps_{\overline{p}}}`$ is obtained simply by its cyclic sum as
$`_{s_ps_{\overline{p}}}`$ $`=`$ $`A_{s_ps_{\overline{p}}}(s_{12},\mathrm{cos}\theta _{12})+A_{s_ps_{\overline{p}}}(s_{23},\mathrm{cos}\theta _{23})+A_{s_ps_{\overline{p}}}(s_{31},\mathrm{cos}\theta _{31}).`$ (2.4)
Cross section is given by
$`d\sigma `$ $``$ $`{\displaystyle _{2m_\pi }^{2m_pm_\pi }}𝑑\sqrt{s_{12}}{\displaystyle \frac{\sqrt{s_{12}}}{\pi }}{\displaystyle \frac{|𝐩|}{8\pi m_p}}{\displaystyle \frac{|𝐪|}{8\pi \sqrt{s_{12}}}}{\displaystyle _1^1}𝑑\mathrm{cos}\theta \overline{||^2},`$
$`\mathrm{where}`$ (2.5)
$`\overline{||^2}(1/4){\displaystyle \underset{s_p,s_{\overline{p}}}{}}|_{s_ps_{\overline{p}}}|^2=(1/2)|_+|^2,`$
$`s_{12}+s_{23}+s_{13}=4m_p^2+3m_\pi ^2,`$
$`s_{23}=4m_p(|𝐩||𝐪|/\sqrt{s_{12}})\mathrm{cos}\theta +m_\pi ^2+2m_pE_3,`$
$`E_3=\sqrt{m_\pi ^2+𝐩^2}=(4m_p^2s_{12}+m_\pi ^2)/4m_p.`$
## 3 Results of Analysis
We analyze the experimental data, the $`\pi ^0\pi ^0`$ mass spectra and angular distributions around $`K\overline{K}`$-threshold and at 1.5 GeV, which are published in the paper by Crystal Barrel collaboration, using our formulas Eq.(2.5). We take into consideration as the physical particles $`f_0=\sigma ,f_0(980),f_0(1370)`$, $`f_0(1500)`$, and $`f_2=f_2(1275)`$,$`f_2(1565)`$.
We take into account the $`\pi \pi `$ and $`K\overline{K}`$ couplings of the relevant resonances, where we consider the effects of all the inelastic channels are represented by the $`K\overline{K}`$ coupling.
The result of the fit is shown in Fig. 1.
The mass and width of $`\sigma `$ obtained are $`m_\sigma =540\stackrel{+36}{29}`$MeV and $`\mathrm{\Gamma }_\sigma =385\stackrel{+64}{80}`$MeV (error corresponding to the 5$`\sigma `$-deviation). The reduced $`\chi ^2`$ is given by $`746.9/(28130)=2.98`$. Respective contributions to the $`\chi ^2`$ from the mass and angular distributions(a.d.) around $`K\overline{K}`$ and at 1.5GeV are 406.5(mass), 148.1(a.d. around $`K\overline{K}`$) and 192.2(a.d. at 1.5GeV).<sup>\***)</sup><sup>\***)</sup>\***) The obtained $`\chi ^2`$ value may be compared with the one by the original fit: We have estimated their corresponding $`\chi ^2`$ by reading the deviations of their fit from the experimental points and obtained the numerical values as $`\chi ^2/(N_{data}N_{param})=835.4/(28225)=3.25.`$ Respective contributions to the $`\chi ^2`$ are 432.8(mass), 113.4(a.d. around $`K\overline{K}`$) and 289.3(a.d. at 1.5GeV). Almost the similar but slightly improved fit is obtained in the present VMW method. However, note that their parameters were determined by analyzing the original two-dimensional Dalitz plot directly. They reported that $`\chi ^2`$ in this case is $`\chi ^2/N_F=2028/(133834)=1.6`$. <sup>\****)</sup><sup>\****)</sup>\****) Number of data points in the figures given in ref.? is 282=82(mass)+100(a.d. around $`K\overline{K}`$)+100(a.d. at 1.5GeV). In our analysis the first one point close to the threshold of mass spectra is removed since this point is at 279MeV, where $`\pi ^0\pi ^0`$ channel is open but $`\pi ^+\pi ^{}`$ channel is almost closed, and accordingly a special treatment of $`\pi \pi `$ widths of resonances is required. The properties of all the relevant resonances used in the fit are given in Table 1.
We have also tried the fit without considering the $`K\overline{K}`$ couplings of the resonances. The result of our fit with the $`K\overline{K}`$ couplings is almost the same as the one without the $`K\overline{K}`$ couplings; and in order to determine the $`K\overline{K}`$ couplings and the other ones, it is necessary to analyze the data on the corresponding channels directly.
In order to see the effect of $`\sigma `$-meson production in our fit, in Fig. 1 the spectra given by setting $`r_\sigma =0`$ are also given by dashed lines. Effect of $`\sigma `$-production is seen to be crucially important in reproducing the structure of mass spectra below 1 GeV.
We have also tried the fit without introducing the $`\sigma `$ Breit-Wigner amplitude. In this case the broad peak structure below 1 GeV in the mass spectra are only roughly reproduced by the combinatorial background coming from the higher mass $`2`$-$`\pi ^0`$ resonances due to the statistics property of $`3\pi ^0`$ system. The corresponding $`\chi ^2`$ is 3112/(281-26)=12.2, which is much worse than our best fit with $`\sigma `$ meson. This seems to give a strong evidence for $`\sigma `$-existence.
## 4 Comparison with other analyses
Several extensive analyses (including the original one) on the relevant experimental data by Crystal Barrel collaboration have been thus far done. However, all the analyses seem to be done, more or less, under the influence of “universality argument.”
According to this argument, all the $`\pi \pi `$ production amplitudes $``$ must be parame- trized through the “universal” $`\pi \pi `$ scattering amplitude $`𝒯`$ as
$``$ $`=`$ $`\alpha (s)𝒯,`$ (4.1)
with slowly varying real function $`\alpha (s)`$ which corresponds to the $`\pi \pi `$ production couplings of the production channel and is process-dependent. This equation is believed to be based on the unitarity or final state interaction(FSI) theorem.
However, it has been pointed out that, in order to apply the FSI theorem presently, after knowing quark physics, we must make a special attention on the bases of representation of $`𝒯`$ and $``$. Before knowing quark physics, the $`S`$-matrix of strong interaction $`S_{\mathrm{str}}`$ was represented in terms of only stable particles such as $`\pi `$ and $`N`$. However, now $`S_{\mathrm{str}}`$ should be described by the interaction Hamiltonian $`_{I,\mathrm{str}}`$ among the bare states, the color neutral stable bound states of quarks and/or antiquarks. The observed $`\pi \pi `$ resonances, such as $`\sigma `$ or $`f_0(980)`$, and all the resonant states must be treated equally to the stable states, such as one pion or two pion states, as complete set describing the hadron world. They have mutually independent $`\pi \pi `$( and/or $`K\overline{K}`$) production couplings, in principle. Accordingly, Eq.(4.1) in its original form has proved not to be correct.
In the pioneering work by Aker et al. of Crystal Barrel(CB) collaboration, the $`I=0`$ $`S`$ wave production amplitude $`_S(s_{12})`$ (in our notation) is taken as
$`_S(s_{12})`$ $`=`$ $`\alpha {\displaystyle \frac{1}{\rho _1(s_{12})}}e^{i\delta _S^{I=0}(s_{12})}\mathrm{sin}\delta _S^{I=0}(s_{12})=\alpha 𝒯`$ (4.2)
with the constant $`\pi \pi `$ coupling $`\alpha `$, basing on the original “universality argument” Eq.(4.1). The symmetric production amplitude is obtained by taking the cyclic sum as Eq.(2.4), while the scattering amplitude $`𝒯`$ is picked up from the reference by Au, Morgan, Pennington(AMP) in ?, where the conventional $`𝒦`$ matrix analysis was done for the CERN-Munich $`\pi \pi `$ scattering phase shift $`\delta _S^{I=0}`$.
In a series of works by Anisovich, Sarantsev, Bugg and Zou et al. being done in the line of this thought, the relation of “the universality argument” $`=\alpha (s)𝒯`$ with real $`\alpha (s)`$ function is argued to be not correct, since of the possible effect of the strong phases(basing on the detailed consideration of their origins, for example, triangle singularities). They used $`N/D`$ formalism and mentioned that in production processes the complex $`N`$ function $`N^{}(s)`$, being independent of $`\pi \pi `$ scattering $`N`$ function $`N(s)`$, is necessary. A great variety of modification is allowed for the parametrization of $`N^{}(s)`$, and they try to fit the spectra of $`p\overline{p}`$ annihilation with three tentative forms of $`N^{}(s)`$ with complex couplings. Here, they introduce the scalar($`f_0(1365)`$ and $`f_0(1520)`$) and tensor($`f_2(1270)`$ and $`f_2(1560)`$) Breit-Wigner amplitudes above 1.1 GeV with complex production couplings. In our interpretation, their analysis is, so to speak, two-fold: Below 1.1 GeV it was done without taking into account of the freedom of production couplings of resonances by following the viewpoint of “the universality argument,” while above 1.1 GeV their method is equivalent to the VMW method, where this freedom is explicitly introduced.
In all the above analyses the pole positions of $``$ matrix in low energy region below $``$ 1 GeV are determined only through the analyses of CERM-Munich $`\delta _S^{I=0}`$ by $`𝒦`$ matrix method, although there are many varieties of parametrization methods of $`𝒦`$. This is also the case in the original analysis by Amsler et al. of CB collaboration. In this analysis the $`\pi \pi `$ scattering and production amplitudes are given, respectively, in the $`𝒦`$ matrix representation as: $`𝒯=𝒦/(1i\rho 𝒦),=𝒫/(1i\rho 𝒦)`$, where $`𝒦`$-matrix and $`𝒫`$-matrix are taken in pole-dominative form;
$`𝒦={\displaystyle \underset{\alpha }{}}{\displaystyle \frac{g_\alpha ^2}{(m_\alpha ^2s)}}+c_{\mathrm{BG}},𝒫={\displaystyle \underset{\alpha }{}}{\displaystyle \frac{e^{i\theta _\alpha }\xi _\alpha g_\alpha }{(m_\alpha ^2s)}}.`$ (4.3)
Here the summation $`\alpha `$ is taken for the $`𝒦`$-matrix states, which are related to the physical states corresponding to the poles of $`𝒯`$ (or $``$). In the case with no production phases, $`\theta _\alpha =0`$, the $``$ and $`𝒯`$ have the same phase, coming from the common factor $`1/(1i\rho 𝒦)`$, and the final state interaction theorem is satisfied. Actually in their analysis this phase was set to be the experimental scattering phase shift $`\delta _\mathrm{S}^{\mathrm{I}=0}`$. However, in the above $`𝒦`$-matrix parametrization, when $`s`$ is close to $`m_\alpha ^2`$, $`𝒦`$ diverges and the phase must take the value $`90^{}(+n`$$`\times `$$`180^{})`$. This gives a very strong constraint for the value of $`m_\alpha `$. The experimental $`\delta _\mathrm{S}^{\mathrm{I}=0}`$ passes through $`90^{}`$ at about $`\sqrt{s}900`$ MeV, and so the $`m_\alpha `$ becomes $`m_\alpha 900`$ MeV , which is much larger than $`m_\sigma (`$600MeV). Thus, no existence of light $`\sigma `$ is implicitly assumed from the beginning. This situation is common in all the analyses thus far made.
In our method, the $`\delta _\mathrm{S}^{\mathrm{I}=0}`$ is analyzed by introducing the repulsive background phase shift $`\delta _{\mathrm{BG}}`$, which is required from chiral symmetry. The scattering $`S`$-matrix and correspondingly the $`𝒦`$-matrix are parametrized by
$`S=S^{\mathrm{Res}}S^{\mathrm{BG}},𝒦={\displaystyle \frac{𝒦^{\mathrm{Res}}+𝒦_{\mathrm{BG}}}{1\rho ^2𝒦^{\mathrm{Res}}𝒦_{\mathrm{BG}}}}.`$ (4.4)
The $`𝒦^{\mathrm{Res}}`$ in denominator of $`𝒦`$ removes the poles of $`𝒦^{\mathrm{Res}}=_\alpha g_\alpha ^2/(m_\alpha ^2s)`$ in the numerator in the total $`𝒦`$-matrix and we can take the light $`m_\alpha 600`$MeV.<sup>\*****)</sup><sup>\*****)</sup>\*****) On the other hand, the background matrix, $`c_{BG}`$, in the conventional $`𝒦`$-matrix cannot describe the global phase motion corresponding to $`\delta _{\mathrm{BG}}`$ in our method, and the small value of $`m_\alpha `$ is not permissible.
Correspondingly, the $``$ is represented by <sup>\******)</sup><sup>\******)</sup>\******)Here we neglect the possible effect of non-resonant $`3\pi ^0`$ production.
$`={\displaystyle \frac{𝒫^{\mathrm{Res}}}{1i\rho 𝒦^{\mathrm{Res}}}}e^{i\delta _{\mathrm{BG}}},`$ (4.5)
which satisfies the final state interaction theorem. This $``$ is able to be rewritten into the form, applied in VMW method, in the physical state representation. In our approach, whether the light $`\sigma `$ meson exists or not is determined directly from the experimental data themselves, as was done in §3.
## 5 Conclusion
Through the results of analyses given above we may conclude that the effects of production of light $`\sigma `$ meson are clearly shown. The numerical values of mass and width of $`\sigma `$ are obtained as $`m_\sigma =540\stackrel{+36}{29}`$MeV, $`\mathrm{\Gamma }_\sigma =385\stackrel{+64}{80}`$MeV, which are consistent with those obtained in our phase shift analysis ($`(m_\sigma ,\mathrm{\Gamma }_\sigma )`$=$`(535675,385\pm 70)`$MeV). However, the effect of $`\sigma `$ in this process is, in principle, not able to be discriminated from those of higher mass resonances, such as $`f_0(1370)`$, due to the effects coming from statistics property of $`3\pi ^0`$ system. In order to avoid this, it is desirable to analyze also the process, $`\overline{p}n\pi ^0\pi ^0\pi ^{}`$, through the similar method.
Finally it should be noted that the experimental data of the spectra applied in this paper were obtained by reading out the corresponding figures of ref. ? and incomplete. The excellent reproduction of the data encourages us to study more in details. It is desirable to reanalyze directly the experimental data of the Dalitz plot.
|
warning/0005/hep-th0005160.html
|
ar5iv
|
text
|
# (Anti-)Instantons and the Atiyah-Hitchin Manifold
## 1 Introduction
The advent of the now mundane dualities of supersymmetric field and string theories has made it possible to obtain a wealth of non-perturbatively exact results for various couplings in the low-energy effective action of these theories, all of them severely constrained by supersymmetry. In the seminal case of the Coulomb phase of $`N=2`$ four-dimensional gauge theories , the holomorphicity of the prepotential together with electric-magnetic duality was sufficient to fix the dynamics of the vector-multiplets at two-derivative order for all values of the (dimensionally transmuted) gauge coupling $`\mathrm{\Lambda }^{b_0}=\mu ^{b_0}\mathrm{exp}(8\pi ^2/g_{YM}^2)`$, and the non-perturbative corrections of order $`\mathrm{\Lambda }^{nb_0}e^{in\theta }`$ to the one-loop prepotential were identified as the contribution from $`n`$ Yang-Mills instantons. In a similar manner, exact higher derivative couplings in type I and type II string theories have been obtained from the requirements of S-duality and harmonicity (see for a review and references), and shown to encapsulate the contributions from bound states of arbitrary number $`n`$ of D-instantons, plus its complex-conjugate series of anti-instanton contributions.
While holomorphy or harmonicity give powerful constraints on the half-BPS–saturated couplings of vector-multiplets of $`N=2`$ supersymmetry (and presumably any multiplets of higher supersymmetry), the couplings of hypermultiplets are by no means as well understood, even though they occur with the same amount of supersymmetry. The reason is that the hyperkähler or quaternionic constraints on the moduli space of these fields lack a concise and tractable formulation (except perhaps for the twistor methods, which are holomorphic with respect to an auxiliary variable; see for a review). Indeed, there are to date no explicitly known non-homogeneous quaternionic manifolds, and very few explicit examples of hyperkähler manifolds, all of them one-dimensional (in quaternionic units) and with a large number of isometries. Among them are the (multi) Eguchi-Hanson and Taub-NUT gravitational instantons, with asymptotic geometry $`^4/_k`$ and $`^3\times S^1`$ respectively, which both possess a triholomorphic $`U(1)`$ isometry (times $`SO(3)`$ in the single instanton case; see for a review). Hence, they can be obtained from a harmonic function, which is a sum of a finite number of pole terms in both cases (plus a constant for Taub-NUT). By combining infinite series of poles, it is possible to generate a non-perturbative behaviour, and indeed it was shown by Ooguri and Vafa that such a space describes the quantum-corrected moduli space of the conifold hypermultiplet . In the weak coupling (or asymptotic) limit, one recovers a series of $`n`$ D-instanton effects, together with its complex-conjugate series of anti-instanton effects, as in the Type IIB case above.
The only explicit example of four-dimensional hyperkähler space without triholomorphic isometry (but with an $`SO(3)`$ group of “rotational symmetries” not rotating the three Kähler forms) is the Atiyah-Hitchin manifold, first introduced in the context of the moduli space of two BPS monopoles in a 3+1 $`SU(2)`$ gauge theory with a Higgs field in the adjoint . This space is one-dimensional (after omitting a trivial center-of-mass factor of $`^3\times S^1`$), and consists of three relative positions of the monopoles (measured in units of the $`W`$-boson mass <sup>2</sup><sup>2</sup>2We normalize the kinetic term of the Higgs field to $`(\varphi )^2/g_{YM}^2`$. $`M_W=\varphi `$) and a relative phase $`\sigma `$. The metric on that space controls the slow motion of the two monopoles , which primarily interact through the exchange of massless photons and scalars at long distance . In this regime, the moduli space reduces to a Taub-NUT space, albeit with a negative mass parameter, and hence a singularity at finite distance $`r=2`$ between the monopoles on the order of the Higgs vev. The metric is independent of the gauge coupling, there are therefore no quantum corrections to this motion, whether perturbative or not. There are however corrections to the long-range interaction due to the exchange of massive $`W`$-bosons and Higgs field, which dominates when they come close to each other and resolves the singularity. The exact metric was derived in (see also ) on the basis of $`SO(3)`$ isometry and self-duality (which expresses hyperkählerity in 1 dimension), in terms of elliptic functions. The deviation from the Taub-NUT limit is exponentially small in the distance between the monopoles, and is most easily expressed as a deviation to particular components of the Riemann tensor which we shall make precise later,
$`R_{(0)}`$ $`=`$ $`{\displaystyle \frac{4}{(r2)^3}}+\mathrm{}`$ (1)
$`R_{(+2)}`$ $`=`$ $`{\displaystyle \frac{8(3+9r6r^2+r^3)}{(r2)^2}}e^{r+i\sigma }+\mathrm{}`$ (2)
$`R_{(2)}`$ $`=`$ $`{\displaystyle \frac{8(3+9r6r^2+r^3)}{(r2)^2}}e^{ri\sigma }+\mathrm{}`$ (3)
where the dots denote subleading corrections of order $`e^{2r}`$. The exponential terms in this expression can be interpreted as the semi-classical effect of the Euclidean worldline of a massive $`W`$-boson stretching between the two monopoles. In the following, we will be mostly interested in the structure of the subleading terms in this expansion, displayed in (18) below, which will reveal the interplay between instantons and anti-instantons.
While the above occurrence of the Atiyah-Hitchin manifold was purely classical, the same manifold arises in many other instances in string or field theory, where the radial parameter $`r`$ takes another meaning, and in particular can have a coupling-dependent scale. The long distance expansion of the monopole problem then becomes a weak coupling expansion, and the exponentially small corrections can be truly identified with instanton effects. It will be our goal to interpret these effects in the various cases where the Atiyah-Hitchin manifold provides the exact answer, in the hope of drawing lessons for cases where an explicit answer is missing. This program has already been carried out at the one-instanton level in the context of three-dimensional gauge theories : our goal is to extend this study to all higher order non-perturbative contributions, and to other settings where the same effects appear.
The plan of this paper is as follows. We will first review the various instances of the Atiyah-Hitchin manifold in supersymmetric gauge theories, brane constructions and string backgrounds, find the relevant instanton configurations and identify the parameters $`r`$ and $`\sigma `$ in these settings. In Section 3, we will revisit the Atiyah-Hitchin metric in a way that allows us to easily extract the series of exponential corrections to the Taub-NUT limit, and identify precisely which instanton configurations contribute to which components of the metric. In Section 4, we shall justify our claim that the exponential corrections arise as contributions from instanton–anti-instanton bound states, and discuss the consequences of this phenomenon for the general question of non-perturbative corrections to hyperkähler manifolds. Some computational details pertaining to Section 3 are relegated to the appendices.
## 2 Atiyah-Hitchin Manifold, a Festival
In this section, we would like to review some of the many instances of the Atiyah-Hitchin manifold in field or string-theoretic situations, where it provides an exact resummation of all quantum corrections. Being the moduli space of two $`SU(2)`$ monopoles, it naturally arises whenever we embed monopole solutions in string theory. For a recent discussion on this point see . By duality, it also appears in many other situations where the relevance of monopoles is not immediately obvious. Finally, being a hyperkähler manifold, it appears in many backgrounds with high degree of supersymmetry. Our aim here will be to understand the source of non-perturbative effects, and in particular identify the weight of the semi-classical configurations in terms of the monopole variables (up to factors of 2 and $`\pi `$, which take care of themselves). In the course of our discussion, we will also mention the occurrence of higher monopole moduli spaces. Even though they are not as well understood as the two $`SU(2)`$ monopole case, there is yet a considerable amount of knowledge about them which can be carried over to these dual situations.
### 2.1 Brane lifting of the monopole problem
Much insight into the dynamics of gauge theories has been gained by embedding them into string theory. The monopole problem is no exception, and can be given a simple brane realization. In the limit of far separation, a $`k`$-monopole configuration of an $`SU(2)`$ gauge theory is represented by $`k`$ oriented D-strings stretched between two D3-branes of Type IIB theory. That a D-string ending on a D3-brane acts as magnetic source follows from the worldvolume anomalous coupling $`FB_{RR}`$ on the D3-brane worldvolume . The four scalars associated to each monopole correspond to the three spatial coordinates of the D1-brane on the D3-brane, together with a fourth scalar $`\sigma `$ measuring the zero-mode of the U(1) gauge field $`A`$ on the stretched D1-brane world-line.
Monopoles of $`SU(N)`$ can be similarly represented as D-strings stretching between different foils of a stack of $`N`$ D3-branes. Monopoles in 3+1 dimensions can also be lifted to higher extended objects in higher dimensions, or reduced to instantons in 3 dimensions, and so the configuration generalizes to $`k`$ D$`p`$-branes stretched between $`N`$ D$`(p+2)`$-branes, for $`p=0\mathrm{}6`$. One virtue of this description is that it makes obvious Nahm’s construction of monopoles, by switching the perspective from the D3-brane to the D1-brane worldvolume: the matrices appearing in Nahm’s equations simply describe the fluctuations of the transverse positions of the D1-branes stretched between the D3-branes .
Another virtue of this representation is that it gives a simple geometric representation of the non-perturbative contributions appearing in (1): the value of the Higgs field $`\varphi `$ being related to the distance $`L`$ between the two D3-branes through $`\varphi =L/l_s^2`$, the weight $`e^{\varphi r}`$ simply corresponds to the action of an Euclidean fundamental string stretched between the two D1-branes and D3-branes as depicted on Figure 1. We shall refer to this configuration as a worldsheet instanton. Since the fundamental string stretched between two D3 is the massive $`W`$-boson of the gauge theory, this rephrases our previous statement that the exponential corrections in (1) come from Euclidean worldlines of $`W`$-bosons stretching between the two monopoles. The imaginary part $`e^{i\sigma }=e^{i(\sigma _1\sigma _2)}`$ of the instanton weight comes from the electric coupling of the boundaries of the fundamental string worldsheet to the U(1) gauge fields $`\sigma _1`$ and $`\sigma _2`$ living on the two D1-branes.
### 2.2 Monopoles and three-dimensional gauge theories
The above brane configuration for $`p=3`$ can be S-dualized to a configuration of $`k`$ D3-branes stretched between $`N`$ NS5-branes. This system was studied in detail in as a way of deriving the conjectured equivalence between moduli spaces of monopoles and the Coulomb branch of three-dimensional gauge theories with 8 supersymmetries . Indeed, the theory living on the D3-branes is effectively at distances larger than the 5-brane separation a three-dimensional gauge theory with three Higgs fields in the adjoint (the extra four Higgs fields living on the D3-branes in vacuo are projected out by the boundary conditions imposed by the NS5-branes). In the Coulomb phase, the three-dimensional gauge fields can be dualized into pseudoscalars $`\sigma _i`$, which combine with the three Higgs fields to make $`k`$ hypermultiplets taking value in some hyperkähler manifold. For $`(N,k)=(2,2)`$, it was shown on symmetry grounds that the only possibility is the Atiyah-Hitchin manifold . The fact that the moduli space of two $`SU(2)`$ monopoles appears as the moduli space of the three-dimensional gauge theory is naturally explained from the brane point of view, by simply switching the perspective from the NS5 to the D3-brane worldvolume . More generally, the quantum corrected moduli space of the three-dimensional $`SU(k)`$ gauge theory is identified with the moduli space of $`k`$ monopoles in $`SU(2)`$ , and a similar equivalence also holds between moduli spaces of monopoles in ADE gauge groups and Coulomb branches of quiver three-dimensional gauge theories .
From the point of view of the three-dimensional gauge theory, the $`1/(r2)^3`$ contribution to the Riemann tensor $`R_{(0)}`$ appears as a one-loop effect, while the non-perturbative corrections arise as instanton effects from monopoles. This can easily be seen by following the instanton configuration on Figure 1 under S-duality: it becomes an Euclidean D1-brane stretched between both the pair of D3 branes and the pair of NS5-branes, as first discussed in (see Figure 2). From the point of view of the D3-brane, this is an Euclidean monopole, and hence an instanton in the three-dimensional gauge theory. Its classical action is given by $`rL/(g_sl_s^2)`$, where $`L=\varphi l_s^2`$ is the distance between the D3-branes, and $`r`$ the distance between the NS5-branes. This can be rewritten in terms of the gauge theory variables as $`\varphi /g_{YM}^2`$, which is the appropriate weight for an instanton of the three-dimensional gauge theory. As shown by Polyakov , instanton contributions should in addition be weighted by a term $`e^{i\sigma }`$ where $`\sigma `$ is the dual of the gauge field in three dimensions, or equivalently the fourth component of the dual magnetic gauge field $`A_m`$ in four dimensions. Indeed, this imaginary part of the action naturally arises from a magnetic coupling $`A_m`$ on the D1-brane boundary, dual to the more familiar electric coupling $`A`$ on the fundamental string boundary.
If the moduli space of the $`N=4`$ three-dimensional gauge theory is really the Atiyah-Hitchin manifold, it should be possible to recover the exponential correction in (1) from a one-instanton computation. This was carried out successfully in for the case without matter, and in for the case with one hypermultiplet, which is conjectured to be described by a double cover of the Atiyah-Hitchin manifold . We will be interested in the extension of these considerations to higher order in the instanton expansion.
### 2.3 Orientifold Eight-planes with NS5-branes
Let us now go back to the case of the D$`p`$-D$`(p+2)`$ system, for $`p=6`$, now embedded in Type I. Enhanced $`SU(2)`$ gauge symmetry occurs when two D8-branes become coincident, but also when the string coupling on one of the $`O8^{}`$ orientifold planes becomes infinite. The $`SU(2)`$ symmetry is Higgsed at finite string coupling, and the massive $`W`$-bosons are provided by D0-branes stuck on the orientifold plane. It should also be possible to describe $`SU(2)`$ monopoles on an $`O8^{}`$ orientifold at finite coupling. This problem was discussed in , where it was shown that the monopoles are NS5-branes stuck on the orientifold plane (see Figure 3 left). These have a relative three- dimensional distance with a phase angle corresponding to the distance between the NS branes in the eleventh direction (string coupling direction). These four scalars form a hypermultiplet which is again parameterizing the AH manifold . The monopoles interact by exchange of $`W`$-bosons, hence by D0-brane instantons stretched between the two NS5-branes. The non-perturbative corrections to the moduli space are thus suppressed by $`\mathrm{exp}[r/(g_sl_s)]`$, where $`r`$ is the distance between the NS5-branes. This effect is non-perturbative in the string coupling.
### 2.4 Heterotic string on an ALE space
One can now T-dualize this configuration to a Type I background and further S-dualize to a Heterotic string background. The NS5-branes turn into an $`^4/_2`$ singularity and the corresponding background is the Heterotic string on $`^{5,1}\times ALE`$, where $`ALE`$ is the Eguchi-Hanson manifold which resolves the orbifold singularity. The hypermultiplet controlling the size and B-flux of the blown-up two-cycle can be argued to take value again in Atiyah-Hitchin manifold . The origin of the exponential corrections can be traced by following the sequence of dualities, and correspond to world-sheet instantons of the fundamental Heterotic string. This is in agreement with the non-renormalization property of the hypermultiplet moduli space in Heterotic theories with 8 supersymmetries. The exponential corrections are of order $`\mathrm{exp}(A/l_s^2+iB)`$, where $`A`$ denotes the area of the two-sphere in the blown-up ALE space and $`B`$ is the B-flux on that cycle. Since these corrections occur purely at string tree-level, it should be possible to recover the exact Atiyah-Hitchin metric in a sigma-model computation. This is however not easy due to the singular nature of the conformal field theory at hand. Note that the identification of Heterotic hypermultiplet spaces and monopole moduli spaces has been generalized in .
### 2.5 D6-branes on $`K3`$
In a recent series of papers a yet different type of realization of the Atiyah-Hitchin manifold was considered. One looks at a collection of $`N`$ D6 branes wrapped on a smooth $`K3`$. The low energy dynamics of these branes is given by a 2+1 dimensional gauge theory with 8 supercharges and gauge group $`U(N)`$. For the special case of $`N=2`$ this configuration leads to the gauge theory which was discussed in subsection 2.2 and hence the moduli of the wrapped D6 branes parameterize the Atiyah-Hitchin manifold. The gauge coupling of the 2+1 dimensional theory can be computed to be $`1/g_3^2=(V_{K3}l_s^4)/(g_sl_s^3)`$, where the negative term comes from extrinsic curvature terms on the D6-brane. It vanishes at $`V_{K3}=l_s^4`$, which is a point of $`SU(2)`$ enhanced symmetry in target space (as seen from duality with the Heterotic side for example). The D6-branes can be seen as the monopoles of this broken gauge symmetry in 5+1 dimensions, as they become tensionless at the enhanced symmetry point. The Higgs vev of the three-dimensional gauge theory on the D6-brane is related to the distance $`r`$ between the 2 D6-branes by $`\varphi =r/l_s^2`$. The expansion parameter controlling the corrections to the D6-brane moduli space is then $`r(V_{K3}l_s^4)/(g_sl_s^5)`$. This is also the action of Euclidean D4-branes wrapped on $`K3`$ and stretching between the D6-branes, which are therefore the relevant instanton configurations. This is consistent with the identification of the wrapped D4-branes as the $`W`$-bosons of the target-space $`SU(2)`$ enhanced symmetry. This also identifies the Higgs vev of the target-space gauge symmetry as $`(V_{K3}l_s^4)/(g_sl_s^5)`$. In fact, the statement made in is slightly different, since it is concerned with the moduli space of a probe D6-brane in the background of a large number $`N`$ of D6-branes creating a repulson singularity. The claim is that the corrected moduli space is again the Atiyah-Hitchin manifold. The latter then appears as a four-dimensional submanifold of the moduli space of a large number $`N`$ of $`SU(2)`$ monopoles.
### 2.6 Atiyah-Hitchin Backgrounds
An interesting occurrence which is different from all examples discussed above is a string/M-theory background on a manifold which contains the Atiyah-Hitchin manifold, $`M_{AH}`$. A simple example for such a background is given by M-theory on $`^{6,1}\times M_{AH}`$. This background is known to be the strong coupling dual of Type IIA string theory on $`^{6,1}\times ^3/I_3\times \mathrm{\Omega }\times (1)^{F_L}`$ or in its more familiar form, an $`O6^{}`$ plane . One may argue that D2-brane probes in the vicinity of an orientifold 6-plane behave as monopoles of an Higgsed $`SU(2)`$ gauge symmetry (see Figure 3, right).
Indeed, the worldvolume theory on a pair of D2-branes is a pure three-dimensional $`N=4`$ $`SU(2)`$ theory, and hence has the Atiyah-Hitchin manifold as its moduli space. It should therefore be the case that the singular O6-plane be resolved in the strong string coupling limit into a smooth manifold, the Atiyah-Hitchin manifold itself. In order to relate the Higgs vev to the string parameters, let us consider the gauge theory on the D2-branes. The three dimensional gauge coupling is given by the usual $`1/g_3^2=l_s/g_s`$. Denote by $`r/2`$ the distance of the D2 brane to the $`O6^{}`$ plane, measured by the $`SU(2)`$ vev $`\varphi _{D2}=r/l_s^2`$. The exponential corrections are given by Euclidean D0-branes which are stretched in between the D2 brane and its image with an expansion parameter $`\varphi _{D2}/g_3^2=r/g_sl_s`$. This implies that the $`SU(2)`$ gauge theory of which the D2-branes are monopoles will have a Higgs vev $`\varphi _{O6}=1/g_sl_s`$. The enhanced $`SU(2)`$ symmetry therefore occurs at scale $`\varphi _{O6}`$, which is also the radius of the eleventh dimension. The $`W`$-bosons of the enhanced $`SU(2)`$ symmetry are the D0-branes, i.e. the momentum modes of the graviton on the compact eleventh dimension.
## 3 Atiyah-Hitchin revisited
Having recalled a few occurrences of the Atiyah-Hitchin manifold in string and field theory, and identified what type of non-perturbative corrections it purports to resum, we now would like to extract the precise form of these corrections beyond the one-instanton effect that was displayed in (1). For convenience, we shall use the monopole terminology, but the other cases can be obtained by simply reinterpreting the meaning of the $`r`$ and $`\sigma `$ coordinates.
### 3.1 The Atiyah-Hitchin metric and modular forms
The metric found by Atiyah and Hitchin was originally expressed in terms of elliptic functions, whose appearance is hardly surprising given the fact that the algebraic curve underlying the Nahm equations has genus 1 in the two-monopole case. The asymptotic expansion of elliptic functions is most easily obtained after expressing them in terms of Jacobi Theta functions and other Eisenstein series, whose $`q`$-expansion is well known. It turns out that the metric can be expressed very concisely in that form, as we now briefly show <sup>3</sup><sup>3</sup>3After obtaining this result, we were informed by I. Bakas that it had already appeared in the mathematical literature . . We follow the notations and conventions of for the Atiyah-Hitchin metric, and of for modular forms.
The general $`SO(3)`$ invariant four-dimensional metric can be chosen in the Bianchi IX form,
$$ds^2=(abc)^2dt^2+a^2\sigma _1^2+b^2\sigma _2^2+c^2\sigma _3^2$$
(4)
where $`\sigma _{1,2,3}`$ are a basis of $`SO(3)`$-invariant one-forms fulfilling the algebra
$$d\sigma _i=\frac{1}{2}ϵ_{ijk}\sigma _j\sigma _k.$$
(5)
An explicit representation for these one-forms is given in terms of the Euler parameterisation of $`SU(2)`$ as
$`\sigma _1`$ $`=`$ $`\mathrm{sin}\psi d\theta \mathrm{cos}\psi \mathrm{sin}\theta d\varphi ,`$
$`\sigma _2`$ $`=`$ $`\mathrm{cos}\psi d\theta \mathrm{sin}\psi \mathrm{sin}\theta d\varphi ,`$ (6)
$`\sigma _3`$ $`=`$ $`d\psi +\mathrm{cos}\theta d\varphi `$
The ranges of $`\theta `$, $`\varphi `$ and $`\psi `$ are $`[0,\pi ]`$, $`[0,2\pi ]`$ and $`[0,2\pi ]`$ respectively, up to identifications to be discussed below. Requiring the curvature to be self-dual puts three constraints on the coefficients $`a,b,c`$,
$$\frac{a^{}}{abc}=\frac{b^2+c^2a^2}{2bc}\lambda ,\text{etc }$$
(7)
plus the two others obtained by cyclic permutation of $`a,b,c`$. Here, the prime denotes differentiation with respect to the radial parameter $`t`$ (see below equation 3.1 for a relation between $`r`$ and $`t`$.) and $`\lambda `$ is an integration constant which is set to 1 in the Atiyah-Hitchin case. Following , we define $`w_1=bc,w_2=ca,w_3=ab`$ to rewrite the differential system as
$$d(w_1+w_2)/dt=2w_1w_2,\text{etc}$$
(8)
known as the Halphen system. Now we observe that Jacobi Theta functions give a simple solution of that system, since they fulfill the modular identity
$$\frac{d}{dt}\left(\frac{\vartheta _3^{}}{\vartheta _3}+\frac{\vartheta _4^{}}{\vartheta _4}\right)=\frac{2}{\pi }\frac{\vartheta _3^{}}{\vartheta _3}\frac{\vartheta _4^{}}{\vartheta _4}$$
(9)
where the complex modulus of the Theta function is $`\tau =it`$ <sup>4</sup><sup>4</sup>4The relation of modular forms to Halphen-like differential systems has been discussed in .. This implies that a solution of (8) can be chosen as
$`w_1`$ $`=`$ $`{\displaystyle \frac{2}{\pi }}{\displaystyle \frac{\vartheta _2^{}}{\vartheta _2}}={\displaystyle \frac{\pi }{6}}\left(E_2+\vartheta _3^4+\vartheta _4^4\right)`$
$`w_2`$ $`=`$ $`{\displaystyle \frac{2}{\pi }}{\displaystyle \frac{\vartheta _3^{}}{\vartheta _3}}={\displaystyle \frac{\pi }{6}}\left(E_2+\vartheta _2^4\vartheta _4^4\right),`$ (10)
$`w_3`$ $`=`$ $`{\displaystyle \frac{2}{\pi }}{\displaystyle \frac{\vartheta _4^{}}{\vartheta _4}}={\displaystyle \frac{\pi }{6}}\left(E_2\vartheta _2^4\vartheta _3^4\right),`$
which will be shown to satisfy the appropriate boundary conditions. The second equality on the righthand side follows from a standard equality involving the holomorphic Eisenstein series $`E_2`$. Note that $`w_1<0,w_2<0`$ and $`w_3>0`$, or equivalently $`a>0,b>0,c<0`$. The relation to Atiyah and Hitchin’s original formulae is detailed in Appendix A.
The above elliptic functions involve two different asymptotic regimes, $`\tau i\mathrm{}`$ and $`\tau 0`$, corresponding to the coincident limit and the large separation limit of the monopoles, respectively. Indeed, as $`t\mathrm{}`$, $`w_2`$ and $`w_3`$ become equal and the metric takes the form
$$ds^2=4\pi ^2e^{2\pi t}\left(\pi ^2dt^2+4e^{2\pi t}\sigma _1^2\right)+\frac{\pi }{2}\left(\sigma _2^2+\sigma _3^2\right)$$
(11)
up to exponentially small corrections. Changing variables to $`\stackrel{~}{u}=e^{\pi t}`$, this is recognized as a bolt singularity, which is a mere coordinate singularity if we take the quotient by the symmetry $`I_1:\theta \pi \theta ,\varphi \varphi +\pi ,\psi \psi `$. The resulting space is a double cover of the Atiyah-Hitchin manifold, the latter being obtained after modding out in addition by $`I_3:\psi \psi +\pi `$.
In the $`t0`$ limit, it is necessary to perform a modular transformation $`w_i(t)=\frac{1}{t}\frac{1}{t^2}w_j(1/t)`$, where $`j=3`$ for $`i=1`$; $`j=1`$ for $`i=3`$; and $`j=2`$ for $`i=2`$. This yields
$`w_1(t)`$ $`=`$ $`{\displaystyle \frac{1}{t}}{\displaystyle \frac{\pi }{t^2}}(4q+8q^2+\mathrm{})r/\pi `$
$`w_2(t)`$ $`=`$ $`{\displaystyle \frac{1}{t}}+{\displaystyle \frac{\pi }{t^2}}(4q8q^2+\mathrm{})r/\pi `$ (12)
$`w_3(t)`$ $`=`$ $`{\displaystyle \frac{1}{t}}+{\displaystyle \frac{\pi }{t^2}}({\displaystyle \frac{1}{2}}+4q^2+\mathrm{})r^2(12/r)/(2\pi )`$
where $`q=e^{\pi /t}`$ and we defined $`r=\pi /t`$. The first term in $`1/t`$ arises because of the anomalous modular property of the Eisenstein series $`E_2`$. The Atiyah-Hitchin metric (4) thus reduces to
$$ds^2=\frac{1}{2\pi }\left[\left(1\frac{2}{r}\right)\left(dr^2+r^2(\sigma _1^2+\sigma _2^2)\right)+\frac{4}{1\frac{2}{r}}\sigma _3^2\right]$$
(13)
where we recognize the Taub-NUT metric with mass parameter $`1`$. In particular, the asymptotic geometry is $`^3\times S_1`$ (mod $`_2`$) and $`r`$ can be identified as the distance between the monopoles, in units of the $`W`$-boson mass. The angle $`\psi `$ parameterizes the circle $`S_1`$ and is the coordinate that we called $`\sigma `$ before; giving momentum in that direction amounts to giving electric charge to the monopole, turning it into a dyon.
### 3.2 Four-fermion terms and curvature
We now would like to extract the exponential corrections from the metric described above. For this purpose, we find it convenient to use the language of the three-dimensional $`SU(2)`$ gauge theory with 8 supercharges. As described in Subsection 2.1, the three Higgs scalars combine with the pseudoscalar $`\sigma `$ dual to the gauge field in three dimensions to make an hypermultiplet taking values in the Atiyah-Hitchin manifold. The $`N=4`$ gauge theory has a symmetry group $`SU(2)_H\times SU(2)_V\times SU(2)_E`$ which is the product of the R-symmetry $`SU(2)_H`$ already present in the six-dimensional $`N=1`$ theory, the R-symmetry $`SU(2)_V`$ coming from compactification from 6 to 3 dimensions, and the Euclidean group in three dimensions. The fermions transform as $`(2,2,2)`$ and the bosons as $`(1,3+1,1)`$, so that $`SU(2)_V`$ is broken to a $`U(1)_V`$ subgroup by the Higgs vev. We choose to align the Higgs field along the “vertical” direction, $`\theta =0`$.
Given that the theory has 8 supersymmetries, the first quantum corrections arise in the metric of the scalars, or the four-fermion interactions which are related to the former by supersymmetry. It is thus convenient to concentrate on the four-fermion terms, which are contracted with the Riemann tensor of the bosonic moduli space. The antisymmetric product of four fermions transforms as
$$(5,1,1)+(1,5,1)+(1,1,5)+(3,3,3)+(3,3,1)+(3,1,3)+(1,3,3)+(1,1,1)$$
out of which we must keep the singlets under the Euclidean group $`SU(2)_E`$ and the R-symmetry group $`SU(2)_H`$, since the scalars are neutral under them. This only leaves the $`(1,5,1)+(1,1,1)`$ component.
On the other hand, the Riemann tensor of a hyperkähler manifold transforms as $`(1,5,1)`$. Indeed, since the Riemann tensor is self-dual, the independent components are $`R_{i;j}:=R_{0i0j}`$ which make a symmetric tensor of $`SU(2)`$, and its trace is zero by the cyclic property of the Riemann tensor. It is therefore contracted in the effective action with the $`(1,5,1)`$ part of the 4 fermion product only. Since the Higgs vev breaks $`SU(2)_V`$ to $`U(1)_V`$, we can split the fluctuations of the three scalars into a complex field $`w=x+iy`$ of $`U(1)_V`$ charge $`++=1`$, its complex conjugate $`\overline{w}=xiy`$ of $`U(1)_V`$ charge $`=1`$, and a real scalar $`z`$ of charge $`+=0`$. Taking into account the change of basis, the Riemann tensor decomposes into
$`R_{(2)}`$ $`=`$ $`R_{++;++}=e^{i\sigma }(R_{1;1}R_{2;2}+2iR_{1;2})`$
$`R_{(1)}`$ $`=`$ $`R_{++;+}=e^{i\sigma /2}(R_{1;3}+iR_{2;3})`$
$`R_{(0)}`$ $`=`$ $`R_{++;}R_{+;+}=R_{1;1}+R_{2;2}R_{3;3}`$ (14)
$`R_{(1)}`$ $`=`$ $`R_{;+}=e^{i\sigma /2}(R_{1;3}iR_{2;3})`$
$`R_{(2)}`$ $`=`$ $`R_;=e^{i\sigma }(R_{1;1}R_{2;2}2iR_{1;2})`$
These components hence appear in the effective action contracted with the fermion quadrilinears in such a way that the $`SU(2)_H\times U(1)_V\times SU(2)_E`$ quantum numbers are trivial. In particular, since $`\sigma `$ enters in the theory only through exponential effects $`e^{in\sigma }`$, with $`n`$ integer, it should be the case that $`R_{(1)}=R_{(1)}=0`$ exactly, and $`R_{(2)}=R_{(2)}=0`$ in perturbation theory. We shall shortly see that this indeed holds.
### 3.3 The non-perturbative expansion of the Riemann tensor
Having identified the precise components of the Riemann tensor to which instanton effects contribute, we can now proceed to evaluate them using the modular expression of the Atiyah-Hitchin metric. We use the orthonormal basis $`e_0=abcdt,e_i=a_i\sigma _i`$. In terms of the parameters $`a,b,c`$ entering (4), we find (see Appendix B and )
$$R_{1010}=\frac{a^{}}{a(abc)^2},R_{2020}=\frac{b^{}}{b(abc)^2},R_{3030}=\frac{c^{}}{c(abc)^2},R_{i0j0}=0\text{if}ij.$$
(15)
Using the differential system (7), this can be re-expressed in terms of the $`w`$’s as
$$R_{1010}=\frac{1}{w_2w_3}\frac{d}{dt}\left(\frac{(w_1w_3)^2+(w_1w_2)^2(w_2w_3)^2}{w_1^2w_2w_3}\right)\text{etc.}$$
(16)
and expressed in terms of modular forms using (3.1),
$`R_{1010}`$ $`=`$ $`{\displaystyle \frac{1}{(E_2+e_3)(E_2+e_4)}}\times `$
$`{\displaystyle \frac{d}{dt}}\left[{\displaystyle \frac{E_2^4+4E_2^3e_22E_2^2(e_2^2+2e_3e_4)4E_2E_6+e_2^4(e_3e_4)^24e_2E_6}{(E_2+e_2)^2(E_2+e_3)(E_2+e_4)}}\right].`$
The modular forms $`e_2,e_3,e_4`$ are the roots of the polynomial $`x^33E_4x2E_6`$ and are defined in Appendix A. This result, although not very enlightening, has the virtue of showing that $`R_{1;1}`$ is a modular form invariant under a $`\mathrm{\Gamma }^{}(2)`$ subgroup of $`Sl(2,)`$ exchanging $`e_3`$ and $`e_4`$ but leaving $`e_2`$ invariant. Similarly $`R_{2;2}`$ is invariant under $`\mathrm{\Gamma }^0(2)`$ and $`R_{3;3}`$ under $`\mathrm{\Gamma }^+(2)`$, while general modular transformations permute these groups. An important consequence is that $`R_{1;1}`$ and $`R_{2;2}`$ have the same $`q`$-expansion up to alternating signs. Hence the expansion of $`R_{(0)}`$ only involves even powers, while the expansion of $`R_{(\pm 2)}`$ only involves odd powers. As we shall see shortly, this is an important consistency check on the interpretation of the result as coming from instanton–anti-instantons bound states.
More precisely, keeping into account the anomalous modular transformation of $`w_i(t)`$, we find that the Riemann tensor has a large distance expansion
$`R_{(0)}`$ $`=`$ $`{\displaystyle \frac{4}{(r2)^3}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{P_{2n}(r)}{(r2)^{n+3}}}q^{2n}`$ (18)
$`R_{(\pm 2)}`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{P_{2n+1}(r)}{(r2)^{n+2}}}q^{2n+1}e^{\pm i\sigma }`$ (19)
where $`q=e^r`$, and $`P_n(r)`$ is a polynomial of order $`3([n/2]+1)`$ in $`r`$ with alternating integer coefficients. The first terms are easily computed using Mathematica,
$`P_1(r)`$ $`=`$ $`8(3+9r6r^2+r^3)`$
$`P_2(r)`$ $`=`$ $`32(654r+123r^2124r^3+66r^418r^5+2r^6)`$
$`P_3(r)`$ $`=`$ $`32(6117r+459r^2702r^3+495r^4162r^5+20r^6)`$
$`P_4(r)`$ $`=`$ $`64(6+333r2811r^2+9023r^314928r^4+14352r^58392r^6+2952r^7`$
$`576r^8+48r^9)`$
The leading $`r`$ power can be extracted at each order in $`q`$, yielding
$`R_{(0)}`$ $`=`$ $`{\displaystyle \frac{2^2}{r^3}}2^6r^2q^232^{10}r^4q^432^{15}r^6q^652^{19}r^8q^8372^{22}r^{10}q^{10}+\mathrm{}`$
$`R_{(\pm 2)}`$ $`=`$ $`e^{\pm i\sigma }\left(2^3rq+52^7r^3q^3+52^{13}r^5q^5+112^{16}r^7q^7+572^{19}r^9q^9+\mathrm{}\right)`$
Several remarks are in order about these results.
* First, as anticipated in the last section, the perturbative correction only occurs in $`R_{(0)}`$, as it should be since $`U(1)_V`$ is a symmetry preserved by perturbation theory.
* Second and most importantly, exponential corrections come with an arbitrary power $`n`$ of the semiclassical weight $`q=e^r`$, but only with zero-th power of $`e^{i\sigma }`$ for $`n`$ even, or first power for $`n`$ odd. This suggests that the correct interpretation of the non-perturbative exact result is rather a sum of bound states of $`n`$ instantons and $`\overline{n}`$ anti-instantons, with overall topological charge $`n\overline{n}=0`$ for $`R_{(0)}`$ and $`\pm 1`$ for $`R_{(\pm 2)}`$, as appropriate for charge conservation.
* Third, each power of $`q`$ comes in with an additional power of $`(2r)^{3/2}/(r2)^{1/2}`$, which indicates that the action of the semiclassical configuration is
$$S_{cl}=(n+\overline{n})\left(r\frac{1}{2}\mathrm{log}\frac{(2r)^3}{r2}\right)+i(n\overline{n})\sigma $$
(20)
This result agrees at large $`r`$ with the one-instanton result of , who found a contribution proportional to $`\varphi e^{8\pi ^2\varphi /g_3^2}`$, where the prefactor $`\varphi `$ originates from the one-loop determinant in the instanton background. Our result implies higher loop corrections in the one-instanton background. Besides, it predicts higher order contributions from bound states of instantons and anti-instantons.
* We may have considered instead the near-coincident limit $`r0`$. The Riemann tensor also admits a $`q`$-expansion in that regime, with $`q=e^{1/r}`$. In fact, this expansion only involves powers of $`q`$ and not of $`r`$, since no modular transformation is required, and all coefficients are integer numbers. We do not know of a semi-classical interpretation of this expansion.
This concludes our dissection of the Atiyah-Hitchin metric. The form of the expansion (18) therefore strongly suggests the interpretation of the non-perturbative corrections as contributions of instanton–anti-instanton bound states. We will now try to give further support for this unorthodox claim.
## 4 Discussion
Half-BPS saturated couplings in supersymmetric string or field theories are commonly thought to satisfy some sort of non-renormalization theorem, restricting them to receive contributions from a limited order in perturbation theory (usually one-loop), as well as exponential corrections coming from half-BPS instantons. This is a well-established fact in a few particular cases, including the prepotential of $`N=2`$ gauge theories in 4 dimensions, or some higher-derivative terms in string theory. In more general cases, this expectation is based on a simple zero-mode counting argument: an $`n`$-fermion vertex in the low energy effective action can only receive corrections from instanton configurations with less than $`n`$ fermionic zero-modes, and hence breaking at most $`n`$ supersymmetries. This crude argument would seem to rule out contributions from bound states of $`N`$ BPS instantons, which possess $`nN`$ zero-modes, but $`n(N1)`$ of them are usually lifted by the fermionic interactions in the action governing the collective coordinates. Non-BPS instantons (and in particular superpositions of instantons and anti-instantons) break more supersymmetries, and hence would seem to have too many fermionic zero-modes to make any contribution to the half-BPS couplings.
In the case at hand, we are interested in a four-fermion coupling in a theory with eight supersymmetries. In the language of the three-dimensional gauge theory, the relevant instanton configurations are ’t Hooft-Polyakov monopoles, which have four fermionic zero-modes. In the particular case where the three Higgs fields are aligned (which is automatic in the $`SU(2)`$ case), the $`N`$-monopole configurations are in fact exact solutions (in contrast to Yang-Mills instantons in four dimensional gauge theories with a Higgs vev), and hence the $`4N`$ fermionic zero-modes cannot be lifted. This argument thus predicts contributions from only one BPS monopole to the four-fermion coupling . Our expansion (18) clearly points to a flaw in this line of thought. Indeed, we seem to find contributions from arbitrary numbers of instantons and anti-instantons at the same time, with a net instanton number $`n\overline{n}=0,\pm 1`$ <sup>5</sup><sup>5</sup>5A crucial check on this interpretation is provided by the agreement between the net instanton number and the parity of the semi-classical weight of the instanton correction.. The argument above still correctly predicts the net instanton number, at least when we focus on the particular four-fermion vertex $`\psi _+\psi _+\psi _+\psi _+`$. This is in fact a simple consequence of $`U(1)_V`$ charge conservation, which is violated by $`2(n\overline{n})`$ units in a classical background with topological charge $`(n\overline{n})`$. An instanton–anti-instanton configuration on the other hand has vanishing topological and hence $`U(1)_V`$ charge, and therefore can be added without disturbing charge conservation.
The argument based on fermionic zero-mode counting can also be evaded. An instanton–anti-instanton configuration does break all supersymmetries, but the action of the supercharges on the instanton configuration does not generate zero-modes of the Dirac operator, since the configuration is not an exact saddle point of the action in the first place. It is in the limit of far separation only, which implies that the fermionic determinant vanishes in the large distance regime in the space of collective coordinates. At finite distance, the number of exact zero-modes is given by the index theorem, which yields $`4|n\overline{n}|`$ independent of $`n`$ and $`\overline{n}`$ separately <sup>6</sup><sup>6</sup>6The instanton field configuration being not self-dual anymore, the index theorem only counts the difference between zero-modes of different chiralities, but one does not expect more zero-modes than the minimum number $`4|n\overline{n}|`$.. For $`|n\overline{n}|=1`$, those are saturated by the four-fermion vertex. We are thus left to evaluate
$$A=_{_{(n,\overline{n})}}\mathrm{Pf}^{}(D)(\stackrel{}{det}(^2))^1e^{S_{cl}+i\sigma (n\overline{n})}$$
(21)
where $`_{(n,\overline{n})}`$ denotes the space of bosonic collective coordinates of the $`(n,\overline{n})`$ instanton–anti-instanton configuration, $`\mathrm{Pf}^{}(D)`$ the Pfaffian with the $`4|n\overline{n}|`$ zero-modes deleted, and $`det^{}(^2)`$ the bosonic fluctuation determinant with zero-modes deleted. The space of bosonic collective coordinates is not a completely well-defined notion, but makes sense in a dilute gas approximation. The integrand vanishes in the limit of far separation, but can give a finite value from the bulk of the moduli space. In fact, our analysis of the Atiyah-Hitchin metric predicts the value of the integral (21) for any $`|n\overline{n}|=0,1`$. It is interesting to note that the total weight in (20) is simply $`(n+\overline{n})`$ times the weight of an individual instanton (which does receive a logarithmic correction from the quantum fluctuations around it). The effects of the interactions between instantons and anti-instantons are encoded in the polynomial factors appearing in (18), and it would be interesting to analyze them in more detail. In particular, the fact that the coefficients are integer, albeit warranted by the underlying modular symmetry, suggests some topological or counting interpretation of these coefficients.
The physical picture that emerges from this discussion is therefore that the exponential corrections to the monopole moduli space arise from semi-classical configurations of $`2n+1`$ fundamental string world-sheets connecting the two D-strings in the D3-brane setting, or $`2n+1`$ Euclidean world-lines of $`W`$-bosons linking the two monopoles. As depicted in Figure 4, $`n+1`$ of them are oriented in one way and $`n`$ in the other, so that the total charge cancels (taking into account the net polarization induced by the $`\psi _+\psi _+\psi _+\psi _+`$ vertex). In the three-dimension-al gauge theory language, these are $`n+1`$ instantons and $`n`$ anti-instantons occurring at arbitrary Euclidean time.
We can ask how this picture generalizes to more than two monopoles, where only the large distance behaviour corresponding to massless exchange is known . Having surmounted the mental barrier of supersymmetry, we suggest that the semi-classical configurations controlling the $`n`$-SU(2) monopole dynamics are given by strings connecting monopoles in charge-conserving configurations. This includes two-particle interactions as in Figure 4, but also higher point interactions as depicted in Figure 5. Note that in the three-monopole case, the tantalizing $`Y`$-shaped configuration is not a minimum action configuration, since, due to charge conservation, it is really two $`Y`$’s with opposite orientations, which prefer to relax into two counter-rotating triangles. Of course, these semi-classical contributions are sub-leading with respect to the power corrections around the two-monopole interactions due to the presence of a third. These have been discussed in . It would be interesting to test our prediction against explicit results for multi-monopoles moduli spaces, as obtained for instance in .
It is also natural to ask if it is possible to precisely compute the higher-instanton effects in (18) from first principles. Equation (21) is first-principled, but seemingly amenable only to numerical computation. The embedding as a D1-D3 system seems more tractable, but would require some understanding of stacks of fundamental strings with opposite orientations. The NS5-brane setting may also offer some interesting insight, since the instanton is a D-string for which there exists a second quantized description allowing to consider stacks of them (see for instance ). The best bet may actually be the Heterotic setting, since the instanton corrections simply arise in that case from world-sheet instantons in the sigma-model $`^4/_2`$. The latter is unfortunately little understood due to the unresolved singularity. We can however speculate on a possible solution: since the Atiyah-Hitchin metric arises in the decoupling limit of the singularity, one may try to approximate the $`^4/_2`$ singularity by a $`CP^1`$ sigma-model describing the vanishing two-sphere. Interestingly enough, Cecotti and Vafa have found long ago similar instanton–anti-instanton contributions to the topological-anti-topological fusion coefficients in the $`(2,2)`$ supersymmetric $`CP^1`$ sigma model (the only example so far of such instanton–anti-instanton effects to our knowledge) <sup>7</sup><sup>7</sup>7We thank C. Vafa for bringing this work to our attention.. Even more tantalizingly, the $`t\overline{t}`$ equation controlling the $`CP^1`$ fusion coefficients is nothing but the $`SU(2)`$ Toda equation , while the Atiyah-Hitchin is also known to be controlled by an $`SU(\mathrm{})`$ Toda equation (see for a review). The study of $`(4,0)`$ supersymmetric sigma-models may therefore shed an interesting new light on the corrections to hypermultiplet manifolds.
###### Acknowledgments.
We are grateful to S. Cherkis, M. Douglas, R. Gopakumar, M. Gutperle, K. Hori, D. Tong, C. Vafa for valuable discussions, and especially to I. Bakas for correspondence and useful guidance into the literature. We would also like to thank the organizers of the workshop TMR 2000 in Tel Aviv, 7-11 January 2000, where this work was initiated, for their kind invitation and financial support. B.P. is supported in part by DOE grant DE-FG02-91ER-40654. A. H. is partially supported by the DOE under grant no. DE-FC02-94ER40818, by an A. P. Sloan Foundation Fellowship and by a DOE OJI award.
## Appendix A From elliptic functions to Theta functions
Atiyah and Hitchin linearize (8) by introducing a solution of the auxiliary equation ,
$$\frac{d^2u}{d\theta ^2}+\frac{1}{4\mathrm{sin}^2\theta }u=0$$
(22)
where $`dt/d\theta =1/u^2`$. Then $`w_{1,2,3}`$ are given by
$`w_1`$ $`=`$ $`u{\displaystyle \frac{du}{d\theta }}{\displaystyle \frac{1}{2\mathrm{sin}\theta }}u^2`$
$`w_2`$ $`=`$ $`u{\displaystyle \frac{du}{d\theta }}+{\displaystyle \frac{1}{2\mathrm{tan}\theta }}u^2`$ (23)
$`w_3`$ $`=`$ $`u{\displaystyle \frac{du}{d\theta }}+{\displaystyle \frac{1}{2\mathrm{sin}\theta }}u^2`$
The solution of (22) satisfying the appropriate limits is given by a complete elliptic integral of the first kind,
$$u=\sqrt{2\mathrm{sin}\theta }K\left(\mathrm{sin}\frac{\theta }{2}\right)$$
(24)
where
$$K(k)=_0^{\pi /2}\frac{d\varphi }{\sqrt{1k^2\mathrm{sin}^2\varphi }}$$
(25)
The region of infinite separation corresponds to $`k\mathrm{}`$.
In order to see the equivalence with our solution (3.1), recall that the elliptic integral (25) is naturally associated to an elliptic curve with complex modulus $`\tau =iK/K^{}`$, where $`K^{}=K((1k^2)^{1/2})=K(k^{})`$ with $`k^{}=\mathrm{cos}(\theta /2)`$. The relation between $`\tau `$ and $`K,K^{}`$ can be rewritten as
$$k=\frac{\vartheta _2^2}{\vartheta _3^2},k^{}=\frac{\vartheta _4^2}{\vartheta _3^2},K=\frac{\pi }{2}\vartheta _3^2$$
(26)
Expressing the derivative of the elliptic integral $`dK/dk=E/(kk^{}_{}{}^{}2)K/k`$ in terms of the elliptic integral of the second kind $`E(k)`$, related to Jacobi Theta functions by
$$E=\frac{\vartheta _3^4+\vartheta _4^4}{3\vartheta _3^4}K+\frac{\pi ^2}{12K}E_2(\tau )=K+\frac{\pi ^2}{12K}(E_2+e_4)$$
(27)
where we follow the conventions of for modular forms. we find that the solutions (A) can be rewritten as
$`w_1`$ $`=`$ $`{\displaystyle \frac{\pi }{6}}\left(E_2+\vartheta _3^4+\vartheta _4^4\right)`$
$`w_2`$ $`=`$ $`{\displaystyle \frac{\pi }{6}}\left(E_2+\vartheta _2^4\vartheta _4^4\right)`$ (28)
$`w_3`$ $`=`$ $`{\displaystyle \frac{\pi }{6}}\left(E_2\vartheta _2^4\vartheta _3^4\right)`$
which is precisely the same as in (3.1). The line element $`d\theta `$ can be simply related to $`dt`$ as $`d\theta =\pi \vartheta _2^2\vartheta _4^2dt`$, which agrees with the previous relation $`dt/d\theta =1/u^2`$. It is also useful to introduce the modular forms
$$e_2=\vartheta _3^4+\vartheta _4^4,e_3=\vartheta _2^4\vartheta _4^4,e_4=\vartheta _2^4\vartheta _3^4$$
(29)
which are the roots of the Weierstrass polynomial $`x^33E_4x2E_6`$ and are permuted under modular transformations.
## Appendix B Riemann tensor of the Bianchi IX ansatz
Let us now compute the curvature of the Atiyah-Hitchin metric, in the orthonormal basis $`e_0=abcdt,e_i=a_i\sigma _i`$. The Levi-Civita connection is easily found to be
$$\omega _{i0}=a_i^{}/(abc)\sigma _i,\omega _{jk}=ϵ_{ijk}(a_k^{}/abc+1)\sigma _k$$
(30)
and we can compute the curvature through $`R=d\omega +[\omega ,\omega ]`$,
$$R_{10}=\frac{d}{dt}\left(\frac{a^{}}{abc}\right)dt\sigma _1+\frac{2abca^{}(a^2+b^2c^2)b^{}c(a^2+c^2b^2)bc^{}}{2a^2b^2c^2}\sigma _2\sigma _3$$
$`R_{23}`$ $`=`$ $`{\displaystyle \frac{d}{dt}}\left({\displaystyle \frac{b^2+c^2a^2}{2bc}}\right)dt\sigma _1`$
$`{\displaystyle \frac{4b^{}c^{}+bc(a^2+b^2c^2)(a^2+c^2b^2)2a^2bc(b^2+c^2a^2)}{4(abc)^2}}\sigma _2\sigma _3`$
In particular, a necessary condition for the metric to be self-dual is
$$\frac{a^{}}{abc}=\frac{b^2+c^2a^2}{2bc}\lambda $$
(32)
which agrees with (7) for $`\lambda =1`$. Defining $`\widehat{a}=a^{}/(abc)`$, we can rewrite this as
$`R_{10}`$ $`=`$ $`R_{23}=\widehat{a}^{}dt\sigma _1+(\widehat{a}\widehat{b}\widehat{c}2\widehat{b}\widehat{c})\sigma _2\sigma _3`$
$`R_{20}`$ $`=`$ $`R_{31}=\widehat{b}^{}dt\sigma _2+(\widehat{b}\widehat{c}\widehat{a}2\widehat{c}\widehat{a})\sigma _3\sigma _1`$ (33)
$`R_{30}`$ $`=`$ $`R_{12}=\widehat{c}^{}dt\sigma _3+(\widehat{c}\widehat{a}\widehat{b}2\widehat{a}\widehat{b})\sigma _1\sigma _2`$
Using the identity
$$\frac{\widehat{a}^{}}{aabc}=\frac{\widehat{a}\widehat{b}\widehat{c}2\widehat{b}\widehat{c}}{bc}$$
(34)
we see that the Riemann tensor has the correct symmetry property.
|
warning/0005/hep-th0005277.html
|
ar5iv
|
text
|
# 1 Introduction
## 1 Introduction
Calogero-Moser models are one-dimensional dynamical systems with long range interactions having a remarkable property that they are integrable at both classical and quantum levels. In fact the integrability or more precisely the triangularity of the quantum Hamiltonian was first discovered by Calogero for the model with inverse square potential plus a confining harmonic force and by Sutherland for the particles on a circle with the inverse square potential. Later classical integrability of the models in terms of Lax pairs was proved by Moser . Olshanetsky and Perelomov showed that these models were based on $`A_r`$ root systems and generalisations of the models based on other root systems including the non-crystallographic ones were introduced .
In this paper we discuss quantum Calogero-Moser models with degenerate potentials, that is the rational with/without harmonic force, the hyperbolic and the trigonometric potentials based on all root systems. We demonstrate, based on previous works of universal Lax pairs for classical and quantum models , that various results known for the quantum $`A_r`$ models can be generalised to the models based on any root systems as well. They are (i) Construction of a complete set of quantum conserved quantities in terms of quantum Lax pairs and other methods. (ii) Universal proof of Liouville integrability for the rational, hyperbolic and trigonometric potential models. Namely, the quantum conserved quantities commute among themselves. (iii) Triangularity of the quantum Hamiltonian is demonstrated explicitly for all the models. In other words, the Hamiltonian is shown, in certain bases, to be decomposed into a sum of finite dimensional triangular matrices. Thus any eigenvalue equation can be solved by finite steps of linear algebraic processes only. This also gives the entire discrete spectrum of the models. As unique eigenfunctions of the Hamiltonian, generalisation of Jack polynomials and multivariable Laguerre (Hermite) polynomials are defined for all root systems. (iv) Equivalence of the quantum Lax pair method and that of so-called differential-reflection (Dunkl) operators is demonstrated. (v) For rational models with harmonic confining force, an algebraic construction of all excited states in terms of creation (annihilation) operators is achieved.
For the $`A_r`$ models, the Lax pairs, conserved quantities and their involution were discussed by many authors with varied degrees of completeness and rigour, see for example, , . The point (iv) was shown by Wadati and collaborators and point (v) was initiated by Perelomov and developed by Brink and collaborators and Wadati and collaborators . A rather different approach by Heckman and Opdam to Calogero-Moser models with degenerate potentials based on any root systems should also be mentioned in this connection.
For the general background and the motivations of this series of papers and the physical applications of the Calogero-Moser models with various potentials to lower dimensional physics, ranging from solid state to particle physics and supersymmetric Yang-Mills theory, we refer to our previous papers and references therein.
This paper is organised as follows. In section two, quantum Calogero-Moser Hamiltonian with degenerate potentials is introduced as a factorised form (2.5). Connection with root systems and the Coxeter invariance is emphasised. Some rudimentary facts of the root systems and reflections are summarised in Appendix A. A universal Coxeter invariant ground state wavefunction together with its energy eigenvalue are presented. In section three we show that all the excited states are also Coxeter invariant and that the Hamiltonian is triangular in certain bases. Complete sets of quantum conserved quantities are derived from quantum Lax operator $`L`$ in section four. Instead of the trace, the total sum of $`L^n`$ is conserved. That is Ts$`(L^n)=_{\mu ,\nu }(L^n)_{\mu \nu }`$, in which $``$ is a representation space of the Coxeter group. The details of the complete set for each root system are given in Appendix B. In section five the creation and annihilation operators for the rational models with harmonic force are derived. In section six, the equivalence of the Lax pair operator formalism and the so-called differential-reflection (Dunkl) operators is demonstrated and the quantum conserved quantities are expressed in terms of the latter. In section seven an algebraic construction of excited states in terms of the differential-reflection (Dunkl) operators for rational models with harmonic force is presented. The complete sets of explicit eigenfunctions for the rank two models are derived in terms of separation of variables based on the Coxeter invariant variables. Section eight gives a universal proof of the Liouville integrability for models with rational (without the confining force), hyperbolic and trigonometric potentials. For rational models with harmonic force, the involution is demonstrated for those based on classical root systems. A simple use of the quantum Lax pair with spectral parameter is mentioned. The final section is for summary and comments.
## 2 Quantum Calogero-Moser Models
In this section we briefly introduce the quantum Calogero-Moser models along with appropriate notation and background for the main body of this paper. A Calogero-Moser model is a Hamiltonian system associated with a root system $`\mathrm{\Delta }`$ of rank $`r`$, which is a set of vectors in $`𝐑^r`$ with its standard inner product. A brief review of the properties of the root systems and the associated reflections will be found in the Appendix A.
### 2.1 Factorised Hamiltonian
The dynamical variables of the Calogero-Moser model are the coordinates $`\{q_j\}`$ and their canonically conjugate momenta $`\{p_j\}`$, with the canonical commutation relations:
$$[q_j,p_k]=i\delta _{jk},[q_j,q_k]=[p_j,p_k]=0,j,k=1,\mathrm{},r.$$
(2.1)
These will be denoted by vectors in $`𝐑^r`$
$$q=(q_1,\mathrm{},q_r),p=(p_1,\mathrm{},p_r).$$
(2.2)
The momentum operator $`p_j`$ acts as
$$p_j=i\frac{}{q_j},j=1,\mathrm{},r.$$
As for the interactions we consider only the degenerate potentials, that is the rational (with/without harmonic force), hyperbolic and trigonometric potentials:
$$V(\rho q)=\{\begin{array}{ccc}\hfill 1/(\rho q)^2,& \text{type I},\hfill & \\ \hfill a^2/\mathrm{sinh}^2a(\rho q),& \text{type II},\hfill & \rho \mathrm{\Delta },\hfill \\ \hfill a^2/\mathrm{sin}^2a(\rho q),& \text{type III},\hfill & \end{array}$$
(2.3)
in which $`a`$ is an arbitrary real positive constant, determining the period of the trigonometric potentials. They imply integrability for all of the Calogero-Moser models based on the crystallographic root systems. Those models based on the non-crystallographic root systems, the dihedral group $`I_2(m)`$, $`H_3`$, and $`H_4`$, are integrable only for the rational potential. The rational potential models are also integrable if a confining harmonic potential
$$\frac{1}{2}\omega ^2q^2,\omega >0,\text{type V},$$
(2.4)
is added to the Hamiltonian. Since we will discuss the universal properties and solutions applicable to all the interaction types as well as those for specific interaction potentials, let us adopt the conventional nomenclature for them. We call the models with rational, hyperbolic, trigonometric and rational with harmonic force the type I, II, III and V models, respectively. (Type IV models have elliptic potentials which we will not discuss in this paper.)
The Hamiltonian for the quantum Calogero-Moser model can be written in a ‘factorised form’:
$``$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{j=1}{\overset{r}{}}}\left(p_ji{\displaystyle \frac{W}{q_j}}\right)\left(p_j+i{\displaystyle \frac{W}{q_j}}\right),`$ (2.5)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{j=1}{\overset{r}{}}}\left(p_j^2+\left({\displaystyle \frac{W}{q_j}}\right)^2\right)+{\displaystyle \frac{1}{2}}{\displaystyle \underset{j=1}{\overset{r}{}}}{\displaystyle \frac{^2W}{q_j^2}},`$ (2.6)
$`=`$ $`{\displaystyle \frac{1}{2}}p^2+{\displaystyle \frac{1}{2}}{\displaystyle \underset{\rho \mathrm{\Delta }_+}{}}g_{|\rho |}(g_{|\rho |}1)|\rho |^2V(\rho q)+({\displaystyle \frac{\omega ^2}{2}}q^2)_0.`$ (2.7)
It should be noted that the above factorised Hamiltonian (2.7) consists of an operator part $`\widehat{}`$, which is the Hamiltonian in the usual definition, and a constant $`_0`$ which is the ground state energy to be discussed later:
$``$ $`=`$ $`\widehat{}_0,`$ (2.8)
$`\widehat{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}p^2+{\displaystyle \frac{1}{2}}{\displaystyle \underset{\rho \mathrm{\Delta }_+}{}}g_{|\rho |}(g_{|\rho |}1)|\rho |^2V(\rho q)+({\displaystyle \frac{\omega ^2}{2}}q^2).`$ (2.9)
The real positive coupling constants $`g_{|\rho |}`$ are defined on orbits of the corresponding Coxeter group, i.e. they are identical for roots in the same orbit. That is, for the simple Lie algebra cases one coupling constant $`g_{|\rho |}=g`$ for all roots in simply-laced models and two independent coupling constants, $`g_{|\rho |}=g_L`$ for long roots and $`g_{|\rho |}=g_S`$ for short roots in non-simply laced models. For the $`I_2(m)`$ models, there is one coupling if $`m`$ is odd, and two independent ones if $`m`$ is even. Let us call them $`g_e`$ for even roots and $`g_o`$ for odd roots. Throughout this paper we consider the coupling constants at generic values. We parametrise the positive roots of $`I_2(m)`$ as
$$\rho _j=(\mathrm{cos}((j1)\pi /m),\mathrm{sin}((j1)\pi /m)),j=1,\mathrm{},m.$$
(2.10)
The $`H_3`$ and $`H_4`$ models have one coupling constant $`g_{|\rho |}=g`$, since these root systems are simply-laced. It should be noted that the operator part of the Hamiltonian $`\widehat{}`$ is strictly positive for $`g_{|\rho |}1`$.
The simplest way to introduce the factorised form is through supersymmetry , in which function $`W`$ is called a superpotential:
$$W(q)=\underset{\rho \mathrm{\Delta }_+}{}g_{|\rho |}\mathrm{ln}|w(\rho q)|+(\frac{\omega }{2}q^2),g_{|\rho |}>0,\omega >0.$$
(2.11)
The potential $`V(u)`$ (2.3) and the function $`w(u)`$ are related by
$`y(u)`$ $``$ $`{\displaystyle \frac{d}{du}}x(u),{\displaystyle \frac{dw(u)}{du}}/w(u)x(u),`$ (2.12)
$`V(u)`$ $`=`$ $`y(u)=x^2(u)+a^2\times \{\begin{array}{cc}\hfill 0& \text{rational},\hfill \\ \hfill 1& \text{hyperbolic},\hfill \\ \hfill 1& \text{trigonometric}.\hfill \end{array}`$ (2.16)
The following Table I gives these functions for each potential:
| potential | type | $`w(u)`$ | $`x(u)`$ | $`y(u)`$ |
| --- | --- | --- | --- | --- |
| rational | I & V | $`u`$ | $`1/u`$ | -$`1/u^2`$ |
| hyperbolic | II | $`\mathrm{sinh}au`$ | $`a\mathrm{coth}au`$ | -$`a^2/\mathrm{sinh}^2au`$ |
| trigonometric | III | $`\mathrm{sin}au`$ | $`a\mathrm{cot}au`$ | -$`a^2/\mathrm{sin}^2au`$ |
Table I: Functions appearing in the Lax pair and superpotential.
For proofs that the factorised Hamiltonian (2.6) actually gives the quantum Hamiltonian (2.7) for all the root systems and potentials see . It is easy to verify that for any potential $`V(u)`$, the Hamiltonian is invariant under reflection of the phase space variables in the hyperplane perpendicular to any root
$$(s_\alpha (p),s_\alpha (q))=(p,q),\alpha \mathrm{\Delta }$$
(2.17)
with $`s_\alpha `$ defined by (A.2).
Some remarks are in order. For all of the root systems and for any choice of potential (2.3), the Calogero-Moser model has a hard repulsive potential $`1/(\alpha q)^2`$ near the reflection hyperplane $`H_\alpha =\{q𝐑^r,\alpha q=0\}`$. The strength of the singularity is given by the coupling constant $`g_{|\alpha |}(g_{|\alpha |}1)`$ which is independent of the choice of the normalisation of the roots. (Thus for rational models with/without harmonic force there is equivalence, $`A_2I_2(3)`$, $`B_2I_2(4)`$, $`G_2I_2(6)`$, $`B_rC_rBC_r`$.) This determines the form of the ground state wavefunction, as we will see in subsection 2.2. This repulsive potential is classically insurmountable. Thus the motion is always confined within one Weyl chamber. This feature allows us to constrain the configuration space to the principal Weyl chamber ($`\mathrm{\Pi }`$: set of simple roots, see Appendix A)
$$PW=\{q𝐑^r|\alpha q>0,\alpha \mathrm{\Pi }\},$$
(2.18)
without loss of generality. In the case of the trigonometric potential, the configuration space is further limited due to the periodicity of the potential to
$$PW_T=\{q𝐑^r|\alpha q>0,\alpha \mathrm{\Pi },\alpha _hq<\pi /a\},$$
(2.19)
where $`\alpha _h`$ is the highest root.
The fact that the classical motions are confined in the Weyl chamber (alcove) $`PW(PW_T)`$ does not necessarily mean that the corresponding quantum wavefunctions vanish identically outside of the region. On the contrary, as we will see soon, the ground state wavefunction (subsection 2.2) and all the other excited states wavefunctions (section 3) are Coxeter invariant, reflecting the Coxeter invariance of the Hamiltonian (2.17). In the early years of Calogero-Moser models in which those based on the $`A_r`$ root system were mainly discussed, these Coxeter invariant solutions were considered as totally symmetric states of bosonic systems. We will not, however, adopt this interpretation, for in the models based on the other root systems the reflection is not the same as particle interchange. The quantum theory we are discussing is the so-called first quantised theory. That is, the notions of identical particles and the associated statistics are non-existent.
### 2.2 Ground state wavefunction and energy
One merit of the factorised Hamiltonian (2.5) is the ease of derivation of the ground state wavefunction and of the Hamiltonian (3.3) derived by the similarity transformation in terms of the ground state wavefunction. Supersymmetric formulation of the Calogero-Moser models provides a natural setting for the introduction of the factorised Hamiltonian. The universal ground state wavefunction is
$$\mathrm{\Phi }_0(q)=e^{W(q)}=\underset{\rho \mathrm{\Delta }_+}{}|w(\rho q)|^{g_{|\rho |}}e^{\frac{\omega }{2}q^2}.$$
(2.20)
The exponential factor $`e^{\frac{\omega }{2}q^2}`$ exists only for the rational potential case with the harmonic confining force. It is easy to see that it is an eigenstate of the Hamiltonian (2.5) with zero eigenvalue:
$$\mathrm{\Phi }_0(q)=\frac{1}{2}\underset{j=1}{\overset{r}{}}\left(p_ji\frac{W}{q_j}\right)\left(p_j+i\frac{W}{q_j}\right)\mathrm{\Phi }_0(q)=0,$$
(2.21)
since it satisfies
$$\left(p_j+i\frac{W}{q_j}\right)e^{W(q)}=0,j=1,\mathrm{},r.$$
(2.22)
By using the decomposition of the factorised Hamiltonian into the operator Hamiltonian (2.9) and a constant, we obtain
$$\widehat{}e^W\left(\frac{1}{2}p^2+\frac{1}{2}\underset{\rho \mathrm{\Delta }_+}{}g_{|\rho |}(g_{|\rho |}1)|\rho |^2V(\rho q)+(\frac{\omega ^2}{2}q^2)\right)e^W=_0e^W.$$
(2.23)
In other words, the above solution (2.20) provides an eigenstate of the Hamiltonian operator $`\widehat{}`$ with energy $`_0`$. The fact that it is a ground state (for type I, III and V) can be easily shown within the framework of the supersymmetric model thanks to the positivity of the supersymmetric Hamiltonian. It should be stressed that $`_0`$ is determined purely algebraically , without really applying the operator on the left hand side of (2.23) to the wavefunction. This type of ground states has been known for some time. It is derived by various methods, see for example , and also by using supersymmetric quantum mechanics for the models based on classical root systems .
The ground state energy for the rational potential cases are
$$_0=\{\begin{array}{ccc}0& & \text{type I},\hfill \\ \omega \left(\frac{r}{2}+_{\rho \mathrm{\Delta }_+}g_{|\rho |}\right)& & \text{type V}.\hfill \end{array}$$
(2.24)
The same for the hyperbolic and trigonometric potential cases are
$$_0=2a^2\varrho ^2\times \{\begin{array}{cc}1& \text{hyperbolic},\hfill \\ 1& \text{trigonometric},\hfill \end{array}$$
(2.25)
in which
$$\varrho =\frac{1}{2}\underset{\rho \mathrm{\Delta }_+}{}g_{|\rho |}\rho $$
(2.26)
can be considered as a ‘deformed Weyl vector’ . Again these formulas are universal. That is they apply to all of the Calogero-Moser models based on any root systems. A negative $`_0`$ for the obviously positive Hamiltonian of the hyperbolic potential model indicates that the interpretation of $`e^W`$ as an eigenfunction is not correct. This function diverges as $`|q|\mathrm{}`$ for the hyperbolic and the rational potential cases, destroying the hermiticity of the Hamiltonian. Obviously we have
$$_{PW(PW_T)}e^{2W(q)}𝑑q=\{\begin{array}{cc}\mathrm{}& \text{: type I and II},\hfill \\ \text{finite}& \text{: type III and V},\hfill \end{array}$$
(2.27)
in which $`PW`$ and $`PW_T`$ denote that the integration is over the regions defined in (2.18) and (2.19). Naturally, most existing results in quantum Calogero-Moser models are for the models with trigonometric potential and the rational potential with harmonic force which have normalisable states and discrete spectra.
It should be remarked that the domain of the universal ground state wavefunction $`\mathrm{\Phi }_0`$ could be considered as the entire $`𝐑^r`$ space except for the points on the reflection hyperplanes, that is the disjoint union of all the Weyl chambers (alcoves), instead of the initial Weyl chamber/alcove ($`PW`$, $`PW_T`$) in which the classical motions are restricted due to the singular potential. In fact, $`\mathrm{\Phi }_0`$ and $`W`$ are characterised as Coxeter invariant:
$$\stackrel{ˇ}{s}_\rho \mathrm{\Phi }_0=\mathrm{\Phi }_0,\stackrel{ˇ}{s}_\rho W=W,\rho \mathrm{\Delta },$$
(2.28)
in which $`\stackrel{ˇ}{s}_\rho `$ is the representation of the reflection in the function space. For an arbitrary function $`f`$ of $`q`$, its action is defined by
$$(\stackrel{ˇ}{s}_\rho f)(q)=f(s_\rho (q)).$$
(2.29)
This definition can be generalised to the entire Coxeter group $`G_\mathrm{\Delta }`$: for an arbitrary element $`g`$ of $`G_\mathrm{\Delta }`$, $`\stackrel{ˇ}{g}`$ is defined by:
$$(\stackrel{ˇ}{g}f)(q)=f(g^1(q)),gG_\mathrm{\Delta }.$$
(2.30)
In the rest of this paper we discuss mainly the type III and V models which have normalisable states and discrete spectra.
## 3 Coxeter invariant excited states, triangularity and spectrum
In this section we show that all the excited states wavefunctions are Coxeter invariant, too. In other words, the Fock space consists of Coxeter invariant functions only. With the knowledge of the ground state wavefunction $`e^W`$, the other states of the Calogero-Moser models can be easily obtained as eigenfunctions of a differential operator $`\stackrel{~}{}`$ obtained from $``$ by a similarity transformation:
$`\stackrel{~}{}`$ $`=`$ $`e^We^W`$ (3.1)
$`=`$ $`e^W\left({\displaystyle \frac{1}{2}}p^2+{\displaystyle \frac{1}{2}}{\displaystyle \underset{\rho \mathrm{\Delta }_+}{}}g_{|\rho |}(g_{|\rho |}1)|\rho |^2V(\rho q)+({\displaystyle \frac{\omega ^2}{2}}q^2)_0\right)e^W,`$
$$\stackrel{~}{}\mathrm{\Psi }_\lambda =\lambda \mathrm{\Psi }_\lambda \mathrm{\Psi }_\lambda e^W=\lambda \mathrm{\Psi }_\lambda e^W.$$
(3.2)
Thanks to the factorised form of the Hamiltonian $``$ (2.5), (2.6), the transformed Hamiltonian $`\stackrel{~}{}`$ takes a simple form:
$$\stackrel{~}{}=\frac{1}{2}\underset{j=1}{\overset{r}{}}\left(\frac{^2}{q_j^2}+2\frac{W}{q_j}\frac{}{q_j}\right).$$
(3.3)
The Coxeter invariance of $`W`$ implies those of $``$ and $`\stackrel{~}{}`$:
$$\stackrel{ˇ}{s}_\rho \stackrel{ˇ}{s}_\rho =,\stackrel{ˇ}{s}_\rho \stackrel{~}{}\stackrel{ˇ}{s}_\rho =\stackrel{~}{},\rho \mathrm{\Delta }.$$
(3.4)
For type I and III models we introduce proper bases of Fock space consisting of Coxeter invariant functions and show that the above Hamiltonian $`\stackrel{~}{}`$ (3.3) is triangular in these bases. This establishes the integrability of the type I and III models universally <sup>1</sup><sup>1</sup>1Triangularity of the $`A_r`$ type V and III Hamiltonians was noted in the original papers of Calogero and Sutherland . That of rank two models in the Coxeter invariant bases was shown in . and also gives the entire spectrum of the Hamiltonian, see (3.12), (3.13) and (3.44).
### 3.1 Rational potential with harmonic force
First, let us determine the structure of the set of eigenfunctions of the transformed Hamiltonian $`\stackrel{~}{}`$, for the type V models:
$$\stackrel{~}{}=\omega q\frac{}{q}\frac{1}{2}\underset{j=1}{\overset{r}{}}\frac{^2}{q_j^2}\underset{\rho \mathrm{\Delta }_+}{}\frac{g_{|\rho |}}{\rho q}\rho \frac{}{q}.$$
(3.5)
Obviously a constant and $`\omega q^2_0/\omega `$ are its eigenfunctions with eigenvalue 0 and $`2\omega `$, respectively. Let us suppose that a polynomial $`P(q)`$ is an eigenfunction of $`\stackrel{~}{}`$:
$$\stackrel{~}{}P(q)=\lambda P(q).$$
(3.6)
Due to the Coxeter invariance of $`\stackrel{~}{}`$ (3.4), we know that $`\stackrel{ˇ}{s}_\rho P`$ together with the difference
$$Q=(1\stackrel{ˇ}{s}_\rho )P$$
are also eigenfunctions with the same eigenvalue:
$$\stackrel{~}{}Q(q)=\lambda Q(q),$$
(3.7)
if the latter is not identically zero. Since $`Q`$ is a polynomial which is odd under reflection $`\stackrel{ˇ}{s}_\rho `$
$$\stackrel{ˇ}{s}_\rho Q(q)=Q(q),$$
it can be factorised as
$$Q(q)=(\rho q)^{2n+1}\stackrel{~}{Q}(q),\stackrel{~}{Q}|_{\rho q=0}0,$$
(3.8)
with a non-negative integer $`n`$ and a polynomial $`\stackrel{~}{Q}`$. By substituting (3.8) into (3.7) and using the explicit form of $`\stackrel{~}{}`$ near the reflection hyperplane $`\rho q=0`$, we obtain
$$(\rho q)^{2n1}(2n+1)(n+g_{|\rho |})\rho ^2\stackrel{~}{Q}+𝒪[(\rho q)^{2n}]=\lambda (\rho q)^{2n+1}\stackrel{~}{Q},$$
(3.9)
which would imply the vanishing of $`\stackrel{~}{Q}`$ on the reflection hyperplane
$$\stackrel{~}{Q}|_{\rho q=0}=0,$$
an obvious contradiction. Thus we are led to the conclusion that the eigenfunctions are Coxeter invariant polynomials and that the Hamiltonian $``$ (3.5) maps a Coxeter invariant polynomial to another.
An obvious basis in the space of Coxeter invariant polynomials is the homogeneous polynomials of various degrees. This basis has a natural order given by the degree. For a given degree the space of homogeneous Coxeter invariant polynomials is finite-dimensional. The explicit form of $`\stackrel{~}{}`$ (3.5) shows that it is lower triangular in this basis and the diagonal elements are $`\omega \times degree`$ as given by the first term. Independent Coxeter invariant polynomials exist at the degrees $`f_j`$ listed in Table II:
$$f_j=1+e_j,j=1,\mathrm{},r,$$
(3.10)
in which $`\{e_j\}`$, $`j=1,\mathrm{},r`$, are the exponents of $`\mathrm{\Delta }`$. Let us denote them by
$$z_1(q),\mathrm{},z_r(q);z_j(\kappa q)=\kappa ^{f_j}z_j(q).$$
(3.11)
| $`\mathrm{\Delta }`$ | $`f_j=1+e_j`$ | $`\mathrm{\Delta }`$ | $`f_j=1+e_j`$ |
| --- | --- | --- | --- |
| $`A_r`$ | $`2,3,4,\mathrm{},r+1`$ | $`E_8`$ | $`2,8,12,14,18,20,24,30`$ |
| $`B_r`$ | $`2,4,6,\mathrm{},2r`$ | $`F_4`$ | $`2,6,8,12`$ |
| $`C_r`$ | $`2,4,6,\mathrm{},2r`$ | $`G_2`$ | $`2,6`$ |
| $`D_r`$ | $`2,4,\mathrm{},2r2;r`$ | $`I_2(m)`$ | $`2,m`$ |
| $`E_6`$ | $`2,5,6,8,9,12`$ | $`H_3`$ | $`2,6,10`$ |
| $`E_7`$ | $`2,6,8,10,12,14,18`$ | $`H_4`$ | $`2,12,20,30`$ |
Table II: The degrees $`f_j`$ in which independent Coxeter invariant polynomials exist.
Thus we arrive at:
the quantum Calogero-Moser models with the rational potential and the harmonic confining force is algebraically solvable for any (crystallographic and non-crystallographic) root system $`\mathrm{\Delta }`$. The spectrum of the operator Hamiltonian $`\widehat{}`$ is
$$\omega N+_0,$$
(3.12)
with a non-negative integer $`N`$ which can be expressed as
$$N=\underset{j=1}{\overset{r}{}}n_jf_j,n_j𝐙_+,$$
(3.13)
and the degeneracy of the above eigenvalue (3.12) is the number of different solutions of (3.13) for given $`N`$. This is generalisation of Calogero’s original argument for the $`A_r`$ model to the models based on arbitrary root systems. Now let us denote by $`\stackrel{}{N}`$ the set of non-negative integers in (3.13):
$$\stackrel{}{N}=(n_1,n_2,\mathrm{},n_r),$$
(3.14)
and by $`\varphi _\stackrel{}{N}(q)`$ the homogeneous Coxeter invariant polynomial of determined by $`\stackrel{}{N}`$ and the above basis $`\{z_j\}`$ (3.11):
$$\varphi _\stackrel{}{N}(q)=\underset{j=1}{\overset{r}{}}z_j^{n_j}(q).$$
(3.15)
As shown above, there exists a unique eigenstate $`\psi _\stackrel{}{N}(q)`$ for each $`\varphi _\stackrel{}{N}(q)`$:
$`\psi _\stackrel{}{N}(q)`$ $`=`$ $`\varphi _\stackrel{}{N}(q)+{\displaystyle \underset{\stackrel{}{N}^{}<\stackrel{}{N}}{}}d_\stackrel{}{N}^{}\varphi _\stackrel{}{N}^{}(q),d_\stackrel{}{N}^{}:const,`$ (3.16)
$`\stackrel{~}{}\psi _\stackrel{}{N}(q)`$ $`=`$ $`\omega N\psi _\stackrel{}{N}(q).`$ (3.17)
It satisfies the orthogonality relation
$$(\psi _\stackrel{}{N},\varphi _\stackrel{}{N}^{})=0,\stackrel{}{N}^{}<\stackrel{}{N},$$
(3.18)
with respect to the inner product in $`PW`$:
$$(\psi ,\phi )=_{PW}\psi ^{}(q)\phi (q)e^{2W(q)}𝑑q.$$
(3.19)
These polynomials $`\{\psi _\stackrel{}{N}(q)\}`$ are generalisations of the multivariable Laguerre (Hermite) polynomials known for the $`A_r`$ ($`B_r`$, $`D_r`$) root systems to arbitrary root systems.
Some remarks are in order.
1. There is no Coxeter invariant linear function in $`q`$. The quadratic invariant polynomial $`q^2=qq`$ exists in all the root systems. This corresponds to the universal fact that $`f_1=2,(e_1=1)`$ for all the root systems. Moreover, this is related to the fact that a special sub-series of the excited states with $`N=2n_1=n_1f_1`$, $`n_j=0,`$ $`j2`$, can be expressed universally in terms of Laguerre polynomials in $`q^2`$. This will be discussed at the end of section 5 and in subsection 7.2.
2. The other Coxeter invariants corresponding to the degrees $`f_2,\mathrm{},f_r`$ could be interpreted as special ‘angular’ variables of a unit sphere $`S^{r1}`$ ($`q^2=1`$), with the first Coxeter invariant $`\sqrt{q^2}`$ being the radial coordinate. These would provide proper variables for describing solutions. Solutions in terms of separation of variables are in general possible only for the simplest cases, that is the rank two models, which will be demonstrated in subsection 7.3.
3. For $`\mathrm{\Delta }=A_1`$, the simplest root system of rank one, the Hamiltonian $`\stackrel{~}{}`$ can be rewritten in terms of a Coxeter invariant variable $`u=\omega q^2`$ as:
$$\stackrel{~}{}=\omega q\frac{d}{dq}\frac{1}{2}\frac{d^2}{dq^2}\frac{g}{q}\frac{d}{dq}=2\omega \left\{u\frac{d^2}{du^2}+(g+\frac{1}{2}u)\frac{d}{du}\right\}.$$
(3.20)
The Laguerre polynomial satisfying the differential equation
$$\left\{u\frac{d^2}{du^2}+(g+\frac{1}{2}u)\frac{d}{du}+n\right\}L_n^{(g\frac{1}{2})}(u)=0,$$
(3.21)
provides an eigenfunction with eigenvalue $`2\omega n`$, which corresponds to the eigenvalue $`2\omega n+_0`$ of $`\widehat{}`$. This is a well-known result.
4. Triangularity of the type I models is also obvious from the above argument.
### 3.2 Trigonometric potential
Here we consider those root systems associated with Lie algebras. In order to determine the excited states of the type III models, we have to consider the periodicity. The superpotential $`W`$ and the Hamiltonian $``$ are invariant under the following translation:
$$W(q^{})=W(q),(p,q^{})=(p,q),q^{}=q+l^{}\pi /a,$$
(3.22)
in which $`l^{}`$ is an element of the dual weight lattice, that is
$$l^{}=\underset{j=1}{\overset{r}{}}l_j\frac{2}{\alpha _j^2}\lambda _j,l_j𝐙,\alpha _j\mathrm{\Pi },\alpha _j^{}\lambda _k=\delta _{jk}.$$
(3.23)
As is well known in quantum mechanics with periodic potentials, the wavefunctions diagonalising the translation operators are expressed as
$$e^{2ia\mu q}\left(\underset{\alpha L(\mathrm{\Delta })}{}b_\alpha e^{2ia\alpha q}\right)e^W,b_\alpha :const,L(\mathrm{\Delta }):\text{root lattice},$$
(3.24)
in which a vector $`\mu 𝐑^r`$ is as yet unspecified. In other words, up to the overall phase factor $`e^{2ia\mu q}`$, this is a Fourier expansion in terms of the simple roots.
Let $`P_T(q)`$ be a polynomial in $`e^{\pm 2ia\alpha _jq}`$, $`\alpha _j\mathrm{\Pi }`$ and suppose that a function $`\varphi (q)`$
$$\varphi (q)=e^{2ia\mu q}P_T(q),\mu 𝐑^r,$$
(3.25)
is an eigenfunction of $`\stackrel{~}{}`$:
$$\stackrel{~}{}\varphi (q)=\lambda \varphi (q),$$
(3.26)
in which the explicit form of $`\stackrel{~}{}`$ is given by
$$\stackrel{~}{}=\frac{1}{2}\underset{j=1}{\overset{r}{}}\frac{^2}{q_j^2}a\underset{\rho \mathrm{\Delta }_+}{}g_{|\rho |}\mathrm{cot}(a\rho q)\rho \frac{}{q}.$$
(3.27)
As above, due to the Coxeter invariance of $`\stackrel{~}{}`$ (3.4), we know that $`\stackrel{ˇ}{s}_\rho \varphi `$ together with the difference
$$\phi =(1\stackrel{ˇ}{s}_\rho )\varphi $$
are also eigenfunctions with the same eigenvalue:
$$\stackrel{~}{}\phi (q)=\lambda \phi (q).$$
(3.28)
Since $`\phi `$ is odd under reflection $`\stackrel{ˇ}{s}_\rho `$
$$\stackrel{ˇ}{s}_\rho \phi (q)=\phi (q),$$
it can be expressed as
$$\phi (q)=(\rho q)^{2n+1}\stackrel{~}{\phi }(q)+𝒪[(\rho q)^{2n+3}],\stackrel{~}{\phi }|_{\rho q=0}0,$$
(3.29)
in a neighbourhood of the reflection hyperplane $`\rho q=0`$. In this neighbourhood, the singularity structure of $`\stackrel{~}{}`$ for the trigonometric potential is the same as that of $`\stackrel{~}{}`$ for the rational potential discussed in the previous subsection. Thus we obtain, as before, a contradiction $`\stackrel{~}{\phi }|_{\rho q=0}=0`$. In other words, the eigenfunction $`\varphi `$ (3.25) must be Coxeter invariant. This in turn requires that the unspecified vector $`\mu `$ in (3.25) to be an element of the weight lattice
$$\mu \mathrm{\Lambda }(\mathrm{\Delta }).$$
(3.30)
Since
$$\stackrel{ˇ}{s}_\rho \varphi (q)=e^{2ias_\rho (\mu )q}\stackrel{ˇ}{s}_\rho P_T(q)=e^{2ia(\mu q\rho ^{}\mu \rho q)}\stackrel{ˇ}{s}_\rho P_T(q),$$
the following condition is necessary, but not sufficient,
$$\rho ^{}\mu 𝐙,\rho \mathrm{\Delta }$$
(3.31)
for the Coxeter invariance of $`\varphi `$. Thus we arrive at (3.30).
Let us introduce a basis for the Coxeter invariant functions of the form (3.25). Let $`\lambda `$ be a dominant weight
$$\lambda =\underset{j=1}{\overset{r}{}}m_j\lambda _j,m_j𝐙_+,$$
(3.32)
and $`W_\lambda `$ be the orbit of $`\lambda `$ by the action of the Weyl group:
$$W_\lambda =\{\mu \mathrm{\Lambda }(\mathrm{\Delta })|\mu =g(\lambda ),gG_\mathrm{\Delta }\}.$$
(3.33)
We define
$$\varphi _\lambda (q)\underset{\mu W_\lambda }{}e^{2ia\mu q},$$
(3.34)
which is Coxeter invariant. The set of functions $`\{\varphi _\lambda \}`$ has an order $`>`$:
$$|\lambda |^2>|\lambda ^{}|^2\varphi _\lambda >\varphi _\lambda ^{}.$$
(3.35)
Next we show that $`\stackrel{~}{}`$ is lower triangular in this basis. By using (3.27) we obtain
$$\stackrel{~}{}\varphi _\lambda =2a^2\lambda ^2\varphi _\lambda 2ia^2\underset{\rho \mathrm{\Delta }_+}{}\underset{\mu W_\lambda }{}g_{|\rho |}\mathrm{cot}(a\rho q)(\rho \mu )e^{2ia\mu q}.$$
(3.36)
First let us fix one positive root $`\rho `$ and a weight $`\mu `$ in $`W_\lambda `$ such that $`\rho \mu 0`$. Then
$$\mu ^{}s_\rho (\mu )=\mu (\rho ^{}\mu )\rho W_\lambda ,\rho \mu ^{}=\rho \mu .$$
(3.37)
Without loss of generality we may assume
$$\rho ^{}\mu =k>0,k𝐙.$$
(3.38)
The contribution of the pair $`(\mu ,\mu ^{})`$ in the summation of (3.36) reads
$`|\rho \mu |e^{2ai\mu q}(1e^{2aik\rho q})\mathrm{cot}(a\rho q)`$ (3.39)
$`=`$ $`i|\rho \mu |\left(e^{2ai\mu q}+e^{2ai\mu ^{}q}+2{\displaystyle \underset{j=1}{\overset{k1}{}}}e^{2ai(\mu j\rho )q}\right),`$
which is the generalisation of Sutherland’s fundamental identity eq(15) in to arbitrary root systems. The summation in the expression correspond to $`\varphi _\lambda ^{}`$ with $`\lambda ^{}`$ being lower than $`\lambda `$. Thus (3.36) reads
$$\stackrel{~}{}\varphi _\lambda =2a^2\lambda ^2\varphi _\lambda +2a^2\underset{\rho \mathrm{\Delta }_+}{}\underset{\mu W_\lambda }{}g_{|\rho |}|\rho \mu |e^{2ia\mu q}+\underset{|\lambda ^{}|<|\lambda |}{}c_\lambda ^{}\varphi _\lambda ^{},$$
(3.40)
in which $`\{c_\lambda ^{}\}`$’s are constants. It is easy to see that ($`\mu =g(\lambda )`$, $`gG_\mathrm{\Delta }`$)
$$\underset{\rho \mathrm{\Delta }_+}{}g_{|\rho |}|\rho \mu |=\underset{\rho \mathrm{\Delta }_+}{}g_{|\rho |}|g(\rho )\lambda |=(\underset{\rho \mathrm{\Delta }_+}{}g_{|\rho |}\rho )\lambda =2\varrho \lambda ,$$
(3.41)
which is independent of $`\mu `$. Thus we have demonstrated the triangularity of $`\stackrel{~}{}`$:
$$\stackrel{~}{}\varphi _\lambda =2a^2(\lambda ^2+2\varrho \lambda )\varphi _\lambda +\underset{|\lambda ^{}|<|\lambda |}{}c_\lambda ^{}\varphi _\lambda ^{},$$
(3.42)
or that of $`\widehat{}`$
$$\widehat{}\varphi _\lambda e^W=2a^2(\lambda +\varrho )^2\varphi _\lambda e^W+\underset{|\lambda ^{}|<|\lambda |}{}c_\lambda ^{}\varphi _\lambda ^{}e^W,$$
(3.43)
with the eigenvalue
$$2a^2(\lambda +\varrho )^2.$$
(3.44)
In other words, for each dominant weight $`\lambda `$ there exists an eigenstate of $`\stackrel{~}{}`$ with eigenvalue proportional to $`\lambda (\lambda +2\varrho )`$. Let us denote this eigenfunction by $`\psi _\lambda (q)`$:
$`\psi _\lambda (q)`$ $`=`$ $`\varphi _\lambda (q)+{\displaystyle \underset{|\lambda |^{}<|\lambda |}{}}d_\lambda ^{}\varphi _\lambda ^{}(q),d_\lambda ^{}:const,`$ (3.45)
$`\stackrel{~}{}\psi _\lambda (q)`$ $`=`$ $`2a^2\lambda (\lambda +2\varrho )\psi _\lambda (q),`$ (3.46)
and call it a generalised Jack polynomial -. It satisfies the orthogonality relation
$$(\psi _\lambda ,\varphi _\lambda ^{})=0,|\lambda |^{}<|\lambda |,$$
(3.47)
with respect to the inner product in $`PW_T`$:
$$(\psi ,\phi )=_{PW_T}\psi ^{}(q)\phi (q)e^{2W(q)}𝑑q.$$
(3.48)
In the $`A_r`$ model, specifying a dominant weight $`\lambda `$ is the same as giving a Young diagram which designates a Jack polynomial. It should be emphasised, however, that $`\{\psi _\lambda \}`$ are not identical to the Jack polynomials even for the $`A_r`$ root systems, because of different treatments of the center of mass coordinates. Detailed properties of these polynomials for various root systems will be published elsewhere.
Thus we arrive at:
the quantum Calogero-Moser models with the trigonometric potential are algebraically solvable for any crystallographic root system $`\mathrm{\Delta }`$. The spectrum of the Hamiltonian $`\widehat{}`$ is given by (3.44) in which $`\lambda `$ is an arbitrary dominant weight. This is generalisation of Sutherland’s original argument to the models based on arbitrary root systems.
Some remarks are in order.
1. The weights $`\mu `$ appearing in the lower order terms $`\{\varphi _\lambda ^{}\}`$’s are those weights contained in the Lie algebra representation belonging to the highest weight $`\lambda `$.
2. As a simple corollary we find that for a minimal weight $`\lambda `$,
$$\psi _\lambda (q)=\varphi _\lambda (q)=\underset{\mu W_\lambda }{}e^{2ia\mu q}$$
is an eigenfunction of $`\stackrel{~}{}`$. A minimal representation consists of a single Weyl orbit and all of its weights $`\mu `$ satisfy $`\rho ^{}\mu =0,\pm 1`$, $`\rho \mathrm{\Delta }`$.
3. If $`\lambda =\alpha _h`$, the highest root of a simply-laced root system, $`W_\lambda `$ is the set of roots itself. Then the lower order terms are constants only. We find that
$$\psi _{\alpha _h}(q)=2\underset{\rho \mathrm{\Delta }_+}{}\left(\mathrm{cos}(2a\rho q)+g\rho ^2/\alpha _h(\alpha _h+2\varrho )\right)$$
is an eigenstate of $`\stackrel{~}{}`$.
4. If $`\lambda =\alpha _{Sh}`$, the highest short root of a non-simply laced root system, $`W_\lambda `$ is the set of short roots itself. The lower order terms are constants, too. Similarly as above, we find that
$$\psi _{\alpha _{Sh}}(q)=2\underset{\rho \mathrm{\Delta }_{L}^{}{}_{+}{}^{}}{}\mathrm{cos}(2a\rho q)+2\underset{\rho \mathrm{\Delta }_{S}^{}{}_{+}{}^{}}{}\left(\mathrm{cos}(2a\rho q)+g_S\rho _S^2/\alpha _{Sh}(\alpha _{Sh}+2\varrho )\right)$$
is an eigenstate of $`\stackrel{~}{}`$. Here $`\mathrm{\Delta }_{L(S)}`$ is the set of long (short) roots.
5. If $`\lambda W_\lambda `$ then there is another set of functions containing the weight $`\lambda `$ which belongs to the same eigenvalue.
6. The Coxeter invariant trigonometric polynomials specified by the fundamental weights $`\{\lambda _j\}`$
$$\varphi _{\lambda _j}(q)=\underset{\mu W_{\lambda _j}}{}e^{2ia\mu q},\lambda _j:\text{fundamental weight},j=1,\mathrm{},r$$
(3.49)
are expected to play the role of the fundamental variables .
7. Let us consider the well-known case $`\mathrm{\Delta }=A_1`$, the simplest root system of rank one. By rewriting the Hamiltonian $`\stackrel{~}{}`$ in terms of the Coxeter invariant variable $`z=\mathrm{cos}(a\rho q)`$, we obtain
$$\widehat{}=\frac{1}{2}\frac{d^2}{dq^2}ag\rho \mathrm{cot}(a\rho q)\frac{d}{dq}=\frac{1}{2}a^2|\rho |^2\left\{(1z^2)\frac{d^2}{dz^2}(1+2g)z\frac{d}{dz}\right\}.$$
(3.50)
The Gegenbauer polynomials , a special case of Jacobi polynomials $`P_n^{(\alpha ,\beta )}`$ provide eigenfunctions:
$$P_n^{(g\frac{1}{2},g\frac{1}{2})}\left(\mathrm{cos}(a\rho q)\right),=a^2|\rho |^2(n+g)^2/2,n𝐙_+.$$
(3.51)
The Jacobi polynomial $`P_{\mathrm{}}^{(\alpha ,\beta )}(z)`$ satisfies differential equation
$$\left\{(1z^2)\frac{d^2}{dz^2}+\beta \alpha (2+\alpha +\beta )z\frac{d}{dz}+\mathrm{}(\mathrm{}+\alpha +\beta +1)\right\}P_{\mathrm{}}^{(\alpha ,\beta )}(z)=0.$$
(3.52)
Here we follow the notation of . They form orthogonal polynomials with weight $`e^{2W}=|\mathrm{sin}(a\rho q)|^{2g}`$ in the interval $`q[0,\pi /a\rho ]`$, (2.19). The Gegenbauer polynomials have a definite parity, $`(1)^n`$, reflecting the periodicity. The fundamental period is $`\pi /a\rho `$, the length of the interval itself. The (odd) even degree ones corresponding to (half-odd) integer spin representations are (anti-) periodic.
8. Triangularity of type II models follows from the same algebraic reasoning.
## 4 Quantum Lax pair and quantum conserved quantities
Historically, Lax pairs for Calogero-Moser models were presented in terms of Lie algebra representations , in particular, the vector representation of the $`A_r`$ models. However, the invariance of Calogero-Moser models is that of Coxeter group but not that of the associated Lie algebra, which does not exist for the non-crystallographic root systems. Thus the universal and Coxeter covariant Lax pairs are given in terms of the representations of the Coxeter group.
### 4.1 General case
Here we recapitulate the essence of the quantum Lax pair operators for the Calogero-Moser models with degenerate potentials and without spectral parameter. The case with spectral parameter will be discussed briefly in subsection 8.3 in connection with the proof of involution of the quantum conserved quantities for $`A_r`$ model . The quantum Lax pair in this subsection applies to all the degenerate potential cases except for the case of the rational potential with the harmonic force, which will be treated separately in subsection 4.2. For details and a full exposition, see . The Lax operators without spectral parameter are
$`L(p,q)`$ $`=`$ $`p\widehat{H}+X(q),X(q)=i{\displaystyle \underset{\rho \mathrm{\Delta }_+}{}}g_{|\rho |}(\rho \widehat{H})x(\rho q)\widehat{s}_\rho ,`$ (4.1)
$`M(q)`$ $`=`$ $`{\displaystyle \frac{i}{2}}{\displaystyle \underset{\rho \mathrm{\Delta }_+}{}}g_{|\rho |}|\rho |^2y(\rho q)\widehat{s}_\rho {\displaystyle \frac{i}{2}}{\displaystyle \underset{\rho \mathrm{\Delta }_+}{}}g_{|\rho |}|\rho |^2y(\rho q)\times I,`$ (4.2)
in which $`I`$ is the identity operator and $`\{\widehat{s}_\alpha ,\alpha \mathrm{\Delta }\}`$ are the reflection operators of the root system. In contrast with $`\{\stackrel{ˇ}{s}_\alpha \}`$ operators (2.29) which act in function space, $`\{\widehat{s}_\alpha \}`$ act on a set of $`𝐑^r`$ vectors $`=\{\mu ^{(k)}𝐑^r,k=1,\mathrm{},d\}`$, permuting them under the action of the reflection group. The vectors in $``$ form a basis for the representation space $`𝐕`$ of dimension $`d`$. The operator $`M`$ satisfies the relation
$$\underset{\mu }{}M_{\mu \nu }=\underset{\nu }{}M_{\mu \nu }=0,$$
(4.3)
which is essential for deriving quantum conserved quantities. The matrix elements of the operators $`\{\widehat{s}_\alpha ,\alpha \mathrm{\Delta }\}`$ and $`\{\widehat{H}_j,j=1,\mathrm{},r\}`$ are defined as follows:
$$(\widehat{s}_\rho )_{\mu \nu }=\delta _{\mu ,s_\rho (\nu )}=\delta _{\nu ,s_\rho (\mu )},(\widehat{H}_j)_{\mu \nu }=\mu _j\delta _{\mu \nu },\rho \mathrm{\Delta },\mu ,\nu .$$
(4.4)
The form of the function $`x`$ depends on the chosen potential, and the function $`y`$ are defined by (2.12), (2.16). Note that these relations are only valid for the degenerate potentials (2.3).
The underlying idea of the Lax operator $`L`$, (4.1), is quite simple. As seen from (4.8), $`L`$ is a “square root” of the Hamiltonian. Thus one part of $`L`$ contains $`p`$ which is not associated with roots and another part contains $`x(\rho q)`$, a “square root” of the potential $`V(\rho q)`$, which being associated with a root $`\rho `$ is therefore accompanied by the reflection operator $`\widehat{s}_\rho `$. Another explanation is the factorised Hamiltonian $``$ (2.5). We obtain, roughly speaking, $`L\sqrt{}p+i\frac{W}{q}\widehat{s}`$ and the property of reflection $`\widehat{s}^2=1`$ explains the sign change in the first term in (2.5).
It is straightforward to show that the quantum Lax equation
$$\frac{d}{dt}L=i[,L]=[L,M],$$
(4.5)
is equivalent to the quantum equations of motion derived from the Hamiltonian (2.7). From this it follows:
$$\frac{d}{dt}(L^n)_{\mu \nu }=i[,(L^n)_{\mu \nu }]=[L^n,M]_{\mu \nu }=\underset{\lambda }{}\left((L^n)_{\mu \lambda }M_{\lambda \nu }M_{\mu \lambda }(L^n)_{\lambda \nu }\right),n=1,\mathrm{}.$$
Thanks to the property of the $`M`$ operator (4.3):
$$\underset{\mu }{}M_{\mu \nu }=\underset{\nu }{}M_{\mu \nu }=0,$$
we obtain quantum conserved quantities as the total sum $`(\text{Ts})`$ of all the matrix elements of $`L^n`$ <sup>2</sup><sup>2</sup>2This type of conserved quantities is known for $`A_r`$ models .:
$$Q_n=\text{Ts}(L^n)\underset{\mu ,\nu }{}(L^n)_{\mu \nu },[,Q_n]=0,n=1,\mathrm{}.$$
(4.6)
Independent conserved quantities appear at such power $`n`$ that
$$n=1+exponent$$
(4.7)
of each root system. These are the degrees at which independent Coxeter invariant polynomials exist. There are $`r`$ exponents for each root system $`\mathrm{\Delta }`$ of rank $`r`$. Thus we have $`r`$ independent conserved quantities in Calogero-Moser models. We list in Table II these powers for each root system. In particular, the power 2 is universal to all the root systems and the quantum Hamiltonian (2.7) is given by
$$=\frac{1}{2C_{}}\text{Ts}(L^2)+const,$$
(4.8)
where the constant $`C_{}`$ is the quadratic Casimir invariant, which depends on the representation. It is defined by
$$\text{Tr}(\widehat{H}_j\widehat{H}_k)\underset{\mu }{}(\widehat{H}_j\widehat{H}_k)_{\mu \mu }=\underset{\mu }{}\mu _j\mu _k=C_{}\delta _{jk}.$$
(4.9)
Some remarks are in order.
1. The Lax pair is Coxeter covariant:
$$L(s_\rho (p),s_\rho (q))_{\mu \nu }=L(p,q)_{\mu ^{}\nu ^{}},M\left(s_\rho (q)\right)_{\mu \nu }=M(q)_{\mu ^{}\nu ^{}},\mu ^{}s_\rho (\mu ),\nu ^{}s_\rho (\nu ),$$
(4.10)
which ensures the Coxeter invariance of the conserved quantities.
2. Lax pairs can be written down in various representations and the quantum conserved quantities $`Q_n`$ do depend on the representations, in general. If necessary, we denote by $`Q_n^{}`$ the explicit representation dependence.
3. The availability of plural representations of the Lax pair and the conserved quantities is essential for the completeness of the set of conserved quantities as polynomials of the momentum operators. For example, let us consider the case of $`D_r`$ with even $`r`$, which has two independent conserved quantities at power $`r`$, see Table II. At least two different representations of the Lax pair are necessary in order to represent them in the form of (4.6). Those based on the vector, and the (anti)-spinor representations give two independent conserved quantities. For $`D_4`$ case, we obtain
$$Q_4^v=2\underset{j=1}{\overset{4}{}}p_j^4,Q_4^sQ_4^a=24\underset{j=1}{\overset{4}{}}p_j,$$
(4.11)
in which $`v`$, $`s`$ and $`a`$ stand for the vector, spinor and anti-spinor representations and we have set $`g=0`$ for simplicity. If two conserved quantities are independent for zero coupling constants, surely they are so at non-vanishing couplings. Here we have used an explicit parametrisation of the $`D_r`$ root system:
$$D_r\text{root system :}\mathrm{\Delta }=\{\pm e_j\pm e_k,j,k=1,\mathrm{},r|e_j𝐑^r,e_je_k=\delta _{jk}\}.$$
(4.12)
4. If a representation $``$ contains a vector $`\mu `$ and its negative $`\mu `$ at the same time, then we have Ts$`(L^{odd})=0`$. In such a case the corresponding Lie algebra representations are called real. In order to construct the odd power conserved quantities appearing in $`A_r`$ for all $`r`$, $`D_r`$ for odd $`r`$, $`E_6`$ and $`I(m)`$ for odd $`m`$, we need a Lax pair in non-real representations. For $`A_r`$ all the fundamental representations corresponding to the fundamental weights $`\lambda _j`$, $`j=1,\mathrm{},r`$ except for the middle one $`\lambda _{(r+1)/2}`$ for odd $`r`$ are non-real. For $`D_r`$ with odd $`r`$, the spinor and anti-spinor representations are non-real. For $`E_6`$ the 27 and $`\overline{\mathrm{𝟐𝟕}}`$ are non-real. $`I_2(m)`$ is the symmetry group of a regular $`m`$-sided polygon. The set of $`m`$ vectors with ‘half” angles of the roots (see (2.10)) to be denoted by $`V_m`$ (B.3), provides a non-real representation when $`m`$ is odd.
5. In Appendix B we list for each root system how the full set of independent conserved quantities are obtained by choosing proper representations of the Lax pair.
### 4.2 Rational potential with harmonic force
The quantum Lax pair for the type V models needs a separate formulation. The explicit form of the Hamiltonian is
$$=\frac{1}{2}p^2+\frac{1}{2}\omega ^2q^2+\frac{1}{2}\underset{\rho \mathrm{\Delta }_+}{}g_{|\rho |}(g_{|\rho |}1)\frac{|\rho |^2}{(\rho q)^2}_0.$$
(4.13)
The canonical equations of motion are equivalent to the following Lax equations for $`L^\pm `$:
$$\frac{d}{dt}L^\pm =i[,L^\pm ]=[L^\pm ,M]\pm i\omega L^\pm ,$$
(4.14)
in which (see section 4 of ) $`M`$ is the same as before (4.2), and $`L^\pm `$ and $`Q`$ are defined by
$$L^\pm =L\pm i\omega Q,Q=q\widehat{H},$$
(4.15)
with $`L`$, $`\widehat{H}`$ as earlier (4.1), (4.4). If we define hermitian operators $`_1`$ and $`_2`$ by
$$_1=L^+L^{},_2=L^{}L^+,$$
(4.16)
they satisfy Lax-like equations
$$\dot{}_k=[_k,M],k=1,2.$$
(4.17)
From these we can construct conserved quantities
$$\text{Ts}(_j^n),j=1,2,n=1,2,\mathrm{},$$
(4.18)
as before. Such quantum conserved quantities have been previously reported for models based on $`A_r`$ root systems . It should be remarked that Ts$`(_2^n)`$ is no longer the same as Ts$`(_1^n)`$ due to quantum corrections. It is elementary to check that the first conserved quantities give the Hamiltonian (4.13)
$$\text{Ts}(_1)=\text{Ts}(_2)+const.$$
(4.19)
This then completes the presentation of the quantum Lax pairs and quantum conserved quantities for all of the quantum Calogero-Moser models with non-elliptic potentials.
## 5 Algebraic construction of excited states I
In this section we show that all the excited states of the type V Calogero-Moser models can be constructed algebraically. Later in section 7 we show the same results in terms of the $`\mathrm{}`$ operators to be introduced in section 6. The main result is surprisingly simple and can be stated universally:
Corresponding to each partition of an integer $`N`$ which specify the energy level (3.12) into the sum of the degrees of Coxeter invariant polynomials (3.13), we have an eigenstate of the Hamiltonian $`\widehat{}`$ with eigenvalue $`\omega N+_0`$:
$$\underset{j=1}{\overset{r}{}}(B_{f_j}^+)^{n_j}e^W,N=\underset{j=1}{\overset{r}{}}n_jf_j,n_j𝐙_+,$$
(5.1)
in which the integers $`\{f_j\}`$, $`j=1,\mathrm{},r`$ are listed in Table II. They exhaust all the excited states. In other words the above states give the complete basis of the Fock space. The creation operators $`B_{f_j}^+`$ and the corresponding annihilation operators <sup>3</sup><sup>3</sup>3We adopt the notation by Olshanetsky and Perelomov . $`B_{f_j}^{}`$ are defined in terms of the Lax operators $`L^\pm `$ (4.15) as follows:
$$B_{f_j}^\pm =\text{Ts}(L^\pm )^{f_j},j=1,\mathrm{},r.$$
(5.2)
They are hermitian conjugate to each other
$$(B_{f_j}^\pm )^{}=B_{f_j}^{}$$
(5.3)
with respect to the standard hermitian inner product of the states defined in $`PW`$:
$$(\psi ,\phi )=_{PW}\psi ^{}(q)\phi (q)𝑑q.$$
(5.4)
We will show later in section 6, (6.15) that the creation (annihilation) operators commute among themselves:
$$[B_k^+,B_l^+]=[B_k^{},B_l^{}]=0,k,l\{f_j|j=1,\mathrm{},r\},$$
(5.5)
so that the state (5.1) does not depend on the order of the creation.
The proof is very simple. By using (4.14) we obtain
$$\frac{d}{dt}(L^\pm )^n=i[,(L^\pm )^n]=[(L^\pm )^n,M]\pm in\omega (L^\pm )^n,$$
(5.6)
from which
$$[,B_n^\pm ]=\pm n\omega B_n^\pm ,$$
(5.7)
follows after taking the total sum. This simply says that $`B_n^\pm `$ creates (annihilates) a state having energy $`n\omega `$. In other words we have
$$\widehat{}\underset{j=1}{\overset{r}{}}(B_{f_j}^+)^{n_j}e^W=\left(_0+\omega \underset{k=1}{\overset{r}{}}n_kf_k\right)\underset{j=1}{\overset{r}{}}(B_{f_j}^+)^{n_j}e^W.$$
Moreover, it is trivial to show that
$$\underset{\nu }{}(L^{})_{\mu \nu }e^W=\left(p\mu i\omega q\mu +i\underset{\rho \mathrm{\Delta }_+}{}\frac{\rho \mu }{\rho q}\right)e^W=\mu \left(p+i\frac{W}{q}\right)e^W=0,$$
(5.8)
which implies that the ground state is annihilated by all the annihilation operators
$$B_{f_j}^{}e^W=0,j=1,\mathrm{},r.$$
(5.9)
Some remarks are in order.
1. In most cases the energy levels are highly degenerate. The above basis is neither orthogonal nor normalised.
2. The independence of the creation-annihilation operators can also be shown in a similar way to that of the conserved quantities. As with the conserved quantities, plural representations are necessary to define the full set of creation-annihilation operators in some models. This aspect will be discussed in later sections in connection with the $`\mathrm{}`$ operators.
3. Reflecting the universality of the first exponent, $`f_1=2`$, the creation and annihilation operators of the least quanta, $`2\omega `$, exist in all the models. They form an $`sl(2,𝐑)`$ algebra together with the Hamiltonian $`\widehat{}`$:
$$[\widehat{},b_2^\pm ]=\pm 2\omega b_2^\pm ,[b_2^+,b_2^{}]=\omega ^1\widehat{},$$
(5.10)
in which $`b_2^\pm `$ are normalised forms of $`B_2^\pm `$:
$$b_2^\pm =\underset{\mu ,\nu }{}(L^\pm )_{\mu \nu }^2/(4\omega C_{}).$$
(5.11)
The $`sl(2,𝐑)`$ algebra was discussed by many authors (see, for example, and others) in connection with the models based on classical root systems. We will show later in subsection 7.2 that the states created by $`B_2^+`$ ($`b_2^+`$) only can be expressed by the Laguerre polynomial:
$$(b_2^+)^ne^W=n!L_n^{(\stackrel{~}{}_01)}(\omega q^2)e^W,\stackrel{~}{}_0_0/\omega .$$
(5.12)
It is trivial to verify that $`L_n^{(\stackrel{~}{}_01)}(\omega q^2)`$ is an eigenfunction of $`\stackrel{~}{}`$ (3.5)
$$\stackrel{~}{}L_n^{(\stackrel{~}{}_01)}(\omega q^2)=2n\omega L_n^{(\stackrel{~}{}_01)}(\omega q^2).$$
(5.13)
The normalisation of the state
$$(b_2^+)^ne^W^2=n!𝒩_0/\mathrm{\Gamma }(n+\stackrel{~}{}_0),𝒩_0e^W^2\mathrm{\Gamma }(\stackrel{~}{}_0),$$
(5.14)
is also dictated by the $`sl(2,𝐑)`$ relations. The Laguerre polynomial wavefunctions appear as ‘radial’ wavefunctions in all the cases . This will be shown explicitly for for the rank two models given in subsection 7.3.
4. As is emphasised by Perelomov and Gambardella the $`sl(2,𝐑)`$ algebra and the corresponding Laguerre wavefunctions are more universal than Calogero-Moser models. They arise when the potentials are homogeneous functions in $`q`$ of degree $`2`$ with the confining harmonic force.
5. The operators $`\{Q_n\}`$ and $`\{B_n^\pm \}`$ do not form a Lie algebra. They satisfy interesting non-linear relations, for example,
$$[[B_n^+,b_2^{}],b_2^+]=nB_n^+,[[B_n^{},b_2^+],b_2^{}]=nB_n^{}.$$
(5.15)
This tells, for example, that although $`B_n^+`$ and $`b_2^+`$ create different units of quanta $`n`$ and $`2`$, they are not independent
$$[B_n^+,b_2^{}]0[B_n^{},b_2^+].$$
Clarification of the algebraic structure (5.15) for each root system is wanted.
## 6 $`\mathrm{}`$ operators
In this section we will show the equivalence of the quantum conserved quantities obtained in the Lax operator formalism of section 4 and those derived in the ‘commuting differential operators’ formalism initiated by Dunkl and followed by many authors. Again the equivalence is universal, applicable to the models based on any root systems. We propose to call the operators in the latter approach simply ‘$`\mathrm{}`$ operators’, since they are essentially the same as the $`L`$ operator in the Lax pair formalism and that they are not mutually commuting, as we will show presently, when the interaction potentials are trigonometric (hyperbolic), (6.14). Although these two formalisms are formally equivalent, the $`\mathrm{}`$ operator formalism has many advantages over the Lax pair one. Roughly speaking, the ‘vector-like’ objects $`\mathrm{}_\mu `$’s are easier to handle than the matrix $`L_{\mu \nu }`$.
Let us fix a representation $``$ of the Coxeter group $`G_\mathrm{\Delta }`$ and define for each element $`\mu `$ the following differential-reflection operator
$$\mathrm{}_\mu =\mathrm{}\mu =p\mu +i\underset{\rho \mathrm{\Delta }_+}{}g_{|\rho |}(\rho \mu )x(\rho q)\stackrel{ˇ}{s}_\rho ,\mu .$$
(6.1)
It is linear in $`\mu `$ and Coxeter covariant
$$\mathrm{}_{\mu +\nu }=\mathrm{}_\mu +\mathrm{}_\nu ,\stackrel{ˇ}{s}_\rho \mathrm{}_\mu \stackrel{ˇ}{s}_\rho =\mathrm{}_{s_\rho (\mu )},\rho \mathrm{\Delta }.$$
(6.2)
They are hermitian operators, $`\mathrm{}_\mu ^{}=\mathrm{}_\mu `$, with respect to the standard inner product for the states (5.4).
It is straightforward to show that the quantum conserved quantities $`Q_n`$ derived in the previous section (4.6) can be expressed as polynomials in the $`\mathrm{}`$ operators as follows:
$$Q_n\psi =\underset{\mu ,\nu }{}(L^n)_{\mu \nu }\psi =(\underset{\mu }{}\mathrm{}_\mu ^n)\psi ,$$
(6.3)
in which $`\psi `$ is an arbitrary Coxeter invariant state, $`\stackrel{ˇ}{s}_\rho \psi =\psi `$. This also illustrates the Coxeter invariance of $`Q_n`$ clearly, since $`\stackrel{ˇ}{s}_\rho (_\mu \mathrm{}_\mu ^n)\stackrel{ˇ}{s}_\rho =_\mu \mathrm{}_{s_\rho (\mu )}^n=_\mu \mathrm{}_\mu ^n`$. For $`n=1`$ it is trivial, since
$`{\displaystyle \underset{\nu }{}}(L)_{\mu \nu }\psi `$ $`=`$ $`\left(p\mu +i{\displaystyle \underset{\rho \mathrm{\Delta }_+}{}}g_{|\rho |}(\rho \mu )x(\rho q){\displaystyle \underset{\nu }{}}(\widehat{s}_\rho )_{\mu \nu }\right)\psi `$ (6.4)
$`=`$ $`\left(p\mu +i{\displaystyle \underset{\rho \mathrm{\Delta }_+}{}}g_{|\rho |}(\rho \mu )x(\rho q)\stackrel{ˇ}{s}_\rho \right)\psi =\mathrm{}_\mu \psi ,`$
in which $`_\nu (\widehat{s}_\rho )_{\mu \nu }=1`$ and $`\stackrel{ˇ}{s}_\rho \psi =\psi `$ are used. Let us assume that
$$\underset{\nu }{}(L^n)_{\mu \nu }\psi =\mathrm{}_\mu ^n\psi ,$$
(6.5)
is correct, then we obtain
$`{\displaystyle \underset{\nu }{}}(L^{n+1})_{\mu \nu }\psi ={\displaystyle \underset{\lambda ,\nu }{}}L_{\mu \lambda }(L^n)_{\lambda \nu }\psi ={\displaystyle \underset{\lambda }{}}L_{\mu \lambda }\mathrm{}_\lambda ^n\psi ,`$
$`=`$ $`{\displaystyle \underset{\lambda }{}}\left(p\mu \delta _{\mu \lambda }+i{\displaystyle \underset{\rho \mathrm{\Delta }_+}{}}g_{|\rho |}(\rho \mu )x(\rho q)(\widehat{s}_\rho )_{\mu \lambda }\right)\mathrm{}_\lambda ^n\psi .`$
In the second summation only such $`\lambda `$ as $`\lambda =s_\rho (\mu )`$ contributes and we find
$$\mathrm{}_{s_\rho (\mu )}^n\psi =(\stackrel{ˇ}{s}_\rho \mathrm{}_\mu ^n\stackrel{ˇ}{s}_\rho )\psi =\stackrel{ˇ}{s}_\rho \mathrm{}_\mu ^n\psi .$$
Thus we arrive at
$$\underset{\nu }{}(L^{n+1})_{\mu \nu }\psi =\mathrm{}_\mu ^{n+1}\psi ,$$
(6.6)
and the equivalence of the two expressions of the conserved quantity (6.3) is proved.
Commutation relations among $`\mathrm{}`$ operators can be evaluated in a similar manner as those appearing in the Lax pair , that is, by decomposing the roots into two-dimensional sub-root systems. We obtain
$$[\mathrm{}_\mu ,\mathrm{}_\nu ]=a^2\underset{\rho ,\sigma \mathrm{\Delta }_+}{}g_{|\rho |}g_{|\sigma |}(\rho \mu )(\sigma \nu )[\stackrel{ˇ}{s}_\rho ,\stackrel{ˇ}{s}_\sigma ]\times \{\begin{array}{cc}\hfill 0& \text{rational},\hfill \\ \hfill 1& \text{hyperbolic},\hfill \\ \hfill 1& \text{trigonometric}.\hfill \end{array}$$
(6.7)
One important use of the $`\mathrm{}`$ operators is the proof of involution of quantum conserved quantities. For type I models Heckman gave a universal proof based on the commutation relation (6.7):
$$[Q_n,Q_m]\psi =\underset{\mu ,\nu }{}[\mathrm{}_\mu ^n,\mathrm{}_\nu ^m]\psi =0,\text{rational model}.$$
(6.8)
This was the motivation for the introduction of the commuting differential-reflection operators by Dunkl, . In fact, Dunkl’s and Heckman’s operators were the similarity transformation of $`\mathrm{}_\mu `$ by the ground state wavefunction $`e^W`$:
$$\stackrel{~}{\mathrm{}}_\mu =e^W\mathrm{}_\mu e^W=p\mu +i\underset{\rho \mathrm{\Delta }_+}{}g_{|\rho |}\frac{(\rho \mu )}{(\rho q)}(\stackrel{ˇ}{s}_\rho 1).$$
(6.9)
As for type V models, we define $`\mathrm{}^\pm `$ corresponding to $`L^\pm `$ (4.15):
$$\mathrm{}_\mu ^\pm =\mathrm{}^\pm \mu =p\mu \pm i\omega (q\mu )+i\underset{\rho \mathrm{\Delta }_+}{}g_{|\rho |}\frac{(\rho \mu )}{(\rho q)}\stackrel{ˇ}{s}_\rho ,\mu .$$
(6.10)
They are linear in $`\mu `$, Coxeter covariant and hermitian conjugate of each other with respect to the standard inner product (5.4):
$$\stackrel{ˇ}{s}_\rho \mathrm{}_\mu ^\pm \stackrel{ˇ}{s}_\rho =\mathrm{}_{s_\rho (\mu )}^\pm ,(\mathrm{}_\mu ^\pm )^{}=\mathrm{}_\mu ^{}.$$
(6.11)
The conserved quantities are expressed as polynomials in $`\mathrm{}^\pm `$ operators:
$`\text{Ts}(_1^n)\psi `$ $`=`$ $`{\displaystyle \underset{\mu ,\nu }{}}(L^+L^{})_{\mu \nu }^n\psi ={\displaystyle \underset{\mu }{}}(\mathrm{}_\mu ^+\mathrm{}_\mu ^{})^n\psi ,`$ (6.12)
$`\text{Ts}(_2^n)\psi `$ $`=`$ $`{\displaystyle \underset{\mu ,\nu }{}}(L^{}L^+)_{\mu \nu }^n\psi ={\displaystyle \underset{\mu }{}}(\mathrm{}_\mu ^{}\mathrm{}_\mu ^+)^n\psi .`$
Likewise the creation and annihilation operators $`B_n^\pm `$ (5.2) are expressed as
$$B_n^\pm \psi =\text{Ts}(L^\pm )^n\psi =\underset{\mu ,\nu }{}(L^\pm )_{\mu \nu }\psi =\underset{\mu }{}(\mathrm{}_\mu ^\pm )^n\psi .$$
(6.13)
The commutation relations among $`\mathrm{}^\pm `$ operators are easy to evaluate, since $`\mathrm{}`$ operators commute in the rational potential models (6.7):
$$[\mathrm{}_\mu ^+,\mathrm{}_\nu ^+]=[\mathrm{}_\mu ^{},\mathrm{}_\nu ^{}]=0,[\mathrm{}_\mu ^{},\mathrm{}_\nu ^+]=2\omega \left(\mu \nu +\underset{\rho \mathrm{\Delta }_+}{}g_{|\rho |}(\rho \mu )(\rho ^{}\nu )\stackrel{ˇ}{s}_\rho \right).$$
(6.14)
From these it follows that the creation (annihilation) operators $`B_n^\pm `$ do commute among themselves:
$$[B_n^+,B_m^+]\psi =[B_n^{},B_m^{}]\psi =0.$$
(6.15)
It is also clear that $`\mathrm{}_\mu ^\pm /\sqrt{2\omega }`$ are the ‘deformation’ of the creation (annihilation) operators of the ordinary multicomponent harmonic oscillators. In fact we have
$$\mathrm{}_\mu ^+e^W=2i\omega (\mu q)e^W\mathrm{and}\mathrm{}_\mu ^{}e^W=0.$$
(6.16)
In the next section we present an alternative scheme of algebraic construction of excited states of type V models by pursuing the analogy that $`\mathrm{}^\pm `$ are the creation and annihilation operators of the unit quantum. This method was applied to the $`A_r`$ models by Brink et. al and others .
## 7 Algebraic construction of excited states II
### 7.1 Operator solution of the triangular Hamiltonian
In subsection 3.1, we have shown that an eigenfunction of $``$ with eigenvalue $`N\omega `$ is given by
$$\left(P_N(q)+\stackrel{~}{P}_{N2}(q)\right)e^W,$$
(7.1)
in which $`P_N(q)`$ is a Coxeter invariant polynomial in $`q`$ of homogeneous degree $`N`$ and $`\stackrel{~}{P}_{N2}(q)`$ is a Coxeter invariant polynomial in $`q`$ of degree $`N2`$ and lower. The non-leading polynomial $`\stackrel{~}{P}_{N2}(q)`$ is completely determined by the leading one $`P_N(q)`$ due to the triangularity. This solution can be written in an operator form as follows.
Suppose $`P_N(q)`$ is expressed as
$$P_N(q)=\underset{\{\mu \}}{}c_{\{\mu \}}(q\mu _1)\mathrm{}(q\mu _N),\mu _j,c_{\{\mu \}}:const.$$
(7.2)
We obtain a Coxeter invariant polynomial in the creation operators $`\mathrm{}^+`$ by replacing $`q\mu `$ by $`\mathrm{}_\mu ^+/(2i\omega )`$:
$$P_N(q)\frac{1}{(2i\omega )^N}P_N(\mathrm{}^+).$$
This creates the above eigenfunction of $``$ from the ground state:
$$\frac{1}{(2i\omega )^N}P_N(\mathrm{}^+)e^W=\left(P_N(q)+\stackrel{~}{P}_{N2}(q)\right)e^W.$$
(7.3)
The proof is again elementary. By using the commutation relations among $`\mathrm{}^\pm `$ operators it is straightforward to derive the explicit expression of the Hamiltonian in terms of $`\mathrm{}^\pm `$:
$$=\frac{1}{2𝒞_{}}\underset{\mu }{}\mathrm{}_\mu ^+\mathrm{}_\mu ^{}+\underset{\rho \mathrm{\Delta }_+}{}g_{|\rho |}\left(\omega +\frac{1}{2}\frac{|\rho |^2}{(\rho q)^2}\right)(\stackrel{ˇ}{s}_\rho 1),$$
(7.4)
in which the second term vanishes upon acting on a Coxeter invariant state. Next we obtain
$$\frac{1}{2𝒞_{}}\underset{\mu }{}[\mathrm{}_\mu ^+\mathrm{}_\mu ^{},\mathrm{}_\nu ^\pm ]=[\mathrm{}_\nu ^\pm ,S]\pm \omega \mathrm{}_\nu ^\pm ,S\underset{\rho \mathrm{\Delta }_+}{}g_{|\rho |}\stackrel{ˇ}{s}_\rho ,$$
(7.5)
which is an $`\mathrm{}`$ operator version of (4.14). Since a commutator is a derivation, we obtain
$$\frac{1}{2𝒞_{}}\underset{\mu }{}[\mathrm{}_\mu ^+\mathrm{}_\mu ^{},P_N(\mathrm{}^+)]=[P_N(\mathrm{}^+),S]+N\omega P_N(\mathrm{}^+),$$
(7.6)
in which the first term in r.h.s. vanishes due to the Coxeter invariance of $`P_N`$. Thus we arrive at the desired commutation relation
$$[,P_N(\mathrm{}^+)]=N\omega P_N(\mathrm{}^+)+\underset{\rho \mathrm{\Delta }_+}{}g_{|\rho |}[\frac{1}{2}\frac{|\rho |^2}{(\rho q)^2},P_N(\mathrm{}^+)](\stackrel{ˇ}{s}_\rho 1),$$
(7.7)
and the eigenvalue equation
$$P_N(\mathrm{}^+)e^W=N\omega P_N(\mathrm{}^+)e^W.$$
(7.8)
Since the action of the creation operators on the ground state is
$$\mathrm{}_{\mu _1}^+\mathrm{}\mathrm{}_{\mu _N}^+e^W=[(2i\omega )^N(q\mu _1)\mathrm{}(q\mu _N)+\text{lower powers of }q]e^W,$$
(7.9)
our assertion (7.3) is proved. It should be stressed that in this formalism the Coxeter invariance of the polynomial $`P`$ is important but not how it is obtained.
Like the above Hamiltonian (7.4), the $`\mathrm{}`$ operator formulas of higher conserved quantities (6.12) contain extra terms:
$$\text{Ts}(_1^n)=\underset{\mu ,\nu }{}(L^+L^{})_{\mu \nu }^n=\underset{\mu }{}(\mathrm{}_\mu ^+\mathrm{}_\mu ^{})^n+VT.$$
(7.10)
Here $`VT`$ stands for vanishing terms when they act on a Coxeter invariant state. The same is true for most formulas derived in section 6.
### 7.2 States Created by $`B_2^+`$
Here we derive the explicit forms of the subseries of eigenstates obtained by multiple applications of the least quanta creation operator $`B_2^+`$ (5.2), or its normalised form $`b_2^+`$ (5.11). It is convenient to work with the similarity transformed operator
$$\stackrel{~}{b}_2^+=e^Wb_2^+e^W=\frac{1}{4\omega C_{}}\underset{\mu }{}(\stackrel{~}{\mathrm{}}_\mu ^+)^2+VT,$$
(7.11)
in which
$$\stackrel{~}{\mathrm{}}_\mu ^+=p\mu +2i\omega (q\mu )+i\underset{\rho \mathrm{\Delta }_+}{}\frac{\rho \mu }{\rho q}(\stackrel{ˇ}{s}_\rho 1).$$
(7.12)
Let $`f(u)`$ be an arbitrary function of $`u\omega q^2`$, then it is Coxeter invariant. We find
$$\stackrel{~}{\mathrm{}}_\mu ^+f(u)=2i\omega (q\mu )(1\frac{d}{du})f(u),u\omega q^2,$$
(7.13)
and
$$\stackrel{~}{b}_2^+f(u)=\left[u\left(1\frac{d}{du}\right)^2\stackrel{~}{}_0\left(1\frac{d}{du}\right)\right]f(u).$$
(7.14)
Since $`\stackrel{~}{b}_2^+1=\stackrel{~}{}_0u=L_1^{(\stackrel{~}{}_01)}(u)`$, we assume
$$(\stackrel{~}{b}_2^+)^n1=n!L_n^{(\stackrel{~}{}_01)}(u).$$
(7.15)
By using the Laguerre differential equation (3.21) and the recurrence formulas of the Laguerre polynomial $`L_n^{(\alpha )}(u)`$,
$`u{\displaystyle \frac{d}{du}}L_n^{(\alpha )}(u)=nL_n^{(\alpha )}(u)(n+\alpha )L_{n1}^{(\alpha )}(u),`$ (7.16)
$`nL_n^{(\alpha )}(u)+(u2n\alpha +1)L_{n1}^{(\alpha )}(u)+(n+\alpha 1)L_{n2}^{(\alpha )}(u)=0,`$ (7.17)
we can show
$$\left[u(1\frac{d}{du})^2\stackrel{~}{}_0(1\frac{d}{du})\right]L_n^{(\stackrel{~}{}_01)}(u)=(n+1)L_{n+1}^{(\stackrel{~}{}_01)}(u).$$
(7.18)
Thus the induction is proved and we arrive at (5.12). The orthogonality of the states
$$((B_2^+)^ne^W,(B_2^+)^me^W)=0,nm$$
(7.19)
can be easily understood as the $`du`$ part of the measure
$$e^{2W}d^rq=e^uu^{\stackrel{~}{}_01}dud\mathrm{\Omega },d\mathrm{\Omega }:\text{angular part},$$
is the proper weight function for the Laguerre polynomial $`L_n^{(\stackrel{~}{}_01)}(u)`$.
### 7.3 Explicit solutions of the rank two models
For rank two models, the Liouville integrability, or the involution of conserved quantities is automatically satisfied since the second conserved quantity is already obtained. For rank two type V models, the complete set of orthogonal wavefunctions can be written down explicitly in terms of separation of variables by using the Coxeter invariant polynomials. These are based on the dihedral root systems $`I_2(m)`$, with $`A_2I_2(3)`$ , $`B_2I_2(4)`$ and $`G_2I_2(6)`$ . The Coxeter invariant polynomials exist at degree 2, i.e. $`q^2`$ and $`m`$ which is
$$\underset{j=1}{\overset{m}{}}(v_jq),$$
(7.20)
where $`\{v_j\}`$ is a set of vectors given in (B.3). If we introduce the two-dimensional polar coordinates system <sup>4</sup><sup>4</sup>4We believe no confusion arises here, between the radial coordinate variable $`r`$ and the rank of the root system $`r`$, which in this case is 2 of $`I_2(m)`$. for $`q`$
$$q=r(\mathrm{sin}\theta ,\mathrm{cos}\theta ),$$
(7.21)
then the principal Weyl chamber is
$$PW:0<r^2<\mathrm{},0<\theta <\pi /m.$$
(7.22)
The two Coxeter invariant variables read:
$$q^2=r^2,\underset{j=1}{\overset{m}{}}(v_jq)=2(\frac{r}{2})^m\mathrm{cos}m\theta ,$$
(7.23)
and the latter variable varies the full range, $`1<\mathrm{cos}m\theta <1`$ in the $`PW`$. Thus solving the eigenvalue equation for $`\stackrel{~}{}`$ (3.6) by separation of variables in the polar coordinate system is compatible with Coxeter invariance. We adopt as two independent variables
$$u\omega r^2,z\mathrm{cos}m\theta .$$
(7.24)
The solutions consist of a Gegenbauer (Jacobi) polynomial in $`\mathrm{cos}m\theta `$ times a Laguerre polynomial in $`\omega r^2`$. The former we have encountered in the $`A_1`$ Sutherland problem, subsection 3.2 and the latter in the $`A_1`$ Calogero problem subsections 3.1 and 7.2.
Let us demonstrate this for odd $`m`$ with a single coupling constant and for even $`m`$ with two independent coupling constants, in parallel. In terms of the Coxeter invariant variables (7.24) the $`I_2(m)`$ Hamiltonians take surprisingly simple forms:
$`\stackrel{~}{}`$ $`=`$ $`\omega r{\displaystyle \frac{}{r}}{\displaystyle \frac{1}{2r}}{\displaystyle \frac{}{r}}\left(r{\displaystyle \frac{}{r}}\right){\displaystyle \frac{1}{2r^2}}\left[{\displaystyle \frac{^2}{\theta ^2}}+m\left\{{\displaystyle \genfrac{}{}{0pt}{}{2g\mathrm{cot}m\theta }{g_0\mathrm{tan}\frac{m\theta }{2}+g_e\mathrm{cot}\frac{m\theta }{2}}}\right\}{\displaystyle \frac{}{\theta }}\right]`$
$`=`$ $`2\omega \left[u{\displaystyle \frac{^2}{u^2}}+(1u){\displaystyle \frac{}{u}}\right]{\displaystyle \frac{\omega m^2}{2u}}\left[(1z^2){\displaystyle \frac{^2}{z^2}}+\left\{{\displaystyle \genfrac{}{}{0pt}{}{0}{g_og_e}}\right\}\left\{{\displaystyle \genfrac{}{}{0pt}{}{1+2g}{1+g_e+g_o}}\right\}z{\displaystyle \frac{}{z}}\right].`$
The $`z`$ part admits polynomial solutions
$`\left[(1z^2){\displaystyle \frac{d^2}{dz^2}}+\left\{{\displaystyle \genfrac{}{}{0pt}{}{0}{g_og_e}}\right\}\left\{{\displaystyle \genfrac{}{}{0pt}{}{1+2g}{1+g_e+g_o}}\right\}z{\displaystyle \frac{d}{dz}}\right]`$ $`P_{\mathrm{}}^{\left\{\genfrac{}{}{0pt}{}{(g\frac{1}{2},g\frac{1}{2})}{(g_o\frac{1}{2},g_e\frac{1}{2})}\right\}}(z)`$
$`=\mathrm{}\left(\mathrm{}+\left\{{\displaystyle \genfrac{}{}{0pt}{}{2g}{g_e+g_o}}\right\}\right)`$ $`P_{\mathrm{}}^{\left\{\genfrac{}{}{0pt}{}{(g\frac{1}{2},g\frac{1}{2})}{(g_o\frac{1}{2},g_e\frac{1}{2})}\right\}}(z),`$ (7.26)
in which $`\mathrm{}`$ is the degree of the polynomial. After substituting them, the radial part of the Hamiltonian $`\stackrel{~}{}_r`$ reads
$$\stackrel{~}{}_r=2\omega \left[u\frac{d^2}{du^2}+(1u)\frac{d}{du}\frac{m^2}{4u}\mathrm{}\left(\mathrm{}+\left\{\genfrac{}{}{0pt}{}{2g}{g_e+g_o}\right\}\right)\right].$$
(7.27)
By similarity transformation in terms of $`u^{m\mathrm{}/2}r^m\mathrm{}`$, which is the radial part of the highest term of the polynomial $`P_{\mathrm{}}^{(\alpha ,\beta )}(r^m\mathrm{cos}m\theta )`$, it reads
$$u^{m\mathrm{}/2}\stackrel{~}{}_ru^{m\mathrm{}/2}=2\omega \left[u\frac{d^2}{du^2}+\left(m\left(\left\{\genfrac{}{}{0pt}{}{\mathrm{}+g}{\mathrm{}+\frac{1}{2}(g_e+g_o)}\right\}\right)+1u\right)\frac{d}{du}\frac{m\mathrm{}}{2}\right].$$
(7.28)
This is the main part of the differential equation for the Laguerre polynomial (3.21):
$$\left[u\frac{d^2}{du^2}+\left(m\mathrm{}+\stackrel{~}{}_0u\right)\frac{d}{du}n\right]L_n^{(m\mathrm{}+\stackrel{~}{}_01)}(u)=0,$$
in which the indices $`m(\mathrm{}+g)`$ and $`m(\mathrm{}+(g_o+g_e)/2)`$ can be written in a unified way as $`m\mathrm{}+\stackrel{~}{}_01`$. Thus the eigenstates of the Hamiltonian are obtained:
$`\stackrel{~}{}_ru^{m\mathrm{}/2}L_n^{(m\mathrm{}+\stackrel{~}{}_01)}(u)`$ $`=`$ $`\omega (2n+m\mathrm{})u^{m\mathrm{}/2}L_n^{(m\mathrm{}+\stackrel{~}{}_01)}(u),`$ (7.29)
$`\widehat{}u^{m\mathrm{}/2}L_n^{(m\mathrm{}+\stackrel{~}{}_01)}(u)`$ $`P_{\mathrm{}}^{\left\{\genfrac{}{}{0pt}{}{(g\frac{1}{2},g\frac{1}{2})}{(g_o\frac{1}{2},g_e\frac{1}{2})}\right\}}(z)`$
$`=`$ $`(\omega (2n+m\mathrm{})+_o)u^{m\mathrm{}/2}L_n^{(m\mathrm{}+\stackrel{~}{}_01)}(u)P_{\mathrm{}}^{\left\{\genfrac{}{}{0pt}{}{(g\frac{1}{2},g\frac{1}{2})}{(g_o\frac{1}{2},g_e\frac{1}{2})}\right\}}(z).`$
It is instructive to note that the Hamiltonians $`\widehat{}`$ look also simple:
$$\widehat{}=\omega r\frac{}{r}\frac{1}{2r}\frac{}{r}\left(r\frac{}{r}\right)\frac{1}{2r^2}\frac{^2}{\theta ^2}\frac{m^2}{2r^2}\left\{\begin{array}{c}\frac{g(g1)}{\mathrm{sin}^2m\theta }\\ \frac{g_o(g_o1)}{4\mathrm{cos}^2\frac{m\theta }{2}}+\frac{g_e(g_e1)}{4\mathrm{sin}^2\frac{m\theta }{2}}\end{array}\right\}.$$
(7.31)
Olshanetsky and Perelomov obtained the above solutions starting from these formulas.
## 8 Involution of conserved quantities
### 8.1 Universal proof of involution of quantum conserved quantities for type I, II and III models
Here we present a proof of involution of quantum conserved quantities $`\{Q_n\}`$ derived from the universal Lax pair in subsection 4.1 for type I, II and III models. The proof is applicable to all models based on any root systems. Though a universal proof of involution for type I models is given by Heckman as recapitulated in section 6, we believe the universal proof applicable to type II and III models as well is new. It depends on a theorem by Olshanetsky and Perelomov . Our own contribution is that we have provided a universal Lax pair and conserved quantities satisfying all the requirements of the theorem.
Liouville’s theorem states the complete integrability as the existence of an involutive set of conserved quantities as many as the degrees of freedom. We have already given conserved quantities $`\{Q_n\}`$ (4.6) independent and as many as the degrees of freedom (see Appendix B). They have the following properties:
1. Coxeter invariance
$$Q_n(s_\rho (p),s_\rho (q))=Q_n(p,q),\rho \mathrm{\Delta }.$$
(8.1)
2. $`Q_n(p,q)`$ is a homogeneous polynomial of degree $`n`$ in variables $`(p_1,\mathrm{},p_r,x(\rho q))`$.
3. Scaling property for those of type I models:
$${}_{}{}^{I}Q_{n}^{}(\kappa ^1p,\kappa q)=\kappa ^n{}_{}{}^{I}Q_{n}^{}(p,q)$$
(8.2)
as a consequence of the above point.
4. For type II and III models, the asymptotic behaviour near the origin:
$$Q_n(p,q)={}_{}{}^{I}Q_{n}^{}(p,q)(1+𝒪(|q|)),\text{for}|q|0.$$
(8.3)
We need to show the vanishing of
$$J_{lm}[Q_l,Q_m],$$
(8.4)
which is a polynomial in $`\{p\}`$ of degree $`s`$
$$s<l+m.$$
(8.5)
Let us decompose $`J_{lm}`$ into the leading part and the rest:
$$J_{lm}=J_{lm}^0+J_{lm}^{rest},J_{lm}^0=c^{j_1,\mathrm{},j_s}(q)p_{j_1}\mathrm{}p_{j_s}$$
(8.6)
and $`J_{lm}^{rest}`$ is a polynomial in $`\{p\}`$ of degree less than $`s`$. From Jacobi identity and conservation $`[,Q_{l(m)}]=0`$, we obtain
$$[,J_{lm}]=0.$$
(8.7)
Considering the explicit form of the Hamiltonian (2.7) ($`\omega =0)`$, the leading (i.e. of degree $`s+1`$ in $`\{p\}`$) part of $`[,J_{lm}]`$ comes only from the free part
$$[p^2,J_{lm}^0]$$
and it vanishes if the following conditions are satisfied:
$$\underset{\sigma }{}\frac{}{q_t}c^{k_1,\mathrm{},k_s}(q)=0,$$
(8.8)
where the sum is taken over all permutations of indices $`\sigma (t,k_1,\mathrm{},k_s)=(j_1,\mathrm{},j_{s+1})`$. In it is proved (Lemma 2.5, p. 407) that the system (8.8) has only polynomial solutions. Then Olshanetsky and Perelomov argue that for type I models the scaling property tells that $`c^{k_1,\mathrm{},k_s}(\kappa q)=\kappa ^{slm}c^{k_1,\mathrm{},k_s}(q)`$. Since $`s<l+m`$ (8.5), it follows that the only polynomial solution satisfying the condition is the null polynomial. Thus we obtain $`c^{j_1,\mathrm{},j_s}(q)=0`$ $`J_{lm}^0=0`$ and $`J_{lm}=0`$. The same results follow for type II and III models by considering the asymptotic behaviour for $`|q|0`$. Thus the involution of all the conserved quantities $`\{Q_n\}`$ is proved. This result also implies the involution of classical conserved quantities by taking the classical limit ($`\mathrm{}0`$).
### 8.2 Rational models with the harmonic confining force
In this subsection we show the involution of quantum conserved quantities for the type V models based on the root systems of classical Lie algebras. The method is a straightforward generalisation of the one developed by Polychronakos on the $`A_r`$ model. This is made possible by the availability of the universal Lax pair formalism , in particular the root type and minimal type Lax pairs. We apply it to the models based on $`B_r`$ and $`D_r`$ root systems. For the rational potential, the $`B_r`$, $`C_r`$ and $`BC_r`$ models are equivalent. Let us choose $`=\{\pm e_j𝐑^r|e_je_k=\delta _{jk}\}`$ as the representation space of the Coxeter group consisting of orthogonal vectors and their negatives. They are the set of short roots of $`B_r`$ and the set of vector weights of $`D_r`$ in the parametrisation of roots given in (B.1) and (4.12), respectively.
The conserved quantities in the $`\mathrm{}`$ operator form are given by
$$Q_n={}_{}{}^{\mathrm{}}Q_{n}^{}+VT,{}_{}{}^{\mathrm{}}Q_{n}^{}\underset{j=1}{\overset{r}{}}(\mathrm{}_j^+\mathrm{}_j^{})^n,n=1,\mathrm{},r,$$
(8.9)
in which we abbreviate $`\mathrm{}_{e_j}^\pm `$ as $`\mathrm{}_j^\pm `$. In this case the commutation relation among $`\mathrm{}`$ operators (6.14) are greatly simplified thanks to the orthogonality of $`\{e_j\}`$’s and the explicit forms of the roots:
$$[\mathrm{}_j^{},\mathrm{}_k^+]=2\omega \stackrel{~}{g}(\stackrel{ˇ}{s}_{jk}\overline{\stackrel{ˇ}{s}}_{jk}),$$
(8.10)
in which
$$\stackrel{~}{g}=\{\begin{array}{cc}g_L& B_r\text{model},\hfill \\ g& D_r\text{model},\hfill \end{array}\stackrel{ˇ}{s}_{jk}\stackrel{ˇ}{s}_{e_je_k},\overline{\stackrel{ˇ}{s}}_{jk}\stackrel{ˇ}{s}_{e_j+e_k}.$$
(8.11)
By repeating them we obtain
$`[\mathrm{}_j^+\mathrm{}_j^{},\mathrm{}_k^+\mathrm{}_k^{}]`$ $`=`$ $`2\omega \stackrel{~}{g}[\mathrm{}_j^+\mathrm{}_j^{},m_{jk}],m_{jk}\stackrel{ˇ}{s}_{jk}+\overline{\stackrel{ˇ}{s}}_{jk}=m_{kj},jk,`$ (8.12)
$`[(\mathrm{}_j^+\mathrm{}_j^{})^n,\mathrm{}_k^+\mathrm{}_k^{}]`$ $`=`$ $`2\omega \stackrel{~}{g}[(\mathrm{}_j^+\mathrm{}_j^{})^n,m_{jk}]=+2\omega \stackrel{~}{g}[(\mathrm{}_k^+\mathrm{}_k^{})^n,m_{jk}],`$ (8.13)
$`[\mathrm{}_j^+\mathrm{}_j^{},(\mathrm{}_k^+\mathrm{}_k^{})^m]`$ $`=`$ $`+2\omega \stackrel{~}{g}[(\mathrm{}_k^+\mathrm{}_k^{})^m,m_{jk}]=2\omega \stackrel{~}{g}[(\mathrm{}_j^+\mathrm{}_j^{})^m,m_{jk}].`$ (8.14)
Here and later the identity $`[(\mathrm{}_k^+\mathrm{}_k^{})^t,m_{jk}]=[(\mathrm{}_j^+\mathrm{}_j^{})^t,m_{jk}]`$ is used repeatedly. Then (8.13) leads to
$$[(\mathrm{}_j^+\mathrm{}_j^{})^n,(\mathrm{}_k^+\mathrm{}_k^{})^m]=+2\omega \stackrel{~}{g}\underset{t=0}{\overset{m1}{}}\left((\mathrm{}_k^+\mathrm{}_k^{})^{t+n}m_{jk}(\mathrm{}_k^+\mathrm{}_k^{})^{mt1}(\mathrm{}_k^+\mathrm{}_k^{})^tm_{jk}(\mathrm{}_k^+\mathrm{}_k^{})^{m+nt1}\right)$$
(8.15)
and (8.14) to
$$[(\mathrm{}_j^+\mathrm{}_j^{})^n,(\mathrm{}_k^+\mathrm{}_k^{})^m]=2\omega \stackrel{~}{g}\underset{t=0}{\overset{n1}{}}\left((\mathrm{}_j^+\mathrm{}_j^{})^{t+m}m_{jk}(\mathrm{}_j^+\mathrm{}_j^{})^{nt1}(\mathrm{}_j^+\mathrm{}_j^{})^tm_{jk}(\mathrm{}_j^+\mathrm{}_j^{})^{m+nt1}\right).$$
(8.16)
Summing over $`j`$ and $`k`$ and adding (8.15) and (8.16) together with the interchange of the dummy indices $`jk`$ in the latter produces
$$2[{}_{}{}^{\mathrm{}}Q_{n}^{},{}_{}{}^{\mathrm{}}Q_{m}^{}]=+2\omega \stackrel{~}{g}\underset{t=0}{\overset{n+m1}{}}\underset{j,k}{}\left((\mathrm{}_k^+\mathrm{}_k^{})^tm_{jk}(\mathrm{}_k^+\mathrm{}_k^{})^{n+mt1}(\mathrm{}_k^+\mathrm{}_k^{})^tm_{jk}(\mathrm{}_k^+\mathrm{}_k^{})^{n+mt1}\right)=0.$$
(8.17)
Thus we obtain
$$[Q_n,Q_m]=0,n,m=1,\mathrm{},r$$
(8.18)
on the Fock space of Coxeter invariant states.
### 8.3 Lax pair with spectral parameter
In the theory of classical Calogero-Moser models, Lax pair with spectral parameter ($`\xi `$) plays an important role, in particular, in elliptic potential models for derivation of spectral curves, etc. In quantum theory, however, the meaning and use of the spectral parameter are yet to be established, partly because of the underdeveloped stage of the quantum models with elliptic potentials. Here we point out a small use of the quantum Lax pair with spectral parameter in the quantum model with trigonometric potential. Namely, it accounts for the useful trick by Polychronakos for the proof of the involution of conserved quantities in trigonometric $`A_r`$ model. (Now we have a universal proof of involution for type I, II and III models, see the previous subsection.)
From the theory of the generalised Lax pair and its quantum version for degenerate potential models , we find that the $`L`$ operator can contain one additional complex parameter $`\xi `$:
$$L^\xi =p\widehat{H}+X^\xi ,X^\xi =i\underset{\rho \mathrm{\Delta }_+}{}g_{|\rho |}(\rho \widehat{H})\left(x(\rho q)x(\rho ^{}\widehat{H}\xi )\right)\widehat{s}_\rho ,$$
(8.19)
in which the function $`x`$ are given in Table I for the degenerate potentials. With the same $`M`$ operator as before, the quantum equations of motion can be written in a matrix form:
$$\frac{d}{dt}L^\xi =i[,L^\xi ]=[L^\xi ,M].$$
(8.20)
In other words, the $`\xi `$ dependent part decouples. This allows us to define a one parameter family of conserved quantities
$${}_{}{}^{\xi }Q_{n}^{}=\text{Ts}(L^\xi )^n,n=1,2,\mathrm{},$$
(8.21)
which turns out to be a $`\xi `$ dependent sum of $`Q_n`$ and the lower order conserved quantities $`Q_m`$, $`m<n`$. A special limit $`\xi i\mathrm{}`$ in the trigonometric models provides a convenient combination which allows easy proof of involution in the $`A_r`$ model. (In the rational model, the limit reduces to the Lax pair without spectral parameter. In the hyperbolic models this limit is ill-defined.) In the rest of this subsection we consider only the trigonometric potential models. Let us denote the limiting $`L^\xi `$ operator by $`L^{\mathrm{}}`$ which reads
$$L^{\mathrm{}}=L+a\underset{\rho \mathrm{\Delta }_+}{}g_{|\rho |}|\rho \widehat{H}|\widehat{s}_\rho $$
(8.22)
and the corresponding $`\mathrm{}`$ operators are given by
$$\mathrm{}_\mu ^{\mathrm{}}=\mathrm{}_\mu +t_\mu ,t_\mu =a\underset{\rho \mathrm{\Delta }_+}{}g_{|\rho |}|\rho \mu |\stackrel{ˇ}{s}_\rho .$$
(8.23)
For the $`A_r`$ model in the vector representation $`=\{\mu _j𝐑^r,j=1,\mathrm{},r+1\}`$ with the standard normalisation of roots $`\rho ^2=2`$, the above expression simplifies to
$$\mathrm{}_j^{\mathrm{}}=\mathrm{}_j+t_j,t_j=ag\underset{kj}{}\stackrel{ˇ}{s}_{jk},$$
(8.24)
in which as before we abbreviate $`\mathrm{}_{\mu _j}`$ as $`\mathrm{}_j`$ and $`\stackrel{ˇ}{s}_{jk}\stackrel{ˇ}{s}_{\mu _j\mu _k}`$. These are the operators introduced in , $`\mathrm{}_j\pi _j`$, $`\mathrm{}_j^{\mathrm{}}\stackrel{~}{\pi }_j`$. The Lax operator with the spectral parameter gives an ‘explanation’ for the rather ad hoc introduction of $`\stackrel{~}{\pi }_j`$. It is straightforward to show
$$[\mathrm{}_j,\mathrm{}_k]=[t_j,t_k],$$
(8.25)
which leads to
$$[\mathrm{}_j^{\mathrm{}},\mathrm{}_k^{\mathrm{}}]=[\mathrm{}_j^{\mathrm{}},t_k]+[t_j,\mathrm{}_k^{\mathrm{}}]=2ag[\mathrm{}_j^{\mathrm{}},\stackrel{ˇ}{s}_{jk}]=2ag[\mathrm{}_k^{\mathrm{}},\stackrel{ˇ}{s}_{jk}].$$
(8.26)
This has the same structure as (8.12) in the previous subsection. By repeating the same argument we arrive at
$$[{}_{}{}^{\mathrm{}}Q_{n}^{},{}_{}{}^{\mathrm{}}Q_{m}^{}]=0,n,m=1,\mathrm{},r,{}_{}{}^{\mathrm{}}Q_{n}^{}=\underset{j=1}{\overset{r}{}}(\mathrm{}_j^{\mathrm{}})^n,$$
(8.27)
which then imply the involution of the conserved quantities obtained from the original $`L`$ operator
$$[Q_n,Q_m]=0,n,m=1,\mathrm{},r.$$
(8.28)
## 9 Summary and comments
We have discussed various issues related to quantum integrability of Calogero-Moser models based on all root systems. These are construction of quantum conserved quantities and a unified proof of their involution, the relationship between the Lax pair and the differential-reflection (Dunkl) operators formalisms, construction of excited states by creation operators, etc. They are mainly generalisations of the results known for the models based on $`A_r`$ root systems, which are shown to apply to the models based on any root systems. There are some interesting works discussing the integrability issues of the models based on other classical root systems and the exceptional ones including the non-crystallographic models , -.
Here we list some comments on interesting issues which are not treated in this paper. The structure and properties of the eigenfunctions of the trigonometric potential models, which are generalisations of the Jack polynomials -, will be discussed in future publications. A comprehensive treatment of Liouville integrability of rational models with harmonic force is wanted. Our starting point, the factorised Hamiltonian (2.5) for degenerate potential models, is closely related with supersymmetry and shape invariance . Further investigation in this direction is a future problem. It is a great challenge to formulate various aspects of quantum Calogero-Moser models with elliptic potentials; Lax pair, the differential-reflection operators , conserved quantities, supersymmetry and excited states wavefunctions.
## Acknowledgements
We thank K. Takasaki, S. Odake and P. Ghosh for useful discussion. This work is partially supported by the Grant-in-aid from the Ministry of Education, Science and Culture, priority area (#707) “Supersymmetry and unified theory of elementary particles”. S. P. K. and A. J. P. are supported by the Japan Society for the Promotion of Science.
## Appendix A: Root Systems
In this Appendix we recapitulate the rudimentary facts of the root systems and reflections to be used in the main text. The set of roots $`\mathrm{\Delta }`$ is invariant under reflections in the hyperplane perpendicular to each vector in $`\mathrm{\Delta }`$. In other words,
$$s_\alpha (\beta )\mathrm{\Delta },\alpha ,\beta \mathrm{\Delta },$$
(A.1)
where
$$s_\alpha (\beta )=\beta (\alpha ^{}\beta )\alpha ,\alpha ^{}2\alpha /|\alpha |^2.$$
(A.2)
The set of reflections $`\{s_\alpha ,\alpha \mathrm{\Delta }\}`$ generates a group $`G_\mathrm{\Delta }`$, known as a Coxeter group, or finite reflection group. The orbit of $`\beta \mathrm{\Delta }`$ is the set of root vectors resulting from the action of the Coxeter group on it. The set of positive roots $`\mathrm{\Delta }_+`$ may be defined in terms of a vector $`U𝐑^r`$, with $`\alpha U0,\alpha \mathrm{\Delta }`$, as those roots $`\alpha \mathrm{\Delta }`$ such that $`\alpha U>0`$. Given $`\mathrm{\Delta }_+`$, there is a unique set of $`r`$ simple roots $`\mathrm{\Pi }=\{\alpha _j,j=1,\mathrm{},r\}`$ defined such that they span the root space and the coefficients $`\{a_j\}`$ in $`\beta =_{j=1}^ra_j\alpha _j`$ for $`\beta \mathrm{\Delta }_+`$ are all non-negative. The highest root $`\alpha _h`$, for which $`_{j=1}^ra_j`$ is maximal, is then also determined uniquely. The subset of reflections $`\{s_\alpha ,\alpha \mathrm{\Pi }\}`$ in fact generates the Coxeter group $`G_\mathrm{\Delta }`$. The products of $`s_\alpha `$, with $`\alpha \mathrm{\Pi }`$, are subject solely to the relations
$$(s_\alpha s_\beta )^{m(\alpha ,\beta )}=1,\alpha ,\beta \mathrm{\Pi }.$$
(A.3)
The interpretation is that $`s_\alpha s_\beta `$ is a rotation in some plane by $`2\pi /m(\alpha ,\beta )`$. The set of positive integers $`m(\alpha ,\beta )`$ (with $`m(\alpha ,\alpha )=1,\alpha \mathrm{\Pi }`$) uniquely specify the Coxeter group. The weight lattice $`\mathrm{\Lambda }(\mathrm{\Delta })`$ is defined as the $`𝐙`$-span of the fundamental weights $`\{\lambda _j\}`$, $`j=1,\mathrm{},r`$, defined by
$$\alpha _j^{}\lambda _k=\delta _{jk},\alpha _j\mathrm{\Pi }.$$
(A.4)
The root systems for finite reflection groups may be divided into two types: crystallographic and non-crystallographic. Crystallographic root systems satisfy the additional condition
$$\alpha ^{}\beta 𝐙,\alpha ,\beta \mathrm{\Delta },$$
(A.5)
which implies that the $`𝐙`$-span of $`\mathrm{\Pi }`$ is a lattice in $`𝐑^r`$ and contains all roots in $`\mathrm{\Delta }`$. We call this the root lattice, which is denoted by $`L(\mathrm{\Delta })`$. These root systems are associated with simple Lie algebras: {$`A_r,r1\}`$, $`\{B_r,r2\}`$, $`\{C_r,r2\}`$, $`\{D_r,r4\}`$, $`E_6`$, $`E_7`$, $`E_8`$, $`F_4`$ and $`G_2`$. The Coxeter groups for these root systems are called Weyl groups. The remaining non-crystallographic root systems are $`H_3`$, $`H_4`$, whose Coxeter groups are the symmetry groups of the icosahedron and four-dimensional 600-cell, respectively, and the dihedral group of order $`2m`$, $`\{I_2(m),m4\}`$.
## Appendix B: Conserved quantities
Here we list for each root system how the full set of independent conserved quantities are obtained by choosing proper representations of the Lax pair. We choose those of the lowest dimensionality for the convenience of practical calculation. Of course there are many other choices of representations giving equally good sets of conserved quantities. The independence of the conserved quantities can be easily verified by considering the free limit: $`g_{|\rho |}0`$.
1. $`A_r`$: For all powers, the vector representation ($`r+1`$ dimensions) is enough.
2. $`B_r`$: For all powers, the representation consisting of short roots $`\{\pm e_j:j=1,\mathrm{},r\}`$, ($`2r`$ dimensions) is enough. Here we adopt the following explicit parametrisation of the $`B_r`$ root system:
$$B_r\text{root system :}\mathrm{\Delta }=\{\pm e_j\pm e_k,\pm e_j,j,k=1,\mathrm{},r|e_j𝐑^r,e_je_k=\delta _{jk}\}.$$
(B.1)
3. $`C_r`$: For all powers, the representation consisting of long roots $`\{\pm 2e_j:j=1,\mathrm{},r\}`$, ($`2r`$ dimensions) is enough. The following parametrisation of the root system is used:
$$C_r\text{root system :}\mathrm{\Delta }=\{\pm e_j\pm e_k,\pm 2e_j,j,k=1,\mathrm{},r|e_j𝐑^r,e_je_k=\delta _{jk}\}.$$
(B.2)
4. $`D_r`$: For all even powers, the vector representation ($`2r`$ dimensions) is enough. For the additional one occurring at power $`r`$, the (anti)-spinor representation ($`2^{r1}`$ dimensions) would be necessary. They are minimal representations.
5. $`E_6`$: For all powers, the 27 (or $`\overline{\mathrm{𝟐𝟕}}`$) dimensional representation of the Lie algebra is enough. They are minimal representations.
6. $`E_7`$: For all powers, the 56 dimensional representation of the Lie algebra is enough. This is a minimal representation.
7. $`E_8`$: For all powers, the 240 dimensional representation consisting of all the roots is enough. This is not the same as the adjoint representation of the Lie algebra.
8. $`F_4`$: For all powers, either of the 24 dimensional representation consisting of all the long roots or the short roots is enough. These are not Lie algebra representations.
9. $`G_2`$: For all powers, either of the 6 dimensional representations consisting of all the long roots or the short roots is enough. These are not Lie algebra representations.
10. $`I_2(m)`$: For both powers 2 and $`m`$, the representation consisting of $`V_m`$ is enough. Here, $`V_m`$ is the set of vectors with ‘half” angles of the roots (see (2.10)) given by
$$V_m=\{v_j=(\mathrm{cos}((2j1)\pi /2m),\mathrm{sin}((2j1)\pi /2m))𝐑^2|j=1,\mathrm{},m\}.$$
(B.3)
11. $`H_3`$: For all powers, the representation consisting of all the 30 roots is enough.
12. $`H_4`$: For all powers, the representation consisting of all the 120 roots is enough.
|
warning/0005/cond-mat0005469.html
|
ar5iv
|
text
|
# Schrödinger cat state of a Bose-Einstein condensate in a double-well potential
## I Introduction
Since the first observations of Bose-Einstein condensation in dilute alkali-metal atomic gases the ultra-cold atomic gases have stimulated significant theoretical and experimental interest . The scientific progress has been rapid and examples of recent experiments include the development of accurate detection methods , the state preparation of topological structures , and the applications in nonlinear atom optics . Due to the macroscopic quantum coherence Bose-Einstein condensates (BECs) could possibly be also used in the future as a test for the foundations of quantum mechanics. One particularly puzzling and controversial issue has been the existence of macroscopic quantum superposition states in many-particle quantum systems. In this paper we propose a method of creating the Schrödinger cat states of different atom occupation numbers in a weakly interacting BEC confined in a double-well potential.
The existence of the superpositions of macroscopically distinguishable states in BECs has been addressed by several authors . The superposition state may arise as the ground state of a coherently coupled BEC in a double-well potential . Under certain conditions it could be reached as a result of a unitary time evolution . Previously, we proposed a method of creating Schrödinger cat states in BECs by means of scattering light from two BECs moving with opposite velocities . The nonunitary evolution due to the detections of scattered photons drives the condensates to macroscopic quantum superposition states. In this paper we show that a continuous quantum measurement process could also drive a trapped coherently coupled BEC in a double-well potential to a Schrödinger cat state. The advantage of the proposed scheme is that the BEC is almost stationary and trapped. Moreover, as a result of the back-action of quantum measurement process the superposition state could be reached rapidly unlike in a slow unitary evolution, which may be very sensitive to decoherence.
The paper is organized as follows: We begin in Sec. II A by introducing the unitary system Hamiltonian. The scattering of light and the measurement geometry is described in Sec. II B. In Sec. II C we study the dynamics of the open quantum system in terms of stochastic trajectories of state vectors. The results of the numerical simulations are presented in Sec. III. Finally, a few concluding remarks are made in Sec. IV.
## II System dynamics
### A Unitary evolution
We consider the evolution of a BEC in a double-well potential in a two-mode approximation. Macroscopic quantum coherence of BECs results in coherent quantum tunneling of atoms between the two modes representing ‘two BECs’. This is analogous to the coherent tunneling of Cooper pairs in a Josephson junction . To obtain the system Hamiltonian in the two-mode approximation for the unitary evolution of the BEC we approximate the total field operator by the two lowest quantum modes $`\psi (𝐫)\psi _b(𝐫)b+\psi _c(𝐫)c`$, where $`\psi _b`$ and $`\psi _c`$ stand for the local mode solutions of the individual wells with small spatial overlap. The corresponding annihilation operators are denoted by $`b`$ and $`c`$. The Hamiltonian in the two-mode approximation reads :
$$\frac{H_S}{\mathrm{}}=\xi b^{}b+\mathrm{\Omega }(b^{}c+c^{}b)+\kappa [(b^{})^2b^2+(c^{})^2c^2].$$
(1)
Here $`\xi `$ is the energy difference between the modes. The tunneling between the two wells is described by $`\mathrm{\Omega }`$, which is proportional to the overlap of the spatial mode function of the opposite wells. The short-ranged two-body interaction strength is obtained from $`\kappa =2\pi a\mathrm{}/m|\psi _b(𝐫)|^4`$, where $`a`$ and $`m`$ denote the scattering length and the atomic mass, respectively. For simplicity, here we have assumed that $`|\psi _b(𝐫)|^4=|\psi _c(𝐫)|^4`$. A necessary condition for the validity of the single-mode approximation in a harmonic trap is that the oscillation energy of the atoms does not dominate over the mode energy spacing of the trap.
According to the Josephson effect, the atom numbers of the BECs determined by the Hamiltonian (1) may oscillate even if the number of atoms in each well is initially equal. Due to the nonlinear self-interaction the number oscillations also exhibit collapses and revivals. These have been studied numerically in Ref. . We may also obtain a simple analytical description by solving the dynamics in the rotating wave approximation in the limit $`\mathrm{\Omega }N\kappa `$ as described in Ref. . Here $`N`$ denotes the total number of atoms. In particular, we may solve the number of atoms $`N_bb^{}b`$ in well $`b`$. We consider a coherent state in the both wells as an initial state. Then we obtain
$$N_b=\frac{N}{2}\left[1+e^{N[\mathrm{cos}(\kappa t)1]}(\sqrt{1\beta ^2}\mathrm{cos}\eta \beta \mathrm{sin}\phi \mathrm{sin}\eta )\right],$$
(2)
with $`\eta N\beta \mathrm{sin}(\kappa t)\mathrm{cos}\phi 2\mathrm{\Omega }t`$. Here all the operators on the right-hand side have been evaluated at $`t=0`$. It is useful to define the real expectation values $`\beta `$ and $`\phi `$ in the following way:
$$\beta e^{i\phi }\frac{2}{N}b^{}c.$$
(3)
For a coherent state with equal atom numbers in the two wells we obtain the visibility $`\beta =1`$. The relative phase between the wells is $`\phi `$. For a number state there is no phase information and $`\beta =0`$. For unequal atom numbers the maximum visibility is $`\beta _{\mathrm{max}}=2(N_bN_c)^{1/2}/N`$. We see that the number of atoms in Eq. (2) may oscillate in the case of initially equal atom numbers $`\beta =1`$. The amplitude of sinusoidal oscillations, representing the macroscopic coherence, collapses. For instance, for $`\phi =\pi /2`$ and $`\beta =1`$ we may obtain the short time decay by considering the time scales $`N\kappa t1\mathrm{\Omega }t`$. Then the decay of the oscillations has the form $`\mathrm{exp}(N\kappa ^2t^2/2)`$. This is the rate of the phase diffusion and it may be interpreted as the width of the relative phase $`\mathrm{\Delta }\phi ^2(t)N\kappa ^2t^2\kappa ^2t^2/\mathrm{\Delta }\phi ^2(0)`$. Perhaps surprisingly the functional dependence of the width in this case is the same as in the case of two uncoupled BECs .
If the phase is unknown we may obtain the ensemble average by integrating over the relative phase $`\phi `$ in Eq. (2):
$`N_b`$ $`=`$ $`{\displaystyle \frac{N}{2}}\{1+e^{N[\mathrm{cos}(\kappa t)1]}\sqrt{1\beta ^2}`$ (5)
$`\times \mathrm{cos}(2\mathrm{\Omega }t)J_0[N\beta \mathrm{sin}(\kappa t)]\},`$
where $`J_0`$ is the 0th order Bessel function. If the both wells have initially equal number of atoms, $`\beta =1`$, the atom numbers do not oscillate. For unequal atom numbers the oscillations collapse at the first zero of the Bessel function $`t2.4/(N\beta \kappa )`$ for $`(N_bN_c)^{1/2}1`$.
### B Quantum measurement process
The time evolution of the system is nonunitary, when we include the effect of quantum measurement process. We consider the nondestructive measurement of the number of atoms in the both wells by means of shining coherent light beams through the atom clouds. The scattered light beams are combined by a 50-50 beam splitter. We display the measurement setup in Fig. 1.
We assume that the incoming light fields are detuned far from the atomic resonance. For instance, if the shape of the gas is flat and the light is shone through a thin dimension, the multiple scattering is negligible and the sample can be considered optically thin. A BEC atom scatters back to the BEC via coherent spontaneous scattering, stimulated by a large number of atoms in the BEC. Coherently scattered photons are emitted into a narrow cone in the forward direction. By spontaneous scattering we mean that the emission is not stimulated by light, although it is stimulated by atoms. The decay into noncondensate center-of-mass states is also stimulated by the Bose-Einstein statistics. However, at very low temperatures this stimulation is much weaker because most of the particles are in the BEC. As a first approximation we ignore the scattering from and to the noncondensate modes. Then the measurement is nondestructive in the sense that BEC atoms in the modes $`b`$ and $`c`$ scatter back to the same modes $`b`$ and $`c`$. Because the overlap of the mode functions of the different wells is assumed to be small, the scattering between the two wells is ignored.
The detection rate of photons on the detectors is the intensity of the scattered light $`I(𝐫)`$ integrated over the scattering directions divided by the energy of a photon $`\mathrm{}ck`$. Here $`k`$ and $`c`$ stand for the wave number and the velocity of light. We obtain the detection rate at the channel $`j`$:
$$\gamma _j=\frac{1}{\mathrm{}ck}𝑑\mathrm{\Omega }_{\widehat{𝐧}}r^2I_j(𝐫)=2\mathrm{\Gamma }C_j^{}C_j.$$
(6)
The photon annihilation operator at the output channel $`j`$ of the beam splitter is denoted by $`C_j`$. For a symmetric measurement geometry we obtain
$$C_1=\frac{1}{\sqrt{2}}(b^{}bc^{}c),C_2=\frac{1}{\sqrt{2}}(b^{}b+c^{}c).$$
(7)
Because the total number of atoms is assumed to be conserved, the operator $`\widehat{N}b^{}b+c^{}c`$ contributes to the measurements only through a constant phase shift. Therefore, we may ignore the effect of the scattering channel 2 on the dynamics.
The scattered intensity may be written in terms of the positive frequency component of the scattered electric field $`𝐄^+(𝐫)`$
$$I(𝐫)=2cϵ_0𝐄^{}(𝐫)𝐄^+(𝐫),$$
(8)
Here $`ϵ_0`$ denotes the permittivity of the vacuum.
We assume that the driving electric fields may be approximated by plane waves $`𝐄_d^+(𝐫)=\widehat{𝐞}e^{i(𝐤𝐫\mathrm{\Omega }t)}/2`$. In the limit of large atom-light detuning $`\mathrm{\Delta }`$ we use the first Born approximation and write the electric fields in the far radiation zone ($`kr1`$). Then the scattered field from the well $`b`$ has the following form :
$`𝐄^+(\widehat{𝐧}r)`$ $`=`$ $`{\displaystyle \frac{k^2e^{ikr}}{4\pi ϵ_0\mathrm{\Delta }r}}\widehat{𝐧}\times (\widehat{𝐧}\times 𝐝)`$ (10)
$`{\displaystyle d^3r^{}e^{i(𝐤k\widehat{𝐧})𝐫^{}}|\psi _b(𝐫^{})|^2b^{}b}.`$
Here we have defined the Rabi frequency $``$ of the atomic dipole matrix element $`𝐝`$ by $`d/(2\mathrm{})`$. We also assumed that $`𝐝\widehat{𝐞}=d`$. In the limit that the characteristic length scale $`\mathrm{}`$ of the BECs is much larger than the inverse of the wave number of the incoming light $`\mathrm{}1/k`$, the momentum of the scattered photon is approximately conserved, and we obtain in Eq. (10):
$$d^3r^{}e^{i(𝐤k\widehat{𝐧})𝐫^{}}|\psi _b(𝐫^{})|^2\delta (k\widehat{𝐧}𝐤)$$
In this simple case the scattering rate $`\mathrm{\Gamma }`$ may be easily evaluated:
$$\mathrm{\Gamma }=\frac{3\gamma ^2}{8\pi \mathrm{\Delta }^2},$$
(11)
Here $`\gamma =d^2k^3/(6\pi \mathrm{}ϵ_0)`$ denotes the optical linewidth of the atom.
### C Stochastic Schrödinger equation
The dissipation of energy from the quantum system of macroscopic light fields and the BEC in a double-well potential is described by the coupling to a zero temperature reservoir of vacuum modes, resulting in a spontaneous emission linewidth for the atoms. The dynamics of the continuous quantum measurement process may be unraveled into stochastic trajectories of state vectors . The procedure consists of the evolution of the system with a non-Hermitian Hamiltonian $`H_{\mathrm{eff}}`$, and randomly decided quantum ‘jumps’. In our case the quantum jumps correspond to the detections of spontaneously emitted photons. The system evolution is thus conditioned on the outcome of a measurement. The non-Hermitian Hamiltonian has the following form:
$$H_{\mathrm{eff}}=H_Si\mathrm{}\mathrm{\Gamma }C_1^{}C_1,$$
(12)
where the unitary system Hamiltonian $`H_S`$ is determined by Eq. (1).
Equation (12) corresponds to the modification of the state of the system associated with a zero detection result for scattered photons. Because the output is being continuously monitored, we gain information about the system even if no photons have been detected.
The Hamiltonian $`H_{\mathrm{eff}}`$ determines the evolution of the state vector $`\psi _{\mathrm{sys}}(t)`$. If the wave function $`\psi _{\mathrm{sys}}(t)`$ is normalized, the probability that a photon from the output channel 1 of the beamsplitter is detected during the time interval $`[t,t+\delta t]`$ is
$$P(t)=2\mathrm{\Gamma }\psi _{\mathrm{sys}}(t)|C_1^{}C_1|\psi _{\mathrm{sys}}(t)\delta t.$$
(13)
We implement the simulation algorithm as follows: At the time $`t_0`$ we generate a quasi-random number $`ϵ`$ which is uniformly distributed between 0 and 1. We assume that the state vector $`\psi _{\mathrm{sys}}(t_0)`$ at the time $`t_0`$ is normalized. Then we evolve the state vector by the non-Hermitian Hamiltonian $`H_{\mathrm{eff}}`$ iteratively for finite time steps $`\mathrm{\Delta }t`$. At each time step $`n`$ we compare $`ϵ`$ to the reduced norm of the wave function, until
$$\psi _{\mathrm{sys}}(t_0+n\mathrm{\Delta }t)|\psi _{\mathrm{sys}}(t_0+n\mathrm{\Delta }t)<ϵ,$$
when the detection of a photon occurs. If the photon has been observed during the time step $`tt+\mathrm{\Delta }t`$ we take the new wave function at $`t+\mathrm{\Delta }t`$ to be
$$|\psi _{\mathrm{sys}}(t+\mathrm{\Delta }t)=\sqrt{2\mathrm{\Gamma }}C_1|\psi _{\mathrm{sys}}(t),$$
(14)
which is then normalized.
## III Numerical results
We simulate the effect of the system Hamiltonian $`H_S`$ and the quantum measurement process of scattered photons by means of the stochastic Schrödinger equation. For simplicity, we set the total number of atoms to be reasonably small $`N=200`$. We start from a slightly asymmetric initial state with the two modes in number states $`N_b=102`$ and $`N_c=98`$. We choose $`\mathrm{\Gamma }/\mathrm{\Omega }=5\times 10^6`$.
After just a few detected photons we observe the emergence of two well-separated amplitude maxima in the occupation number of atoms in one of the two wells. These correspond to a macroscopic number state superposition or a Schrödinger cat state. Because the total number of atoms is assumed to be conserved, the atom numbers in the two wells are entangled and we have a Bell-type of superposition state. In Fig. 2 we display the absolute value of the wave function $`|\psi _b|`$ in mode $`b`$ in number state basis in a single run at two different times. In this case the nonlinearity vanishes $`\kappa =0`$ and $`\xi /\mathrm{\Omega }=0.1`$. We clearly recognize the two distinct peaks in the number distribution. For instance, the peaks in Fig. 2 (b) are centered at $`N_b10`$ and $`N_b190`$ corresponding to maxima at $`N_c190`$ and $`N_c10`$, respectively. In Fig. 3 we show the absolute value of the wave function for a different run with $`N\kappa /\mathrm{\Omega }=0.2`$ and $`\xi /\mathrm{\Omega }=0.001`$.
We also describe the state of the BEC in terms of the quasiprobability $`Q`$ distribution. For the number state distribution of atoms $`|\psi _b=_nc_n|n`$ in mode $`b`$ we obtain :
$$Q(\alpha )=\frac{|\alpha |\psi _b|^2}{\pi }=\frac{e^{|\alpha |^2}}{\pi }\left|\underset{n=0}{\overset{N}{}}\frac{\alpha ^nc_n^{}}{\sqrt{n!}}\right|^2.$$
(15)
The $`Q`$ function represents the phase-space distribution. The amplitude and phase quadratures are denoted by $`X`$ and $`Y`$. In polar coordinates the radius in the $`xy`$ plane is equal to $`N_b^{1/2}`$ and the polar angle is the relative phase between the atoms in the two wells. In Fig. 4 we show the $`Q`$ function distribution of the number state superposition displayed in Fig. 2(a).
It is interesting to emphasize that the measurement of the number of atoms in only one of the wells affects the system dynamics quite differently. In that case the number state distribution remains well localized and approximately approaches a coherent state . Even though we start from a number state with no phase information, the detections of spontaneously scattered photons establish a macroscopic coherence or the off-diagonal long-range order (ODLRO) between the atoms in the two separate wells. This is similar to establishing the coherence between two BECs as a result of the counting of atoms . However, in the present case the continuous measurement process drives the system to a Schrödinger cat state and the ODLRO remains small. We may describe the visibility of the macroscopic coherence between the two wells by the real parameter $`\beta `$ defined in Eq. (2). We show the relative visibility $`\beta _r\beta /\beta _{\mathrm{max}}`$ and the number of atoms in well $`b`$ as a function of the number of measurements for $`\kappa =0`$ in Fig. 5 and for $`N\kappa /\mathrm{\Omega }=0.2`$ in Fig. 6. Due to the emergence of the superposition state the visibility remains below one. The measurement process of the scattered photons significantly complicates the dynamics of the number of atoms predicted by the unitary time evolution of Eq. (1).
## IV Final remarks
We studied the generation of the macroscopic superposition states or the Schrödinger cat states of a BEC in a double-well potential. The Schrödinger cat state was shown to emerge as a result of the continuous quantum measurement process of scattered photons. The particular detection geometry increses the fluctuations of the relative atom number between the two wells. Therefore the superposition states are more stable in the detection process. The proposed setup is an open quantum system and the creation of the Schrödinger cat state in this case is not based on reaching the ground state of a BEC in a double-well potential . The advantage over previously proposed open systems schemes is that the BEC is stably trapped and the superposition state for a small BEC could be created by scattering only a few photons.
In the present discussion we ignored the effect of decoherence . The interaction of the BEC with its environment results in the decoherence of the superposition states. We can identify several sources of decoherence. Decoherence by amplitude damping or by phase damping has been estimated in Ref. . The inelastic two-body and three-body collisions between the condensate atoms and the noncondensate atoms change the number of condensate atoms and introduce amplitude damping. The phase damping corresponds, e.g., to elastic collisions between the condensate and noncondensate atoms in which case the number of BEC atoms is conserved. If the number of atoms in a BEC is not large, the scattering between the condensate and noncondensate atom fractions may not be negligible. This also introduces amplitude decoherence. Additional sources of decoherence may be, e.g., the imperfect detection of the scattered photons and the fluctuations of the magnetic trap. In Ref. it was proposed that the decoherence rate of a BEC could be dramatically reduced by symmetrization of the environment and by changing the geometry of the trapping potential to reduce the size of the thermal cloud. Moreover, the continuous measurement process increases the information about the system and therefore it could also reduce the decoherence rate.
## Acknowledgements
We are indebted to the late Prof. Walls for his support, inspiration, and encouragement during his last years. This work was financially supported by the EC through the TMR Network ERBFMRXCT96-0066.
|
warning/0005/cond-mat0005222.html
|
ar5iv
|
text
|
# Two-Magnon States in Cu(NO3)₂⋅2.5D2O using Inelastic Neutron Scattering.
## Abstract
We report measurements of the two-magnon states in a dimerized antiferromagnetic chain material, copper nitrate (Cu(NO<sub>3</sub>)$`{}_{2}{}^{}2.5`$D<sub>2</sub>O). Using inelastic neutron scattering, we have studied the one- and two-magnon excitation spectra in a large single crystal of this material. We compare this new data with perturbative expansions of the alternating Heisenberg chain and find good agreement with these calculations. The data may also show evidence for the recently proposed $`S=1`$ two-magnon bound state (Phys. Rev. B54, R9624 (1996)).
PACS numbers: 75.10.Jm, 75.40.Gb, 78.70.Nx
Competition between hopping and binding effects in elementary excitations is a general feature of low dimensional hard core systems. Such effects are predicted to be manifest in the structure factor of multiparticle continua, accessible using neutron and Raman scattering, as bound modes and enhancement of continuum scattering . Experimental systems of particular interest in this respect are the $`S=1/2`$ alternating Heisenberg chains (AHC) and the new spin ladder systems which are predicted to have bound modes below the two-magnon excitation continua. To investigate this phenomenon we present an experimental study of the two-magnon states in a near ideal example of an AHC compound, using inelastic neutron scattering.
The $`S=1/2`$ AHC spin Hamiltonian is
$$H=\underset{i}{}J\stackrel{}{S}_{2i1}\stackrel{}{S}_{2i}+\alpha J\stackrel{}{S}_{2i}\stackrel{}{S}_{2i+1}$$
(1)
where $`J>0`$ is the intradimer coupling, and $`\alpha J`$ ($`0\alpha 1`$) the interdimer one, which alternate between chain sites $`i`$ . Computational methods are very effective at calculating perturbative properties of (1), and ground and low-lying excited state wavefunctions are given to $`𝒪(\alpha ^5)`$ in , together with many experimentally important quantities to $`𝒪(\alpha ^9)`$. Higher order expansions for selected quantities have recently been reported . Here we compare these results to Cu(NO<sub>3</sub>)$`{}_{2}{}^{}2.5`$D<sub>2</sub>O, or CN for short.
The magnetic properties of CN are well characterized. It is monoclinic ($`I12/c1`$ ), with low temperature lattice parameters $`a=16.1`$, $`b=4.9`$, $`c=15.8`$ Å and $`\beta =92.9^{}`$, with spin $`S=1/2`$ magnetic Cu<sup>2+</sup> ions. The dominant magnetic exchange integral $`J`$ is between pairs of spins forming dimers. Dimers separated by $`𝐮`$ are coupled together by exchanges $`J_𝐮^{}`$. Only the $`J_{[\frac{1}{2},\pm \frac{1}{2},\frac{1}{2}]}^{}`$ exchange paths are of appreciable strength, giving two sets of alternating Heisenberg chains (AHCs) running in the $`[111]`$ and $`[1\overline{1}1]`$ directions which repeat every $`𝐮_0=[111]/2`$ and $`𝐮_0^{}=[1\overline{1}1]/2`$ respectively.
Bulk magnetic measurements give information on the gap, exchange and ground state energies: Applied magnetic fields induce spin flop ordering in CN above a critical field $`H_{c1}27`$ kOe, with a transition to full alignment at $`H_{c2}43`$ kOe. Because the orbital moment is quenched, and demagnetization effects are negligible, $`H_{c1}`$ directly gives the excitation gap, $`\mathrm{\Delta }=0.378\pm .007`$ meV and $`H_{c2}`$ gives the sum of exchange couplings $`J+_𝐮J_𝐮^{}=0.580\pm .007`$ meV. Normally it is not possible to measure the ground state energy of quantum antiferromagnets but high field magnetization techniques in our case make this achievable.
The ground state energy-per-spin $`e_0`$ can be found from the low temperature isothermal magnetization $`M(H)`$ using $`e_0e_fSH_{c2}+_0^{H_{c2}}M(H)𝑑H`$ where the fully aligned energy-per-spin is $`e_f=S^2/2(J+_𝐮J_𝐮^{})=H_{c2}/8=0.0725\pm .001`$ meV. Using the 270 mK data of Diederix et al. in Figure (3) of (measured using proton resonance) to determine the integral over magnetization gives an experimental ground state energy-per-spin $`e_0=0.174\pm .004`$ meV. This is essentially the $`T=0`$ result, as the gap activation energy corresponds to 4.4 K. To estimate thermodynamic properties we approximate the sum of the exchanges by the single coupling $`\alpha J=_𝐮J_𝐮^{}`$ of equation (1). Using the $`𝒪(\alpha ^9)`$ expansions for $`\mathrm{\Delta }(\alpha )`$ and $`e_0`$ gives $`J=0.455\pm .002`$ meV and $`\alpha =0.277\pm .006`$; in agreement with the results of and , $`J=0.45`$ meV and $`\alpha =0.27`$. Our calculated values of the thermodynamic parameters $`J+_𝐮J_𝐮^{}=0.581`$ meV, $`\mathrm{\Delta }=0.379`$ meV and $`e_0=0.172`$ meV agree within error with the experimental values.
The neutron scattering structure factor $`𝒮(𝐐,\omega )`$ probes the ground $`|G`$ and excited states $`|E`$ of a magnetic system through the matrix element $`|E|S^\alpha (𝐐)|G|^2`$ of the Fourier transformed spin operator $`S^\alpha (𝐐)`$ $`(\alpha =x,y,z)`$. We measured $`𝒮(𝐐,\omega )`$ using inelastic neutron scattering from two deuterated single crystals of CN with a total mass of 14.1 g, using the SPINS cold neutron triple-axis spectrometer at the NIST Center for Neutron Research. The substitution of D for H reduces incoherent scattering of neutrons and does not significantly change the magnetic properties of the material. The sample was mounted with $`(h,0,l)`$ as the scattering plane in a pumped <sup>3</sup>He cryostat at a base temperature of 300 mK.
The spectrometer was set up with $`80^{}`$ before the sample as the only collimation. A vertically focused pyrolytic graphite PG(002) monochromator and a horizontally focused analyzer array composed of eleven independently rotatable PG(002) blades were employed. A cooled Be filter before the sample removed higher-order contamination from the beam. Measurements were made with fixed final energy $`E_f=2.5`$ meV by scanning incident energy at various reduced wavevector transfers along the chain, $`\stackrel{~}{q}=𝐐𝐮_0`$. The wide angular acceptance ($`14^{}`$) of the analyzer dominated the instrumental resolution making it highly elongated perpendicular to the scattered wavevector in the scattering plane. Scan trajectories were chosen to maintain the final wavevector $`𝐤_f`$ along the $`(101)`$ direction so as to integrate over nondispersive directions while maintaining good resolution in $`\stackrel{~}{q}`$.
Because the ground state is a singlet it cannot be probed directly by neutrons, but its composition is reflected through spin matrix elements to triplet excited states. Neutron scattering matrix elements to the $`S=1`$ one magnon states have been calculated to $`𝒪(\alpha ^5)`$ : the leading order scattering process is from the bare dimer component of the ground state, and an $`\alpha /2\mathrm{cos}(\stackrel{~}{q})`$ component in the one magnon structure factor arises from an $`O(\alpha )`$ two-dimer excitation in the ground state. Similarly, transitions between various components of the full ground and excited states can be identified through distinctive modulations of the $`𝐐`$ dependence of $`𝒮(𝐐,\omega )`$.
Figure (1) shows scans in energy in CN: Panel (a) shows a scan at the antiferromagnetic zone-center, $`\stackrel{~}{q}=2\pi `$, taken at $`T=300`$ mK. Strong elastic scattering from incoherent nuclear processes is clearly seen as well as a one magnon peak at 0.4 meV, close to the dimer energy $`J=0.45`$ meV. The non-magnetic background (dashed line) was modelled by a Gaussian (incoherent) component and a power-times-Lorentzian (broad, quasielastic) component. A second magnetic peak appears at about 0.9 meV, double the dimer energy. Panels (b) and (c) show this peak with the nonmagnetic background subtracted at $`\stackrel{~}{q}=2\pi `$ and $`\stackrel{~}{q}=3\pi `$, respectively. This feature is considerably weaker than the one-magnon scattering and narrows at the zone-boundary, panel (c); this behavior is consistent with two-magnon scattering.
Figure (2) shows excitations calculated to $`𝒪(\alpha )`$ using degenerate perturbation theory, including states up to two-magnon, for $`J=0.45`$ meV and $`\alpha =0.27`$. The calculation is similar to that for the two-soliton continuum in the XXZ Ising chain in , except that a dimer excited state basis is used and the one magnon band is included . An interesting feature of the spectrum is the existence of an $`S=1`$ two-magnon bound state for a range of $`\stackrel{~}{q}`$ around $`\stackrel{~}{q}(2n+1)\pi `$ where $`n`$ is an integer. This bound state is due to the attraction between adjacent excited dimers with total spin 1 (and 0). It exists only over a limited range of $`\stackrel{~}{q}`$ around the bandwidth minimum in the two-magnon continuum because of the nature of the hopping matrix elements. Near the two-magnon bandwidth minimum the hopping element is small, and binding occurs. Far from this minimum the gain in hopping energy of excited dimers is dominant and no $`S=1`$ bound mode exists. An interesting consequence of the competition is that an $`S=0`$ bound mode (not visible to neutron scattering, but can be observed using Raman scattering) should exist below the continuum for all $`\stackrel{~}{q}`$, because it has a larger attractive interaction between adjacent dimers.
The strength of production of the $`S=1`$ bound mode in neutron scattering depends on the spin operator matrix element to this state. Perturbation analysis in $`\alpha `$ shows that there is no zeroth order coupling of the bare dimer ground state to the two-magnon bound state and continuum through the neutron scattering matrix element. The leading perturbative contribution to $`𝒮(𝐐,\omega )`$ appears at $`𝒪(\alpha ^2)`$, and is due to a transition from the $`𝒪(\alpha )`$ two-excited-dimer component of the ground state and to the $`𝒪(\alpha )`$ one-excited-dimer component of the bound state. The matrix element of these two basis components has a complicated and characteristic wavevector dependence.
Figure (2) shows the energy corresponding to a weighted average of scattering as a thick grey line. It is notable that at $`\stackrel{~}{q}=3\pi `$ to a good approximation the neutrons couple only to the bound mode, so that nearly all the scattering weight is in it, not the continuum. The calculated neutron scattering intensity from the bound state is 2% of the one-magnon intensity which agrees with the data in Figure (1).
The one- and two-magnon scattering at 300 mK was scanned from $`\stackrel{~}{q}=\pi `$ to $`5\pi `$ in steps of $`\pi /4`$. The background subtracted data are plotted in the upper panel of Figure (3). A calculation of the magnetic scattering based on the perturbation calculation described above, including the dimer envelope function , is shown in the lower panel of Figure (3). The calculation is directly comparable with the data in the upper panel of Figure (3), and at a qualitative level there is good agreement with experiment.
A quantitative comparison between theory and data is shown in Figure (4). The measured positions of one- and two-magnon peaks are plotted in the left panel. Energies, widths and intensities for each peak were extracted by least-squares fitting of Gaussians. Considerable dispersion of the one-magnon modes is evident. Measurements of the one magnon dispersion were previously fitted using an $`𝒪(\alpha )`$ model of chains with an additional weak interdimer coupling $`J_{[1/200]}^{}`$ and $`J_{[001/2]}^{}`$ and .
The dispersion relation predicted by this model gives a good account of our data, Figure (4). The two-magnon peak is broader than experimental resolution, and the extracted positions (grey filled circles) are nearly dispersionless and are located at the calculated weighted average energies (grey band) replotted from Figure (2). Interchain effects are effectively integrated over in the two-magnon scattering, and this results in line broadening rather than shifts in energy; we thus have implicitly included interchain coupling effects in our definition of $`\alpha =0.27`$ for (1).
The right-hand panels of Figure (4) show one and two magnon integrated intensities extracted from the fit. The one-magnon intensity (lower panel) is well described by the theory (solid line), and shows a simple oscillation. Because the one-magnon intensity is modulated by the dimer envelope function in $`𝒮(𝐐,\omega )`$ the intensity should go to zero near $`\stackrel{~}{q}=4\pi `$. The residual intensity comes from secondary elastic scattering from incoherent processes, and the theory (solid line) is corrected for this. The two-magnon intensity (upper panel) shows a more complicated $`\stackrel{~}{q}`$ dependence, which is due to the various basis transitions that contribute to the coupling of the ground to excited states. Although the $`𝒪(\alpha )`$ result looks qualitatively similar to the data, it underestimates the scattering at $`\stackrel{~}{q}\frac{9}{2}\pi `$ and overestimates it at $`2\pi `$, which may indicate that higher order terms in the scattering amplitude are important. It is notable that the two-magnon intensity is very strongly dependent on the spatial arrangement of magnetic ions.
The binding energy of the $`S=1`$ state is predicted to be $`E_B=J\left(\frac{1}{4}\alpha \frac{13}{32}\alpha ^2\right)=0.017`$ meV for CN. Scattering around $`\stackrel{~}{q}=3\pi `$ is centered at $`0.852\pm .007`$ meV, which gives $`E_B=0.03\pm .02`$ meV. Although the energy and intensity around $`\stackrel{~}{q}=3\pi `$ lend support to binding around this bandwidth minimum, the experimental error means this does not constitute definitive proof of the effect in CN. The much better energy discrimination of time-of-flight (TOF) neutron spectrometers could provide this by resolving the bound mode from the continuum. Previously, TOF techniques have proven successful in the study of similar binding effects at the bandwidth minimum of the two-soliton continuum scattering of the $`S=1/2`$ XXZ Ising chain material CsCoCl<sub>3</sub> . Unlike the AHC, binding only occurs when extra terms in the Hamiltonian, such as exchange mixing or next-nearest neighbor coupling , are present. However limited neutron fluxes may make such measurements difficult for CN .
In conclusion, we have used inelastic neutron scattering to investigate the ground and excited states of the near-ideal alternating Heisenberg chain material Cu(NO<sub>3</sub>)$`{}_{2}{}^{}2.5`$D<sub>2</sub>O. Our measurements are consistent with predictions of this model for several magnetic properties of this system, including the ground state energy, one- and two-magnon excitation spectra and intensities, and possibly the existence of a two-magnon bound state. Much experimental work remains to be done to establish the phenomenology of binding in isotropic 1D systems.
DAT wishes to thank Drs B. Lebech, R. Hazell, and D. McMorrow for their help and Risø for generous support. This work was partly supported by Oak Ridge National Laboratory, managed by UT-Battelle, LLC, for the US Dept. of Energy under contract DE-AC05-00OR22725. The NSF supported work at SPINS through DMR-9423101 and work at JHU through DMR-9453362 and DMR-9801742. DHR acknowledges the generous support of the David and Lucile Packard Foundation.
|
warning/0005/cond-mat0005162.html
|
ar5iv
|
text
|
# Aging and Non-Linear Glassy Dynamics in a Mean-Field Model
## I Introduction
Thermal equilibrium is the situation where all fast processes have already taken place while slow processes have not yet started happening . At the opposite, systems with slow dynamics are characterised by a broad distribution of relaxation times, ranging from the microscopic scale ($`10^{12}`$ s) to the macroscopic one (hours or days). For instance, glassy systems have an equilibration time, either infinite, or much longer than the laboratory time scale. These systems reveal their out-of-equilibrium state in phenomena like aging or non-linear response.
In these systems, the microscopic time scale is not the only relevant one, and much slower processes also take place. The slow dynamics is generally attributed to the presence of thermally activated barrier crossing in the configuration space, but others mechanisms, such as the so-called “entropic barriers” may also contribute . Glassy dynamics is observed when the relaxation time $`\tau `$ becomes larger than the laboratory typical time scale, as it is the case for supercooled liquids .
The aging behaviour of spin glasses has been thoroughly investigated . The thermoremanent magnetization of field cooled samples shows a strong waiting-time dependence (where the waiting-time $`t_w`$ is the time interval between the temperature quench and the measurement). These systems have an a-priori infinite internal relaxation time, and the late stage of the relaxation is instead controlled by the waiting-time itself. Moreover, field cooled (waiting for $`t_w`$ and switching off the field) and zero field cooled (waiting for $`t_w`$ and switching on the field) show a remarkable complementarity of the magnetization curves . While out-of-equilibrium, as indicated by its significant waiting-time dependence, the response of the system is linear in the applied field, provided this one is weak enough.
Glassy dynamics is also observed in the dissipative dynamics of high-Tc superconductors. Supraconducting samples with quenched disorder, at magnetic field and temperature large enough, offer a significant resistance to a flowing dc current, due to the thermal motion of the flux lines. A transition line is believed to separate an ohmic regime (the vortex liquid) from a true superconducting state (the vortex glass) . In the latter, and in the limit of a vanishingly small current $`j`$, the dissipation occurs by activation of “bundles” of flux lines over pinning energy barriers. According to the scaling theory of the vortex glass, the typical time needed for such a move, $`\tau (j)`$, diverges exponentially fast as $`j`$ tend to zero . In this situation, the response (the voltage) is a non-linear function of the driving force (the current). The system is out-of-equilibrium, because of a constant rate dissipation, but stationary, at variance with the spin glass aging. The relaxation time $`\tau (j)`$ which would be infinite in the absence of driving force, is regularized by any small but constant $`j`$, and inversely related to the magnitude of $`j`$.
Important and related issues are the supercooled liquids dynamics and the rheology of soft glassy materials. In the first case, the mode-coupling approach predicts an increase of the structural relaxation time $`\tau `$ upon cooling . Aging have been found during the early stage of molecular dynamics studies of the Lennard-Jones fluid . On the other hand, a constant shear rate flow seems to be able to change the value of $`\tau `$, resulting in a shear rate dependent viscosity, e.g. a non-linear response of the fluid . More generally, this shear-thinning behaviour is well known in the context of soft-matter rheology, and a phenomenological treatment of this phenomenon based upon glassy dynamics has been proposed . These are situations where aging and non-linear response certainly coexist as manifestations of a more general glassy dynamics.
Among the existing theoretical approaches on glassy systems, the mean-field dynamics is a very promising one. It was already used to suggest that the presence of dissipative forces generically prevents the aging phenomenon . In this framework, the out-of-equilibrium character of a system is made precise by the existence of a generalised fluctuation dissipation theorem, related in turn to the entropy creation rate .
A model of particular interest is the mean-field dynamics of a particle with a quenched pinning potential. Isolated, the particle presents an aging behaviour with a logarithmic growth of the time correlation functions . In the presence of a time-independent driving force, the dynamics is believed to be stationary, with a power law dependence of the particle’s velocity in the applied force .
Recently, the author presented a geometrical description of the aging and linear response regime of this model, at zero temperature . This approach is extended, in this paper, to the non-linear stationary regime. As a result, we find that aging, linear response dynamics on the one hand, and stationary, non-linear response dynamics on the other hand, are indeed dual manifestations of a single out-of-equilibrium state. The constant force is found to interrupt efficiently the aging relaxational dynamics, and to control the characteristic times, which in turn control the effective friction coefficient in the stationary regime. The resulting velocity-force characteristics is $`vF^4`$, while the cross-over time between aging and stationary regime is $`t_FF^3`$. These predictions are confronted to the numerical integration of the mean-field equations. We finally suggest a scaling behaviour for the correlation function of this model which, according to our numerical findings, interpolate smoothly between the two different regimes, demonstrating their common origin.
## II Mean field equations and the Horner result
We focus on the zero temperature relaxational dynamics of a particle in a quenched random gaussian potential . The particle evolves in a $`N`$-dimensional space, under the simultaneous effect of a pinning force $`\mathbf{}V`$ and a constant force $``$, and the equation of motion for the vector position $`𝐱(t)`$ is :
$$\dot{𝐱}(t)=\mathbf{}V(𝐱(t))+.$$
(1)
The potential $`V(𝐱)`$ is a quenched disorder, chosen from a gaussian distribution, with correlations (the overline stands for the average over the quenched disorder) :
$$\overline{V(𝐱)V(𝐱^{})}=N\mathrm{exp}\left(\frac{𝐱𝐱^{}^2}{N}\right);\overline{V(𝐱)}=0.$$
(2)
This form ensures a meaningful $`N\mathrm{}`$ limit, in which each coordinate $`𝐱_i(t)`$, or gradient component $`_iV(𝐱)`$, remains of order one, while the norms $`𝐱(t)`$, $`\mathbf{}V`$ scale like $`N^{1/2}`$. The force is directed along the direction 1.
The thermodynamic limit $`N\mathrm{}`$ is taken first, which makes the zero temperature dynamics non trivial . In this mean-field limit, the relaxation process is completely described by the displacement $`u`$, the response function $`r`$ and the correlation functions $`b`$ and $`d`$, $`i\stackrel{~}{x}`$ being the Martin-Siggia-Rose auxiliary time .
$`u(t)`$ $`=`$ $`N^{1/2}\overline{x_1(t)};`$ (3)
$`r(t,t^{})`$ $`=`$ $`N^1{\displaystyle \underset{j=1}{\overset{N}{}}}\overline{x_j(t)i\stackrel{~}{x}_j(t^{})};`$ (4)
$`b(t,t^{})`$ $`=`$ $`N^1{\displaystyle \underset{j=2}{\overset{N}{}}}\overline{(x_j(t)x_j(t^{}))^2};`$ (5)
$`d(t,t^{})`$ $`=`$ $`N^1{\displaystyle \underset{j=1}{\overset{N}{}}}\overline{(x_j(t)x_j(t^{}))^2};`$ (6)
$`=`$ $`b(t,t^{})+[u(t)u(t^{})]^2.`$ (7)
The Dyson equations for $`r,b,d,u`$ are a closed system of coupled integro-differential equations, which contains the equations of as a particular case :
$`_tr(t,t^{})`$ $`=`$ $`\delta (tt^{})`$ (9)
$`4{\displaystyle _0^t}\text{d}s\mathrm{exp}(d(t,s))r(t,s)[r(t,t^{})r(s,t^{})];`$
$`_tb(t,t^{})`$ $`=`$ $`(2T)4{\displaystyle _0^t}\text{d}s\mathrm{exp}(d(t,s))[r(t,s)r(t^{},s)]`$ (11)
$`4{\displaystyle _0^t}\text{d}s\mathrm{exp}(d(t,s))r(t,s)[b(t,s)+b(t,t^{})b(s,t^{})];`$
$`_tu(t)`$ $`=`$ $`F4{\displaystyle _0^t}\text{d}s\mathrm{exp}(d(t,s))r(t,s)[u(t)u(s)].`$ (12)
The temperature term $`(2T)`$ is actually zero in our case. It is also convenient to define the integrated response $``$ and the energy $``$:
$`(t,t^{})`$ $`=`$ $`{\displaystyle _t^{}^t}\text{d}sr(t,s);`$ (14)
$`(t)`$ $`=`$ $`2{\displaystyle _0^t}\text{d}s\mathrm{exp}(d(t,s))r(t,s).`$ (15)
In a seminal paper, Horner described the stationary state reached by the system when driven with a finite force $`F`$, in the case of short range, power-law correlations . We have found that the system (LABEL:eq:Dyson-2temps:u) does indeed lead to a stationary situation, that we study in detail in this paper. One must mention however that the stationary state reached by the particle depends on how the system is prepared, in the same way as thermalised initial conditions can prevent aging in the $`p`$-spin case . The stationary state is found only if the system is quenched from a high enough temperature . Results for a similar driven system have also been recently published .
Let us summarise the main properties of the stationary solution found by Horner . The correlation functions are time-translationally invariant (TTI) : $`r(t,t^{})=R(tt^{})`$, $`B(t,t^{})=B(tt^{})`$ while the displacement goes linearly with time $`u(t)=vt`$. The system (LABEL:eq:Dyson-2temps:u) becomes a set of non-causal equations to be solved self-consistently. A non trivial feature of this solution is the emergence of characteristic time scales dependent on the velocity. With the notations of , $`t_p(v)`$ is the characteristic time for breaking the fluctuation dissipation theorem, while $`t_a(v)`$ controls the main “$`\alpha `$” relaxation of the correlation function $`B(t)`$. These characteristic times play a very similar role than the time scales $`t_f`$ and $`t_b`$ respectively, introduced in , and we identify subsequently $`t_ft_p`$, $`t_bt_a`$. In the long time regime, $`tt_b`$:
$$B(t)=q+\widehat{B}\left(\frac{t}{t_b(v)}\right),$$
(16)
with,
$`t_f(v)`$ $``$ $`v^{(\eta 1)\zeta },`$ (17)
$`t_b(v)`$ $``$ $`v^{\eta 1},\mathrm{\hspace{0.33em}0}<\eta ,\zeta <1.`$ (18)
The exponents depend (in a complex way) on the correlator (2. $`\widehat{B}`$ is a scaling function discussed in appendix B, and $`q`$ is the “plateau value” of $`B(t)`$, equal to 0 in the zero temperature limit. Meanwhile, the fluctuation dissipation theorem, obeyed for $`tt_f`$, is violated around $`tt_f`$, and becomes :
$$\text{d}B(t)/\text{d}t=2\overline{T}R(t);tt_f.$$
(19)
The effective temperature $`\overline{T}`$ and the plateau value $`q`$ are identical to those obtained in the aging case .
The velocity-force characteristics is given by (LABEL:eq:Dyson-2temps:u), and in the limit of small velocities,
$$vF/t_b(v).$$
(20)
The time $`t_b(v)`$ plays the role of an effective friction coefficient, controlled by the velocity. The $`vF`$ characteristics is a power law.
$$vF^{1/\eta }.$$
(21)
Let us mention for completeness the presence of a third time scale, called $`t_a^{}`$ in , defined by $`B(t_a^{})=v^2t_{}^{}{}_{a}{}^{2}`$. As we consider an exponential correlator, we have $`\mathrm{exp}(B(t)v^2t^2)\mathrm{exp}(B(t))\times \mathrm{exp}(v^2t^2)`$ and in our case $`t_a^{}v^1`$. We believe that apart from this point, the results of all qualitatively apply to the exponential correlator, and anyway $`t_a^{}`$ does not play a direct role in the dynamics of the short range correlated models.
## III Geometrical description of the driven stationary dynamics
In was proposed a geometrical description of the relaxational dynamics of the particle. This approach makes use of a comoving frame, defined by the eigendirections of the hessian matrix $`\mathbf{}\mathbf{}V(𝐱)`$ at the precise point $`𝐱(t)`$ where the particle stands. This frame is made of $`N`$ vectors $`\{𝐞_i\}`$, each one eigenvector of the hessian $`\mathbf{}\mathbf{}V(𝐱)`$. The distribution of the corresponding eigenvalues is a semi-circle of radius $`4`$, shifted towards the positive values, and such that the lowest one is equal to $`\mathrm{SS}`$. Each eigenvector $`𝐞_i`$ has an eigenvalue $`\lambda _i\mathrm{SS}`$, and the density of states of the $`\lambda _i`$ is :
$$\rho (\lambda )=(8\pi )^1\sqrt{\lambda (8\lambda )}.$$
(22)
The quantity $`\mathrm{SS}`$ is positive, and depends linearly on the energy of the system, e.g. $`\mathrm{SS}(t)=4+2V(𝐱(t))/N=4+2(t)`$. In the aging case, $`\mathrm{SS}(t)`$ is a time-dependent function, while in the stationary case, $`\mathrm{SS}`$ is constant.
One projects the instantaneous velocity $`\dot{𝐱}`$ is the above frame such that:
$$\dot{𝐱}=\underset{i=1}{\overset{N}{}}\gamma _i𝐞_i$$
(23)
Because the spacing of the eigenvalues is of order $`1/N`$, the set of $`\lambda _i`$ becomes dense, and one replaces the discrete sum over the index $`i`$ by a continuous one, involving the semi-circular density of eigenvalues $`\rho (\lambda )`$.
$$\dot{𝐱}^2=_0^8\text{d}\lambda \rho (\lambda )g(\lambda ,t)$$
(24)
The distribution $`g(\lambda ,t)`$ represents the mean value of the component $`𝐱_i^2`$, locally averaged over the indices $`i`$ such that $`\lambda _i\lambda `$. We have justified in the following self-similar form for $`g(\lambda ,t)`$:
$`𝐱_i^2`$ $``$ $`g(\lambda _i,t)`$ (25)
$`g(\lambda ,t)`$ $`=`$ $`\mathrm{SS}(t)\widehat{G}\left(\lambda /\mathrm{SS}(t)\right),`$ (26)
where $`\lambda `$ stands for any direction with a curvature of the potential equal to $`\lambda \mathrm{SS}`$ . The prefactor $`\mathrm{SS}`$ in front of the distribution comes in fact from an assumption about the value of the exponent $`\kappa `$ governing the power law decay of $`(t)`$ and $`\mathrm{SS}(t)t^\kappa `$, which we believe to be $`2/3`$. This assumption is supported by our numerical results.
The characteristic times $`t_f`$ and $`t_b`$ are controlled by $`\mathrm{SS}`$, and scale like :
$`t_f`$ $``$ $`\mathrm{SS}^1;`$ (27)
$`t_b`$ $``$ $`\mathrm{SS}^{3/2};`$ (28)
making the instantaneous velocity equal to :
$$\dot{u}(t)=F\mathrm{SS}^{3/2}=F/t_b(t).$$
(29)
When a constant force $`F`$ is applied, the dynamics changes from an aging linear-response behaviour to a stationary regime . At short times, the displacement $`u(t)`$ is proportional to the integrated response $`F(t,0)`$, while at long times, it becomes equal to $`vt`$.
One expects the stationary regime to take over the aging regime when the dynamics is dominated by the external force $``$ rather than by the gradient $`\mathbf{}V(𝐱(t))`$. This happens at a time $`t_F`$, inversely related to the magnitude of the force. Our numerical definition of $`t_F`$, is the time where the slope of the asymptotic curve $`u(t)/Fvt/F`$ is equal to the slope of the (logarithmic) integrated response $`(t,0)`$, as shown on Figure (1) with $`F=0.3`$.
In our situation, the dynamics is controlled by the value of $`\mathrm{SS}`$; inverse friction and diffusion coefficients are both proportional to $`\mathrm{SS}^{3/2}`$. In the aging case, $`\mathrm{SS}(t)`$ tends to zero as a power law, and the dynamics of the system is slower and slower. A look at Figure (2) however shows that in the presence of a force, $`\mathrm{SS}`$ does not go to zero, but to a finite value $`\mathrm{SS}(F)`$, controlled by $`F`$, and inversely related to the magnitude of $`F`$. The same is true for the energy $`(t)=2+\mathrm{SS}(t)/2`$, which stands higher than in the absence of driving force. Both diffusivity and mobility are kept finite thanks to a non-zero driving force. What is needed is to compute $`\mathrm{SS}(F)`$. For this purpose, one assumes that the self-similar form (26) is still valid in the stationary regime. A justification is provided in appendix A.
The zero temperature relaxation equation is:
$$\dot{𝐱}_i=_iV(𝐱)+_i,$$
(30)
while the energy obeys :
$`\dot{}(t)`$ $`=`$ $`1/N{\displaystyle \underset{i}{}}_iV(𝐱(t))\dot{𝐱}_i(t),`$ (31)
$`\dot{}(t)`$ $`=`$ $`\dot{𝐱}^2(t)/N+F\dot{u}(t).`$ (32)
One uses now the distribution $`g(\lambda ,t)`$ of the instantaneous velocity components $`\dot{𝐱}_i^2`$, the density of eigenvalues $`\rho (\lambda )`$ and finds :
$$\dot{}(t)+F\dot{u}(t)=\text{d}\lambda \rho (\lambda )g(\lambda ,t).$$
(33)
In the stationary regime, $`\dot{}(t)=0`$ and $`\dot{u}(t)=v`$. The equation (33) reduces to a balance between the mechanical power given by the force, and a kind of intrinsic dissipation $`(\dot{𝐱}^2)`$.
$$Fv=\text{d}\lambda \rho (\lambda )g(\lambda ).$$
(34)
Assuming that $`g`$ is still equal to $`\mathrm{SS}\widehat{G}(\lambda /\mathrm{SS})`$ (cf appendix A), one gets :
$$Fv\mathrm{SS}^{5/2}.$$
(35)
From (35) and (29), one finally finds $`\mathrm{SS}`$ as a function of the force,
$$\mathrm{SS}F^2,$$
(36)
the resulting velocity force characteristics,
$$vF^4,$$
(37)
and the force and velocity dependence of the time scales :
$`t_f`$ $``$ $`F^2v^{1/2};`$ (38)
$`t_b`$ $``$ $`F^3v^{3/4}.`$ (39)
These results are in full qualitative agreement with the findings of Horner . The main relaxation time $`t_b`$ does not scale as $`v^1`$, as could be expected from a simple dimensional analysis, but is shorter, such that $`lim_{v0}vt_b(v)=0`$.
One determines the cross-over time $`t_F`$ by a matching argument. In the linear response regime, $`\mathrm{SS}(t)`$ decreases as $`t^{2/3}`$, as the force acts only as a weak perturbation. The linear response breaks down when the perturbation modifies the nature of the relaxation itself. This happens when $`\mathrm{SS}(t_F)`$ reaches the order of magnitude of its limit value $`\mathrm{SS}(F)=F^2`$ (equation 36), leading to $`t_F^{2/3}=F^2`$, or:
$$t_F=F^3.$$
(40)
Physically, this means that a typical coordinate $`f_i`$ of the force $``$, along a downhill direction $`i`$, is of the same order of magnitude than the gradient of the potential $`_iV`$ , or the instantaneous velocity $`\dot{𝐱}_i`$ . From equation (26) and $`f_iF`$, one gets $`F^2f_i^2\dot{𝐱}_i^2\mathrm{SS}`$, in agreement with (36).
Let us mention that a qualitatively similar cross-over has been observed in the simulated dynamics of a driven polymer, in the presence of quenched disorder .
## IV Numerical results
We present numerical results which support the findings of the previous section.
Figure (3) shows $`\mathrm{SS}(F)`$ versus $`F`$, in log coordinates, and in regular coordinates (inset), for $`F=0.05,0.1,0.2,0.3,0.4,0.5`$ and $`F=0.6`$. The squares are the values obtained with $`t_{max}=200`$ ($`h=0.1`$), and the full curve with $`t_{max}=400`$ ($`h=0.2`$). One sees that the three first values are not well converged. If we except them, the overall shape of the curve is concave (downward curvature). The slope of the tangent curve between the arrows gives an exponent equal to $`1.81`$ ($`h=0.1`$) and $`1.87`$ ($`h=0.2`$). Because the curve is concave, we believe that these values are a lower bound for the real exponent, compatible with our prediction $`2`$.
Figure (4) shows $`v(F)`$ versus $`F`$, in log coordinates, and in regular coordinates (inset). The squares are the values obtained with $`t_{max}=200`$ ($`h=0.1`$), and the crosses with $`t_{max}=400`$ ($`h=0.2`$). As for Figure (3), the three first values are not well converged. If we except them, the overall shape of the curve is again concave. The slope of the tangent curve between the arrows gives an exponent equal to $`3.73`$ ($`h=0.1`$) and $`3.82`$ ($`h=0.2`$). Repeating the above argument, these values are a lower bound for the real exponent, compatible with 4.
A plot of $`t_F`$ vs $`F`$ is reported on Figure (5). Again the system has not reached its asymptotic regime as far as the three first values $`F0.2`$ are concerned. This can be checked by looking at the first derivative $`\dot{u}(t)`$ which must be constant when $`t`$ reaches the upper limit of the time window, here $`t=400`$. The fitted value on the straight part of the graph, in logarithmic coordinates, gives an exponent $`2.72`$ instead of $`3`$. Again the true asymptotic limit $`F0`$ is out of reach, due to our limited computer facilities.
These numerical results are not good enough to prove the exactness of the equations (36) and (37). However they provide lower bounds which constraint the exponents to be larger than $`1.8`$ for $`\mathrm{SS}`$, and larger than $`3.8`$ for $`v`$. On the other hand, if we assume that we are close enough to the asymptotic regime where equations (36) and (37) apply, one expects the real exponents to be not too much different from the above numerical values. In this respect, we think that the numerics is in agreement with our findings. As far as the cross-over time is concerned, the numerical exponent is $`2.72`$ instead of $`3`$. A larger time window would certainly improve the agreement.
## V A unified description of the out-of-equilibrium regimes
In the isolated aging regime, at zero temperature, the correlation function obeys, as a particular case of equation (B1) of appendix B, as shown in :
$$b(t,t^{})=\mathrm{ln}\left(\frac{h(t)}{h(t^{})}\right),$$
(41)
where the parametrisation function is related to the time-scale $`t_b`$ by :
$$t_b(t)=h(t)/h^{}(t)$$
(42)
As the time scale $`t_b`$ is proportional to $`\mathrm{SS}^{3/2}`$, we have :
$`{\displaystyle \frac{h(t)}{h(t^{})}}`$ $`=`$ $`\mathrm{exp}\left[C{\displaystyle _t^{}^t}\text{d}s\mathrm{SS}^{3/2}\right],`$ (43)
$`b(t,t^{})`$ $`=`$ $`C{\displaystyle _t^{}^t}\text{d}s\mathrm{SS}^{3/2}.`$ (44)
On the other hand, equation (29) leads immediately to
$$\frac{u(t)u(t^{})}{F}=C^{}_t^{}^t\text{d}s\mathrm{SS}^{3/2}.$$
(45)
Now, one observes that the scaling form (16) resemble to (41), (B1), with $`q=0`$. We prove in appendix B that the scaling function of the aging regime and the driven regime are indeed equal, and thus $`\widehat{B}(x)=x`$ in (16) (strictly speaking, $`\widehat{B}(x)`$ is only proportional to $`x`$, but one can choose $`t_b`$ such as $`\widehat{B}(x)=x`$). The equations (44) and (45) make sense in the aging regime as well as in the stationary regime.
The integral $`_t^{}^t\text{d}s\mathrm{SS}^{3/2}`$ is the effective time variable for the system, interpolating smoothly between $`\mathrm{ln}(t)`$ (aging, linear response regime) and $`t/t_b(F)`$ (stationary regime) while $`S^{3/2}`$ is an effective age, growing like the waiting time, in the aging regime, and bounded in the stationary regime. Interestingly, a similar effective age has been used in the context of the stick-slip motion in dry friction experiments .
The prediction for (44) and (45) is checked by plotting $`b(t,t^{})`$ vs $`[u(t)u(t^{})]/F`$, shown on Figures (6) and (7). One expects $`b(t,t^{})`$ and $`[u(t)u(t^{})]/F`$ to be proportional, both in the aging and stationary regimes, provided equations (44, 45) hold, which is the case for a time separation $`tt^{}`$ large enough.
On Figure (6), $`b(t,t^{})`$ is plotted against $`[u(t)u(t^{})]/F`$ for $`F=0.1`$ (crosses, squares and diamonds) in the linear response regime and for $`F=0.5`$ (continuous lines) in the non-linear regime. The force is zero till $`t=t^{}`$, and then switched on; $`t^{}`$ takes the value 0, 20 and 40. As far as $`F=0.5`$ is concerned, the transition from linear to non-linear regime is not visible on this curve, and in any case very smooth. The slope of the curve defines the effective temperature $`2\overline{T}`$, equal to the ratio $`C/C^{}`$ in equations (44) and (45). The effective temperature thus makes sense in both linear and non-linear regimes.
As the force is switched on at $`t^{}`$, there is a short-time “elastic” displacement. This is how the directions with a positive curvature respond to the new static constraint, and this corresponds to the short horizontal step at the origin, seen on Figures (6) (inset) and (7). The finite slope part of the curve corresponds to the slow wandering motion of the particle in the energy landscape, in the regime where equations (44) and (45) apply. Thus, we conclude that Figure (6) support the proportionality of $`u(t)u(t^{})`$ and $`b(t,t^{})`$, once the short time regime has been taken into account.
A close look near the origin of the graph (inset of Figure (6) shows that the $`F=0.5`$ curve is slightly shifted from $`F=0.1`$, but parallel to it. This shift goes rapidly to zero as $`F0`$. The shift is presumably there because $`0.5`$ is already a large value of the force, leading to a departure from the ideal curve corresponding to $`F1`$.
Figure (7) is the same as Figure (6) for $`F=0.01`$, $`F=0.1`$ and $`F=0.5`$, for three values of $`t^{}`$, 0, 20 and 40, and gives additional details on the short time response of the particle. Again, the horizontal part of the curves corresponds to the short-time displacement (“elastic” or reversible) while the finite slope regime corresponds to the slow motion in the energy landscape (“plastic” or irreversible).
## VI Conclusion
In this paper, we have proposed a consistent picture for the stationary driven dynamics, in the mean field approximation and zero temperature limit, of a particle in a quenched, exponentially correlated, random potential.
The velocity $`(v)`$\- force $`(F)`$ relation is a power law $`v=F^4`$, while the main relaxation time scales as $`t_bv^{3/4}`$. The product $`vt_b`$ tends to zero as $`v`$ vanishes. These findings are consistent with earlier work . The driving force is found to generate a relaxation time smaller than the “dimensional” time scale $`v^1`$, which is probably a generic feature of the mean-field short-range correlated potentials.
If the force $`F`$ is small enough, a linear response around the aging regime is found, up to a time $`t_F`$, scaling as $`F^3`$. A plot of the displacement $`(u(t)u(t^{}))/F`$ vs the correlation $`b(t,t^{})`$ shows no sign of discontinuity, when the linear response regime is replaced with the non-linear stationary regime. We interpret it by saying that, when a small force is applied, the dynamical properties of the system (mobility, diffusivity) are controlled by the effective age $`\mathrm{SS}^{3/2}`$. The quantity $`\mathrm{SS}^{3/2}`$ is proportional to the number of negative eigenvalues in the spectrum of the hessian of the hamiltonian. The effective age is proportional to the waiting time in the aging regime, and finite in the stationary case.
The effective time $`^t\mathrm{SS}^{3/2}\text{d}s`$, closely related to the correlation function $`b(t,t_0)`$, grows logarithmically with $`t`$ in the aging regime, and linearly with $`t`$ in the stationary regime. The effective temperature $`\overline{T}`$ generalising the fluctuation dissipation theorem, remains unchanged in the non-linear regime. However, the geometrical meaning of $`\overline{T}`$, if any, is still unknown.
Future work will determine to what extent are the present features generic from other short range correlated models, and finite dimensional models. Even though such a power law dependence of the characteristic times in the driving force is not observed in realistic systems, the qualitative behaviour presented in this study –cross-over between linear to non-linear regime, coexistence of aging and non-linear stationary dynamics–, could indeed be a very generic situation.
### Acknowledgements
I especially thank L.Cugliandolo and J.Kurchan for having lent me their numerical code, and S.Scheidl, J.P Bouchaud, J.Kurchan, M.Mézard and A.Cavagna for discussions on this field. I thank D.Feinberg for suggestions and criticisms about the manuscript. I warmly thank the hospitality of the Department of Physics, IISc, Bangalore, where a part of the writing has been done.
## A The energy balance
Let $`\gamma _i(t)`$ be the coordinates of the instantaneous velocity $`\dot{𝐱}(t)`$ in the comoving frame $`\{𝐞_i(t)\}`$ (23). When $`=0`$, this definition is equivalent to say that $`\gamma _i`$ is the coordinate of $`\mathbf{}V`$. Using the local average defined in , one finds:
$`\dot{𝐱}^2(t)`$ $`=`$ $`{\displaystyle \underset{i}{}}\dot{𝐱}_i^2(t)={\displaystyle \underset{i}{}}\gamma _i^2(t),`$ (A1)
$`=`$ $`N{\displaystyle \text{d}\lambda \rho (\lambda )g(\lambda ,t)}.`$ (A2)
The derivative of $`\dot{𝐱}^2(t)`$ reads :
$`_t\dot{𝐱}^2(t)`$ $`=`$ $`{\displaystyle \underset{j}{\overset{N}{}}}_t(_jV+_j)^2`$ (A3)
$`=`$ $`2{\displaystyle \underset{jk}{}}_{jk}V\dot{𝐱}_j\dot{𝐱}_k`$ (A4)
$`=`$ $`2N{\displaystyle \text{d}\lambda \rho (\lambda )(\lambda \mathrm{SS})g(\lambda ,t)}.`$ (A5)
We deduce that, in the stationary situation, for all $`\mathrm{SS}`$,
$$\text{d}\lambda \rho (\lambda )\lambda g(\lambda )=\mathrm{SS}\text{d}\lambda \rho (\lambda )g(\lambda ),$$
(A6)
which is in favour of a scaling form $`g(\lambda )=\mathrm{\Gamma }\widehat{G}(\lambda /\mathrm{SS})`$.
As $`_iV(𝐱(t))=\dot{𝐱}_i(t)+_i`$, the equation for $`\dot{}(t)`$ is :
$`\dot{}(t)`$ $`=`$ $`N^1{\displaystyle \underset{j}{}}_jV\dot{𝐱}_j`$ (A7)
$`=`$ $`N^1{\displaystyle \underset{j}{}}\{\dot{𝐱}_j^2(t)+_j\dot{𝐱}_j(t)\}`$ (A8)
The product $`N^1_j_j\dot{𝐱}_j`$ is by construction equal to $`F\dot{u}(t)`$. Thus, (this is equation 33):
$$\dot{}(t)=\text{d}\lambda \rho (\lambda )g(\lambda ,t)+F\dot{u}(t).$$
(A9)
The energy balance (35), and the factorised form of $`g(\lambda )`$ imply in the stationary regime :
$$\mathrm{\Gamma }\mathrm{SS}^{3/2}F^2\mathrm{SS}^{3/2}$$
(A10)
However the relation between $`\mathrm{SS}`$ and $`\mathrm{\Gamma }`$ remains undetermined by the present argument. For the sake of simplicity, we can suppose that the equality $`\mathrm{SS}=\mathrm{\Gamma }`$, true in the aging regime, remains true in the stationary regime. This assumption is in fact equivalent to a matching argument, when the distribution $`g(\lambda ,t)=\mathrm{SS}(t)\widehat{G}(\lambda /\mathrm{SS}(t))`$, crosses over the distribution $`g(\lambda )=F^2\widehat{G}(\lambda /\mathrm{SS}(F))`$ around $`t=t_F`$. The matching of $`g(\lambda ,t)`$ and $`g(\lambda )`$ leads to the identification $`\mathrm{SS}(F)=F^2`$. One cannot rule out, rigorously, more complicated behaviours, which could lead to a different velocity-force characteristics. The assumption $`\mathrm{\Gamma }=\mathrm{SS}`$ is just the most natural one.
## B The scaling form of the correlation function
In the isolated situation, the correlation function in the aging regime reads, for any finite temperature $`T`$ :
$$b(t,t^{})=q+\stackrel{~}{B}\left[\mathrm{ln}\left(\frac{h(t)}{h(t^{})}\right)\right]$$
(B1)
An general equation for $`\stackrel{~}{B}(u)`$ is obtained in reference (equation 6.22, with the opposite sign convention for $`f`$), and reads:
$`0`$ $`=`$ $`\stackrel{~}{B}(u)f^{\prime \prime }(q)f^{}(q+\stackrel{~}{B}(u))+f^{}(q)`$ (B3)
$`+{\displaystyle \frac{2\chi q}{T}}f^{\prime \prime }(q){\displaystyle _0^u}\text{d}u^{}\stackrel{~}{B}^{}(u^{})f^{\prime \prime }(q+\stackrel{~}{B}(u))\stackrel{~}{B}(uu^{})`$
whose solution is $`\stackrel{~}{B}(u)=C^{st}\times u`$, leading to (41), with $`q=0`$. The function $`f=\mathrm{exp}(x)`$ stands for the correlator (2), $`T`$ for the temperature, $`q`$ for plateau value of the correlation function $`b`$, and $`\chi `$ for the fluctuation dissipation violation parameter (see equation B7 below).
On the other hand, in the stationary regime, the equation for $`b(t,t^{})=B(tt^{})=q+\overline{B}(tt^{})`$ and $`R(tt^{})=r(t,t^{})`$ is (equation (2.9) in reference , again with the opposite sign convention for $`f`$):
$`_tB(t)`$ $`=`$ $`2T\left({\displaystyle _0^{\mathrm{}}}\text{d}s4f^{\prime \prime }(B(s)+v^2s^2)R(s)\right)B(t)+{\displaystyle _0^t}\text{d}s4f^{\prime \prime }(B(s)+v^2s^2)R(s)B(ts)`$ (B6)
$`+{\displaystyle _0^{\mathrm{}}}\text{d}s\{(4f^{}(B(t+s)+v^2(t+s)^2)4f^{}(B(s)+v^2s^2))R(s)`$
$`+(4f^{\prime \prime }(B(t+s)+v^2(t+s)^2)R(t+s)4f^{\prime \prime }(B(s)+v^2s^2)R(s))B(s)\}`$
One knows that the main relaxation scale $`t_b`$ is much smaller than $`v^1`$, and asymptotically, $`lim_{t0}vt_b=0`$. The above integrals can be safely cut beyond a cut-off $`\mathrm{\Lambda }`$ such that $`t_b\mathrm{\Lambda }v^1`$. The contributions $`_\mathrm{\Lambda }^{\mathrm{}}`$ are negligible because the relaxation of $`B(t)`$ has already taken place, while in the integrals $`_0^\mathrm{\Lambda }`$, the term $`v^2s^2`$ can be neglected compared with $`B(s)`$ in the argument of the correlators $`f^{}`$ and $`f^{\prime \prime }`$.
One introduces the quasi fluctuation dissipation parameter $`X`$, defined by:
$$X(\overline{B}(t))\text{d}B(t)/\text{d}t=R(t).$$
(B7)
$`X`$ is equal to its equilibrium value $`1/2T`$ if $`\overline{B}<0`$ and to $`\chi `$ if $`\overline{B}>0`$. Equation (B6) becomes:
$`_tB(t)`$ $`=`$ $`2T+\left({\displaystyle _0^\mathrm{\Lambda }}\text{d}s4f^{\prime \prime }(q+\overline{B}(s))X(\overline{B}(s))\text{d}\overline{B}(s)/\text{d}s\right)B(t){\displaystyle _0^t}4f^{\prime \prime }(\overline{B}(s))X(\overline{B}(s))\text{d}\overline{B}(s)/\text{d}sB(ts)`$ (B10)
$`{\displaystyle _0^\mathrm{\Lambda }}\text{d}s\{(4f^{}(q+\overline{B}(t+s))4f^{}(q+\overline{B}(s)))X(\overline{B}(s))\text{d}\overline{B}(s)/\text{d}s`$
$`+(4f^{\prime \prime }(q+\overline{B}(t+s))X(\overline{B}(t+s))\text{d}\overline{B}(t+s)/\text{d}s4f^{\prime \prime }(q+\overline{B}(s))X(\overline{B}(s))\text{d}\overline{B}(s)/\text{d}s)\}B(t)`$
Each integral $`_a^b`$ has to be split to take into account the short time quasi-equilibrium regime and the long time regime. As the time scale $`t_f`$ separates these two regimes, one writes $`_a^b=_a^{a+t_f}+_{a+t_f}^{bt_f}+_{bt_f}^b`$. The parameter $`X`$ is then set to $`1/2T`$ or $`\chi `$ accordingly, and most of the integrals can be reduced to boundary terms. One neglects the time derivative $`_tB(t)`$ in the asymptotic long-time regime, and the result is:
$`0`$ $`=`$ $`2T+\left(4\chi {\displaystyle _{t_f0}^\mathrm{\Lambda }\mathrm{}}f^{\prime \prime }(q+\overline{B})\text{d}\overline{B}\right)(q+\overline{B}(t))`$ (B13)
$`4q\left(\chi +{\displaystyle \frac{1}{2T}}\right)\times \left(f^{}(q)f^{}(q+\overline{B}(t))\right)`$
$`4\chi {\displaystyle _{t_f0}^t}\text{d}s\text{d}\overline{B}(s)/\text{d}sf^{\prime \prime }(q+\overline{B}(s))(q+\overline{B}(ts))`$
By using $`lim_t\mathrm{}\overline{B}(t)=\mathrm{}`$, $`q^2f^{\prime \prime }(q)=T^2`$ and $`4\chi f^{\prime \prime }(\overline{B})\text{d}\overline{B}=2T/q`$, the equation (B13) for $`\overline{B}`$ coincides exactly with (B3). As the equation (B13) is invariant upon time dilatations, $`\widehat{B}(u)=\overline{B}(t/t_b)`$ is a solution of (B3) and without loss of generality, one has:
$$\widehat{B}(u)=\stackrel{~}{B}(u)=u,$$
(B14)
which is the announced result.
Captions
FIGURE 1. Determination of the cross-over time $`t_F`$, defined as the time where the slope of the integrated response $`(t,0)`$ is equal to the velocity $`1/Flim_t\mathrm{}\dot{u}(t)`$.
FIGURE 2. Family of curves $`(t)+2`$ for increasing forces, ranging from $`F=0.05`$ to 0.5. The limit value $`lim_t\mathrm{}(t)+2`$ is a monotonically increasing function of $`F`$, equal to $`\mathrm{SS}(F)/2`$. The effective mobility and diffusivity are directly related to $`\mathrm{SS}(F)`$. The system stays above the marginal states, in a region with a finite extensive number of downhill directions.
FIGURE 3 The parameter $`\mathrm{SS}`$ as a function of the force, for $`F=0.1,0.2,0.3,0.4,0.5`$ and $`0.6`$, in log coordinates, and normal coordinates (inset). The boxes stand for a run, up to a time $`t=200`$ while the straight line corresponds to $`t=400`$. Whenever the boxes differ from the line, the value is not converged. See text for details.
FIGURE 4. The velocity $`v`$ as a function of the force $`F`$, in log coordinates, and normal coordinates (inset). Same remark as for Figure (4).
FIGURE 5. The time $`t_F`$ as a function of the force $`F`$ in log coordinates. Inset : $`t_F`$ as a function of $`F`$. The three first values are not accurate ($`t_F`$ larger than our maximum time). The fitted exponent of the straight part is $`2.72`$ instead of $`3`$; $`2.72`$ is a lower bound for the real value.
FIGURE 6. The correlation $`b(t,t^{})`$ vs the displacement $`[u(t)u(t^{})]/F`$, for $`F=0.1`$ and $`0.5`$. The force is switched on at $`t^{}`$, successively equal to $`0,20`$ and $`40`$. Inset: the short-time behaviour. See text for details.
FIGURE 7. Same as Figure (6), with $`F=0.01`$, $`F=0.1`$ and $`F=0.5`$.
|
warning/0005/hep-th0005123.html
|
ar5iv
|
text
|
# KUNS-1664hep-th/0005123 Open membranes in a constant 𝐶-field background and noncommutative boundary strings
## 1 Introduction
It is surprising that, although it seems that noncommutative geometry is quite a pure mathematical object, noncommutativity does emerge in some definite limits of string theory. For instance, matrix theory compactified on tori gives Yang-Mills theory on noncommutative tori; the quantization of open strings on a D-brane with a background $`B`$-field leads this D-brane world-volume to become noncommutative; the twisted version of the reduced large-$`N`$ Super Yang-Mills model originally considered as a constructive definition of type IIB superstring can be interpreted as noncommutative Yang-Mills theory, and so on.
Recent development on string dualities reveals that M-theory rules nonperturbative features of superstring theories. It is natural to ask what is noncommutativity in M-theory. We do not know so much about M-theory. M-theory leads to eleven dimensional supergravity at the low-energy limit, and M-theory compactified on a circle becomes type IIA superstring by taking the limit for the radius of the circle to become zero. Moreover M-theory contains the two-dimensional extended object, M2-brane, as the fundamental component. Matrix theory proposed by Banks, Fischler, Shenker and Susskind is considered as describing some (or complete as they state originally) degrees of freedom of M-theory. This matrix theory does show noncommutativity in some cases commented above. We can expect naturally that noncommutativity can emerge in M-theory.
On the other hand, a supersymmetric two-dimensional extended object, called supermembrane, is interesting in its connection to superstrings. A quantum extension of supermembrane is expected to give a definition of M-theory. Especially, it is well known that supermembrane in eleven dimensions can consistently couple to eleven dimensional supergravity as its backgrounds. Thus, we have a natural question here; how does supermembrane theory show noncommutativity? It is a very meaningful question in two reasons. First, since we expect that supermembrane is a definition of M-theory, we also expect that supermembrane theory has noncommutativity in a definite limit or a background. Secondly, we wonder what is noncommutativity in more than two-dimensional extended objects. To clear this second point, let us compare it with the string case. In string theory, the end of open strings becomes noncommutative and a D-brane world-volume on which open strings can end has noncommutative geometry. Then, let us consider an open membrane which has one-dimensional boundary and focus on the behavior of these boundaries. Here, we face a conceptual jump. In string theory, open string ends are “points” and on a D-brane world-volume points do not commute with each other, while in membrane case, we find that its boundaries are “strings” and noncommutativity means one-dimensional strings do not commute with each other. Thus, we can learn a new feature of noncommutative geometry by studying membrane noncommutativity. A primitive analysis was carried out in .
In string theory, we can find noncommutativity by quantizing open strings in background NS-NS fields. Some authors have applied the Dirac procedure to boundary conditions. This method is very transparent and can be easily extended to other systems. We attempt to investigate an open membrane in a background three-form field in this way. It is well known that to investigate membrane theory has severe difficulties, for example, non-linearity of world-volume theory, non-renormalizability of three-dimensional sigma model, and so on. Thus, we must take an appropriate approximation, as explained later.
Our plan of investigation is as follows. In seeing the noncommutativity, supersymmetry was not essential in the string case. We drop the fermionic parts and consider a bosonic membrane. We start with a bosonic open membrane in a constant gauge field background. Since we should take our bosonic membrane as a toy model of eleven dimensional supermembrane, we restrict the background fields to the massless bosonic fields of eleven dimensional supergravity, the metric $`g_{\mu \nu }`$ and the three-form tensor field $`C_{\mu \nu \rho }`$. We consider only a bosonic background and drop the fermionic field, the gravitino $`\chi ^\mu `$. Without introducing a two-form gauge field, there can not exist open membranes by gauge-invariance. Also in supermembrane case, we can not introduce an open supermembrane without braking all the supersymmetries in flat Minkowski space-time. However we can formulate a supersymmetric open supermembrane when there exists a “topological defect” as a background . These defects are interpreted as, for instance, M5-brane, “end of the world” 9-plane in Hořava-Witten’s sense, etc. We shall introduce fixed $`p`$-branes in this bosonic case. We assume our open membranes are bounded to these “boundary planes,” and there is a two-form field, to which open membrane boundaries can couple, on these planes<sup>1</sup><sup>1</sup>1In , an open membrane probe was used to derive the equations of motion of boundary $`M5`$-branes.. In these settings, we calculate the Dirac brackets and confirm noncommutativity on these boundary planes. Our calculation is only to second order in $`C`$ and not exact.
This paper is organized as follows. In section 2, we propose our setup. We consider a bosonic open membrane in a constant $`C`$-field background. We suppose that one direction of the target space is compactified to a circle, another direction is compactified to an interval and there exist two fixed planes at the boundaries of this direction. We fix the reparametrization invariance of the world-volume with a static gauge and simplify the action by taking a limit. Equations of motion and boundary conditions are found, we go on to the canonical formalism and impose the boundary conditions as constraints. In section 3, we solve the constraints with an approximation. We take the radius of the compactification circle to be very large and the distance between the boundary planes to be infinitesimally small. In section 4, we calculate the Dirac brackets and confirm the noncommutativity on the boundary planes. Section 5 is served to discussions and remarks. In appendix A, we review the application of Dirac’s procedure for constrained systems to the boundary constraints in the string case.
## 2 An open membrane in a constant $`C`$-field
Let us consider an open membrane in the background of a constant three-form tensor field $`C_{\mu \nu \rho }`$. We suppose that our membrane topology is cylindrical and the background is eleven dimensional, compactified to $`\text{R}^{9p}\times M^p\times S^1\times I`$, where $`M^p`$ is a $`p`$ dimensional flat Minkowski space-time and $`I`$ is an interval with a finite length<sup>2</sup><sup>2</sup>2Conventions of indices are as follows. $`\mu ,\nu ,\mathrm{}`$ are eleven dimensional suffices and $`i,j,\mathrm{}`$ represent the spatial directions of the $`p`$-brane world-volume. Membrane world-volume indices are $`\alpha ,\beta ,\mathrm{}`$ and $`a,b`$ are world-volume spatial indices, $`a,b=1,2`$.. There exist at the boundaries of $`I`$ two $`p`$-branes on which an open membrane can end, and the $`p`$-branes wrap once around the $`S^1`$. $`\text{R}^{9p}\times I`$ is transverse to these $`p`$-branes. We drop the fermionic part, that is, restrict ourselves to considering a bosonic membrane.
In this case, the action of the membrane is
$$S=Td^3\xi \left\{\sqrt{deth_{\alpha \beta }}+\frac{1}{3!}ϵ^{\alpha \beta \gamma }C_{\mu \nu \rho }_\alpha X^\mu _\beta X^\nu _\gamma X^\rho \right\},$$
(1)
where $`\xi ^\alpha `$ are the world-volume coordinates $`(\tau ,\sigma _1,\sigma _2)`$ and $`h_{\alpha \beta }`$ is the induced metric on the world-volume, $`h_{\alpha \beta }_\alpha X^\mu _\beta X_\mu `$.
First, we fix the gauge freedom of world-volume reparametrization invariance with the static gauge,
$$\{\begin{array}{ccc}X^0=\tau \hfill & & \tau (\mathrm{},\mathrm{})\hfill \\ X^9=\sigma _1L\hfill & & \sigma _1[0,\pi ]\hfill \\ X^{10}=\sigma _2R\hfill & & \sigma _2[0,2\pi ),\hfill \end{array}$$
(2)
and the radius of the compactified direction $`X^{10}`$ is $`R`$,
$$X^{10}X^{10}+2\pi R.$$
(3)
We also compactify the $`X^9`$ direction on an interval. Suppose that there are two “fixed planes” placed at a distance of $`\pi L`$ in the $`X^9`$ direction. Here, $`\pi L`$ is the length of the interval, and the two boundaries of a membrane are bound to each of these “fixed planes”,
$$\mathrm{\Delta }X^9=\pi L.$$
(4)
These “fixed planes” are, for example, regarded as M5-branes in M-theory when $`p=5`$. Since the dimension of the $`p`$-brane is not essential in our analysis, we assume $`p=9`$ from now on.
Under the static gauge condition,
$`deth`$ $`=`$ $`\left|\begin{array}{ccc}1+(\dot{X}^i)^2& \dot{X}^i_1X^i& \dot{X}^i_2X^i\\ \dot{X}^i_1X^i& L^2+(_1X^i)^2& _1X^i_2X^i\\ \dot{X}^i_2X^i& _1X^i_2X^i& R^2+(_2X^i)^2\end{array}\right|`$ (8)
$`=`$ $`L^2R^2+L^2R^2(\dot{X}^i)^2R^2(_1X^i)^2L^2(_2X^i)^2+𝒪\left((X)^4\right),`$ (9)
and we get the first part of the action (Dirac part) as
$$S_\mathrm{D}=Td^3\xi \left[1+\frac{1}{2}(\dot{X}^i)^2\frac{1}{2}(_1X^i)^2\frac{1}{2}(_2X^i)^2+𝒪\left((X)^4\right)\right],$$
(10)
where we have made a rescaling, $`L\sigma _1\sigma _1,R\sigma _2\sigma _2`$.
Next, we go on to consider the $`C`$-field part,
$$S_C=_\mathrm{\Sigma }C_{[3]},$$
(11)
where $`\mathrm{\Sigma }`$ is the world-volume of a membrane. At the beginning, note that our action (1) is not gauge-invariant for an open membrane. So as to make an open membrane gauge-invariant, we introduce a two-form gauge field $`B`$ coupled to the boundaries of a membrane,
$$S_B=_\mathrm{\Sigma }B_{[2]},$$
(12)
which transforms as $`BB\mathrm{\Lambda }`$ under the $`C`$-field gauge transformation, $`CC+d\mathrm{\Lambda }`$, where $`\mathrm{\Lambda }`$ is a two-form field. Here, this $`B`$-field is on the boundary planes and has the field strength $`FdB`$ on these planes. Gauge-invariance requires that $`C`$ and $`F`$ always appear with the form of $`C+F`$, so the constant $`C`$-field leads to a constant field strength $`F`$ on the boundary planes. Then, we gauge away $`F`$ and only consider the effects of the $`C`$-field. Moreover, we suppose that the $`C`$-field is not only constant but also “magnetic”, that is, their non-zero components are only $`C_{ijk}`$. Finally, the $`C`$-field part of the action is
$$S_C=Td^3\xi C_{ijk}\dot{X}^i_1X^j_2X^k,$$
(13)
where we have made a rescaling $`C(LR)^1C`$.
A part of difficulties of membrane theory comes from its non-linearity of world-volume theory. Here, to avoid it, we take the limit $`\alpha \mathrm{}`$,
$`T`$ $`\alpha ^2T,`$
$`X`$ $`{\displaystyle \frac{1}{\alpha }}X,`$
$`C`$ $`\alpha C,`$
and also drop the constant term of the Dirac part. This limit means that the self-interactions of the world-volume theory are weak compared to the interactions with the background gauge fields. Finally, we get the effective action as follows,
$$S^{\mathrm{eff}}=Td^3\xi \left[\frac{1}{2}\left\{(\dot{X}^i)^2(_1X^i)^2(_2X^i)^2\right\}C_{ijk}\dot{X}^i_1X^j_2X^k\right],$$
(14)
where the ranges of the world-volume coordinates are
$`\sigma _1`$ $`[0,\pi L],`$ (15)
$`\sigma _2`$ $`[0,2\pi R),`$ (16)
and the area of the membrane is $`2\pi ^2LR`$.
To find the equations of motion and the boundary conditions, we vary the effective action (14),
$`\delta S^{\mathrm{eff}}`$ $`=T{\displaystyle d^3\xi \left[\ddot{X}^i(_1)^2X^i(_2)^2X^i\right]\delta X^i}`$
$`+T{\displaystyle d^3\xi _1\left[\left(_1X^iC_{ijk}\dot{X}^k_2X^j\right)\delta X^i\right]}.`$ (17)
$`\delta S^{\mathrm{eff}}=0`$ leads to the equations of motion,
$$\mathrm{}X^i=0,$$
(18)
where $`\mathrm{}\eta ^{\alpha \beta }_\alpha _\beta =_\tau ^2_1^2_2^2`$, and also leads to the boundary conditions,
$$_1X^iC_{ijk}\dot{X}^j_2X^k|_{\sigma _1=0,\pi L}=0.$$
(19)
The conjugate momenta are
$$P_i=\frac{\delta }{\delta \dot{X}^i}L=T\left(\dot{X}_iC_{ijk}_1X^j_2X^k\right),$$
(20)
so the Hamiltonian is
$`H`$ $`{\displaystyle d^2\sigma \left(\dot{X}^iP_i\right)}`$
$`={\displaystyle \frac{T}{2}}{\displaystyle d^2\sigma \left[\left(\frac{P^i}{T}+C_{ijk}_1X^j_2X^k\right)^2+(_1X^i)^2+(_2X^i)^2\right]}.`$ (21)
To follow the calculations in the string case, we regard the boundary conditions as primary constraints,
$$\varphi _1^i=_1X^iC_{ijk}\left(\frac{P^j}{T}+C_{jlm}_1X^l_2X^m\right)_2X^k|_{\sigma _1=0,\pi L}0.$$
(22)
Poisson brackets are ordinarily defined as
$$\begin{array}{c}\{X^i(\sigma _1,\sigma _2),P_j(\sigma _1^{},\sigma _2^{})\}=\delta _j^i\delta ^2(\sigma \sigma ^{}),\hfill \\ \{X^i,X^j\}=\{P_i,P_j\}=0.\hfill \end{array}$$
(23)
Using these, we get the equations of motion,
$$\dot{X}^i\{X^i(\sigma ),H\}=\frac{P^i}{T}+C_{ijk}_1X^j_2X^k,$$
(24)
and
$`\dot{P}^i\{P_i(\sigma ),H\}`$ $`=T\left\{\ddot{X}^iC_{ijk}\left(_1\dot{X}^j_2X^k+_1X^j_2\dot{X}^k\right)\right\}`$
$`=T[C_{ijk}(_2X^j_1({\displaystyle \frac{P^k}{T}}+C_{klm}_1X^l_2X^m)`$
$`_1X^j_2({\displaystyle \frac{P^k}{T}}+C_{klm}_1X^l_2X^m))+\mathrm{\Delta }X^i],`$ (25)
where Laplacian $`\mathrm{\Delta }`$ is defined as $`_1^2+_2^2`$ and dot means $`\tau `$ derivative.
For simplicity, we set $`T=1`$. We can recover $`T`$ by replacing $`P`$ with $`P/T`$.
## 3 Solving constraints
The method described in appendix A leads us to find the Dirac brackets of the membrane in the constant $`C`$-field. First, we consider the consistency conditions of the constraints
$$\dot{\varphi }\{\varphi ,H_\mathrm{T}\}0,$$
(26)
and find an infinite chain of secondary constraints as follows
$`\varphi _2^i`$ $`\dot{\varphi }_1^i=\{\varphi _1^i,H\}`$
$`=_1\dot{X}^iC_{ijk}\ddot{X}^j_2X^kC_{ijk}\dot{X}^j_2\dot{X}^k,`$
$`\varphi _3^i`$ $`\dot{\varphi }_2^i`$
$`=_1\ddot{X}^iC_{ijk}\left[X^{(3)j}_2X^k+2\ddot{X}^j_2\dot{X}^k+\dot{X}^j_2\ddot{X}^k\right],`$
$`\mathrm{}`$
$`\varphi _{n+1}^i`$ $`\varphi _1^{(n)i}`$
$`=_1X_1^{(n)i}C_{ijk}{\displaystyle \underset{\mathrm{}=0}{\overset{\mathrm{}}{}}}\left(\begin{array}{c}n\\ \mathrm{}\end{array}\right)X^{(n+1\mathrm{})j}_2X^{(\mathrm{})k},`$ (27)
where
$$\varphi ^{(n)i}\frac{^n}{\tau ^n}\varphi ^i.$$
(28)
Note that the equation of motion (24) tells that each secondary constraint has at most $`C^3`$, and all the constraints are second class. Explicit computations show that the first few constraints are given by
$`\varphi _1^i`$ $`=_1X^iC_{ijk}\left(P^j+C_{jlm}_1X^j_2X^k\right)_2X^k|_{\sigma _1=0,\pi L}0,`$ (29)
$`\varphi _2^i`$ $`=_1P^i`$
$`+C_{ijk}\left[_1X^j_1_2X^k_2^2X^j_2X^kP^j_2P^k\right]`$
$`+C_{ijk}C_{jlm}\left[_2P^k_1X^l_2X^m+P^k_2(_1X^l_2X^m)\right]`$
$`C_{ijk}C_{jlm}C_{kop}[_1X^l_2X^m_2(_1X^o_2X^p)]|_{\sigma _1=0,\pi L}0,`$ (30)
$`\varphi _3^i`$ $`=_1\mathrm{\Delta }X^k`$
$`+C_{ijk}[\mathrm{\Delta }P^j_2X^k+2_2P^j\mathrm{\Delta }X^kP^j_2\mathrm{\Delta }X^k]`$
$`+C_{ijk}C_{jlm}[2\mathrm{\Delta }X^k_2(_1X^l_2X^m)`$
$`_2X^k\mathrm{\Delta }(_1X^l_2X^m)_2\mathrm{\Delta }X^k(_1X^l_2X^m)\left]\right|_{\sigma _1=0,\pi L}0.`$ (31)
These constraints look too hard to solve completely unlike the string case. Thus, we shall take an approximation to solve them.
At this stage, we take the limit $`L0`$ and $`R\mathrm{}`$ <sup>3</sup><sup>3</sup>3Note that this limit is a tensionless string limit in Strominger’s sense .. This leads to simplification as follows. For $`\sigma _1`$, we suppose that no oscillations are excited. Hence, after solving the constraints, $`X^i(\tau ,\sigma _1,\sigma _2)`$ and $`P^i(\tau ,\sigma _1,\sigma _2)`$ are determined by their boundary values. And for $`\sigma _2`$, we neglect terms which is of order $`(1/R)^3`$ or higher, which means that we drop the terms involving three derivatives of $`\sigma _2`$ or higher,
$$_2^3X^i=0,_2^2X^i_2X^j=0\text{etc}\mathrm{}.$$
(32)
To solve the constraints, we shall include the effects of the $`C`$-field order by order. At order of $`C^0`$, the boundary conditions are
$$_1X^i|_{\sigma _1=0,\pi L}=0\text{and}X^i(\tau ,\sigma _1,\sigma _2)=X^i(\tau ,\sigma _1,\sigma _2+2\pi R).$$
(33)
Since no oscillations of $`\sigma _1`$ are excited under the $`L0`$ limit, the solution is
$$X^i(\tau ,\sigma _1,\sigma _2)=x_0^{(0)i}(\tau ,\sigma _2).$$
(34)
where the subscript $`0`$ of $`x_0^{(0)i}`$ means we are considering only the zero-mode of $`\sigma _1`$. Since the $`C`$-field background changes the $`\sigma _1`$ boundary conditions, the $`\sigma _1`$ dependence of fields $`X`$ and $`P`$ would be altered:
$$X^i=x_0^{(0)i}(\tau ,\sigma _2)+(\text{corrections which depend also on }\sigma _1\text{ and }C).$$
(35)
Let us calculate the corrections to second order in $`C`$. Consider the expansions of $`X`$ and $`P`$ in terms of $`C`$
$`X_0^i(\tau ,\sigma _1,\sigma _2)=x_0^{(0)i}+x_0^{(1)i}+x_0^{(2)i},`$ (36)
$`P_0^i(\tau ,\sigma _1,\sigma _2)=p_0^{(0)i}+p_0^{(1)i}+p_0^{(2)i},`$ (37)
where $`x_0^{(0)}`$ and $`p_0^{(0)}`$ are functions of $`\tau `$ and $`\sigma _2`$, independent of $`\sigma _1`$ and unconstrained. We substitute them into the constraints (29) and (3). Of order $`C^1`$, we get
$`\varphi _1^i`$ $`=_1x_0^{(1)i}C_{ijk}p_0^{(0)j}_2x_0^{(0)k}|_{\sigma _1=0,\pi L}0,`$
$`\varphi _2^i`$ $`=_1p_0^{(1)i}+C_{ijk}\left(p_0^{(0)j}_2p_0^{(0)k}\right)|_{\sigma _1=0,\pi L}0,`$ (38)
and find solutions at this order as follows
$`x_0^{(1)i}(\tau ,\sigma _1,\sigma _2)=A_0^{(1)i}(\tau ,\sigma _2)+C_{ijk}p_0^{(0)j}_2x_0^{(0)k}\sigma _1,`$ (39)
$`p_0^{(1)i}(\tau ,\sigma _1,\sigma _2)=B_0^{(1)i}(\tau ,\sigma _2)+C_{ijk}p_0^{(0)j}_2p_0^{(0)k}\sigma _1,`$ (40)
where $`A_0`$ and $`B_0`$ in the right hand sides are unconstrained. In succession, the equations of order $`C^2`$ are
$`\varphi _1^i`$ $`=_1x_0^{(2)i}C_{ijk}\left[p_0^{(1)j}_2x_0^{(0)k}+p_0^{(0)j}_2x_0^{(1)k}\right]`$
$`C_{ijk}C_{jlm}\left(p_0^{(0)l}_2p_0^{(0)m}_2x_0^{(0)k}p_0^{(0)l}_2x_0^{(0)m}p_0^{(0)k}\right)\sigma _1,`$ (41)
$`\varphi _2^i`$ $`=_1p_0^{(2)i}C_{ijk}\left[p_0^{(1)j}_2p_0^{(0)k}+p_0^{(0)j}_2p_0^{(1)k}\right]`$
$`C_{ijk}C_{jlm}\left(p_0^{(0)l}_2p_0^{(0)m}_2p_0^{(0)k}p_0^{(0)l}_2p_0^{(0)m}p_0^{(0)k}\right)\sigma _1,`$ (42)
and we find the solutions,
$`x_0^{(2)i}(\tau ,\sigma _1,\sigma _2)`$ $`=A_0^{(2)i}(\tau ,\sigma _2)+C_{ijk}\left[B_0^{(1)j}_2x_0^{(0)k}+p_0^{(0)j}_2A_0^{(1)k}\right]\sigma _1`$
$`+C_{ijk}C_{jlm}\left(p_0^{(0)l}_2p_0^{(0)m}_2x_0^{(0)k}p_0^{(0)l}_2x_0^{(0)m}p_0^{(0)k}\right){\displaystyle \frac{\sigma _1^2}{2}},`$
$`p_0^{(2)i}(\tau ,\sigma _1,\sigma _2)`$ $`=B_0^{(2)i}(\tau ,\sigma _2)+C_{ijk}\left[B_0^{(1)j}_2p_0^{(0)k}+p_0^{(0)j}_2B_0^{(1)k}\right]\sigma _1`$
$`+C_{ijk}C_{jlm}\left(p_0^{(0)l}_2p_0^{(0)m}_2p_0^{(0)k}p_0^{(0)l}_2p_0^{(0)m}p_0^{(0)k}\right){\displaystyle \frac{\sigma _1^2}{2}}.`$ (43)
Putting them together, we find that the $`X^i(\tau ,\sigma _1,\sigma _2)`$ and $`P^i(\tau ,\sigma _1,\sigma _2)`$ are determined by the unconstrained boundary values, $`X_0(\tau ,\sigma _2)=x_0^{(0)}+A_0^{(1)}+A_0^{(2)}`$ and $`P_0(\tau ,\sigma _2)=p_0^{(0)}+B_0^{(1)}+B_0^{(2)}`$ as follows,
$`X^i(\tau ,\sigma _1,\sigma _2)=`$ $`X_0^i+\sigma _1C_{ijk}P_0^j_2X_0^k`$
$`+{\displaystyle \frac{\sigma _1^2}{2}}C_{ijk}C_{jlm}\left[_2X_0^kP_0^l_2P_0^mP_0^k_2(P_0^l_2X_0^m)\right],`$ (44)
$`P^i(\tau ,\sigma _1,\sigma _2)=`$ $`P_0^i+\sigma _1C_{ijk}P_0^j_2P_0^k`$
$`+{\displaystyle \frac{\sigma _1^2}{2}}C_{ijk}C_{jlm}\left[_2P_0^kP_0^l_2P_0^mP_0^k_2(P_0^l_2P_0^m)\right].`$ (45)
One can confirm that these solutions satisfy the remaining constraints by substituting (3) and (3) into the explicit form of $`\varphi _3^i`$ and taking into account the fact that the other higher constraints involve only higher derivative terms of $`\sigma _1`$ and $`\sigma _2`$. Since we get the solutions of the constraints, we can compute the Dirac brackets of $`X`$ and $`P`$ by the method given in appendix A. This is what we shall do in the following section.
## 4 Computing the Dirac brackets
In order to compute the Dirac brackets, we first calculate Lagrange brackets. In this case, Lagrange bracket L is defined as
$`\mathrm{\Omega }=`$ $`2{\displaystyle d^2\sigma 𝑑X^i(\sigma _1,\sigma _2)}dP^i(\sigma _1,\sigma _2)`$
$`=`$ $`{\displaystyle 𝑑x𝑑y𝐋_{xy}^{ij}𝑑\varphi ^i(x)}d\varphi ^j(y),`$ (46)
where we have integrated over $`\sigma _1`$, $`d\varphi =dX_0(\sigma _2)`$ or $`dP_0(\sigma _2)`$, and $`x`$ and $`y`$ denote the $`\sigma _2`$ coordinate. Dirac bracket C is determined by the inverse matrix of this Lagrange brackets, $`\text{C}=\text{L}^1`$. To calculate the Lagrange bracket of this system, we determine the effects of the $`C`$-field order by order, to order $`C^2`$:
$$𝐋=𝐋^{(0)}+𝐋^{(1)}+𝐋^{(2)},$$
(47)
where $`L^{(i)}`$ denotes the terms of order $`C^i`$. Then the Dirac bracket is obtained as
$`𝐂=\text{L}^1=`$ $`𝐋^{(0)1}𝐋^{(0)1}(𝐋^{(1)}+𝐋^{(2)})𝐋^{(0)1}+𝐋^{(0)1}𝐋^{(1)}𝐋^{(0)1}𝐋^{(1)}𝐋^{(0)1}+𝒪(C^3)`$ (48)
$`=`$ $`𝐉𝐉(𝐋^{(1)}+𝐋^{(2)})𝐉+\mathrm{𝐉𝐋}^{(1)}\mathrm{𝐉𝐋}^{(1)}𝐉+𝒪(C^3),`$ (49)
where we have abbreviated $`𝐋^{(0)1}`$ as $`𝐉`$.
Let us start the calculation. In zeroth order in $`C`$, the Lagrange bracket is determined through the symplectic form
$`\mathrm{\Omega }^{[0]}=`$ $`2{\displaystyle 𝑑\sigma ^2𝑑X_0^i}dP_0^i`$
$`=`$ $`2\pi L{\displaystyle 𝑑x𝑑y\delta ^{ij}\delta (xy)𝑑X_0^i(x)}dP_0^j(y).`$ (50)
We get
$$𝐋^{(0)}=\left(\begin{array}{cc}0& L^{(0)}\\ (L^{(0)})^\mathrm{T}& 0\end{array}\right),$$
(51)
where
$$L^{(0)}=\pi L\delta ^{ij}\delta (xy).$$
(52)
The inverse matrix of this $`𝐋^{(0)}`$ is given by
$$𝐉=(𝐋^{(0)})^1=\left(\begin{array}{cc}& J\\ J& \end{array}\right),J=(L^{(0)})^1=\frac{1}{\pi L}\delta ^{ij}\delta (xy),J^\mathrm{T}=J.$$
(53)
At this stage, we can calculate the Dirac bracket at $`C=0`$:
$`\{X_0^i(x),X_0^j(y)\}_{\mathrm{DB}}=`$ $`0,`$ (54)
$`\{P_0^i(x),P_0^j(y)\}_{\mathrm{DB}}=`$ $`0,`$ (55)
$`\{X_0^i(x),P_0^j(y)\}_{\mathrm{DB}}=`$ $`{\displaystyle \frac{1}{\pi L}}\delta ^{ij}\delta (xy).`$ (56)
These are the original Poisson brackets except for the normalization factor.
### Calculations of $`𝒪(C^1)`$
Next, we shall calculate the $`C^1`$ part. This is the first non-trivial result in these calculations. The symplectic form of this order is
$`\mathrm{\Omega }^{[1]}=`$ $`2{\displaystyle }d^2\sigma [\sigma _1C_{ikl}dX_0^i(dP_0^k_2P_0^l+P_0^k_2dP_0^l)`$
$`+\sigma _1C_{ikl}(dP_0^k_2X_0^l+P_0^k_2dX_0^l)dP_0^i]`$
$`=`$ $`(\pi L)^2{\displaystyle }dxdyC_{ijl}[dX_0^i(x)dP_0^j(y)(2C_{ijl}P_0^l(x)_x\delta (xy))`$
$`dP_0^i(x)dP_0^j(y)_xX_0^l\delta (xy)],`$ (57)
and we get
$$𝐋^{(1)}=\left(\begin{array}{cc}0& L^{(1)}\\ (L^{(1)})^\mathrm{T}& l^{(1)}\end{array}\right),$$
(58)
where
$`L^{(1)}=`$ $`(\pi L)^2C_{ijl}P_0^l(x)_x\delta (xy),`$ (59)
$`l^{(1)}=`$ $`(\pi L)^2C_{ijl}_xX_0^l\delta (xy).`$ (60)
At this order, the Dirac bracket is
$`\{X_0^i(x),X_0^j(y)\}_{\mathrm{DB}}=`$ $`C_{ijl}_xX_0^l\delta (xy),`$ (61)
$`\{P_0^i(x),P_0^j(y)\}_{\mathrm{DB}}=`$ $`0,`$ (62)
$`\{X_0^i(x),P_0^j(y)\}_{\mathrm{DB}}=`$ $`{\displaystyle \frac{1}{\pi L}}\delta ^{ij}\delta (xy)C_{ijl}P_0^l(y)\delta ^{}(yx).`$ (63)
One can check that the Jacobi identity holds at this order,
$`\{\{X_0^i(x),P_0^j(y)\},X_0^k(z)\}+\text{(cyclic.)}`$
$`={\displaystyle \frac{1}{\pi L}}C_{ijk}\left(\delta (yz)\delta ^{}(yx)+\delta (yx)\delta ^{}(yz)+\delta (zx)\delta ^{}(zy)\right)`$
$`={\displaystyle \frac{1}{\pi L}}C_{ijk}\left(\delta (yz)\delta ^{}(yx)+\delta (yx)\delta ^{}(yz)+\delta (yx)\delta ^{}(zy)\delta ^{}(zx)\delta (zy)\right)`$
$`=0.`$ (64)
The Jacobi identity for $`\{X,\{X,X\}\}`$ is trivially satisfied at first order in $`C`$. To see how it is non-trivially satisfied, we turn to the calculations of $`C^2`$.
### Calculations of $`𝒪(C^2)`$
The calculations of order $`C^2`$ turn out to be very complicated, so we split the calculations into some parts.
First, we consider the cross terms, $`(C^1\text{ part})(C^1\text{ part})`$ . The symplectic form of this part is
$`\mathrm{\Omega }^{[21]}=`$ $`2{\displaystyle d^2\sigma \sigma _1^2C_{ijk}C_{ilm}(dP_0^j_2X_0^k+P_0^j_2dX_0^k)}(dP_0^l_2P_0^m+P_0^l_2dP_0^m)`$
$`=`$ $`{\displaystyle \frac{2(\pi L)^3}{3}}{\displaystyle d^2\sigma C_{ikl}C_{jml}}`$
$`\times \{dX_0^i(x)dP_0^j(y)_x\left(P_0^k(x)(2_yP_0^m(y)+P_0^m(y)_y)\delta (xy)\right)`$
$`+dP_0^i(x)dP_0^j(y)[{\displaystyle \frac{1}{2}}(X_0^k(x)P_0^m(x)P_0^k(y)X_0^m(y))\delta (xy)`$
$`(X_0^k(x)P_0^m(x)+P_0^k(y)X_0^m(y))\delta ^{}(xy)]\},`$ (65)
so we get
$$𝐋^{[21]}=\left(\begin{array}{cc}& L^{[21]}\\ (L^{[21]})^\mathrm{T}& l^{[21]}\end{array}\right),$$
(66)
where
$`L^{[21]}`$ $`={\displaystyle \frac{(\pi L)^3}{3}}C_{ikl}C_{jml}_x\left(P_0^k(x)\left(2_yP_0^m(y)+P_0^m(y)_y\right)\delta (xy)\right),`$ (67)
$`l^{[21]}`$ $`={\displaystyle \frac{1}{3}}(\pi L)^3C_{ikl}C_{jml}((X_0^k(x)P_0^m(x)P_0^k(y)X_0^m(y))\delta (xy)`$
$`(X_0^k(x)P_0^m(x)+P_0^k(y)X_0^m(y))\delta ^{}(xy)).`$ (68)
Next, we consider the $`(C^0\text{ part})(C^2\text{ part})`$. The symplectic form of this part is
$`\mathrm{\Omega }^{[22]}`$ $`=2{\displaystyle d^2\sigma \sigma _1^2C_{ijk}C_{ilm}}`$
$`\times \{[_2dX_0^kP_0^l_2P_0^m+_2X_0^kdP_0^l_2P_0^m_2X_0^kP_0^m_2dP_0^l`$
$`dP_0^k_2(P_0^l_2X_0^m)P_0^k_2(dP_0^l_2X_0^mP_0^m_2dX_0^l)]dP_0^i`$
$`+dX_0^i[_2dP_0^kP_0^l_2P_0^m_2P_0^kP_0^m_2dP_0^l+_2P_0^kdP_0^l_2P_0^m`$
$`dP_0^k_2(P_0^l_2P_0^m)P_0^k_2(dP_0^l_2P_0^mP_0^m_2dP_0^l)]\}.`$ (69)
Then we find that the $`\mathrm{\Omega }^{[22]}`$ has the form
$$\mathrm{\Omega }^{[22]}=𝑑x𝑑y𝐋^{[22]}𝑑\varphi ^i(x)d\varphi ^j(y),$$
(70)
where
$$𝐋^{[22]}=𝐌+𝐍,$$
(71)
$$(𝐌)_{xy}^{ij}=\left(\begin{array}{cc}& M\\ M^\mathrm{T}& m\end{array}\right),$$
(72)
$$(𝐍)_{xy}^{ij}=\left(\begin{array}{cc}& N\\ N^\mathrm{T}& n\end{array}\right),$$
(73)
and, $`𝐌`$ and $`𝐍`$ correspond to the following tensor structures of $`C^2`$:
$`𝐌`$ $``$ $`C_{ijk}C_{klm},`$
$`𝐍`$ $``$ $`C_{ikl}C_{jml}.`$
The explicit calculations of M and N are shown in appendix B. The results are
$`M`$ $`={\displaystyle \frac{(\pi L)^3}{3}}C_{ijk}C_{klm}\left[P_0^l(x)_xP_0^m(x)\delta ^{}(xy)\right],`$ (74)
$`m`$ $`={\displaystyle \frac{(\pi L)^3}{3}}C_{ijk}C_{klm}_y\left(P_0^l(y)_yX_0^m(y)\right)\delta (xy),`$ (75)
$`N`$ $`={\displaystyle \frac{(\pi L)^3}{3}}C_{ikl}C_{jml}\left[P_0^k(x)P_0^m(x)\delta ^{\prime \prime }(xy)+P_0^m(x)_x\left(P_0^k(x)\delta (xy)\right)\right],`$ (76)
$`n`$ $`={\displaystyle \frac{(\pi L)^3}{3}}C_{ikl}C_{jml}\left[X_0^k(x)P_0^m(x)+X_0^m(y)P_0^k(y)\right]\delta ^{}(xy).`$ (77)
Thus we get the Lagrange brackets to order $`C^2`$. Let us compute the Dirac brackets.
### Computing the Dirac brackets
By (48), we can calculate the Dirac brackets C,
$$\text{C}_{xy}^{ij}=\left(\begin{array}{cc}\{X_0^i(x),X_0^j(y)\}_{\mathrm{DB}}& \{X_0^i(x),P_0^j(y)\}_{\mathrm{DB}}\\ \{P_0^j(y),X_0^i(x)\}_{\mathrm{DB}}& \{P_0^i(x),P_0^j(y)\}_{\mathrm{DB}}\end{array}\right),$$
(78)
as follows:
$`\{X_0^i(x),X_0^j(y)\}_{\mathrm{DB}}`$ $`=J\left(l^{(1)}\right)J+J\left(l^{(2)}\right)JJl^{\left(1\right)}JL^{\left(1\right)}JJ(L^{\left(1\right)})^\mathrm{T}Jl^{\left(1\right)}J`$
$`={\displaystyle \frac{1}{(\pi L)^2}}\left(l^{(1)}\right)_{xy}^{ij}+{\displaystyle \frac{1}{(\pi L)^2}}\left(l^{(2)}\right)_{xy}^{ij}`$
$`+{\displaystyle \frac{1}{(\pi L)^3}}\left\{(l^{\left(1\right)})_{xz}^{il}(L^{\left(1\right)})_{zy}^{lj}+\left((L^{\left(1\right)})^\mathrm{T}\right)_{xz}^{il}(l^{\left(1\right)})_{zy}^{lj}\right\}+𝒪(C^3),`$ (79)
$`\{X_0^i(x),P_0^j(y)\}_{\mathrm{DB}}`$ $`=J+J\left((L^{(1)})^\mathrm{T}\right)J+J\left((L^{(2)})^\mathrm{T}\right)JJ(L^{\left(1\right)})^\mathrm{T}J(L^{\left(1\right)})^\mathrm{T}J`$
$`={\displaystyle \frac{1}{\pi L}}(\mathrm{𝟏})_{xy}^{ij}+{\displaystyle \frac{1}{(\pi L)^2}}\left((L^{(1)})^\mathrm{T}\right)_{xy}^{ij}`$
$`+{\displaystyle \frac{1}{(\pi L)^2}}\left((L^{(2)})^\mathrm{T}\right)_{xy}^{ij}+{\displaystyle \frac{1}{(\pi L)^3}}\left((L^{\left(1\right)})^\mathrm{T}\right)_{xz}^{il}\left((L^{\left(1\right)})^\mathrm{T}\right)_{zy}^{lj}+𝒪(C^3),`$ (80)
$`\{P_0^i(x),P_0^j(y)\}_{\mathrm{DB}}`$ $`=0.`$ (81)
Explicit computation shows
$`\{X_0^i(x),X_0^j(y)\}_{\mathrm{DB}}`$ $`=C_{ijl}X_0^l(x)\delta (xy)`$
$`{\displaystyle \frac{1}{3}}C_{ikl}C_{jml}[(X_0^k(x)P_0^m(x)X_0^m(y)P_0^k(y))\delta (xy)`$
$`+(X_0^k(x)P_0^m(x)+X_0^m(y)P_0^k(y))\delta ^{}(xy)]`$
$`+{\displaystyle \frac{1}{3}}C_{ijk}C_{klm}_y\left(P_0^l(y)X_0^m(y)\right)\delta (xy)+𝒪(C^3),`$ (82)
$`\{X_0^i(x),P_0^j(y)\}_{\mathrm{DB}}`$ $`=\delta ^{ij}\delta (xy)+C_{ijl}P_0^l(y)\delta ^{}(xy)`$
$`{\displaystyle \frac{1}{3}}C_{ikl}C_{jml}[P_0^k(x)P_0^m(x)\delta ^{\prime \prime }(xy)+3P_0^k(x)P_0^m(x)\delta ^{}(xy)`$
$`+(2P_0^k(x)P_0^{\prime \prime m}(x)+P_0^k(x)P_0^m(x))\delta (xy)]`$
$`+{\displaystyle \frac{1}{3}}C_{ijl}C_{lkm}P_0^k(y)P_0^m(y)\delta ^{}(xy)+𝒪(C^3),`$ (83)
where we have rescaled the momenta, $`\pi LP_0^iP_0^i`$. This is because in the limit $`L0`$, the integrated momenta $`\pi LP_0`$ are more naturally assigned to the boundary strings than the original boundary momenta $`P_0`$.
These results mean that the coordinates of the boundary strings of an open membrane in the constant $`C`$-field background show noncommutativity. It is very curious that the commutation relation between $`X^i`$ and $`X^j`$ depends on other components of transverse fields, $`X^k`$.
## 5 Concluding remarks
In the previous section, we have obtained the Dirac brackets of an open membrane in the $`C`$-field background. The result shows that the boundary string has a loop-space noncommutativity.
We can confirm that the Jacobi identity holds at order in $`C^2`$ with these results, though we do not write down the calculation explicitly. Indeed, the satisfaction of Jacobi identity is trivial from the general properties of Poisson bracket, but the cancellations between the terms are not trivial. This indicates the algebra has complicated structures and more transparent understanding of it from the boundary string viewpoint is desirable.
The results presented above are the Dirac brackets between the coordinates and momenta of the boundary. Dirac brackets between the coordinates on the membrane can be calculated by (3), and there exists noncommutativity not only at the boundary but also on the membrane. In string theory, the string coordinates are commutative except at its ends as explained in appendix A, and to show this it is essential to include all the oscillation modes. Thus we also expect that including all the oscillation modes make the membrane coordinates commutative except at its boundary, because the $`C`$-field part of the action (13) is total derivative for a constant $`C`$, and should change the dynamics only at the boundary.
We have done our analysis in a tractable static gauge condition. Light-cone gauge analysis is more interesting in its relationship with BFSS matrix theory and the results of . It is easy to find the light-cone gauge Hamiltonian,
$$H_{\text{LC}}=d^2\sigma \frac{1}{2P^+}\left[(P^i+C_{ijk}_1X^j_2X^k)^2+\frac{T^2}{2}\{X^i,X^j\}^2\right],$$
(84)
the equations of motion
$$\ddot{X}^i+\{X^j,\{X^i,X^j\}\}=0,$$
(85)
and the boundary conditions
$$T_2X^j\{X^i,X^j\}+C_{ijk}_2X^j\dot{X}^k|_{\sigma _1=0,\pi }=0.$$
(86)
However, the chain of the boundary constraints look too complicated to solve in this case even if some approximations are taken. Moreover, when there is a constant $`C`$-field background, we can not apply the matrix regularization method developed in the third paper of . Thus, analysis in this gauge is remaining as a hard but interesting problem. See comments below.
When this work was in the process of typing, we learned that another group has also employed the quantization of an open membrane in a $`C`$-field background, and they have also investigated the decoupling limit as the open string case. Though their line of thought is different from ours, their results seem to be consistent with ours at least in first order in $`C`$. Moreover, their paper has also studied the light-cone coordinate analysis, but their analysis is within the decoupling limit and slightly different from our interests such as membrane regularization related to matrix models.
## Acknowledgments
S. K. is supported in part by the Japan Society for the Promotion of Science under the Predoctoral Research Program. N. S. is supported in part by Grant-in-Aid for Scientific Research (#12740150), and in part by Priority Area: “Supersymmetry and Unified Theory of Elementary Particles” (#707), from Ministry of Education, Science, Sports and Culture.
## Appendix A A brief review of Dirac’s procedure applied to boundary constraints
In string theory, one can find noncommutativity on a D-brane by quantization procedures for open strings with a background $`B`$-field . A transparent way to confirm the noncommutativity of open strings is the Dirac’s procedure applied to boundary conditions . In this appendix, we briefly review this approach. The calculations described here are mainly based on the appendix of the paper by Kawano and Takahashi .
### Dirac’s procedure
First, we survey the ordinary methods for constrained systems following . In singular systems, we face some constraints, primary constraints, between canonical variables. Consistency conditions for these constraints in time evolution sometimes lead to additional constraints, secondary constraints. We must consider the consistency conditions for these new constraints and possibly find new constraints, secondary constraints for secondary constraints, and so on.
Constraints are classified into two classes; the first class constraints that commute with all the other constraints and the second class constraints that do not. The first class constraints are related to the gauge symmetry of the system and we can treat them as second class by gauge fixing. Thus we may assume all the constraints are second class. The singular system is treated with the Dirac bracket defined as
$$\{F,G\}_{\mathrm{DB}}\{F,G\}\{F,\varphi _A\}C^{AB}\{\varphi _B,G\},$$
(87)
where $`C^{AB}=(C^1)^{AB}`$, $`C_{AB}\{\varphi _A,\varphi _B\}`$ and $`\varphi _A`$, $`\varphi _B`$ are second class constraints. This Dirac brackets are Poisson brackets on the constrained surface, so we can determine the time evolution of this constrained system using the Dirac bracket.
### Boundary condition as constraint
According to , we can treat the boundary conditions of an open string as constraints. The consistency conditions of these constraints lead to an infinite chain of secondary constraints, which are all second class. Thus, we can calculate Dirac brackets of this system in principle. However, we must consider the inverse of an $`\mathrm{}\times \mathrm{}`$ matrix $`C_{AB}`$. Surprisingly, we can completely solve this question in the string case.
Let us explain the string case calculations for example. We consider an open string in a constant NS-NS $`B`$-field background. The action of this system is
$$S=\frac{1}{4\pi \alpha ^{}}_\mathrm{\Sigma }d^2\sigma \left[g_{ij}\left(\dot{X}^i\dot{X}^jX^iX^j\right)+2b_{ij}\dot{X}^iX^j\right],$$
(88)
where
$$X^{}\frac{}{\sigma }X,\dot{X}\frac{}{\tau }X,$$
(89)
and $`b_{ij}=2\pi \alpha ^{}B_{ij}`$. Variation of the action leads to the equations of motion and the boundary conditions:
$$^\alpha _\alpha X^i(\tau ,\sigma )=0,$$
(90)
Dirichlet directions: $`\delta X^{i_\mathrm{D}}=0(X^{i_\mathrm{D}}=\text{const.}),`$
Neumann (or Mixed) directions: $`g_{ij}X^j+b_{ij}\dot{X}^j=0\text{at}\sigma =0,\pi ,`$
where mixed directions are named for their mixtures of some directions and we shall only consider below the directions obeying these mixed boundary conditions. We now go on to the canonical formalism. Conjugate momenta are $`2\pi \alpha ^{}P_i(\tau ,\sigma )=\left(g_{ij}\dot{X}^j+b_{ij}X^j\right)`$ and the boundary conditions are taken to be primary constraints of this system,
$$\varphi _i(\sigma )=G_{ij}X^j+2\pi \alpha ^{}b_{ik}g^{kl}P_l,$$
(91)
where $`G_{ij}g_{ij}(bg^1b)_{ij}`$, so called “open string metric”.
The consistency of the constraints in time evolution leads to an infinite chain of secondary constraints:
$$\frac{^{(2n+1)}}{\sigma ^{(2n+1)}}P_i(\sigma )0\text{and}\frac{d^{(2n)}}{d\sigma ^{(2n)}}\varphi _i(\sigma )0.$$
(92)
The solution to these constraints is
$`X^i(\tau ,\sigma )`$ $`={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}X_n^i(\tau )\mathrm{cos}(n\sigma )+\mathrm{\Theta }^{ij}\left[P_{0j}(\tau )\sigma +{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n}}P_{nj}\mathrm{sin}(n\sigma )\right],`$ (93)
$`P_i(\tau ,\sigma )`$ $`={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}P_{ni}(\tau )\mathrm{cos}(n\sigma ).`$ (94)
### Lagrange bracket
One of the easiest way to find the Dirac bracket is to use the Lagrange brackets and this method was used in in a slightly different way.
Lagrange bracket L for variables $`z^\mu =z^\mu (q,p)`$ is defined through the symplectic form
$$\mathrm{\Omega }=2dq^i(z)dp_i(z)=𝐋^{\mu \nu }dz^\mu dz^\nu ,$$
(95)
where $`q`$ and $`p`$ are canonical variables of this system. Explicitly, Lagrange bracket is written as
$$𝐋^{\mu \nu }=\frac{q^i}{z^\mu }\frac{p_i}{z^\nu }\frac{q^i}{z^\nu }\frac{p_i}{z^\mu }.$$
(96)
An important property of this bracket is that this is the inverse matrix of the Poisson bracket,
$$𝐋_{\mu \nu }\{z^\nu ,z^\rho \}=\delta _\mu ^\rho .$$
(97)
To find the relation to the Dirac bracket, let us take the variables as follows
$$\underset{\text{coodinates on the constrained surface}}{\underset{}{z^1,z^2,\mathrm{},z^{2N2m}}},\underset{\text{2m constraints}}{\underset{}{z^{2N2m+1}=\varphi _1,\mathrm{},z^{2N}=\varphi _{2m}}}.$$
(98)
Then we find that the matrix obtained by limiting variables to the first $`(2N2m)`$ ones is the inverse matrix of the Dirac bracket,
$$\underset{\mu ,\nu =1}{\overset{2N2m}{}}𝐋_{\mu \nu }\{z^\nu ,z^\rho \}_{\mathrm{DB}}=\delta _\mu ^\rho .$$
(99)
This means that Dirac bracket is the Poisson bracket on the constrained surface defined through the conditions, $`z^{2N2m+1}=\mathrm{}=z^{2N}=0`$. Thus, we can compute the Dirac bracket by solving the constraints, constructing the Lagrange bracket and taking its inverse.
In string case, Lagrange brackets are defined by
$`\mathrm{\Omega }`$ $`=2{\displaystyle 𝑑\sigma 𝑑X^i(\sigma )}dP_i(\sigma )`$
$`=2\left[\pi dX_0^idP_{0i}+{\displaystyle \frac{\pi }{2}}dX_n^idP_{ni}\mathrm{\Theta }^{ij}{\displaystyle \frac{\pi ^2}{2}}dP_{0i}dP_{0j}\right].`$ (100)
From this, we can determine the Lagrange brackets for every mode of $`X`$ and $`P`$. Taking the inverse, we obtain
$`\{X^i(\sigma ),P_j(\sigma ^{})\}_{\mathrm{DB}}`$ $`=\delta _j^i\left({\displaystyle \frac{1}{\pi }}+{\displaystyle \frac{2}{\pi }}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\mathrm{cos}(n\sigma )\mathrm{cos}(n\sigma ^{})\right)`$
$`\delta _j^i\stackrel{~}{\delta }(\sigma ,\sigma ^{})`$ (101)
$`\{P_i(\sigma ),P_j(\sigma ^{})\}_{\mathrm{DB}}`$ $`=0`$ (102)
$`\{X^i(\sigma ),X^j(\sigma ^{})\}_{\mathrm{DB}}`$ $`=\{\begin{array}{cc}\mathrm{\Theta }^{ij}\hfill & (\sigma =\sigma ^{}=0)\hfill \\ \mathrm{\Theta }^{ij}\hfill & (\sigma =\sigma ^{}=\pi )\hfill \\ 0\hfill & (\text{otherwise})\hfill \end{array}.`$ (106)
This shows noncommutativity of open strings and this equals the result in .
## Appendix B The explicit calculations of Lagrange brackets at second order in $`C`$
In this appendix, we give the explicit calculations of (74), (75), (76) and (77).
First, we calculate the part of $`M`$. This part of the symplectic form is
$`\mathrm{\Omega }_M^{[22]}`$ $`={\displaystyle \frac{(\pi L)^3}{3}}{\displaystyle 𝑑x𝑑yC_{ijk}C_{jlm}}`$
$`\times [dX_0^i(x)dP_0^k(y)[_y\left(P_0^l(y)P_0^m(y)\delta (xy)\right)_y(P_0^l(y)P_0^m(y))\delta (xy)]`$
$`+dX_0^k(x)dP_0^i(y)[_x\left(P_0^l(x)P_0^m(x)\delta (xy)\right)]]`$
$`={\displaystyle \frac{2(\pi L)^3}{3}}{\displaystyle 𝑑x𝑑yC_{ijk}C_{jlm}𝑑X_0^i(x)}dP_0^j(y)\left[P_0^l(x)P_0^m(x)\delta ^{}(xy)\right].`$ (107)
These correspond to $`(2M)_{xy}^{ij}dX_0^i(x)dP_0^j(y)`$, and hence
$$M=\frac{(\pi L)^3}{3}C_{ijk}C_{klm}\left[P_0^l(x)_xP_0^m(x)\delta ^{}(xy)\right].$$
(108)
Next, we consider the $`N`$ part.
$`\mathrm{\Omega }_N^{[22]}`$ $`={\displaystyle \frac{(\pi L)^3}{3}}{\displaystyle 𝑑x𝑑yC_{ilk}C_{ljm}}`$
$`\times \{dX_0^i(x)dP_0^j(y)[P_0^m(y)P_0^k(y)\delta (xy)+_y\left(P_0^m(y)P_0^k(y)\delta (xy)\right)`$
$`+P_0^m(y)_y\left(P_0^k(y)\delta (xy)\right)+_y\left(P_0^m(y)_y\left(P_0^k(y)\delta (xy)\right)\right)]`$
$`+dX_0^j(x)dP_0^i(y)\left[_x\left(P_0^m(x)_x\left(P_0^k(x)\delta (xy)\right)\right)\right]\}`$
$`={\displaystyle \frac{(\pi L)^3}{3}}C_{ikl}C_{jml}{\displaystyle 𝑑x𝑑y𝑑X_0^i(x)}dP_0^j(y)`$
$`\times [P_0^m(x)P_0^k(x)\delta (xy)P_0^m(x)P_0^k(x)\delta ^{}(xy)`$
$`P_0^m(y)P_0^k(x)\delta ^{}(xy)_y(P_0^m(y)P)^k(x)\delta ^{}(xy))`$
$`+_x\left(P_0^k(x)_x\left(P_0^m(x)\delta (xy)\right)\right)],`$ (109)
where in the last term we make $`km`$. Using
$`P_0^k(x)P_0^m(y)\delta ^{}(xy)`$ $`=P_0^k(x)\left(P_0^{\prime \prime m}(x)\delta (xy)+P_0^m(x)\delta ^{}(xy)\right),`$
$`_y\left(P_0^k(x)P_0^m(y)\delta ^{}(xy)\right)`$ $`=P_0^k(x)P_0^m(x)\delta ^{}(xy)+P_0^k(x)P_0^m(x)\delta ^{\prime \prime }(xy),`$
$`_x\left(P_0^k(x)_x\left(P_0^m(x)\delta (xy)\right)\right)`$ $`=P_0^k(x)P_0^m(x)\delta (xy)+P_0^k(x)P_0^m(x)\delta ^{}(xy)`$
$`+P_0^k(x)P_0^{\prime \prime }(x)\delta (xy)+2P_0^k(x)P_0^m(x)\delta ^{}(xy)`$
$`+P_0^k(x)P_0^m(x)\delta ^{\prime \prime }(xy),`$ (110)
we obtain
$`(\text{109})`$ $`={\displaystyle \frac{(\pi L)^3}{3}}C_{ikl}C_{jml}{\displaystyle 𝑑x𝑑y𝑑X_0^i(x)}dP_0^j(y)`$
$`\times [2P_0^k(x)P_0^m(x)\delta ^{\prime \prime }(xy)+2P_0^k(x)P_0^m(x)\delta ^{}(xy)`$
$`+2P_0^k(x)P_0^m(x)\delta (xy)].`$ (111)
Hence,
$`N`$ $`={\displaystyle \frac{(\pi L)^3}{3}}C_{ikl}C_{jml}\left[P_0^k(x)P_0^m(x)\delta ^{\prime \prime }(xy)+P_0^m(x)_x\left(P_0^k(x)\delta (xy)\right)\right].`$ (112)
The $`m`$ part can be calculated as follows.
$`\mathrm{\Omega }_m^{[22]}`$ $`={\displaystyle \frac{(\pi L)^3}{3}}{\displaystyle 𝑑x𝑑yC_{ijk}C_{jlm}𝑑P_0^k(x)}dP_0^i(y)\left(_x\left(P_0^l(x)X_0^m(x)\right)\delta (xy)\right)`$
$`={\displaystyle \frac{(\pi L)^3}{3}}{\displaystyle 𝑑x𝑑yC_{ijk}C_{klm}𝑑P_0^i(x)}dP_0^j(y)\left(_x\left(P_0^l(x)X_0^m(x)\right)\delta (xy)\right),`$ (113)
then
$$m=\frac{(\pi L)^3}{3}C_{ijk}C_{klm}_y\left(P_0^l(y)_yX_0^m(y)\right)\delta (xy).$$
(114)
Finally, we compute the part of $`n`$. Because the result should be antisymmetric under $`\{i,x\}\{j,y\}`$, and $`n`$ is proportional to $`C_{ikl}C_{jml}`$, we only need to consider the antisymmetric part under $`\{k,x\}\{m,y\}`$.
$`\mathrm{\Omega }_n^{[22]}`$ $`={\displaystyle \frac{(\pi L)^3}{3}}{\displaystyle 𝑑x𝑑yC_{ijk}C_{jlm}}`$
$`\times dP_0^l(x)dP_0^i(y)(P_0^m(x)X_0^k(x)\delta (xy)+_x(P_0^m(x)X_0^k(x)\delta (xy))`$
$`+X_0^m(x)_x(P_0^k(x)\delta (xy)))`$
$`={\displaystyle \frac{(\pi L)^3}{3}}{\displaystyle 𝑑x𝑑yC_{ijk}C_{jlm}}`$
$`\times dP_0^l(x)dP_0^i(y)(P_0^m(x)X_0^k(x)\delta (xy)+P_0^k(x)X_0^m(x)\delta (xy)`$
$`+X_0^m(x)P_0^k(x)\delta ^{}(xy)+X_0^k(y)P_0^m(y)\delta ^{}(xy)).`$ (115)
The terms symmetric under $`\{m,x\}\{k,y\}`$ vanish, so
$`\mathrm{\Omega }_n^{[22]}`$ $`={\displaystyle \frac{(\pi L)^3}{3}}{\displaystyle 𝑑x𝑑yC_{ijk}C_{jlm}}`$
$`\times dP_0^l(x)dP_0^i(y)\left(X_0^m(x)P_0^k(x)\delta ^{}(xy)+X_0^k(y)P_0^m(y)\delta ^{}(xy)\right)`$
$`={\displaystyle \frac{(\pi L)^3}{3}}{\displaystyle 𝑑x𝑑yC_{ikl}C_{jml}}`$
$`\times dP_0^i(x)dP_0^j(y)\left(X_0^k(x)P_0^m(x)\delta ^{}(xy)+X_0^m(y)P_0^k(y)\delta ^{}(xy)\right).`$ (116)
Thus
$$n=\frac{(\pi L)^3}{3}C_{ikl}C_{jml}\left[X_0^k(x)P_0^m(x)+X_0^m(y)P_0^k(y)\right]\delta ^{}(xy).$$
(117)
|
warning/0005/astro-ph0005326.html
|
ar5iv
|
text
|
# A Glitch in an Anomalous X-ray Pulsar
## 1 Introduction
An unusual class of X-ray pulsars, the anomalous X-ray pulsars (AXPs), has been puzzling since the discovery of the first such object some 20 years ago (1E 2259+586, Fahlman & Gregory (1981)). AXPs are characterized by spin periods in the range of 5–12 s, steady spin down, X-ray luminosities greatly exceeding their inferred spin-down luminosities, steep X-ray spectra, and lack of evidence for a binary companion, either optically or from Doppler shifts (Mereghetti & Stella (1995); van Paradijs, Taam, & van den Heuvel (1995)). All five known AXPs are located in the Galactic Plane, and two are coincident with supernova remnants (Fahlman & Gregory (1981); Gotthelf & Vasisht (1998)). A sixth AXP candidate is also at the center of a supernova remnant (Gaensler, Gotthelf, & Vasisht (1999)).
Two main models have been suggested to explain the nature of the AXPs. The lack of evidence for companions and their location in the Galactic plane as well as in supernova remnants suggests that AXPs are young, isolated neutron stars. In this case, the steady spin-down, under the assumption that it is due to magnetic dipole braking as in radio pulsars, implies surface dipolar magnetic fields of $`10^{14}10^{15}`$ G. Such fields are similar to those inferred independently in the soft gamma repeaters; both classes of object have therefore been suggested to be “magnetars” (Duncan & Thompson (1992); Thompson & Duncan (1995, 1996); Kouveliotou et al. (1998, 1999)). The large X-ray luminosities of the AXPs in this model may arise from energy from the decay of the large magnetic field (Thompson & Duncan (1996)) or from enhanced thermal emission (Heyl & Hernquist (1997)).
Alternatively, it has been proposed that AXPs are accreting neutron stars, with either (i) a very low-mass companion (Mereghetti & Stella (1995)) or (ii) with no companion, but with accretion disks perhaps made of material leftover after a companion was disrupted (van Paradijs, Taam, & van den Heuvel (1995)), or, for a young neutron star, material remaining from the supernova explosion (Chatterjee, Hernquist, & Narayan (2000); Perna, Hernquist, & Narayan (2000); Alpar (1999)). In this case, the X-ray luminosity is from accretion, and the prolonged spin-down is a result of the pulsars being close to their equilibrium spin period or of them being in an extended “propeller” regime of centrifugal expulsion (Chatterjee et al. 2000, Alpar (1999)).
One way to discriminate among these models is through timing observations. In the magnetar model, timing irregularities and sudden spin-up events, as are seen in the young radio pulsar population, are expected (Thompson & Duncan (1996)), but long episodes of spin-up should not be seen. Also, a long-term periodicity superimposed on the spin-down might be expected due to radiative precession (Melatos (1999)). By contrast, in an accretion scenario, large random torque fluctuations could be expected, as might extended episodes of spin-up (Baykal & Swank (1996); Chakrabarty et al. (1997); Bildsten et al. (1997)).
Past timing observations of AXPs have been hampered by poor sampling, such that multiple interpretations of the same data set were possible (e.g. Usov (1993); Heyl & Hernquist (1999); Melatos (1999)). Kaspi, Chakrabarty & Steinberger (1999) \[hereafter KCS99\] showed that with monthly observations, phase-coherent timing of at least two AXPs (1E 2259+586, 1RXS 1708$``$4009) was possible, demonstrating that the AXPs can be very steady rotators and that such monitoring observations can in principle distinguish among models.
Here we report on continued monitoring of the 11-s AXP 1RXS J170849.0$``$400910 (hereafter 1RXS 1708$``$4009) with the Rossi X-ray Timing Explorer (RXTE). We show that although 1RXS 1708$``$4009 rotated extremely steadily for nearly 2 yr, a sudden spin-up event occurred between two observations at epochs MJD 51446 and 51472. We show that the properties of the event are very similar to the “glitches” seen in young radio pulsars.
## 2 Observations and Results
The RXTE observations described here are a continuation of those reported by KCS99. We refer the reader to that paper for details of the analysis procedure. Briefly, all observations were obtained with the Proportional Counter Array (Jahoda et al. (1996)), with events in the range 2.5–5.4 keV selected to maximize signal-to-noise ratio. Data have been obtained roughly monthly since 1998 January and were reduced using software designed to handle raw spacecraft telemetry packet data. They were binned at 62.5 ms resolution and reduced to the solar system barycenter in barycentric dynamical time using the JPL DE200 solar system ephemeris.
The spacing of the observations was carefully chosen to permit absolute pulse phase determination using standard radio pulsar techniques. The timing ephemeris of KCS99 was the starting point in the continuing analysis, with individual observations folded at the predicted barycentric period. A total of 64 pulse phase bins were used. Folded profiles were cross-correlated in the Fourier domain with a high signal-to-noise ratio average profile in order to determine an average pulse arrival time. Resulting arrival times were then analyzed using the TEMPO pulsar timing software package.<sup>1</sup><sup>1</sup>1http://pulsar.princeton.edu/tempo
The ephemeris given by KCS99, which was determined from 19 observations made in the interval MJD 50826 – 51324 (1998 January 13 – 1999 May 26), continued to predict phase for over 120 days, until MJD 51446 (1999 September 5). This is clear from the pre-glitch timing residuals (see Figure 1) which have RMS 130 ms (0.012$`P`$, where $`P=1/\nu `$ is the pulse period). The subsequent observation, on MJD 51472 (1999 October 21), was not well-predicted, and the following residuals grew steadily (see Figure 1a). For this reason, we initiated a pre-planned series of three closely spaced observations in order to independently determine the new pulse frequency $`\nu `$. All observations from MJD 51472 onward are well modeled by a single $`\nu `$ and $`\dot{\nu }`$. This revised ephemeris has now properly described 9 arrival times obtained over 142 days, with RMS residuals of only 71 ms (0.006$`P`$). Table 1 summarizes the spin parameters before and after the event, where the values are extrapolated to MJD 51459, the midpoint between MJDs 51446 and 51472. Residuals after subtraction of the pre-glitch model from the pre-glitch data and the post-glitch model from the post-glitch data are shown in Figure 1b.
The frequencies given in Table 1 imply that the pulsar suddenly spun up, with fractional frequency change $`|\mathrm{\Delta }\nu /\nu |=(6.2\pm 0.3)\times 10^7`$. Furthermore, following the event, the spin down rate increased in magnitude by $`|\mathrm{\Delta }\dot{\nu }/\dot{\nu }|=(1.38\pm 0.25)\times 10^2`$. In both cases, the uncertainties are derived by combining those of the pre- and post-glitch values in quadrature. These changes are very similar to those observed in the Vela radio pulsar, as well as in other radio pulsars of “adolescent” age (e.g. McKenna & Lyne (1990); Kaspi et al. (1992); Shemar & Lyne (1996); Lyne et al. (1996); Wang et al. (2000); see §3).
Although the timing event is well described by a simple step function model, it can in principle also be described by a continuous model with a single $`\nu ,\dot{\nu }`$ and significant $`\ddot{\nu }`$. However, in this case, the timing residuals show strong systematic trends, including a clear discontinuity at the epoch of the event, and the RMS residual is approximately three times larger than that in the pre-glitch model. Smooth deviations from a simple spin-down law have been observed in many, if not most radio pulsars and are known as “timing noise” for lack of a better term. However, discrete events like the one we have observed for 1RXS 1708$``$4009, especially since they are always observed to be sudden spin-ups, are a distinct phenomenon classified as glitches (see Lyne (1996) for a review). The identification of discrete events as a distinct phenomenon in radio pulsars has grown out of many years of phase-coherent timing observations of hundreds of sources, something unavailable for AXPs. Thus, by the conventional operational definition for glitches in radio pulsars, and by Occam’s Razor, we conclude that the timing event we have observed in 1RXS 1708$``$4009 is indeed a glitch. However, it should be kept in mind that it may instead represent a new phenomenon not seen in radio pulsars. Only continued timing observations will settle this point with certainty.
We detected no change in the 2.5–5.4 keV X-ray flux from the pulsar at the time of the glitch. We set an upper limit on flux variations of $`<`$20% (3$`\sigma `$) of the mean flux. We also detected no statistically significant change in the X-ray pulse profile at the time of the glitch.
## 3 Discussion
The spin-up event we have observed in 1RXS 1708$``$4009 is very similar to the glitches seen in the Vela radio pulsar and other radio pulsars of comparable age, that is, with $`10^4<\tau _c<10^5`$ yr, where characteristic age $`\tau _cP/2\dot{P}`$ (e.g. Shemar & Lyne (1996); Wang et al. (2000)). In such young pulsars, observed glitches are dominated by frequency steps of size $`\mathrm{\Delta }\nu /\nu 10^710^6`$. Furthermore, such glitches frequently show increases in the magnitude of the spin-down rate of order a few percent, sometimes, but not always, with subsequent relaxation back to the pre-glitch value on time scales of several hundred days. The rates of occurrence of such glitches vary from source to source, with some occurring more frequently than once per year (e.g. PSR J1341$``$6220, Kaspi et al. (1992); Wang et al. (2000)), and most (generally the older pulsars, $`\tau _c\stackrel{>}{_{}}50`$ kyr), never having been observed to glitch. All these properties are consistent with those of the spin-up event in 1RXS 1708$``$4009 ($`\tau _c=9`$ kyr), namely the magnitude of the glitch, the change in the slow-down rate, and even, very crudely, the rate of occurrence, once per $`2`$ yr of observation. Some glitching radio pulsars, especially the well-studied Vela pulsar, have also shown significant relaxation on time scales of hours to days (e.g. Chau et al. (1993)). However, such behavior is on too short a time scale to be detectable in our observations of 1RXS 1708$``$4009.
Large glitches in radio pulsars have been ascribed to sudden unpinning of superfluid neutron vortices (Anderson & Itoh (1975); Alpar, Cheng, & Pines (1989); Alpar et al. (1993)). The neutron star spins down under the influence of an external torque which acts on the crust. For radio pulsars, the torque is magnetic dipole braking. Neutron superfluid in the stellar interior, which is not well coupled to the crust, has its angular momentum carried in quantized vortices. The superfluid can spin down by outward motion of these vortices. However, vortex line pinning to crustal nuclei can impede their outward motion. The crust and superfluid components therefore develop a differential angular velocity. Occasionally, sudden unpinning of vortex lines can occur, and the previously decoupled superfluid can spin down, transferring angular momentum to the crust in the process. A spin-up event is therefore observed. The neutron superfluid thus acts as an angular momentum reservoir to fuel glitches. The similarities in the properties of the spin-up event seen in 1RXS 1708$``$4009 to those seen in the Vela-like pulsars suggests that a similar mechanism is at work in 1RXS 1708$``$4009.
In contrast, smaller glitches observed in the younger Crab pulsar are dominated by changes in spin-down rate rather than in pulse frequency (Lyne, Pritchard, & Smith (1993)). These are ascribed to changes in the neutron star ellipticity due to cracking of the crust. The magnitude and frequency of the Vela-like glitches are incompatible with such a model but agree well with the vortex-line unpinning model, in which the fractional angular momentum change per glitch is roughly constant from source to source (Alpar & Baykal (1994)).
From the observed glitch parameters, we can estimate the fraction of the neutron star moment of inertia in neutron superfluid that is not corotating with the crust, $`I_s`$. First, one can show (e.g. Link, Epstein, & Lattimer (1999)) that
$$\frac{I_s}{I_c}\frac{\overline{\nu }}{|\dot{\nu }|}A,$$
(1)
where $`I_c`$ is that of the crust and all other coupled components, $`\overline{\nu }`$ is the average spin frequency over the observing span, and $`A`$ is the activity parameter (McKenna & Lyne (1990)), where
$$A=\frac{1}{t}\underset{i}{}\frac{\mathrm{\Delta }\nu _i}{\nu }.$$
(2)
Here, $`t`$ is the observing span, and the sum is over all observed glitches. As we have observed only one glitch for 1RXS 1708$``$4009, we can only crudely estimate $`A`$, under the assumption that we were not extremely lucky in detecting the glitch, and perhaps also that the small value of $`\ddot{\nu }`$ observed before the glitch (KCS99) was due to relaxation following a glitch that occurred before our observations began (cf. Lyne et al. (1996)). Hence, we take $`t3`$ yr, so $`A5\times 10^7`$ yr$`{}_{}{}^{1}(3\mathrm{yr}/t)`$, and $`I_s/I_c0.01(3\mathrm{yr}/t)`$, similar to that found for many Vela-like pulsars (see Link, Epstein, & Lattimer (1999) and references therein). Alpar et al. (1993) suggested a different estimate for $`I_s`$, namely $`I_s/I_c\mathrm{\Delta }\dot{\nu }/\dot{\nu }`$, where short-term transient contributions to $`\mathrm{\Delta }\dot{\nu }`$ have been omitted. Since we were not sensitive to short time scale transients, our measured $`\mathrm{\Delta }\dot{\nu }`$ can be used directly, and yields $`I_s/I_c0.01`$, consistent with the first estimate.
Thus, the glitch implies that $`I_s`$ in 1RXS 1708$``$4009 is similar to that in the Vela-like pulsars. We note, as pointed out to us by I. Wasserman (personal communication) that this renders models of long time-scale precession in AXPs (Melatos (1999)) unlikely, because of the expected dynamics of the superfluid interior (Shaham (1977); Alpar & Ögelman (1987)).
Ruderman, Zhu & Cheng (1998) have suggested that the origin of the vortex line unpinning events is cracking of the neutron star crust under stresses imposed by outward-moving magnetic flux tubes. These tubes move because they interact with the outward-migrating angular momentum vortex lines as the neutron star spins down. However, this model predicts that glitch activity should be absent in neutron stars having $`P\stackrel{>}{_{}}0.7`$ s because the vortex motion is too slow to cause the necessary stresses. This is in contradiction with the large glitch in the 11-s 1RXS 1708$``$4009. Thus our observations suggests that the Ruderman et al. (1998) model is inapplicable to the glitch in 1RXS 1708$``$4009. Given the similarity of this event to those seen in Vela-like pulsars, this may cast doubt on the relevance of the model to those sources as well.
Usov (1993) and Heyl & Hernquist (1997) argued that spin-down irregularities in other AXPs (1E 2259+586 and 1E 1048.1$``$5937) are also due to glitches. The data they used did not involve phase coherent observations as did ours, and so their conclusions are much less certain. Furthermore, the fractional amplitude of the glitches they inferred are several orders of magnitude larger than what we have observed for 1RXS 1708$``$4009. Given the glitch in 1RXS 1708$``$4009, one might suspect that the previous claims of glitches in other AXPs were correct. However our ongoing observations of AXPs 1E 2259+586 and 1E 1048.1$``$5937 do not support the conclusion that the timing irregularities in those objects are due to sudden spin-up events. A detailed discussion of these sources will be presented elsewhere.
Glitches in AXPs were predicted in the magnetar model (Thompson & Duncan (1996)). These authors argued that crust fracture and superfluid vortex line unpinning play a major role in outbursts of soft-gamma repeaters (SGRs), and are ultimately due to the stresses imposed on the crust by the large magnetic field. However, the glitch alone does not provide proof of the magnetar hypothesis. The origin of neutron-star glitches in the vortex-line unpinning models is independent of the source of the external torque acting on the crust. Rather, it relies upon an angular velocity differential between the crust and that portion of the superfluid that is effectively decoupled from the crust.
It has been argued (Ruderman (1976); Alpar, Nandkumar, & Pines (1985); Ruderman (1991)) that the different nature of the glitches in the very young Crab pulsar ($`\tau _c=1`$ kyr, Lyne et al. 1993) and the absence of glitches in the young PSR B1509$``$58 ($`\tau _c=1.6`$ kyr, Kaspi et al. (1994)) imply that giant glitches do not occur in the youngest pulsars because they have higher internal temperatures, which allow a more plastic flow of vortex lines. However, in the magnetar model, the X-rays are a result of thermal processes, either magnetic field decay (Thompson & Duncan (1996)) or enhanced thermal emission from initial cooling (Heyl & Hernquist (1997)). In either case, the neutron star is very hot. This is supported by the X-ray spectrum of 1RXS 1708$``$4009: it can be fit with power-law and blackbody components (although the latter is not strictly required), which suggest a surface temperature of $`kT0.4`$ keV (Sugizaki et al. (1997)). This is hotter than is observed in any of the Vela-like pulsars, and higher than expected in the very youngest pulsars for all cooling models (Ögelman (1995)), even if the measured temperature is an overestimate of the true surface temperature because of atmospheric effects (e.g. Meyer, Pavlov, & Mészáros (1994)). Thus, in the magnetar model, the glitch in 1RXS 1708$``$4009 argues that the differences in the glitching behavior of the youngest radio pulsars compared to the “adolescent” Vela-like pulsars may not be primarily due to the difference in internal temperature.
In the recently proposed AXP model in which these sources are accreting from disks of material formed after the supernova explosion (Chatterjee et al. 2000, Perna et al. 2000, Alpar 1999), the spin-down rates are a result of accretion torque, which presumably acts only on magnetic field lines anchored in the crust. Thus, glitches might be expected in this model as well. In this case, the frequent glitches seen in Vela-like pulsars might be less influenced by their age than by their relatively large spin-down rates. A prediction of this hypothesis, independent of AXP phenomenology, is that glitches occur in neutron star X-ray binaries, although large fluctuations in accretion torque (e.g. Bildsten et al. (1997)) make them difficult to detect. The low-mass X-ray binary 4U 1626$``$67, in which the spin-down is extremely stable apart from episodes of sudden torque reversal (Chakrabarty et al. (1997)), should be an excellent candidate for the detection of glitches, although none has been seen in $``$5 yr of timing using the BATSE instrument.
Recent deep infrared observations of the field containing the AXP 1E 2259+586 (Hulleman et al. (2000)) have not detected any emission from a putative accretion disk, casting some doubt on the fallback disk model. If 1RXS 1708$``$4009 is indeed isolated, however, the glitch is consistent with the magnetar model, which provides the required external torque via magnetic dipole braking. However, since constraining optical/infrared observations of the 1RXS 1708$``$4009 field have yet to be done, an accretion scenario for this source cannot be ruled out.
Finally, glitches will complicate the determination of braking indexes in AXPs as well as the search for periodic variations in the spin period predicted in the magnetar model due to precession. Nevertheless, continued long-term monitoring of 1RXS 1708$``$4009 is essential for the determination of the amplitude distribution and frequency of its glitches. Similar observations of other AXPs are necessary to determine whether glitch behavior is ubiquitous.
We thank Evan Smith and the RXTE operations team for their skill and support in scheduling the AXP monitoring program. We also thank A. Alpar, R. Duncan, A. Lyne, D. Nice, S. Thorsett and I. Wasserman for useful discussions, and F. Crawford and D. Fox for helpful comments on the manuscript. This work was supported in part by a NASA LTSA grant (NAG5-8063) to VMK.
|
warning/0005/cond-mat0005366.html
|
ar5iv
|
text
|
# Tricritical Points in the Sherrington-Kirkpatrick Model in the Presence of Discrete Random Fields
## Abstract
The infinite-range-interaction Ising spin glass is considered in the presence of an external random magnetic field following a trimodal (three-peak) distribution. Such a distribution corresponds to a bimodal added to a probability $`p_0`$ for a field dilution, in such a way that at each site the field $`h_i`$ obeys $`P(h_i)=p_+\delta (h_ih_0)+p_0\delta (h_i)+p_{}\delta (h_i+h_0)`$. The model is studied through the replica method and phase diagrams are obtained within the replica-symmetry approximation. It is shown that the border of the ferromagnetic phase may present, for conveniently chosen values of $`p_0`$ and $`h_0`$, first-order phase transitions, as well as tricritical points at finite temperatures. Analogous to what happens for the Ising ferromagnet under a trimodal random field, it is verified that the first-order phase transitions are directly related to the dilution in the fields: the extensions of these transitions are reduced for increasing values of $`p_0`$. Whenever the delta function at the origin becomes comparable to those at $`h_i=\pm h_0`$, first-order phase transitions disappear; in fact, the threshold value $`p_0^{}`$, above which all phase transitions are continuous, is calculated analytically as $`p_0^{}=2(e^{3/2}+2)^10.30856`$. The ferromagnetic boundary at zero temperature also exhibits an interesting behavior: for $`0<p_0<p_0^{}`$, a single tricritical point occurs, whereas if $`p_0>p_0^{}`$ the critical frontier is completely continuous; however, for $`p_0=p_0^{}`$, a fourth-order critical point appears. The stability analysis of the replica-symmetric solution is performed and the regions of validity of such a solution are identified; in particular, the Almeida-Thouless line in the plane field versus temperature is shown to depend on the weight $`p_0`$.
Keywords: Spin Glasses, Random Field, Replica Method.
1. Introduction
Among disordered magnets , spin glasses and ferromagnets in the presence of random fields may be singled out as two of the most puzzling and controversial systems in condensed matter physics.
The random-field Ising model (RFIM), introduced by Imry and Ma , has concentrated a lot of interest after the identification of its physical realizations. Probably the most important physical conception of the RFIM comes out to be a diluted Ising antiferromagnet in the presence of a uniform magnetic field . Since then, many diluted antiferromagnets have been investigated, in such a way that systems like $`\mathrm{Fe}_\mathrm{x}\mathrm{Zn}_{1\mathrm{x}}\mathrm{F}_2`$ and $`\mathrm{Fe}_\mathrm{x}\mathrm{Mg}_{1\mathrm{x}}\mathrm{Cl}_2`$ are nowadays considered as standard experimental realizations of the RFIM . From the theoretical point of view, many important ingredients remain unknown. At the mean-field level, it is well known that different probability distributions for the random fields may lead to distinct phase diagrams, e.g., a Gaussian probability distribution yields a continuous ferromagnetic-paramagnetic boundary , whereas for a bimodal distribution, this boundary exhibits a continuous piece at high temperatures ending up at a tricritical point, which is followed by a first-order phase transition at low temperatures . Such a contrast in the mean-field phase diagrams of the RFIM with the bimodal and Gaussian probability distributions has been proven rigorously . Indeed, Aharony argued that whenever an analytic symmetric distribution for the fields presents a minimum at zero field, one should expect a tricritical point and a first-order transition for sufficiently low temperatures. Further studies of the RFIM at the mean-field level have considered a trimodal (three-peak) distribution
$$P(h_i)=p_+\delta (h_ih_0)+p_0\delta (h_i)+p_{}\delta (h_i+h_0),$$
$`(1.1)`$
in its symmetrical form, i.e., $`p_+=p_{}=\frac{1}{2}(1p_0)`$. Such a distribution, which may be interpreted as a bimodal added to a dilution in the fields with probability $`p_0`$ , is expected to mimic better real systems than its bimodal counterpart. It was shown that the field dilution plays an important role in what concerns the presence of the tricritical point: distinct analyses lead to slightly different estimates for the threshold value, above which the tricritical point disappears (whereas the analysis of Mattis shows that the tricritical point vanishes for $`p_0>0.25`$, according to Kaufman et al. such a behavior should occur for $`p_0>0.24`$). Whether the features in the mean-field phase diagrams of the RFIM should prevail on short-range-interaction models, represents a point which has attracted a lot of interest . For the three-dimensional RFIM, recent Monte Carlo simulations detect a jump in the magnetization but no latent heat, for both bimodal and Gaussian distributions, whereas high-temperature series expansions and a zero-temperature scaling analysis find a continuous transition for both distributions. However, in four dimensions the same zero-temperature analysis leads to a first-order transition in the bimodal case and a continuous one for a Gaussian distribution, in agreement with the mean-field predictions. Apart from that, the low-temperature phase of the RFIM, in finite dimensions, may present a nontrivial structure, with a complicated free-energy landscape, as suggested by perturbative analyses .
The Ising spin-glass (ISG) problem became, nowadays, one of the most controversial issues in the physics of disordered magnets. Its mean-field theory, based on the solution of the infinite-range-interaction model, the so-called Sherrington-Kirkpatrick (SK) model , presents a quite nontrivial behavior. The correct low-temperature solution, as proposed by Parisi , consists of a continuous order-parameter function (i.e., an infinite number of order parameters) associated with many low-energy states, a procedure which is usually denominated as replica-symmetry breaking (RSB). Furthermore, a transition in the presence of an external magnetic field, known as the Almeida-Thouless (AT) line , is found in the solution of the SK model: such a line separates a low-temperature region, characterized by RSB, from a high-temperature one, where a simple one-parameter solution, denominated as replica-symmetric (RS) solution, is stable. The validity of the results of the SK model for the description of real (short-range-interaction) systems represents a very polemic question . The rival theory is the droplet model , based on domain-wall renormalization-group arguments for spin glasses . According to the droplet model, the low-temperature phase of any finite-dimensional short-range spin glass should be described in terms of a single thermodynamic state (together, of course, with its corresponding time-reversed counterpart), i.e., essentially a RS-type of solution. Obviously, the droplet model becomes questionable for increasing dimensionalities, where one expects the existence of a finite upper critical dimension – believed to be six for the ISG – above which the mean-field picture should prevail. Recent analyses of short-range ISG on diamond hierarchical lattices (on which the Migdal-Kadanoff renormalization group is exact) has found evidences of the droplet picture ; however, the applicability of such lattices for the description of ISG on Bravais lattices is doubtful . Numerical simulations are very hard to be carried for short-range ISG on a cubic lattice, due to large thermalization times ; as a consequence, no conclusive results in three-dimensional systems are available. However, in four dimensions the critical temperature is much higher, making thermalization easier; in this case, many works claim to have observed some mean-field features .
From the theoretical point of view these two problems (RFIM and ISG), have been, in most of the cases, studied in separate, with a few exceptions . However, many diluted antiferromagnets, like $`\mathrm{Fe}_\mathrm{x}\mathrm{Zn}_{1\mathrm{x}}\mathrm{F}_2`$ and $`\mathrm{Fe}_\mathrm{x}\mathrm{Mg}_{1\mathrm{x}}\mathrm{Cl}_2`$ , are able to exhibit, within certain concentration ranges, random-field, spin-glass or both behaviors. For the $`\mathrm{Fe}_\mathrm{x}\mathrm{Zn}_{1\mathrm{x}}\mathrm{F}_2`$, one gets a RFIM for $`\mathrm{x}0.40`$, an ISG for $`\mathrm{x}0.24`$, whereas for intermediate concentrations ($`0.24\mathrm{x}0.40`$) one may observe both behaviors depending on the magnitude of the applied external magnetic field \[RFIM (ISG) for small (large) magnetic fields\], with a crossover between them; this latter effect was observed in $`\mathrm{Fe}_{0.31}\mathrm{Zn}_{0.69}\mathrm{F}_2`$ . Certainly, such properties are expected to be properly explained only if one considers a model which takes into account both spin-glass and random-field ingredients. Indeed, the crossover observed in $`\mathrm{Fe}_{0.31}\mathrm{Zn}_{0.69}\mathrm{F}_2`$ was also found in the study of the SK model under a Gaussian random field . On the other hand, the study of the SK model in the presence of a bimodal random field produced interesting results, with first-order phase transitions and tricritical points ; such results may be relevant for explaining the first-order phase transitions observed in $`\mathrm{Fe}_\mathrm{x}\mathrm{Mg}_{1\mathrm{x}}\mathrm{Cl}_2`$ .
In the present work we study the SK model in the presence of a random field following a trimodal probability distribution \[see Eq. (1.1)\]. In addition to that, one may interpolate between the bimodal distribution and a behavior which is qualitatively analogous to the Gaussian one, since by monitoring the delta function at the origin, one is able to control the presence of tricritical points. In the next section we define the model and, through the use of the replica method, we find its free-energy density, equations of state and equations for the validity of the RS solution. In section 3 we exhibit and discuss the phase diagrams of the model. Finally, in section 4 we present our conclusions.
2. The Model and Replica Formalism
The mean-field theory of the ISG is usually formulated as a set of $`N`$ spins, each of them interacting with all others \[a total of $`\frac{1}{2}N(N1)`$ interactions\], known as the SK model . The SK model in the presence of an external random magnetic field may be defined in terms of the Hamiltonian ,
$$=\underset{(ij)}{}J_{ij}S_iS_j\underset{i}{}h_iS_i,$$
$`(2.1)`$
where $`S_i=\pm 1`$, with $`i=1,2,\mathrm{},N`$, and the interactions are infinite-range-like, i.e., the sum $`_{(i,j)}`$ applies to all distinct pairs of spins. The coupling constants $`\{J_{ij}\}`$ and the random fields $`\{h_i\}`$ are quenched variables, following independent probability distributions,
$$P(J_{ij})=\left(\frac{N}{2\pi J^2}\right)^{\frac{1}{2}}\mathrm{exp}\left[\frac{N}{2J^2}\left(J_{ij}\frac{J_0}{N}\right)^2\right],$$
$`(2.2)`$
with $`P(h_i)`$ given by Eq. (1.1) ($`p_++p_0+p_{}=1`$). Let us, for the moment, keep the trimodal probability distribution in its general form of Eq. (1.1); later on, we will see that the ferromagnetic boundary does not exist for $`p_+p_{}`$, and so, in such a case, we will be restricted to the symmetrical form $`p_+=p_{}=\frac{1}{2}(1p_0)`$. It should be mentioned that the above randomnesses ($`\{J_{ij}\}`$ and $`\{h_i\}`$) are usually correlated in real systems; herein for the sake of simplicity, we shall consider two independent probability distributions. Therefore, for a given realization of bonds and site-fields, $`(\{J_{ij}\},\{h_i\})`$, one has a corresponding free energy, $`F(\{J_{ij}\},\{h_i\})`$, such that the average over the disorder, $`[]_{J,h}`$, may be performed as independent integrals,
$$[F(\{J_{ij}\},\{h_i\})]_{J,h}=\underset{(ij)}{}[dJ_{ij}P(J_{ij})]\underset{i}{}[dh_iP(h_i)]F(\{J_{ij}\},\{h_i\}).$$
$`(2.3)`$
The usual procedure consists in applying the replica method , in such a way as to get the free energy per spin as,
$`\beta f`$ $`=`$ $`\underset{N\mathrm{}}{lim}{\displaystyle \frac{1}{N}}[\mathrm{ln}Z(\{J_{ij}\},\{h_i\})]_{J,h}`$ (2.4)
$`=`$ $`\underset{N\mathrm{}}{lim}\underset{n0}{lim}{\displaystyle \frac{1}{Nn}}\left([Z^n]_{J,h}1\right),`$ (2.5)
where $`Z^n`$ is the partition function of $`n`$ copies of the system defined in Eq. (2.1) and $`\beta =1/T`$ (we work in units $`k_B=1`$). Standard calculations lead to
$$\beta f=\frac{(\beta J)^2}{4}+\underset{n0}{lim}\frac{1}{n}\mathrm{min}g(m^\alpha ,q^{\alpha \beta }),$$
$`(2.5)`$
where
$$g(m^\alpha ,q^{\alpha \beta })=\frac{\beta J_0}{2}\underset{\alpha }{}(m^\alpha )^2+\frac{(\beta J)^2}{2}\underset{(\alpha \beta )}{}\left(q^{\alpha \beta }\right)^2p_+\mathrm{ln}\mathrm{Tr}_\alpha \mathrm{exp}(_{eff}^+)$$
$$p_0\mathrm{ln}\mathrm{Tr}_\alpha \mathrm{exp}(_{eff}^0)p_{}\mathrm{ln}\mathrm{Tr}_\alpha \mathrm{exp}(_{eff}^{}),$$
$`(2.6\mathrm{a})`$
$$_{eff}^\pm =\beta J_0\underset{\alpha }{}m^\alpha S^\alpha +(\beta J)^2\underset{(\alpha \beta )}{}q^{\alpha \beta }S^\alpha S^\beta \pm \beta h_0\underset{\alpha }{}S^\alpha ,$$
$`(2.6\mathrm{b})`$
$$_{eff}^0=\beta J_0\underset{\alpha }{}m^\alpha S^\alpha +(\beta J)^2\underset{(\alpha \beta )}{}q^{\alpha \beta }S^\alpha S^\beta .$$
$`(2.6\mathrm{c})`$
In the equations above, the sum indexes $`\alpha `$ and $`\beta `$ $`(\alpha ,\beta =1,2,\mathrm{},n)`$ are replica labels and $`_{(\alpha \beta )}`$ denote sums over distinct pairs of replicas.
The extrema of the functional $`g(m^\alpha ,q^{\alpha \beta })`$ give us the equilibrium equations for the magnetization and spin-glass order parameters, respectively,
$`m^\alpha `$ $`=`$ $`p_+S^\alpha _++p_0S^\alpha _0+p_{}S^\alpha _{},`$ (2.7a)
$`q^{\alpha \beta }`$ $`=`$ $`p_+S^\alpha S^\beta _++p_0S^\alpha S^\beta _0+p_{}S^\alpha S^\beta _{}(\alpha \beta ),`$ (2.7b)
where $`_\pm `$ and $`_0`$ refer to thermal averages with respect to the “effective Hamiltonians” $`_{eff}^\pm `$ and $`_{eff}^0`$ in Eqs. (2.6b) and (2.6c), respectively.
If one assumes the replica-symmetry (RS) ansatz ,
$$m^\alpha =m,\alpha ;q^{\alpha \beta }=q,(\alpha \beta ),$$
$`(2.8)`$
the free energy per spin (Eq. (2.5)) and the equilibrium conditions (Eqs. (2.7)) become
$$\beta f=\frac{(\beta J)^2}{4}(1q)^2+\frac{\beta J_0}{2}m^2p_+𝒟z\mathrm{ln}(2\mathrm{cosh}\xi ^+)$$
$$p_0𝒟z\mathrm{ln}(2\mathrm{cosh}\xi ^0)p_{}𝒟z\mathrm{ln}(2\mathrm{cosh}\xi ^{}),$$
$`(2.9)`$
$$m=p_+𝒟z\mathrm{tanh}\xi ^++p_0𝒟z\mathrm{tanh}\xi ^0+p_{}𝒟z\mathrm{tanh}\xi ^{},$$
$`(2.10)`$
$$q=p_+𝒟z\mathrm{tanh}^2\xi ^++p_0𝒟z\mathrm{tanh}^2\xi ^0+p_{}𝒟z\mathrm{tanh}^2\xi ^{},$$
$`(2.11)`$
where
$$𝒟z\mathrm{}=_{\mathrm{}}^{\mathrm{}}\left(\frac{1}{2\pi }\right)^{\frac{1}{2}}𝑑z\mathrm{exp}(z^2/2)\mathrm{},$$
$`(2.12)`$
and
$`\xi ^\pm `$ $`=`$ $`\beta J_0m+\beta Jq^{1/2}z\pm \beta h_0,`$ (2.13a)
$`\xi ^0`$ $`=`$ $`\beta J_0m+\beta Jq^{1/2}z.`$ (2.13b)
Although the spin-glass order parameter (Eq. (2.11)) is always induced by a nonzero random field ($`p_0<1`$), it may still contribute to a nontrivial behavior; this is provided by the instability of the RS solution. Such an instability occurs at the AT line ,
$$\left(\frac{T}{J}\right)^2=p_+𝒟z\mathrm{sech}^4\xi ^++p_0𝒟z\mathrm{sech}^4\xi ^0+p_{}𝒟z\mathrm{sech}^4\xi ^{},$$
$`(2.14)`$
which may be obtained through the simultaneous solution of Eqs. (2.14), (2.10) and (2.11).
In the next section we shall consider the phase diagrams of the model and the regions of instability of the RS solution, worked out from Eqs. (2.9)–(2.14).
3. Results and Discussion
Let us first consider the case $`J_0=0`$; one may easily see that the only nontrivial behavior in this case is given by the AT instability in the plane magnetic field versus temperature, which may now be obtained from the solution of Eqs. (2.11) and (2.14). The integrals involving $`\xi ^{}`$ may be easily transformed through the change of variables $`zz`$, in such a way that the AT line may be obtained by solving the set of equations,
$`\left({\displaystyle \frac{T}{J}}\right)^2`$ $`=`$ $`(1p_0){\displaystyle 𝒟z\mathrm{sech}^4(\beta Jq^{1/2}z+\beta h_0)}+p_0{\displaystyle 𝒟z\mathrm{sech}^4(\beta Jq^{1/2}z)},`$ (3.1a)
$`q`$ $`=`$ $`(1p_0){\displaystyle 𝒟z\mathrm{tanh}^2(\beta Jq^{1/2}z+\beta h_0)}+p_0{\displaystyle 𝒟z\mathrm{tanh}^2(\beta Jq^{1/2}z)}.`$ (3.1c)
It should be pointed out that the equations above are valid for arbitrary values of the weights in the probability distribution of Eq. (1.1), with $`p_++p_{}=1p_0`$; although the AT line changes with field dilution, it is no altered under a field inversion. The AT lines in the plane magnetic field versus temperature are exhibited in Fig. 1, for typical values of $`p_0`$. Clearly, the AT line for the bimodal distribution ($`p_0=0`$) is identical to the one of the SK model in the presence of a uniform magnetic field , due to the property of invariance under field inversion. For $`0<p_0<1`$, one may calculate analytically the behavior of the AT line in the low-field regime ($`TJ`$),
$$1\frac{T}{J}\left[\frac{3(1p_0)}{4}\right]^{1/3}\left(\frac{h_0}{J}\right)^{2/3},$$
$`(3.2)`$
which leads to a slightly modified amplitude, but the same low-field exponent of the standard AT line . If one considers $`p_00`$, the low-temperature behavior of the AT line may be easily calculated,
$$\frac{T}{J}\frac{4}{3}\frac{1}{\sqrt{2\pi }}\left[(1p_0)\mathrm{exp}\left(\frac{h_0^2}{2J^2}\right)+p_0\right],$$
$`(3.3)`$
which exhibits the usual exponential decay , but with a shift towards higher temperatures for increasing values of $`p_0`$. In all other situations, the AT lines were calculated by solving numerically Eqs. (3.1). One notices that for high values of $`p_0`$, the integrals multiplying $`p_0`$ in Eqs. (3.1) contribute significantly, in such a way that the AT lines become slightly independent of $`h_0`$, for $`h_0`$ large enough, as shown in Fig. 1.
From now on, we will be restricted to $`J_0>0`$; in this case, as far as RS is concerned, if $`p_+p_{}`$ Eqs. (2.10) and (2.11) yield nonzero magnetization and spin-glass order parameters, leading to trivial behavior. Therefore, for the rest of this paper we will concentrate on a symmetrical trimodal distribution, i.e., $`p_+=p_{}=\frac{1}{2}(1p_0)`$. In this case, the random field still induces the parameter $`q`$, leading to no spontaneous spin-glass order (like the one found for the SK model in the absence of external field ). Therefore, the only possible phase transition within the RS approximation is the one associated with the magnetization, similarly to what happened in the case of the bimodal distribution . Hence, two phases are possible, namely, the ferromagnetic ($`m0,q0`$) and the independent ($`m=0,q0`$) ones. Although in the RFIM this latter phase is usually denominated of paramagnetic, in the present problem, within the RS approximation, we shall keep the nomenclature independent, for reasons which will become clear soon.
The critical frontier separating these two phases may be found by solving the equilibrium equations, (2.10) and (2.11); in the case of first-order phase transitions, we shall make use of the free-energy per spin \[Eq. (2.9)\] as well. Expanding Eq. (2.10) in powers of $`m`$ one gets,
$$m=A_1(q)m+A_3(q)m^3+A_5(q)m^5+O(m^7),$$
$`(3.4)`$
where the coefficients depend on $`q`$ \[which on its turn, depends on $`m`$ through Eq. (2.11)\]. Expanding Eq. (2.11) in powers of $`m`$,
$$q=q_0+\frac{(\beta J_0)^2\mathrm{\Gamma }}{1(\beta J)^2\mathrm{\Gamma }}m^2+O(m^4),$$
$`(3.5)`$
with
$$\mathrm{\Gamma }=(1p_0)(14\rho _1^++3\rho _2^+)+p_0(14\rho _1^0+3\rho _2^0),$$
$`(3.6)`$
$`\rho _k^+`$ $`=`$ $`{\displaystyle 𝒟z\mathrm{tanh}^{2k}(\beta Jq_0^{1/2}z+\beta h_0)},`$ (3.7a)
$`\rho _k^0`$ $`=`$ $`{\displaystyle 𝒟z\mathrm{tanh}^{2k}(\beta Jq_0^{1/2}z)},`$ (3.7b)
where $`q_0`$ is independent of $`m`$, corresponding to the solution of Eq. (2.11) with $`m=0`$. Substituting the above results into Eq. (3.4), one gets the $`m`$-independent coefficients of the power expansion,
$`A_1^{}`$ $`=`$ $`\beta J_0[1(1p_0)\rho _1^+p_0\rho _1^0],`$ (3.8a)
$`A_3^{}`$ $`=`$ $`{\displaystyle \frac{(\beta J_0)^3}{3}}\left[{\displaystyle \frac{1+2(\beta J)^2\mathrm{\Gamma }}{1(\beta J)^2\mathrm{\Gamma }}}\right]\mathrm{\Gamma },`$ (3.8c)
$`A_5^{}`$ $`=`$ $`\gamma {\displaystyle \frac{(\beta J_0)^5}{30}}\left[{\displaystyle \frac{1+8(\beta J)^2\mathrm{\Gamma }+36(\beta J)^4\mathrm{\Gamma }^2+15(\beta J)^6\mathrm{\Gamma }^3}{1(\beta J)^2\mathrm{\Gamma }}}\right],`$ (3.8e)
where
$$\gamma =(1p_0)(4+34\rho _1^+60\rho _2^++30\rho _3^+)+p_0(4+34\rho _1^060\rho _2^0+30\rho _3^0).$$
$`(3.9)`$
The critical frontier may be determined using standard procedures, as described below.
(i) For continuous phase transitions, $`A_1^{}=1`$ and $`A_3^{}<0`$.
(ii) A first-order phase transition occurs whenever $`A_1^{}=1`$ and $`A_3^{}>0`$; the proper critical frontier should be found, in this case, through a Maxwell construction, i.e., by equating the free energies of the two phases.
(iii) When both types of phase transitions are present, the continuous- and first-order critical frontiers meet at a tricritical point , which defines the limit of validity of the series expansions; beyond the tricritical point the magnetization is discontinuous. The location of such point is determined by setting $`A_1^{}=A_3^{}=0`$, with the condition $`A_5^{}<0`$ satisfied.
In Figs. 2–4 we show three qualitatively distinct ferromagnetic boundaries of the present problem, for a typical value of $`p_0`$ ($`p_0=0.3`$), compared with those of the bimodal probability distribution ($`p_0=0`$). In Fig. 2 there is a single point along the ferromagnetic boundary at which $`A_3^{}=0`$; such a point may not be considered as tricritical, since there is no first-order phase transition. However, for any value of $`h_0`$ greater than those of Fig. 2 \[$`h_0/J=0.9573`$ ($`p_0=0`$) and $`h_0/J=1.53526`$ ($`p_0=0.3`$)\], one gets first-order phase transitions, and at least one tricritical point. In Fig. 3 we show situations where two tricritical points appear along the ferromagnetic boundary; we have verified that, for a fixed value of $`p_0`$, such a behavior occurs within a narrow interval of $`h_0`$. In Fig. 4 a single tricritical point emerges, separating a continuous boundary (high temperatures) from a first-order critical frontier (low temperatures). From such phase diagrams, one notices that the main effect of the field dilution is to push the tricritical points towards lower temperatures, i.e., the temperature range over which the first-order transitions occur decreases.
As mentioned before, although the spin-glass order parameter is always induced by the random field, it may still exhibit interesting behavior, associated with the instability of the RS solution. The AT instabilities, given by the solution of Eqs. (2.10), (2.11) and (2.14) with $`p_+=p_{}=\frac{1}{2}(1p_0)`$, yields two distinct lines in the phase diagrams of Figs. 2–4, depending on whether one is inside the independent phase ($`m=0`$), or in the ferromagnetic ($`m0`$) one. In the former case, the AT line is a straight line (independent of $`J_0`$), whereas in the latter, it presents the usual decrease with temperature for increasing values of $`J_0`$, in such a way that for low temperatures one gets the exponential decays,
$$\frac{T}{J}\frac{4}{3}\frac{1}{\sqrt{2\pi }}\{\frac{1}{2}(1p_0)\mathrm{exp}[\frac{(J_0+h_0)^2}{2J^2}]+p_0\mathrm{exp}[\frac{J_0^2}{2J^2}]$$
$$+\frac{1}{2}(1p_0)\mathrm{exp}[\frac{(J_0h_0)^2}{2J^2}]\}.$$
$`(3.10)`$
Herein we shall adopt the usual criteria for the identification of the regions where RS is stable and those throughout which a RSB procedure is necessary . The two regions with zero magnetization will be associated with the paramagnetic (high temperatures) and spin-glass (low temperatures) phases, whereas those with nonzero magnetization will be associated with the ferromagnetic (high temperatures) and mixed-ferromagnetic (low temperatures). The several phases exhibited in our phase diagrams are identified as:
| Paramagnetic (P) | ($`m=0`$ ; $`q`$ : RS) | ; |
| --- | --- | --- |
| Spin-Glass (SG) | ($`m=0`$ ; $`q`$ : RSB) | ; |
| Ferromagnetic (F) | ($`m0`$ ; $`q`$ : RS) | ; |
| Mixed Ferromagnetic (F) | ($`m0`$ ; $`q`$ : RSB) | . |
It should be mentioned that the present low-temperature results are questionable inside the phases F and SG, due to the instability of the RS solution; in particular the point for $`p_0=0.3`$ where $`A_3^{}=0`$ in Fig. 2, as well as the low-temperature tricritical points of Fig. 3 may completely disappear under a RSB procedure. However the high-temperature tricritical points, like those of Figs. 3 and 4, are inside the region of stability of the RS solution and will persist under more general treatments; we believe that such points are reminiscent of the tricritical point of the bimodal RFIM.
The two AT lines mentioned above usually meet at a continuous ferromagnetic boundary; however, these lines do not match each other across first-order phase transitions : there is a small (but finite) gap between them in Figs. 3 and 4.
Let us now investigate the ferromagnetic boundary at zero temperature; for $`T=0`$ the spin-glass order parameter is trivial ($`q=1`$), in such a way that one gets for the free energy and magnetization,
$`f`$ $`=`$ $`{\displaystyle \frac{J_0}{2}}m^2{\displaystyle \frac{h_0}{2}}(1p_0)\left[\mathrm{erf}\left({\displaystyle \frac{J_0m+h_0}{J\sqrt{2}}}\right)\mathrm{erf}\left({\displaystyle \frac{J_0mh_0}{J\sqrt{2}}}\right)\right]`$ (3.11e)
$`{\displaystyle \frac{J}{\sqrt{2\pi }}}(1p_0)\left\{\mathrm{exp}\left[{\displaystyle \frac{(J_0m+h_0)^2}{2J^2}}\right]+\mathrm{exp}\left[{\displaystyle \frac{(J_0mh_0)^2}{2J^2}}\right]\right\}`$
$`{\displaystyle \frac{2J}{\sqrt{2\pi }}}p_0\left\{\mathrm{exp}\left[{\displaystyle \frac{(J_0m)^2}{2J^2}}\right]\right\},`$
$`m`$ $`=`$ $`{\displaystyle \frac{1}{2}}(1p_0)\left[\mathrm{erf}\left({\displaystyle \frac{J_0m+h_0}{J\sqrt{2}}}\right)+\mathrm{erf}\left({\displaystyle \frac{J_0mh_0}{J\sqrt{2}}}\right)\right]+p_0\mathrm{erf}\left({\displaystyle \frac{J_0m}{J\sqrt{2}}}\right).`$ (3.11g)
Using a similar procedure as the one for finite temperatures, one may expand Eq. (3.11b),
$$m=a_1m+a_3m^3+a_5m^5+O(m^7),$$
$`(3.12)`$
where,
$`a_1`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{\pi }}}{\displaystyle \frac{J_0}{J}}\left[(1p_0)\mathrm{exp}\left({\displaystyle \frac{h_0^2}{2J^2}}\right)+p_0\right],`$ (3.13a)
$`a_3`$ $`=`$ $`{\displaystyle \frac{1}{6}}\sqrt{{\displaystyle \frac{2}{\pi }}}\left({\displaystyle \frac{J_0}{J}}\right)^3\left[(1p_0)\left({\displaystyle \frac{h_0^2}{J^2}}1\right)\mathrm{exp}\left({\displaystyle \frac{h_0^2}{2J^2}}\right)p_0\right],`$ (3.13c)
$`a_5`$ $`=`$ $`{\displaystyle \frac{1}{120}}\sqrt{{\displaystyle \frac{2}{\pi }}}\left({\displaystyle \frac{J_0}{J}}\right)^5\left[(1p_0)\left({\displaystyle \frac{h_0^4}{J^4}}6{\displaystyle \frac{h_0^2}{J^2}}+3\right)\mathrm{exp}\left({\displaystyle \frac{h_0^2}{2J^2}}\right)3p_0\right].`$ (3.13e)
The critical frontier separating the phases F and SG is shown in Fig. 5 for typical values of $`p_0`$. One notices that the effect of the weight $`p_0`$ is to favour the continuous line, along which $`a_1=1`$ with $`a_3<0`$, i.e.,
$$\frac{J_0}{J}=\sqrt{\frac{\pi }{2}}\frac{1}{p_0+(1p_0)\mathrm{exp}(h_0^2/2J^2)},$$
$`(3.14)`$
while decreasing the extension of the first-order transition line. For small values of $`p_0`$ these two lines meet at a tricritical point, obtained by solving the equations $`a_1=1`$, $`a_3=0`$, with the condition $`a_5<0`$; within the analysis for finite temperatures, this corresponds to the situation where the lower-temperature tricritical point (cf. Fig. 3) hits the zero-temperature axis. If $`p_0=0`$ such an effect occurs at
$$\frac{h_0}{J}=1;\frac{J_0}{J}=\sqrt{\frac{\pi e}{2}}2.0664.$$
$`(3.15)`$
We verified that for $`0<p_0<p_0^{}`$ (where $`p_0^{}`$ will be defined below), such a set of equations presents two solutions, although only one of them represents a tricritical point, satisfying $`a_5<0`$. By increasing $`p_0`$ inside this range, we noticed that such solutions get closer and colapse for $`p_0=p_0^{}`$. We calculated analytically $`p_0^{}=2(e^{3/2}+2)^10.30856`$, at which a fourth-order critical point (characterized by $`a_1=a_3=a_5=0`$, with $`a_7<0`$) occurs at
$$\frac{h_0}{J}=\sqrt{3}1.73207;\frac{J_0}{J}=\frac{\sqrt{2\pi }}{6}(e^{3/2}+2)2.70786.$$
$`(3.16)`$
The value $`p_0^{}`$ represents a threshold of $`p_0`$, above which there are no first-order transitions for any temperature $`T0`$. For $`p_0>p_0^{}`$ the second-order critical frontier of Fig. 5 approaches an asymptote for large values of $`h_0`$; indeed, when $`p_01`$ the zero-temperature ferromagnetic boundary approaches a straight line at $`J_0/J=\sqrt{\pi /2}`$ \[see Eq. (3.14)\], characteristic of the SK model in zero field .
It should be mentioned that the finite-temperature vestigial points where $`A_3^{}=0`$, like the ones in Fig. 2, are qualitatively different from the fourth-order critical point found for $`p_0=p_0^{}`$ at zero temperature, even though both situations represent thresholds for the occurrence of tricritical points. In the former case, $`A_5^{}<0`$, whereas in the latter, $`A_5^{}=0`$. In Fig. 6 we exhibit the behavior of the coefficients $`A_3^{}`$ and $`A_5^{}`$, for temperatures along the ferromagnetic frontier, for the case (b) of Fig. 2, i.e., $`p_0=0.3`$ ($`h_0/J=1.53526`$), and $`p_0=p_0^{}`$ ($`h_0/J=\sqrt{3}`$). One clearly sees that the fourth-order critical point only shows up at zero temperature; its parameters, as defined in Eq. (3.16), correspond to the situation where the vestigial point of Fig. 2 collapses with the zero-temperature axis.
If $`0<p_0<p_0^{}`$, it is always possible to obtain first-order phase transitions by conveniently choosing the value of $`h_0`$. In Fig. 7 we exhibit the ranges of $`p_0`$ and $`h_0/J`$ throughout which first-order phase transitions and tricritical points are possible along the ferromagnetic boundary. In region (a), first-order phase transitions are conceivable at finite and zero temperature, with a single tricritical point (at finite temperatures): typical examples are shown in Fig. 4. Throughout a very narrow range \[region (b)\] two tricritical points appear and the first-order phase transition occurs only for finite temperatures: typical examples are exhibited in Fig. 3. The region (b) is delimited by characteristic values of ($`p_0,h_0/J`$): (i) the threshold for $`h_0/J`$ smaller corresponds to the set of points satisfying $`A_3^{}=0`$, but with no first-order phase transition (e.g., the vestigial points shown in Fig. 2); (ii) the delimiter for $`h_0/J`$ larger corresponds to the coordinates of the tricritical points at zero temperature. The vertical line in Fig. 7 is for $`p_0=p_0^{}`$, defining \[together with the delimiter (i) of region (b)\], the range throughout which the ferromagnetic boundary is always continuous \[region (c)\].
4. Conclusion
We have studied the Sherrington-Kirkpatrick spin glass in the presence of random fields $`\{h_i\}`$, following a trimodal (three-peak) probability distribution, which corresponds to a bimodal plus a probability $`p_0`$ for field dilution, i.e., $`P(h_i)=p_+\delta (h_ih_0)+p_0\delta (h_i)+p_{}\delta (h_i+h_0)`$. We have used the replica method and the phase diagrams were obtained within the replica-symmetry approximation. The boundary of the ferromagnetic phase exhibited an interesting behavior, with the presence of first-order phase transitions and tricritical points: within certain ranges for $`p_0`$ and $`h_0`$, a single or two tricritical points were encountered. We have shown that the first-order phase transitions are directly affected by the dilution in the fields, in such a away that the extension of such lines are reduced by increasing $`p_0`$. In fact, there is a threshold value, $`p_0^{}=2(e^{3/2}+2)^10.30856`$, above which the ferromagnetic boundary is always continuous. Such effects may be reminiscent of those occurring within the mean-field theory of the Ising ferromagnet in the presence of trimodal random fields: the single tricritical point that appears in the case of a bimodal distribution is washed way by the presence of the delta at the origin, whenever $`p_0`$ becomes greater than a certain value .
At zero temperature, if $`0<p_0<p_0^{}`$, the ferromagnetic critical frontier exhibits a single tricritical point, with a first-order phase transition at high values of $`h_0`$. By increasing $`p_0`$, the first-order line gets reduced and, for $`p_0=p_0^{}`$, a fourth-order critical point is observed; for $`p_0>p_0^{}`$, the ferromagnetic boundary is always continuous.
Although the spin-glass order parameter is induced by the random field ($`p_0<1`$), it may still contribute to a nontrivial behavior, in what concerns the stability of the replica-symmetric solution. We have calculated the regions of instability of such a solution, leading to the identification of two low-temperature phases, namely, the spin-glass and mixed ferromagnetic ones. Besides that, the Almeida-Thouless line in the plane field versus temperature was shown to depend on the weight $`p_0`$, with different amplitudes (but the same exponent) in the low-field regime, and qualitatively distinct high-field behaviors.
We have verified that whenever the ferromagnetic boundary presents both continuous and first-order transition lines meeting at a single finite-temperature tricritical point, such a point is located inside the region of stability of the replica-symmetric solution, and it will not be removed by a replica-symmetry-breaking procedure. However, when two tricritical points occur along the ferromagnetic boundary, at least one of them (the one at low temperatures) appears inside the unstable region, and its existence may be an artifact of the replica-symmetric solution.
The applicability of the present results in the description of real systems obviously depends on the survival of the mean-field characteristics in the respective short-range-interaction versions of Ising spin glasses and the Ising ferromagnet in the presence of a random field. However, the trimodal distribution employed herein is expected to mimic better real systems than the bimodal distribution itself. Although we are not aware of experimental observations that match with our results, we believe that the diluted antiferromagnet $`\mathrm{Fe}_\mathrm{x}\mathrm{Mg}_{1\mathrm{x}}\mathrm{Cl}_2`$ is a good candidate, since, for conveniently chosen dilutions, it may exhibit first-order phase transitions , as well as a crossover from first- to second-order behavior .
Acknowledgments
We acknowledge E. M. F. Curado for useful discussions. FDN thanks CNPq and Pronex/MCT (Brazilian granting agencies) for partial financial support.
Figure Captions
Fig. 1: The AT lines, for the SK model in the presence of a trimodal random field, in the plane $`h_0`$ versus $`T`$ (in units of $`J`$), for typical values of $`p_0`$.
Fig. 2: Phase diagram $`T`$ versus $`J_0`$ (in units of $`J`$) of the SK model in the presence of a trimodal random field with $`p_0=0.3`$, compared with one of the bimodal case ($`p_0=0`$), for conveniently chosen values of $`h_0`$. (a) $`h_0/J=0.9573`$ ($`p_0=0`$); (b) $`h_0/J=1.53526`$ ($`p_0=0.3`$). The ferromagnetic boundaries are continuous, except for the points where $`A_3^{}=0`$ \[cf. Eq. (3.7b)\], represented by black squares. These choices signal lower bounds for $`h_0`$, above which first-order phase transitions occur. The phase nomenclature is specified in the text, with the low-temperature phases SG and F delimited by AT lines.
Fig. 3: Phase diagram $`T`$ versus $`J_0`$ (in units of $`J`$) of the SK model in the presence of a trimodal random field with $`p_0=0.3`$, compared with one of the bimodal case ($`p_0=0`$), for conveniently chosen values of $`h_0`$, in such a way as to obtain two tricritical points (black circles) along the ferromagnetic boundary. (a) $`h_0/J=0.97`$ ($`p_0=0`$); (b) $`h_0/J=1.558`$ ($`p_0=0.3)`$. The dashed lines stand for first-order phase transitions. The phase nomenclature is the same as in Fig. 2.
Fig. 4: Phase diagram $`T`$ versus $`J_0`$ (in units of $`J`$) of the SK model in the presence of a trimodal random field with $`p_0=0.3`$, compared with one of the bimodal case ($`p_0=0`$), for conveniently chosen values of $`h_0`$, in such a way as to obtain a single tricritical point (black circle) along the ferromagnetic boundary. (a) $`h_0/J=1.02`$ ($`p_0=0`$); (b) $`h_0/J=1.58`$ ($`p_0=0.3`$). The phase nomenclature and line representations are as in Figs. 2 and 3.
Fig. 5: The zero-temperature phase diagram $`h_0`$ versus $`J_0`$ (in units of $`J`$) of the SK model in the presence of a trimodal random field, for typical values of $`p_0`$. If $`0<p_0<p_0^{}`$ one always gets tricritical points (black circles), followed by first-order phase transitions for high values of $`h_0`$. When $`p_0=p_0^{}`$, one gets a fourth-order critical point (represented by a star). Above the threshold value $`p_0^{}=2(e^{3/2}+2)^10.30856`$, the critical frontier separating the phases SG and F is continuous.
Fig. 6: The ordinate represents either the coefficient $`A_3^{}`$ or $`A_5^{}`$ \[Eqs. (3.7b) and (3.7c), respectively\] along the ferromagnetic boundary, for $`p_0=0.3`$ ($`h_0/J=1.53526`$) (dot-dahed lines) and $`p_0=p_0^{}`$ ($`h_0/J=\sqrt{3}`$) (full lines), as a function of temperature. In the former case, $`A_3^{}=0`$ at $`T/J0.25`$ (with $`A_5^{}<0`$), whereas in the latter, $`A_3^{}=A_5^{}=0`$ at $`T=0`$.
Fig. 7: Ranges of $`p_0`$ and $`h_0/J`$ associated with distinct behaviors for the ferromagnetic boundary. (a) First-order phase transitions at finite and zero temperatures, with a single tricritical point at finite temperatures; (b) Two tricritical points with a first-order phase transition for finite temperatures; (c) Continuous phase transitions.
|
warning/0005/hep-ph0005139.html
|
ar5iv
|
text
|
# TTP00-09 HET-BNL-00/8 May 2000 hep-ph/0005139v3 Quartic mass corrections to 𝑅_had at 𝒪(𝛼_𝑠³)
## Summary
The total cross section for the production of massive quarks in electron positron annihilation can be predicted in perturbative QCD. After expansion in $`m^2/s`$ the quartic terms, i.e. those proportional to $`m^4/s^2`$, were calculated up to order $`\alpha _s^3`$ for vector and axial current induced rates. Predictions relevant for charm, bottom and top quark production were presented. The $`\alpha _s^3`$ corrections were shown to be comparable to terms of order $`\alpha _s`$ and $`\alpha _s^2`$. As a consequence, the predictions exhibit a sizeable dependence on the renormalization scale. Adopting instead of the $`\overline{\mathrm{MS}}`$ scheme a framework where the running mass $`\overline{m}(\mu )`$ is replaced by the invariant mass $`\widehat{m}`$, the stability of the prediction is improved and, at the same time, the relative size of the large order terms decreases. Obviously, an improved understanding of the origin of these large corrections would be highly desirable.
Combining these results with the massless prediction and the quadratic mass terms we have demonstrated that the cross section for massive quark production at electron positron colliders is under control in order $`\alpha _s^3`$ from the high energy region down to fairly low energies.
## Remark
The results of this paper are available in Mathematica format at
`http://www-ttp.physik.uni-karlsruhe.de/Progdata/ttp00/ttp00-09/`.
## Acknowledgments
This work was supported by DFG-Forschergruppe “Quantenfeldtheorie, Computeralgebra und Monte-Carlo-Simulation”. R.H. acknowledges support by Landesgraduiertenförderung at the University of Karlsruhe and by Deutsche Forschungsgemeinschaft.
|
warning/0005/hep-ph0005321.html
|
ar5iv
|
text
|
# New model for the neutrino mass matrix
## Figure captions
Figure 1: Two-loops Feynman diagram which generates $`_{e\mu }`$.
Figure 2: Two-loops Feynman diagram which generates $`_{\mu \tau }`$.
Figure 3: One of the three-loops Feynman diagrams which generate $`_{ee}`$.
|
warning/0005/cond-mat0005154.html
|
ar5iv
|
text
|
# Spin-Waves in itinerant ferromagnets
## I Introduction
The most prominent examples for metals with a ferromagnetic order are the elements of the iron group, namely iron, cobalt and nickel. Although the magnetic behaviour of these materials is a well-known phenomenon, there are still many open questions in this field (for a general introduction, see e.g., Refs. , ). It is generally accepted that the basic reason for ferromagnetic order is the interplay between the kinetic energy and the Coulomb-interaction of the electrons. Nevertheless, it is still a matter of debate which kind of minimal model must be used for the description of ferromagnetic materials. Whatever an appropriate Hamilton may be, from a theoretical point of view one would expect that it will be a hard analytical task to find even an approximate solution for such a real many-particle problem.
The simplest theory, which gives an explanation for metallic ferromagnetism is the Hartree-Fock-Stoner theory . A surprising result of this theory is the statement that ferromagnetism occurs in any system provided that the product of Coulomb-interaction and the electrons’ density of states exceeds a certain amount. This statement is the famous Stoner criterion. This criterion usually leads to surprisingly small critical values of the Coulomb-interaction for which ferromagnetism is predicted to occur. At first sight one may consider this as an a posteriori justification of the Hartree-Fock theory, which certainly fails for stronger Coulomb interactions. However, in some simpler model systems like the one-band Hubbard model it is well known that ferromagnetism requires very large Coulomb-interactions if it exists at all. This is a definite contradiction to Hartree-Fock theory, which indicates that the whole Stoner picture may be inadequate for a thorough understanding of itinerant ferromagnetism. Similar objections could be raised against spin-density functional theory which, like Stoner-theory, is based on an effective one-particle description. Despite their conceptual shortcomings, these theories quite successfully describe some properties of the iron-group elements, i.e. the magnetic moment or the shapes of the multisheet Fermi surfaces . Therefore, a competitive strong-coupling theory must meet these apparent successes of spin-density functional theory before it will be taken seriously.
To this end, we recently introduced a variational treatment of multi-band Hubbard models with a general class of Gutzwiller wave functions . These models allow the description of real materials, for example the elements of the iron group. As a first application we studied the ferromagnetic transition in a two-band model. In contradiction to the Hartree-Fock theory we found that ferromagnetism requires quite large Coulomb interactions. In particular, we demonstrated the decisive role of the intra-atomic exchange interaction, which is found to be irrelevant in the Hartree-Fock approach. The Gutzwiller theory shows that finite values of this exchange interaction are essential for metallic ferromagnetism. Based on these results, we suggested that the complex atomic Coulomb-interaction has to be taken seriously in theories on itinerant ferromagnetism. In particular, it appears to be essential to take into account exchange interactions which form local spins according to Hund’s rule. We are presently calculating physical properties of the iron-group elements from our correlated electron approach . First results for Nickel show that our method is able to resolve all major discrepancies between experiment and spin-density functional theory.
Experiments not only provide information about ground-state physics, but also yield insight into dynamical properties of materials. It is found that metallic and insulating ferromagnets behave similar with respect to their low-energy spin excitations. In both cases, inelastic neutron-scattering experiments show pronounced gapless spin-wave excitations. The understanding of these excitations is very important since they govern the magnetic phase transition at finite temperatures. Theoretical methods which allow the determination of spin-wave dispersions in Hubbard models are quite rare. It is obvious that a convincing description of spin excitations can only be obtained starting from a qualified theory for the ground-state. However, almost all earlier theories on spin-wave excitations in ferromagnets are based on effective one-particle theories. Some of these theories lead to surprisingly accurate results compared to experiments. For example, in Ref. a random-phase approximation for iron and nickel was introduced. In other approaches, the spin-wave dispersion is determined via a mapping of the itinerant system to a ferromagnetic Heisenberg model (see, e.g. Ref. ).
In this work, we present a theory for spin-wave excitations in itinerant ferromagnets which is based on the variational ground states in Ref. . Our paper is organized as follows: In Sec. II we introduce the general class of multi-band Hubbard models and the corresponding Gutzwiller wave functions. Our general approach on the spin-wave problem is presented in Sec. III. In Sec. IV we evaluate the spin-wave dispersion for our general class of multi-band Hubbard models. Finally, we apply our results to a model with two degenerate bands in Sec. V. Short conclusions close our presentation. Technical details are deferred to four appendices.
## II Hamiltonian and variational wave function
### A Multi-band Hamiltonian
In this paper we consider the following general class of multi-band Hubbard models,
$$\widehat{H}=\underset{i()j}{}\underset{\sigma ,\sigma ^{}}{}t_{i,j}^{\sigma ,\sigma ^{}}\widehat{c}_{i;\sigma }^+\widehat{c}_{j;\sigma ^{}}^{}+\underset{i}{}\widehat{H}_{i;\text{at}}\widehat{H}_1+\widehat{H}_{\text{at}}.$$
(1)
Here, $`\widehat{c}_{i;\sigma }^+`$ creates an electron with combined spin-orbit index $`\sigma =1,\mathrm{},2N`$ ($`N=5`$ for 3$`d`$ electrons) at the lattice site $`i`$ of a solid. For simplicity, we assume that the orbitals do not belong to the same representation of the respective point-symmetry group. For example, in cubic symmetry this means that there is only one set of $`s,p,e_g`$ and $`t_{2g}`$-orbitals. In this case, one-particle-states $`|\mathrm{\Phi }_0`$, which respect the symmetry of the lattice, lead to vanishing non-diagonal local hopping-terms, i.e.
$$\mathrm{\Phi }_0\left|\widehat{c}_{i,\sigma }^+\widehat{c}_{i,\sigma ^{}}\right|\mathrm{\Phi }_0\delta _{\sigma ,\sigma ^{}}.$$
(2)
This relation simplifies the calculations in this paper, but there is no fundamental obstacle to extend our method to a more general case.
We further assume that the atomic Hamiltonian
$$\widehat{H}_{i;\text{at}}=\underset{\sigma }{}ϵ_\sigma \widehat{n}_\sigma +\underset{\sigma _1,\sigma _2,\sigma _3,\sigma _4}{}𝒰^{\sigma _1,\sigma _2;\sigma _3,\sigma _4}\widehat{c}_{\sigma _1}^+\widehat{c}_{\sigma _2}^+\widehat{c}_{\sigma _3}^{}\widehat{c}_{\sigma _4}^{}.$$
(3)
is site-independent and readily diagonalized,
$`\widehat{H}_{\text{at}}`$ $`=`$ $`{\displaystyle \underset{\mathrm{\Gamma }}{}}E_\mathrm{\Gamma }\widehat{m}_\mathrm{\Gamma },`$ (5)
$`\widehat{m}_\mathrm{\Gamma }`$ $`=`$ $`|\mathrm{\Gamma }\mathrm{\Gamma }|.`$ (6)
Here, we introduced the eigenvalues $`E_\mathrm{\Gamma }`$ and the eigenstates $`|\mathrm{\Gamma }`$ of $`\widehat{H}_{\text{at}}`$. The diagonalization of $`\widehat{H}_{\text{at}}`$ is a standard exercise (see e.g., Ref. ). Knowledge of the states $`|\mathrm{\Gamma }`$ means that we found their expansion
$$|\mathrm{\Gamma }=\underset{I}{}T_{I,\mathrm{\Gamma }}|I,$$
(7)
in the basis of the configuration states $`|I`$. In these states, a definite set of spin-orbit states $`\sigma `$ is occupied
$$|I=|\sigma _1,\sigma _2,\mathrm{}=\widehat{c}_{\sigma _1}^+\widehat{c}_{\sigma _2}^+\mathrm{}|\text{vacuum}(\sigma _1<\sigma _2<\mathrm{}).$$
(8)
For details about the notation, see Ref. , Sec. II.
### B Gutzwiller-wave-function and diagrammatic evaluation
In Ref. we proposed the following wave-function for a variational examination of the Hamiltonian (1):
$$|\mathrm{\Psi }_\text{G}=\widehat{P}_\text{G}|\mathrm{\Phi }_0=\underset{i}{}\widehat{P}_{i;\text{G}}|\mathrm{\Phi }_0.$$
(9)
Here, $`|\mathrm{\Phi }_0`$ is any normalized single-particle product state and the local Gutzwiller projector $`\widehat{P}_{i;\text{G}}`$ is defined as
$$\widehat{P}_{i;\text{G}}=\underset{\mathrm{\Gamma }}{}\lambda _\mathrm{\Gamma }^{\widehat{m}_\mathrm{\Gamma }}=1+\underset{\mathrm{\Gamma }}{}\left(\lambda _\mathrm{\Gamma }1\right)\widehat{m}_\mathrm{\Gamma }.$$
(10)
To simplify our notation, we suppress the spatial indices wherever a misunderstanding is impossible, e.g., we omitted the index $`i`$ on the rhs. of eq. (10). The real variational parameters $`\lambda _\mathrm{\Gamma }`$ may have values between zero and one, where these two limits generate the ground-state both in the uncorrelated ($`\widehat{H}_{\text{at}}=0`$) and the atomic limit ($`\widehat{H}_1=0`$) of our Hamiltonian (1).
In Ref. we showed that the expectation value of the Hamiltonian (1) can be evaluated for the wave-function (9) in the limit of infinite spatial dimensions. In this section we only summarize the main ideas of the diagrammatic derivation which are important for our treatment of the spin-wave-problem in the next chapters. For all details we refer the reader to Ref. .
First, let us consider the norm of the wave-function (9),
$$\mathrm{\Psi }_\text{G}|\mathrm{\Psi }_\text{G}=\underset{i}{}\mathrm{\Phi }_0\left|\widehat{P}_{i;\text{G}}^2\right|\mathrm{\Phi }_0.$$
(11)
The square of the local Gutzwiller-projector can be written as
$$\widehat{P}_{i;\text{G}}^2=1+\underset{I,I^{}(|I|,|I^{}|2)}{}x_{i;I,I^{}}\widehat{n}_{i;I,I^{}}^{\text{HF}},$$
(12)
where we introduced the (local) Hartree-Fock-operators
$`\widehat{n}_{I,I}^{\text{HF}}`$ $`=`$ $`{\displaystyle \underset{\sigma I}{}}\widehat{n}_\sigma ^{\text{HF}},`$ (14)
$`\widehat{n}_\sigma ^{\text{HF}}`$ $`=`$ $`\widehat{n}_\sigma n_\sigma ^0`$ (15)
for $`I=I^{}`$, and
$`\widehat{n}_{I,I^{}}^{\text{HF}}`$ $`=`$ $`\left[{\displaystyle \underset{\sigma J}{}}\widehat{n}_\sigma ^{\text{HF}}\right]\widehat{n}_{I_1,I_2}(J=II^{};I=JI_1;I^{}=JI_2)`$ (16)
$`\widehat{n}_{I_1,I_2}`$ $`=`$ $`{\displaystyle \underset{\sigma _1I_1}{}}\widehat{c}_{\sigma _1}^+{\displaystyle \underset{\sigma _2I_2}{}}\widehat{c}_{\sigma _2}`$ (17)
for $`II^{}`$. The operators $`\widehat{c}_{\sigma _1}^+`$ ($`\widehat{c}_{\sigma _2}`$) in (17) should be placed in an ascending (descending) order. An explicit expression for the coefficients $`x_{i;I,I^{}}`$ in (12) is derived in Appendix (B).
When we apply Wick’s-Theorem to the right-hand-side of eq. (11), all terms are represented by certain diagrams with lines
$$P_{i,j}^{\sigma ,\sigma ^{}}=\mathrm{\Phi }_0\left|\widehat{c}_{i,\sigma }^+\widehat{c}_{j,\sigma ^{}}\right|\mathrm{\Phi }_0$$
(18)
and local vertices $`x_{i;I,I^{}}`$. The special form of the operator $`\widehat{P}_{i;\text{G}}^2`$ in (12) has two essential consequences for the structure of our diagrams. First, the definition of the Hartree-Fock operators together with eq. (2) guarantees that there are no local lines, i.e.,
$$P_{i,i}^{\sigma ,\sigma ^{}}=0.$$
(19)
Second, the constraint $`|I|,|I^{}|2`$ requires that at least four lines meet at every local vertex. When we evaluate the expectation values in
$`\widehat{H}_{i;\text{at}}_{\mathrm{\Psi }_\text{G}}`$ $`=`$ $`{\displaystyle \underset{\mathrm{\Gamma }}{}}E_\mathrm{\Gamma }\widehat{m}_{i;\mathrm{\Gamma }}_{\mathrm{\Psi }_\text{G}}\text{ and}`$ (21)
$`\widehat{H}_1_{\mathrm{\Psi }_\text{G}}`$ $`=`$ $`{\displaystyle \underset{i,j;\sigma ,\sigma ^{}}{}}t_{i,j}^{\sigma ,\sigma ^{}}\widehat{c}_{i;\sigma }^+\widehat{c}_{j;\sigma ^{}}^{}_{\mathrm{\Psi }_\text{G}}`$ (22)
we obtain diagrams, which contain one (for $`\widehat{H}_{i;\text{at}}`$) or two (for $`\widehat{H}_1`$) external vertices for the lattice sites $`i`$ (and $`j`$). If such a diagram possesses at least one internal vertex, we have lattice sites, which are connected by more than two lines. Such diagrams vanish in infinite dimensions and therefore we concluded in Ref. that the expectation values (19) only include diagrams without any internal vertex. Thus, we can write the expectation values $`\widehat{m}_\mathrm{\Gamma }_{\mathrm{\Psi }_\text{G}}=\widehat{m}_{i;\mathrm{\Gamma }}_{\mathrm{\Psi }_\text{G}}`$ in (21) as
$$m_\mathrm{\Gamma }\widehat{m}_\mathrm{\Gamma }_{\mathrm{\Psi }_\text{G}}=\lambda _\mathrm{\Gamma }^2m_\mathrm{\Gamma }^0\lambda _\mathrm{\Gamma }^2\widehat{m}_\mathrm{\Gamma }_0$$
(23)
with uncorrelated expectation values $`\mathrm{}_0\mathrm{}_{\mathrm{\Phi }_0}`$. This relation allows to replace the original variational parameters $`\lambda _\mathrm{\Gamma }`$ by the new parameters $`m_\mathrm{\Gamma }`$. The expectation value for a hopping term in (22) becomes
$`\widehat{c}_{i;\sigma }^+\widehat{c}_{j;\sigma ^{}}^{}_{\mathrm{\Psi }_\text{G}}`$ $`=`$ $`\sqrt{q_\sigma q_\sigma ^{}}\widehat{c}_{i;\sigma }^+\widehat{c}_{j;\sigma ^{}}^{}_0`$ (25)
$`\sqrt{q_\sigma }`$ $``$ $`\sqrt{{\displaystyle \frac{1}{n_\sigma ^0(1n_\sigma ^0)}}}{\displaystyle \underset{\mathrm{\Gamma },\mathrm{\Gamma }^{}}{}}\sqrt{{\displaystyle \frac{m_\mathrm{\Gamma }m_\mathrm{\Gamma }^{}}{m_\mathrm{\Gamma }^0m_\mathrm{\Gamma }^{}^0}}}`$ (27)
$`\times {\displaystyle \underset{I,I^{}(\sigma I,I^{})}{}}f_\sigma ^If_\sigma ^I^{}\sqrt{m_{(I^{}\sigma )}^0m_I^{}^0}T_{\mathrm{\Gamma },(I\sigma )}^+T_{(I^{}\sigma ),\mathrm{\Gamma }}T_{\mathrm{\Gamma }^{},I^{}}^+T_{I,\mathrm{\Gamma }^{}}.`$
Here, the fermionic sign function
$$f_\sigma ^II\sigma |\widehat{c}_\sigma ^+|I$$
(28)
gives a minus (plus) sign if it takes an odd (even) number of anticommutations to shift the operator $`\widehat{c}_\sigma ^+`$ to its proper place in the sequence of electron creation operators in $`|I\sigma `$. Note that the numbers $`q_\sigma `$ in (25) are just the diagonal-elements of the matrix $`q_\sigma ^\sigma ^{}`$ introduced in Ref. , which is diagonal for our symmetry-restricted orbital basis (see Sec. II A).
## III Spin Waves
The theoretical examination of spin-wave excitations requires the analysis of the imaginary part $`\chi _T(\stackrel{}{q},E)`$ of the transversal susceptibility, which is given as the retarded two-particle Greenfunction
$`G_T(\stackrel{}{q},E)`$ $`=`$ $`{\displaystyle \frac{1}{L}}\widehat{S}_\stackrel{}{q}^+;\widehat{S}_\stackrel{}{q}^{}_E`$ (30)
$`=`$ $`{\displaystyle \frac{i}{L}}{\displaystyle _0^{\mathrm{}}}𝑑te^{iEt}\mathrm{\Psi }_0\left|[\widehat{S}_\stackrel{}{q}^+(t),\widehat{S}_\stackrel{}{q}^{}(0)]\right|\mathrm{\Psi }_0.`$ (31)
Here, we introduced the $`\stackrel{}{q}`$-dependent spin-flip operators
$`\widehat{S}_\stackrel{}{q}^+`$ $`=`$ $`{\displaystyle \underset{l}{}}e^{i\stackrel{}{q}\stackrel{}{R}_l}\widehat{S}_l^+={\displaystyle \underset{l,b}{}}e^{i\stackrel{}{q}\stackrel{}{R}_l}\widehat{c}_{l,b,}^+\widehat{c}_{l,b,}`$ (33)
$`\widehat{S}_\stackrel{}{q}^{}`$ $`=`$ $`(\widehat{S}_\stackrel{}{q}^+)^+={\displaystyle \underset{l,b}{}}e^{i\stackrel{}{q}\stackrel{}{R}_l}\widehat{c}_{l,b,}^+\widehat{c}_{l,b,}`$ (34)
in the Heisenberg-picture, where the sum includes all ($`L`$) lattice sites $`l`$ and orbitals $`b`$. The magnetic excitations of the system are represented by poles of the Greenfunction $`G_T(\stackrel{}{q},E)`$ with energies $`E>0`$. For our further analysis we expand the “spin-wave state”
$$|\mathrm{\Psi }_\stackrel{}{q}^0\widehat{S}_\stackrel{}{q}^{}|\mathrm{\Psi }_0$$
(35)
in terms of exact energy-eigenstates
$`|\mathrm{\Psi }_\stackrel{}{q}^0`$ $`=`$ $`{\displaystyle \underset{n}{}}W_n|\mathrm{\Psi }_n,`$ (36)
$`\widehat{H}|\mathrm{\Psi }_n`$ $`=`$ $`E_n|\mathrm{\Psi }_n.`$ (37)
The Lehmann-representation of (30),
$$G_T(\stackrel{}{q},E)=\frac{i}{L}\underset{n}{}\left[\frac{\mathrm{\Psi }_n\widehat{S}_\stackrel{}{q}^{}\mathrm{\Psi }_0^2}{E(E_nE_0)+i\delta }\frac{\mathrm{\Psi }_n\widehat{S}_\stackrel{}{q}^+\mathrm{\Psi }_0^2}{E+(E_nE_0)+i\delta }\right]$$
(38)
shows that there are poles in $`G_T(\stackrel{}{q},E)`$ for the energies $`E_nE_0>0`$ with weights $`\left|W_n\right|^2`$.
In a ferromagnetic system the state $`\mathrm{\Psi }_{\stackrel{}{q}=\stackrel{}{0}}^0`$ is also a ground state of $`\widehat{H}`$, since the operator $`\widehat{S}_{\stackrel{}{q}=\stackrel{}{0}}^{}`$ just flips a spin in the spin-multiplet of the ground-state $`|\mathrm{\Psi }_0`$. Therefore, we can conclude that $`G_T(\stackrel{}{0},E)`$ has one isolated pole for $`EE_0=0`$. Now we consider finite, but small values of $`\stackrel{}{q}`$, and assume that the expansion (36) is still dominated by a narrow distribution of low-energy states. This scenario explains the pronounced peak in $`\chi _T(\stackrel{}{q},E)`$ for small values of $`E`$ and $`\left|\stackrel{}{q}\right|`$, which is seen in experiments and interpreted as a spin-wave excitation (see, e.g., Ref. ). Then, the spin-wave dispersion $`E_\stackrel{}{q}`$ can be identified as the position of this peak, and $`E_\stackrel{}{q}`$ is approximately determined by the first moment of the distribution $`\left|W_n\right|^2`$,
$`E_\stackrel{}{q}`$ $`=`$ $`{\displaystyle \frac{_nE_n\left|W_n\right|^2}{_n\left|W_n\right|^2}}{\displaystyle \frac{\mathrm{\Psi }_0\left|\widehat{H}\right|\mathrm{\Psi }_0}{\mathrm{\Psi }_0\mathrm{\Psi }_0}}`$ (40)
$`=`$ $`{\displaystyle \frac{\mathrm{\Psi }_0\left|\widehat{S}_\stackrel{}{q}^+\widehat{H}\widehat{S}_\stackrel{}{q}^{}\right|\mathrm{\Psi }_0}{\mathrm{\Psi }_0\left|\widehat{S}_\stackrel{}{q}^+\widehat{S}_\stackrel{}{q}^{}\right|\mathrm{\Psi }_0}}{\displaystyle \frac{\mathrm{\Psi }_0\left|\widehat{H}\right|\mathrm{\Psi }_0}{\mathrm{\Psi }_0\mathrm{\Psi }_0}}.`$ (41)
It is still impossible to derive the spin-wave dispersion $`E_\stackrel{}{q}`$ from eq. (III) since we do not know the ground-state $`|\mathrm{\Psi }_0`$ of our multi-band Hamiltonian (1). If we assume, however, that the variational wave function $`|\mathrm{\Psi }_\text{G}`$ is a good approximation for the true ground-state $`|\mathrm{\Psi }_0`$ we may substitute $`|\mathrm{\Psi }_0`$ in eq. (III) by the variational wave function $`|\mathrm{\Psi }_\text{G}`$.
In the next section we will evaluate the “variational” spin-wave dispersion
$$E_\stackrel{}{q}^{var}=\frac{\mathrm{\Psi }_\text{G}\left|\widehat{S}_\stackrel{}{q}^+\widehat{H}\widehat{S}_\stackrel{}{q}^{}\right|\mathrm{\Psi }_\text{G}}{\mathrm{\Psi }_\text{G}\left|\widehat{S}_\stackrel{}{q}^+\widehat{S}_\stackrel{}{q}^{}\right|\mathrm{\Psi }_\text{G}}\frac{\mathrm{\Psi }_\text{G}\left|\widehat{H}\right|\mathrm{\Psi }_\text{G}}{\mathrm{\Psi }_\text{G}\mathrm{\Psi }_\text{G}}$$
(42)
in the limit of large spatial dimensions. It should be noted that this quantity obviously obeys no strict upper-bound properties. Nevertheless, we expect that $`E_\stackrel{}{q}^{\mathrm{exp}}<E_\stackrel{}{q}^{var}`$ is fulfilled since the expectation values (III) includes high-energy states which do not belong to the spin-wave excitation seen in experiments.
In principle, transversal spin-excitations are given as peaks both in $`\chi _T(\stackrel{}{q},E)`$ and $`\chi _T(\stackrel{}{q},E_0E)`$ for energies $`E>E_0`$. In other word, we also had to consider the contributions from the Green function $`\widehat{S}_\stackrel{}{q}^{};\widehat{S}_\stackrel{}{q}^+_E`$ in our calculation. These contributions are identical to the second term in eq. (38) and we could include them by using the proper spin-wave state
$$|\stackrel{~}{\mathrm{\Psi }}_\stackrel{}{q}^0\left(\widehat{S}_\stackrel{}{q}^{}+\widehat{S}_\stackrel{}{q}^+\right)|\mathrm{\Psi }_G$$
(43)
in our variational approach. However, the contributions from the second operator in (43) vanish for $`\stackrel{}{q}=\stackrel{}{0}`$ and may be neglected for small values of $`|\stackrel{}{q}|`$, where spin-wave excitations are observed in experiments. Nevertheless, there is no fundamental obstacle to extend our diagrammatic approach to the more general spin-wave state (43).
## IV Variational Spin-Wave Dispersion
### A General considerations
In order to determine the variational spin-wave dispersion (42) we need to examine the norm of the state $`|\mathrm{\Psi }_\stackrel{}{q}^\text{G}\widehat{S}_\stackrel{}{q}^{}|\mathrm{\Psi }_\text{G},`$
$$N_\stackrel{}{q}\mathrm{\Psi }_\text{G}\left|\widehat{S}_\stackrel{}{q}^+\widehat{S}_\stackrel{}{q}^{}\right|\mathrm{\Psi }_\text{G},$$
(44)
and the expectation values
$`{\displaystyle \frac{\mathrm{\Psi }_\stackrel{}{q}^\text{G}\left|\widehat{H}_{\text{at}}\right|\mathrm{\Psi }_\stackrel{}{q}^\text{G}}{N_\stackrel{}{q}}}`$ $`=`$ $`{\displaystyle \underset{\mathrm{\Gamma }}{}}E_\mathrm{\Gamma }{\displaystyle \underset{i,j,k}{}}e^{i\stackrel{}{q}(\stackrel{}{R}_i\stackrel{}{R}_j)}{\displaystyle \frac{\mathrm{\Psi }_\text{G}\left|\widehat{S}_i^+\widehat{m}_{k;\mathrm{\Gamma }}\widehat{S}_j^{}\right|\mathrm{\Psi }_\text{G}}{N_\stackrel{}{q}}},`$ (46)
$`{\displaystyle \frac{\mathrm{\Psi }_\stackrel{}{q}^\text{G}\left|\widehat{H}_1\right|\mathrm{\Psi }_\stackrel{}{q}^\text{G}}{N_\stackrel{}{q}}}`$ $`=`$ $`{\displaystyle \underset{i,j,k,l}{}}e^{i\stackrel{}{q}(\stackrel{}{R}_i\stackrel{}{R}_j)}{\displaystyle \underset{\sigma _k,\sigma _l}{}}t_{k,l}^{\sigma _k,\sigma _l}{\displaystyle \frac{\mathrm{\Psi }_\text{G}\left|\widehat{S}_i^+\widehat{c}_{k;\sigma _k}^+\widehat{c}_{l;\sigma _l}\widehat{S}_j^{}\right|\mathrm{\Psi }_\text{G}}{N_\stackrel{}{q}}}.`$ (47)
Before we start to evaluate these quantities, it is necessary to discuss two general problems. First, let us consider the norm
$$N_\stackrel{}{q}=\underset{i,j,b,b^{}}{}e^{i\stackrel{}{q}(\stackrel{}{R}_i\stackrel{}{R}_j)}\mathrm{\Psi }_\text{G}\left|\widehat{c}_{i,b,}^+\widehat{c}_{i,b,}\widehat{c}_{j,b^{},}^+\widehat{c}_{j,b^{},}\right|\mathrm{\Psi }_\text{G}$$
(48)
in the special case $`\stackrel{}{q}=\stackrel{}{0}`$, where $`\widehat{S}_{\stackrel{}{q}=\stackrel{}{0}}^{}`$ is just the total spin-flip operator $`\widehat{S}^{}`$. When we assume that $`|\mathrm{\Psi }_\text{G}`$ has the correct spin-symmetry, i.e. it is an eigenstate of
$$\widehat{S}^z=\underset{i}{}\widehat{S}_{i,z}\text{and}\widehat{\stackrel{}{S}}^2=\left(\underset{i}{}\stackrel{}{S}_i\right)^2$$
(49)
with eigenvalues $`S_\text{G}^z`$ and $`S_\text{G}^z(S_\text{G}^z+1)`$, respectively, we obtain
$$N_{\stackrel{}{q}=\stackrel{}{0}}=2S_\text{G}^z\mathrm{\Psi }_\text{G}\mathrm{\Psi }_\text{G}.$$
(50)
Here, we used the well-known equation
$$\widehat{S}^+\widehat{S}^{}=\stackrel{}{S}^2\widehat{S}^z(\widehat{S}^z1)$$
(51)
for spin operators. In general, however, the wave functions $`|\mathrm{\Psi }_\text{G}`$, as defined in eq. (10), do not fulfill this symmetry. Therefore, it is necessary to introduce some additional constraints on our variational parameters $`\lambda _\mathrm{\Gamma }`$ in (10) to guarantee that $`|\mathrm{\Psi }_\text{G}`$ is an eigenstate of $`\widehat{\stackrel{}{S}}^2`$. In Appendix (A) we explain how these relations may be chosen.
The second problem arises from the evaluation of $`\stackrel{}{q}`$-dependent quantities in the limit of large spatial dimensions $`D`$. For example, let us consider the Hartree-Fock case, where $`|\mathrm{\Psi }_\text{G}=|\mathrm{\Phi }_0`$ is a spin-polarized one-particle state with $`n_{}>n_{}`$. We find
$`{\displaystyle \frac{N_\stackrel{}{q}^0}{L}}`$ $`=`$ $`{\displaystyle \frac{1}{L}}{\displaystyle \underset{b,b^{}}{}}{\displaystyle \underset{i,j}{}}e^{i\stackrel{}{q}(\stackrel{}{R}_i\stackrel{}{R}_j)}\mathrm{\Phi }_0\left|\widehat{c}_{i,b,}^+\widehat{c}_{i,b,}\widehat{c}_{j,b^{},}^+\widehat{c}_{j,b^{},}\right|\mathrm{\Phi }_0`$ (53)
$`=`$ $`{\displaystyle \underset{b}{}}\left(n_{b,}^0n_{b,}^0n_{b,}^0+{\displaystyle \frac{1}{L}}{\displaystyle \underset{ij}{}}e^{i\stackrel{}{q}(\stackrel{}{R}_i\stackrel{}{R}_j)}P_{i,j}^{(b,),(b,)}P_{j,i}^{(b,),(b,)}\right),`$ (54)
where, for simplicity, we assumed that the expectation-values $`P_{i,j}^{\sigma ,\sigma ^{}}`$ as defined in (18) are diagonal with respect to the orbitals $`b,b^{}`$. For any finite value of $`\stackrel{}{q}`$ the sum in eq. (54) vanishes as $`1/D`$. This means that the limits $`\stackrel{}{q}0`$ and $`D\mathrm{}`$ do not commute, because
$$\underset{\stackrel{}{q}\stackrel{}{0}}{lim}\underset{D\mathrm{}}{lim}\frac{N_\stackrel{}{q}^0}{L}=\underset{b}{}n_{b,}^0\left(1n_{b,}^0\right)\underset{b}{}\left(n_{b,}^0n_{b,}^0\right)=\underset{D\mathrm{}}{lim}\underset{\stackrel{}{q}\stackrel{}{0}}{lim}\frac{N_\stackrel{}{q}^0}{L}.$$
(55)
To overcome this problem we will evaluate expressions like the sum in (54) using the realistic three dimensional band-structure of our Hamiltonian (1). This leads to results which are continuous in $`\stackrel{}{q}`$ and, consequently, reproduce the limit $`\stackrel{}{q}=\stackrel{}{0}`$ correctly .
The norm of the state (44) for a correlated wave function $`|\mathrm{\Psi }_\text{G}|\mathrm{\Phi }_0`$ will contain diagrams of an arbitrary order $`1/D^n`$. In this paper we will only consider diagrams up to the leading order $`n=1`$. At first sight one may wonder whether or not $`1/D`$-terms need to be included in the ground-state energy expression as well. Fortunately, to order $`1/D`$, these diagrams only lead to a constant shift of the energy expectation-values for $`|\mathrm{\Psi }_\text{G}`$ and $`|\mathrm{\Psi }_\stackrel{}{q}^\text{G}`$. Thus, for our calculation we may neglect the $`1/D`$-contributions to the ground-state energy since we are only interested in the difference between these two energies.
In the next subsection we evaluate the norm (48). The derivation of the expectation values (46) and (47) requires no additional techniques. The cumbersome calculations are deferred to appendices C and D.
### B Evaluation of the norm
The norm of the state (44) can be written as
$`N_\stackrel{}{q}`$ $`=`$ $`{\displaystyle \underset{i}{}}{\displaystyle \underset{l(i)}{}}\left(\widehat{P}_{i;\text{G}}\widehat{S}_i^+\widehat{S}_i^{}\widehat{P}_{i;\text{G}}\right)\widehat{P}_{l;\text{G}}^2_{\mathrm{\Phi }_0}`$ (57)
$`+{\displaystyle \underset{ij}{}}e^{i\stackrel{}{q}(\stackrel{}{R}_i\stackrel{}{R}_j)}{\displaystyle \underset{l(i,j)}{}}\left(\widehat{P}_{i;\text{G}}\widehat{S}_i^+\widehat{P}_{i;\text{G}}\right)\left(\widehat{P}_{j;\text{G}}\widehat{S}_j^{}\widehat{P}_{j;\text{G}}\right)\widehat{P}_{l;\text{G}}^2_{\mathrm{\Phi }_0}.`$
When we use eq. (12) and apply Wick’s theorem, we obtain diagrams with external vertices for the lattice sites $`i`$ and $`j`$ and internal vertices generated by the Hartree-Fock operators $`x_{l;I,I^{}}\widehat{n}_{l;I,I^{}}^{\text{HF}}`$ in $`\widehat{P}_{l;\text{G}}^2`$. In Refs. it was shown that we only have to evaluate the connected diagrams since the unconnected terms just give the norm of the Gutzwiller wave-function $`N_\text{G}\mathrm{\Psi }_\text{G}\mathrm{\Psi }_\text{G}`$. For the evaluation of the first line (57) we use the local relations
$`\widehat{S}^+\widehat{S}^{}`$ $`=`$ $`{\displaystyle \underset{\mathrm{\Gamma }}{}}S_+(\mathrm{\Gamma })\widehat{m}_\mathrm{\Gamma },`$ (59)
$`S_\pm (\mathrm{\Gamma })`$ $``$ $`S(\mathrm{\Gamma })\left[S(\mathrm{\Gamma })+1\right]S_z(\mathrm{\Gamma })\left[S_z(\mathrm{\Gamma })1\right]`$ (60)
which follow from eq. (51). Here, we introduced the total spin $`S(\mathrm{\Gamma })`$ and the spin-component $`S_z(\mathrm{\Gamma })`$ of the atomic eigenstates $`|\mathrm{\Gamma }`$. Hence, the expectation value of $`\widehat{S}_i^+\widehat{S}_i^{}`$ is given as a linear function of the variational parameters $`m_\mathrm{\Gamma }`$ and we may write (57) as
$`{\displaystyle \frac{N_\stackrel{}{q}}{LN_\text{G}}}`$ $`=`$ $`{\displaystyle \underset{\mathrm{\Gamma }}{}}S_{}(\mathrm{\Gamma })m_\mathrm{\Gamma }`$ (63)
$`+{\displaystyle \frac{1}{L}}{\displaystyle \underset{ij}{}}e^{i\stackrel{}{q}(\stackrel{}{R}_i\stackrel{}{R}_j)}{\displaystyle \underset{l(i,j)}{}}\left\{\left(\widehat{P}_{i;\text{G}}\widehat{S}_i^+\widehat{P}_{i;\text{G}}\right)\left(\widehat{P}_{j;\text{G}}\widehat{S}_j^{}\widehat{P}_{j;\text{G}}\right)\widehat{P}_{l;\text{G}}^2\right\}_{\mathrm{\Phi }_0}^c,`$
where $`\left\{\mathrm{}\right\}_{\mathrm{\Phi }_0}^c`$ denotes the application of Wick’s theorem and taking into account only the connected diagrams.
For a further analysis of (63) we introduce indices $`𝔇=(\sigma _1\sigma _2)`$ for pairs of spin-orbit states $`\sigma _i`$ and the basic RPA-diagrams
$$\stackrel{~}{P}_𝔇^𝔇^{}(\stackrel{}{q})=\stackrel{~}{P}_{(\sigma _3\sigma _4)}^{(\sigma _1\sigma _2)}(\stackrel{}{q})\frac{1}{L}\underset{ij}{}e^{i\stackrel{}{q}(\stackrel{}{R}_i\stackrel{}{R}_j)}P_{i,j}^{\sigma _1\sigma _4}P_{j,i}^{\sigma _3\sigma _2}.$$
(64)
$`\stackrel{~}{P}_𝔇^𝔇^{}(\stackrel{}{q})`$ can be evaluated in momentum space as
$$\stackrel{~}{P}_𝔇^𝔇^{}(\stackrel{}{q})=\delta _{\sigma _1}^{\sigma _4}\delta _{\sigma _3}^{\sigma _2}n_{\sigma _1}^0n_{\sigma _2}^0\frac{1}{L}\underset{\stackrel{}{k}}{}n_\stackrel{}{k}^{\sigma _1\sigma _4}n_{\stackrel{}{k}+\stackrel{}{q}}^{\sigma _3\sigma _2}$$
(65)
with expectation values
$$n_\stackrel{}{k}^{\sigma \sigma ^{}}\widehat{c}_{\stackrel{}{k}\sigma }^+\widehat{c}_{\stackrel{}{k}\sigma ^{}}_{\mathrm{\Phi }_0}$$
(66)
and a modified Kronecker-symbol $`\delta _\sigma ^\sigma ^{}=\delta _{\sigma ,\sigma ^{}}`$. Note that, in contrast to the indices $`I`$, the order of the two spin-orbit states in $`𝔇=(\sigma _1\sigma _2)`$ is significant. Here, its first and second element specify a particle which enters or leaves a vertex, respectively.
For large spatial dimensions (i.e., up to the order $`1/D`$), the only contributions in (63) are RPA-type diagrams as shown in Fig (1). The internal vertices $`\stackrel{~}{x}_𝔇^𝔇^{}`$ are generated by the operators $`\widehat{P}_{l;\text{G}}^2`$, which can be expressed in terms of Hartree-Fock operators $`x_{l;I,I^{}}\widehat{n}_{l;I,I^{}}^{\text{HF}}`$ (see eq. (12)). For our RPA-diagrams we only have to consider two-fermion-operators in $`\widehat{n}_{l;I,I^{}}^{\text{HF}}`$. A general vertex with $`n`$ incoming and outgoing lines is determined by the operator
$`{\displaystyle \underset{I_1,I_2(I_1I_2=\mathrm{})}{}}{\displaystyle \underset{J(I_1,I_2J)}{}}x_{JI_1,JI_2}\left\{\mathrm{}\widehat{n}_{JI_1,JI_2}^{\text{HF}}\mathrm{}\right\}_{\mathrm{\Phi }_0}^c`$ (67)
$``$ $`{\displaystyle \underset{I_1,I_2(I_1I_2=\mathrm{})}{}}{\displaystyle \underset{J(I_1,I_2J)}{}}y_{JI_1,JI_2}\left\{\mathrm{}\widehat{n}_{JI_1,JI_2}\mathrm{}\right\}_{\mathrm{\Phi }_0}^c,`$ (68)
with
$`y_{JI_1,JI_2}=x_{JI_1,JI_2}f_{I_1}^Jf_{I_2}^J.`$Thus, the internal vertex $`\stackrel{~}{x}_𝔇^𝔇^{}=\stackrel{~}{x}_{(\sigma _3\sigma _4)}^{(\sigma _1\sigma _2)}`$ at lattice site $`l`$ stems from the term
$$y_{l;(\sigma _1,\sigma _3)(\sigma _2,\sigma _4)}\left\{\mathrm{}\left(\widehat{c}_{l,\sigma _1}^+\widehat{c}_{l,\sigma _2}\right)\left(\widehat{c}_{l,\sigma _3}^+\widehat{c}_{l,\sigma _4}\right)\mathrm{}\right\}_{\mathrm{\Phi }_0}^cf_{\sigma _1}^{\sigma _3}f_{\sigma _2}^{\sigma _4}.$$
(69)
The round brackets indicate that Wick’s theorem may not be applied to operators in different brackets. Thus, we can identify $`\stackrel{~}{x}_𝔇^𝔇^{}`$ as
$$\stackrel{~}{x}_𝔇^𝔇^{}=\stackrel{~}{x}_{(\sigma _3\sigma _4)}^{(\sigma _1\sigma _2)}=f_{\sigma _2}^{\sigma _4}f_{\sigma _1}^{\sigma _3}x_{(\sigma _1,\sigma _3),(\sigma _2,\sigma _4)}.$$
(70)
An explicit expression for $`x_{I,I^{}}`$ is derived in Appendix B.
Our infinite RPA-sum (see Fig 1) for the second term in (63) leads to the following matrix-equation for $`\stackrel{~}{\mathrm{\Omega }}_𝔇^𝔇^{}(\stackrel{}{q})`$,
$`\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})`$ $``$ $`\stackrel{~}{P}(\stackrel{}{q})+\stackrel{~}{P}(\stackrel{}{q})\stackrel{~}{x}\stackrel{~}{P}(\stackrel{}{q})+\stackrel{~}{P}(\stackrel{}{q})\stackrel{~}{x}\stackrel{~}{P}(\stackrel{}{q})\stackrel{~}{x}\stackrel{~}{P}(\stackrel{}{q})+\mathrm{}`$ (72)
$`=`$ $`\stackrel{~}{P}(\stackrel{}{q})+\stackrel{~}{P}(\stackrel{}{q})\stackrel{~}{x}\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q}),`$ (73)
which has the solution
$$\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})=(\widehat{1}\stackrel{~}{P}(\stackrel{}{q})\stackrel{~}{x})^1\stackrel{~}{P}(\stackrel{}{q}).$$
(74)
Here, the dot “$``$”indicates the usual product between two matrices.
Finally, we have to examine the external vertices, generated by the spin-operators $`\widehat{S}_i^+`$ and $`\widehat{S}_j^{}`$ in (63). First, we find for local operators
$`\widehat{P}_\text{G}S^+\widehat{P}_\text{G}`$ $`=`$ $`{\displaystyle \underset{\mathrm{\Gamma }^{},\mathrm{\Gamma }}{}}\lambda _\mathrm{\Gamma }^{}\lambda _\mathrm{\Gamma }|\mathrm{\Gamma }^{}\mathrm{\Gamma }^{}\left|S^+\right|\mathrm{\Gamma }\mathrm{\Gamma }|`$ (76)
$`=`$ $`{\displaystyle \underset{\mathrm{\Gamma }}{}}\lambda _{\mathrm{\Gamma }_+}\lambda _\mathrm{\Gamma }\sqrt{S_+(\mathrm{\Gamma })}|\mathrm{\Gamma }_+\mathrm{\Gamma }|,`$ (77)
where the normalized spin-flip-state
$$|\mathrm{\Gamma }_+\frac{1}{\sqrt{S_+(\mathrm{\Gamma })}}\widehat{S}^+|\mathrm{\Gamma }$$
(78)
was introduced. We may write eq. (77) as
$$\widehat{P}_\text{G}S^+\widehat{P}_\text{G}=\underset{\mathrm{\Gamma }}{}\lambda _{\mathrm{\Gamma }_+}\lambda _\mathrm{\Gamma }\sqrt{S_+(\mathrm{\Gamma })}\underset{I_1,I_2}{}T_{I_1,\mathrm{\Gamma }_+}T_{\mathrm{\Gamma },I_2}^+\widehat{m}_{I_1,I_2}$$
(79)
with
$`\widehat{m}_{I_1,I_2}`$ $`=`$ $`f_I^Jf_I^{}^J\widehat{m}_J^{I,I^{}}\widehat{n}_{I,I^{}},`$ (81)
$`\widehat{m}_J^{I,I^{}}`$ $``$ $`{\displaystyle \underset{\sigma J}{}}\widehat{n}_\sigma {\displaystyle \underset{\sigma \overline{J}\backslash (II^{})}{}}\left(1\widehat{n}_\sigma \right)`$ (82)
and $`J=I_1I_2`$, $`I_1=JI`$, $`I_2=JI^{}`$.
The external vertices in our RPA-diagrams are met by just two lines. Therefore, we only need the following one-particle contribution of the operator $`\widehat{m}_{I_1,I_2}`$ in (79),
$`\left\{\mathrm{}\widehat{m}_{I_1,I_2}\mathrm{}\right\}_{\mathrm{\Phi }_0}^C`$ $``$ $`{\displaystyle \underset{\sigma _1,\sigma _2}{}}𝔙_{_{I_1,I_2}}^{(\sigma _1\sigma _2)}\left\{\mathrm{}\widehat{c}_{\sigma _1}^+\widehat{c}_{\sigma _2}\mathrm{}\right\}_{\mathrm{\Phi }_0}^c`$ (83)
$`𝔙_{_{I_1,I_2}}^{(\sigma _1\sigma _2)}`$ $``$ $`{\displaystyle \underset{J(\sigma _1,\sigma _2J)}{}}f_{\sigma _1}^Jf_{\sigma _2}^Jm_^J^0\left({\displaystyle \underset{\sigma (\sigma _1,\sigma _2)}{}}{\displaystyle \frac{1}{(1n_\sigma ^0)}}\right)\left[\delta _{I_1}^{J\sigma _1}\delta _{I_2,}^{J\sigma _2}\delta _{\sigma _1}^{\sigma _2}\delta _{I_1}^J\delta _{I_2}^J\right]`$ (84)
where we assumed that $`\left|I_1\right|=\left|I_2\right|`$. When we apply this result to the contribution of the external vertex (63), we only need to consider the first term with $`\sigma _1\sigma _2`$, because $`\sigma _1,\sigma _2`$ must have different spins. Thus, the expression for the external vertex, which stems from $`\widehat{S}_i^+`$ becomes
$`S_{(\sigma _1\sigma _2)}^+`$ $`=`$ $`{\displaystyle \underset{\mathrm{\Gamma }}{}}\lambda _{\mathrm{\Gamma }_+}\lambda _\mathrm{\Gamma }\sqrt{S_+(\mathrm{\Gamma })}{\displaystyle \underset{I_1,I_2}{}}T_{I_1,\mathrm{\Gamma }_+}T_{\mathrm{\Gamma },I_2}^+𝔙_{_{I_1,I_2}}^{(\sigma _1\sigma _2)}`$ (86)
$`=`$ $`{\displaystyle \underset{\mathrm{\Gamma }}{}}\lambda _{\mathrm{\Gamma }_+}\lambda _\mathrm{\Gamma }\sqrt{S_+(\mathrm{\Gamma })}{\displaystyle \underset{J(\sigma _1,\sigma _2J)}{}}T_{J\sigma _1,\mathrm{\Gamma }_+}T_{\mathrm{\Gamma },J\sigma _2}^+{\displaystyle \frac{f_{\sigma _1}^Jf_{\sigma _2}^Jm_^J^0}{(1n_{\sigma _1}^0)(1n_{\sigma _2}^0)}}.`$ (87)
The expansion
$`\left\{\mathrm{}\widehat{P}_\text{G}\widehat{S}^{}\widehat{P}_\text{G}\mathrm{}\right\}_{\mathrm{\Phi }_0}^C`$ $``$ $`{\displaystyle \underset{\sigma _1,\sigma _2}{}}\left[S_{(\sigma _1\sigma _2)}^+\left\{\mathrm{}\left(\widehat{c}_{\sigma _1}^+\widehat{c}_{\sigma _2}\right)\mathrm{}\right\}_{\mathrm{\Phi }_0}^c\right]^{}`$ (89)
$`=`$ $`{\displaystyle \underset{\sigma _1,\sigma _2}{}}\left(S_{(\sigma _1\sigma _2)}^+\right)^{}\left\{\mathrm{}\left(\widehat{c}_{\sigma _2}^+\widehat{c}_{\sigma _1}\right)\mathrm{}\right\}_{\mathrm{\Phi }_0}^c,`$ (90)
shows that the vertex-factor for the operator $`\widehat{S}_j^{}`$ in (63) is given as
$$S_{(\sigma _1\sigma _2)}^{}=\left(S_{(\sigma _2\sigma _1)}^+\right)^{}.$$
(91)
When we consider $`S_𝔇^+`$ and $`S_𝔇^{}`$ as components of vectors $`\stackrel{}{S}^+`$ and $`\stackrel{}{S}^{}`$ with respect to the indices $`𝔇`$ we obtain the following final result for the norm $`N_\stackrel{}{q}`$
$`{\displaystyle \frac{N_\stackrel{}{q}}{LN_\text{G}}}`$ $`=`$ $`{\displaystyle \underset{\mathrm{\Gamma }}{}}S_{}(\mathrm{\Gamma })m_\mathrm{\Gamma }+{\displaystyle \underset{𝔇_1,𝔇_2}{}}S_{𝔇_1}^+\stackrel{~}{\mathrm{\Omega }}_{𝔇_2}^{𝔇_1}(\stackrel{}{q})S_{𝔇_2}^{}`$ (93)
$`=`$ $`{\displaystyle \underset{\mathrm{\Gamma }}{}}S_{}(\mathrm{\Gamma })m_\mathrm{\Gamma }+\stackrel{}{S}^+\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\stackrel{}{S}^{}.`$ (94)
Note that the norm $`N_\text{G}`$ in (93) will cancel out when we calculate the expectation values (44).
A similar derivation gives us the expectation values (44) of the kinetic energy and the Coulomb-interaction, see appendices C and D. There, the number of contributing diagrams is much larger (e.g., $`18`$ for the kinetic energy) since we have up to four external vertices. Nevertheless, the numerical evaluation of these terms is a minor technical problem as soon as the wave-function $`|\mathrm{\Psi }_\text{G}`$ has been determined by the minimization of our variational energy expression. First studies show that the application of our variational scheme to iron and nickel represents a solvable numerical task . Further work in this direction is in progress.
## V Application to a two-band model
In this chapter we will present the results for the spin-wave properties in a system with two degenerate $`e_g`$-type orbitals per lattice site. The appearance of ferromagnetism in this model has already been discussed in Ref. , both for the Hartree-Fock and the Gutzwiller theory. Thus, we will only summarize these results here, before we start to consider the spin-wave dispersion. For all details, we refer to Ref. .
### A Ferromagnetic properties
In a system with two degenerate orbitals the general atomic Hamiltonian (3) becomes
$`\widehat{H}_{\text{at}}`$ $`=`$ $`U{\displaystyle \underset{b}{}}\widehat{n}_{b,}\widehat{n}_{b,}+U^{}{\displaystyle \underset{\sigma ,\sigma ^{}}{}}\widehat{n}_{1,\sigma }\widehat{n}_{2,\sigma ^{}}J{\displaystyle \underset{\sigma }{}}\widehat{n}_{1,\sigma }\widehat{n}_{2,\sigma }`$ (96)
$`+J{\displaystyle \underset{\sigma }{}}\widehat{c}_{1,\sigma }^+\widehat{c}_{2,\sigma }^+\widehat{c}_{1,\sigma }^{}\widehat{c}_{2,\sigma }^{}+J_\text{C}\left(\widehat{c}_{1,}^+\widehat{c}_{1,}^+\widehat{c}_{2,}^{}\widehat{c}_{2,}^{}+\widehat{c}_{2,}^+\widehat{c}_{2,}^+\widehat{c}_{1,}^{}\widehat{c}_{1,}^{}\right).`$
In cubic symmetry the Coulomb- and exchange-integrals $`U`$, $`U^{}`$, $`J`$ and $`J_\text{C}`$ are not independent from each other. Instead we have only two free parameters, because the relations $`J=J_\text{C}`$ and $`UU^{}=2J`$ hold. For the determination of the optimal variational wave function (9), we need the $`16`$ eigenstates of the atomic Hamiltonian. The latter can be found in table I of Ref. .
In the one-particle Hamiltonian $`\widehat{H}_1`$ we take into account first and second nearest neighbor hopping matrix elements. Furthermore, we apply the two-center approximation for the hopping matrix elements, which are chosen as $`t_{dd\sigma }^{(1)}=1`$eV, $`t_{dd\sigma }^{(2)}=0.25\text{eV}`$, and $`t_{dd\sigma }^{(1),(2)}:t_{dd\pi }^{(1),(2)}:t_{dd\delta }^{(1),(2)}=1:(0.3):0.1`$ (see Ref. for the notation). The density of states for these parameters (see fig. 1 of Ref. ) has a pronounced peak for the particle-density $`n_\sigma 0.3`$. For this band-filling we observe the strongest tendency to generate a ferromagnetic order, in qualitative agreement with the simple Stoner criterion. Therefore, we consider the spin-wave properties of our system only for this optimum band-filling. The numerical evaluation in Ref. did not respect the global spin-symmetry of the wave-function $`|\mathrm{\Psi }_\text{G}`$. In this work, we include the additional constraints on the variational-parameters $`\lambda _\mathrm{\Gamma }`$ as described in Appendix A, and we obtain almost the same state as in Ref. . In other words, in the two-band-model even the general variational ground-state $`|\mathrm{\Psi }_\text{G}`$, as defined in (9), is a nearly perfect eigenstate of the global spin-operator. However, it is not clear so far, if this statement also holds for a more general multi-band model.
In Ref. we found significant differences between the Hartree-Fock and the Gutzwiller theory. These differences have to be interpreted as a failure of the Hartree-Fock-theory since the variational space of the Hartree-Fock theory is included in our general class of Gutzwiller wave-functions. Although this statements holds for the corresponding RPA-theory as well, this theory is generally considered as the standard method in the context of spin-waves in itinerant ferromagnets.
The ferromagnetic phase diagrams for our model is shown in Fig. 3. In the Hartree-Fock theory ferromagnetism occurs for considerably smaller values of the correlation parameters. The most important difference lies in the role of the interatomic exchange $`J`$. In the Gutzwiller theory a ferromagnetic ground-state exists only for finite values of $`J`$. This is completely different in the Hartree-Fock theory. There, the only relevant quantities are the Stoner-parameter $`I=\frac{U+J}{2}`$ and the density of states at the Fermi-level $`𝔇_0(E_F)`$ which enter the Stoner criterion
$$I𝔇_0(E_F)>1.$$
(97)
This means that even for $`J=0`$ finite values of $`U`$ exist, where the system makes a transition into a ferromagnetic state.
Another important difference between both theories occurs when we compare the condensation energies $`E_{\text{cond}}`$, i.e., the differences between the energies in the paramagnetic and the ferromagnetic ground states. This quantity should have the same order of magnitude as the Curie-temperature $`T_\text{C}`$ in itinerant ferromagnets. In Stoner theory we observe that $`E_{\text{cond}}`$ is typically of the order of $`U`$ and therefore much larger than $`T_\text{C}`$. This is in agreement with the general observation that mean field methods overestimate the stability of magnetic order. On the other hand, in our correlated electron approach we find relatively small values for the condensation energy even for interaction parameters as large as twice the bandwidth.
### B Spin-waves
Fig. 2 shows the spin-wave dispersion in $`(100)`$-direction for four different magnetizations. As the lattice constant of our simple-cubic lattice we chose $`a=2.5\AA `$, the next neighbor distance in Nickel. Our Gutzwiller theory shows that the dispersion strongly depends on the magnetization, especially for small values of $`\stackrel{}{q}`$. The line shows the respective fits for the function
$$E_\stackrel{}{q}=Dq^2(1\beta q^2).$$
(98)
Note that experimental results of the region where quartic corrections become dominant are usually not very accurate, since Stoner excitations reduce the lifetime of spin-waves significantly. The values for the spin-wave stiffness $`D=1.4\mathrm{eV}\mathrm{\AA }^2`$ and $`D=1.2\mathrm{eV}\mathrm{\AA }^2`$ in both cases of almost fully polarized magnetizations, $`m=0.26`$ and $`m=0.28`$, respectively, are of the right order of magnitude for Nickel where $`D=0.43eV\AA ^2`$.
The spin-wave dispersion is almost isotropic, as can be seen for example in the inset of Fig. 2, where the dispersion is shown in the directions $`(100)`$ and $`(110)`$. Such an isotropic behaviour was also observed in experiments on iron and nickel and it is actually somewhat surprising, since the band-structure in these materials is far from being isotropic. However, we should have in mind that even in a metallic system local moments are formed due to the electrons’ correlations. Therefore, spin-excitations may be interpreted as spin-fluctuations in a system with localized spins. In such a system we have a generic isotropy when the exchange coupling is dominated by terms between nearest neighbors.
Our results show that the magnetic excitations in strong itinerant ferromagnets behave very similar to those in systems of localized spins. These low-energy excitations are responsible for the magnetic phase transition which occurs for temperatures much smaller than the typical Fermi energies in itinerant electron systems. This observation is consistent with the small condensation energy in our variational approach. To support this statement let us consider a Heisenberg-model
$`\widehat{H}_S=J{\displaystyle \underset{i,j}{}}\stackrel{}{S}_i\stackrel{}{S}_j`$on a cubic lattice. In this system we have a spin-wave stiffness $`D=2SJa^2`$. The value of the effective local moment in our itinerant system is given as $`S0.6`$ (see Ref. ). Therefore we obtain $`J=\frac{D}{2Sa^2}0.16\mathrm{eV}`$. For an estimate of the Curie-temperature we use the results from quantum Monte-Carlo calculations $`T_C=1.44JS^2`$. In this way we find $`T_C810^2K`$ which is the same order of magnitude as the condensation energy $`E_{\text{cond}}510^2K`$. Thus, we can summarize that our variational approach gives a consistent picture of both of magnetic excitations and ground-state properties in strong itinerant ferromagnets.
When the spin-waves are treated as non-interacting bosons, their contribution to the specific heat and the magnetization $`M(T)`$ may be calculated from $`E_\stackrel{}{q}`$. The first order contribution stems from the quadratic term and is the well-known $`T^{3/2}`$-law. The next order depends on the behaviour of $`E_\stackrel{}{q}`$ for larger values of $`q`$. It is still not clear, wether or not the quartic term in (98) describes the generic feature in the experiments. Some experiments indicate that the second term in $`M(T)`$ is $`\alpha T^2`$, which could be explained with a linear behaviour in $`E_\stackrel{}{q}`$. Our dispersion $`E_\stackrel{}{q}`$ also allows to calculate such temperature-dependent quantities numerically. However, we did not analyze this for our two-band model, because these results could not be compared to experiments.
## VI Summary and Conclusion
In this paper we presented a variational method for the description of spin-wave excitations in itinerant ferromagnets. Our starting point was a general multi-band Hubbard-model which is an appropriate model for elements of the iron group. Earlier work showed that ferromagnetism in metals requires substantial interaction strengths even in systems with orbital degeneracy. Therefore, we suspect that weak-coupling theories are not able to describe the physics of itinerant ferromagnets correctly. Our method is based on a variational study of multi-band Hubbard models with the help of generalized Gutzwiller wave functions. These wave functions yield the exact ground-state both in the uncorrelated and the atomic limit of the Hubbard model. Therefore, we expect them to describe reasonably the ground-state properties for finite values of the correlation parameters.
From a theoretical point of view, spin waves are given as peaks in the imaginary part $`\chi _T(\stackrel{}{q},E)`$ of the transversal spin susceptibility . This quantity can be measured using inelastic neutron scattering. For $`\stackrel{}{q}=\stackrel{}{0}`$ the spin-symmetry of the ferromagnetic ground-state leads to an isolated peak $`\delta (E)`$ in $`\chi _T`$. We assume that this peak is broadened only moderately for finite values of $`\stackrel{}{q}`$. Then, the position of the peak, which is interpreted as the spin-wave dispersion $`E_\stackrel{}{q}`$, can be calculated as a static expectation-value of spin-operators in the ground state of the system. In our approximation we use the variational instead of the true ground state and calculate all expectation values in the limit of large spatial dimensions.
Our results may be evaluated numerically for general multi-band models to describe iron and nickel. However, the main numerical problem is the determination of the optimum variational wave-function for these realistic systems with a non-trivial atomic Hamiltonian. Work in this direction is still in progress. In this work we applied our method to a two-band model. Here, we found a behaviour which is in qualitative agreement with experimental results for strong ferromagnets. The spin-wave dispersion $`E_\stackrel{}{q}`$ is almost isotropic and quadratic for small values of $`\stackrel{}{q}`$. For larger $`\stackrel{}{q}`$ we found quartic corrections, which are also seen in some of the experiments. The values of the spin-wave stiffness have the right order of magnitude compared to experiments. We concluded that the low-lying magnetic excitations in our correlated and itinerant electron system are similar to those in a localized spin systems. When we estimate the Curie-temperature $`T_C`$ from our spin-wave properties and compare it to the condensation energy we find a consistent picture in our variational approach. The ferromagnetic phase transition is driven by spin-waves and the value for $`T_C`$ is therefore much smaller than typical Fermi energies in itinerant electron systems.
## VII Acknowledgement
The author is very grateful to F. Gebhard for a critical reading of the manuscript and many helpful discussions.
## A Global spin-symmetry
The variational spin-wave dispersion (42) is gapless only if the variational wave-function $`|\mathrm{\Psi }_\text{G}`$ is an eigenstate of the global spin-operator (see eq. (49)). This symmetry is not fulfilled for an arbitrary choice of the variational parameters $`\lambda _\mathrm{\Gamma }`$ in (10). In this appendix, we present two ways to implement this symmetry.
The first way to guarantee the global spin symmetry of $`|\mathrm{\Psi }_\text{G}`$ starts from the fact that the one-particle product-state $`|\mathrm{\Phi }_0`$ in (9) may in general be chosen as an eigenstate both of $`\widehat{\stackrel{}{S}}^2`$ and $`\widehat{S}^z`$. In this case, $`|\mathrm{\Psi }_\text{G}`$ is an eigenstate of $`\widehat{S}^z`$ if the states $`|\mathrm{\Gamma }`$ are chosen as eigenstates of the local spin-operator in $`z`$-direction. Then, the correct spin-symmetry of $`|\mathrm{\Psi }_\text{G}`$ is ensured if $`|\mathrm{\Psi }_\text{G}`$ is an eigenstate of $`\widehat{S}^+\widehat{S}^{}`$. It is equivalent to demand that
$$\widehat{S}^+\widehat{S}^{}_{\mathrm{\Psi }_\text{G}}=\frac{2S_z}{N_\text{G}},$$
(A1)
since we may assume that the spin in $`|\mathrm{\Psi }_\text{G}`$ has a maximal component in $`z`$-direction. The left-hand-side of eq. (A1) is just the norm $`N_\stackrel{}{q}`$ for $`\stackrel{}{q}=\stackrel{}{0}`$ which was derived in (93). Thus, (A1) together with (93) leads to one additional condition for the variational-parameters $`\lambda _\mathrm{\Gamma }`$. However, this condition is not very helpful, because equation (93) includes the optimum values $`\lambda _\mathrm{\Gamma }`$. This means that (A1) has to be included in the minimization algorithm, which determines the optimum wave-function $`|\mathrm{\Psi }_\text{G}`$; numerically, this is a very difficult problem.
In the following we propose a second, more feasible strategy for the implementation of the correct spin-symmetry. To this end, we arrange the orbitals on the atoms of our system into groups which carry the index of the respective representation $`D`$ of the point-symmetry group. Then, we define the operators
$$\widehat{N}_D^s=\underset{i,bD}{}\widehat{n}_{i;(bs)},\widehat{M}_D^s=\underset{i,bD}{}\widehat{m}_{i;(bs)},$$
(A2)
for the gross and net number of electrons in orbitals of the representation $`D`$ and with spin $`s`$. Again, we assume that each representation occurs only once. Under this condition, group theoretical arguments show that the following relation holds
$$\underset{bD}{}\widehat{m}_{i;(bs)}=\underset{bD}{}\widehat{n}_{i;(bs)}\underset{\mathrm{\Gamma }(|\mathrm{\Gamma }|2)}{}f_D^s(\mathrm{\Gamma })\widehat{m}_{i;\mathrm{\Gamma }},$$
(A3)
where
$$f_D^s(\mathrm{\Gamma })=\underset{bD}{}\underset{I[(b,s)I]}{}\left|T_{\mathrm{\Gamma },I(b,s)}\right|^2.$$
(A4)
Thus, we can write
$`\widehat{M}_D^s`$ $`=`$ $`\widehat{N}_D^s{\displaystyle \underset{\mathrm{\Gamma }(|\mathrm{\Gamma }|2)}{}}f_D^s(\mathrm{\Gamma })\widehat{M}_\mathrm{\Gamma },`$ (A6)
$`\widehat{M}_\mathrm{\Gamma }`$ $``$ $`{\displaystyle \underset{i}{}}\widehat{m}_{i;\mathrm{\Gamma }},`$ (A7)
and the Gutzwiller-projector (10) becomes
$$\widehat{P}_\text{G}=\underset{D,s}{}\left(\lambda _D^s\right)^{\widehat{N}_D^s}\underset{\mathrm{\Gamma }}{}\stackrel{~}{\lambda }_\mathrm{\Gamma }^{\widehat{M}_\mathrm{\Gamma }}.$$
(A8)
Here, we already assumed that the parameters $`\lambda _{(b,s)}`$ ($`\lambda _D^s`$) are the same for all orbitals, which belong to $`D`$. Further, we introduced
$$\stackrel{~}{\lambda }_\mathrm{\Gamma }=\lambda _\mathrm{\Gamma }\underset{D,s}{}\left(\lambda _D^s\right)^{f_D^s(\mathrm{\Gamma })}.$$
(A9)
Now, we postulate the following conditions which ensure that $`|\mathrm{\Psi }_\text{G}`$ is an eigenstate of the operator $`\widehat{\stackrel{}{S}}^2`$:
* $`\lambda _D^s=\lambda _D^{}^s\lambda _s`$ for all representations $`D,D^{}`$.
* For all states $`|\mathrm{\Gamma }_S^{S_z},|\mathrm{\Gamma }_S^{S_z^{}}`$, which belong to the same spin-multiplet ($`\mathrm{\Gamma }^S`$) with spin $`S`$, we have $`\stackrel{~}{\lambda }_{\mathrm{\Gamma }_S^{S_z}}=\stackrel{~}{\lambda }_{\mathrm{\Gamma }_S^{S_z^{}}}\stackrel{~}{\lambda }_{\mathrm{\Gamma }_S}`$.
For the prove of this statement, we can first conclude from (I) that the state $`|\mathrm{\Phi }_0`$ is an eigenstate of
$$\underset{D,s}{}\left(\lambda _D^s\right)^{\widehat{N}_D^s}=\underset{s}{}\lambda _s^{\widehat{N}_s},$$
(A10)
where $`\widehat{N}_s`$ is the number-operator for electrons with spin $`s`$. Thus, we have
$`\widehat{S}^+\widehat{S}^{}|\mathrm{\Psi }_\text{G}`$ $`=`$ $`\widehat{S}^+\widehat{S}^{}{\displaystyle \underset{s}{}}\lambda _s^{\widehat{N}_s}{\displaystyle \underset{\mathrm{\Gamma }}{}}\stackrel{~}{\lambda }_\mathrm{\Gamma }^{\widehat{M}_\mathrm{\Gamma }}|\mathrm{\Phi }_0`$ (A12)
$``$ $`\widehat{S}^+\widehat{S}^{}{\displaystyle \underset{\mathrm{\Gamma }}{}}\stackrel{~}{\lambda }_\mathrm{\Gamma }^{\widehat{M}_\mathrm{\Gamma }}|\mathrm{\Phi }_0,`$ (A13)
since $`[\widehat{N}_s,\widehat{M}_\mathrm{\Gamma }]=0`$ for all $`s,\mathrm{\Gamma }`$. We introduce the operator
$$\widehat{M}_{\mathrm{\Gamma }_S}\underset{S_z=S}{\overset{S}{}}\widehat{M}_{\mathrm{\Gamma }_S^{S_z}},$$
(A14)
which has the property $`[\widehat{S}^\pm ,\widehat{M}_{\mathrm{\Gamma }_S}]=0`$. Then, condition (II) finally gives
$`\widehat{S}^+\widehat{S}^{}|\mathrm{\Psi }_\text{G}`$ $``$ $`\widehat{S}^+\widehat{S}^{}{\displaystyle \underset{\mathrm{\Gamma }_S}{}}\stackrel{~}{\lambda }_{\mathrm{\Gamma }_S}^{\widehat{M}_{\mathrm{\Gamma }_S}}|\mathrm{\Phi }_0`$ (A15)
$``$ $`{\displaystyle \underset{\mathrm{\Gamma }^S}{}}\stackrel{~}{\lambda }_{\mathrm{\Gamma }_S}^{\widehat{M}_{\mathrm{\Gamma }_S}}\widehat{S}^+\widehat{S}^{}|\mathrm{\Phi }_0|\mathrm{\Psi }_\text{G}`$ (A16)
such that $`|\mathrm{\Psi }_\text{G}`$ is an eigenstate of $`\widehat{\stackrel{}{S}}^2`$.
Note, that condition (II) reduces the number of variational-parameters significantly, since all parameters $`\lambda _\mathrm{\Gamma }`$ for states $`|\mathrm{\Gamma }`$, which belong to the same spin-multiplet $`\mathrm{\Gamma }_S`$ are now determined by just one parameter $`\stackrel{~}{\lambda }_{\mathrm{\Gamma }_S}`$. However, this restriction still allows different occupations of the several $`S_z`$ components, since
$$m_{\mathrm{\Gamma }_S^{S_z}}=\lambda _{\mathrm{\Gamma }_S}m_{\mathrm{\Gamma }_S^{S_z}}^0$$
(A17)
also depends on the one-particle-state $`|\mathrm{\Phi }_0`$.
The one-particle-occupations
$$m_{(bs)}=\lambda _{(bs)}m_{(bs)}^0=\lambda _sm_{(bs)}^0m_D^s,$$
(A18)
with $`bD`$ depends on the quantities $`n_D^sn_{(bs)}`$ and $`m_{\mathrm{\Gamma }_S^{S_z}}`$ (see eq. (A3)),
$$m_D^s=n_D^s\underset{\mathrm{\Gamma }_S(|\mathrm{\Gamma }_S|2)}{}\underset{S_z=S}{\overset{S}{}}f_D^s(\mathrm{\Gamma }_S^{S_z})m_{\mathrm{\Gamma }_S^{S_z}}.$$
(A19)
However, the additional conditions (I) and (II) prevent us from the derivation of an analytical expression for $`m_D^s`$, since the parameters $`m_{\mathrm{\Gamma }_S^{S_z}}`$ depend on $`m_D^s`$ via (A9) and (A17). Therefore, the relation (A19) has to be implemented into our minimization algorithm with the help of appropriate Lagrange parameters.
## B Diagrams and Vertices
### 1 The vertices $`\stackrel{~}{x}_𝔇^𝔇^{}`$ and $`\stackrel{~}{\xi }_𝔇^{𝔇^{}𝔇^{\prime \prime }}`$
In Section (IV) the vertices $`\stackrel{~}{x}_𝔇^𝔇^{}`$ (see eqs. (70)) have been derived in terms of the coefficients $`x_{II^{}}`$, which occur in the expansion (12). Now we will derive an explicit expression for these coefficients. The operator (82) may be written as
$`m_J^{}^{I_1,I_2}`$ $`=`$ $`{\displaystyle \underset{\sigma J^{}}{}}\left(n_\sigma ^0+\widehat{n}_\sigma ^{\text{HF}}\right){\displaystyle \underset{\sigma \overline{J}^{}\backslash (I_1I_2)}{}}\left(1n_\sigma ^0\widehat{n}_\sigma ^{\text{HF}}\right)`$ (B1)
$`=`$ $`{\displaystyle \underset{J(I_1,I_2J)}{}}\left[\left(1\right)^{\left|J\overline{J}^{}\right|}{\displaystyle \underset{\sigma J^{}\mathrm{}J}{}}n_\sigma ^0{\displaystyle \underset{\sigma \overline{J}^{}\mathrm{}\{JI_1I_2\}}{}}(1n_\sigma ^0)\right]\widehat{n}_J^{\text{HF}}.`$ (B2)
Thus, $`\widehat{P}_\text{G}^2`$ in (12) becomes
$`\widehat{P}_\text{G}^2`$ $`=`$ $`{\displaystyle \underset{\mathrm{\Gamma }}{}}\lambda _\mathrm{\Gamma }^2{\displaystyle \underset{I_1,I_2(I_1I_2=\mathrm{})}{}}{\displaystyle \underset{J^{}(I_1,I_2J^{})}{}}T_{J^{}I_1,\mathrm{\Gamma }}T_{\mathrm{\Gamma },J^{}I_2}^+f_{I_1}^J^{}f_{I_2}^J^{}\left({\displaystyle \underset{\sigma I_1I_2}{}}{\displaystyle \frac{1}{1n_\sigma ^0}}\right)`$ (B4)
$`\times {\displaystyle \underset{J(I_1,I_2J)}{}}\left[(1)^{\left|J\overline{J}^{}\right|}{\displaystyle \underset{\sigma J^{}\mathrm{}J}{}}n_\sigma ^0{\displaystyle \underset{\sigma \overline{J}^{}\mathrm{}\{JI_1I_2\}}{}}(1n_\sigma ^0)\right]\widehat{n}_{JI_1,JI_2}^{\text{HF}}.`$
A comparison of the coefficients in (B4) and (12) gives
$`x_{JI_1,JI_2}`$ $`=`$ $`{\displaystyle \underset{J^{}(I_1,I_2J^{})}{}}T_{J^{}I_1,\mathrm{\Gamma }}T_{\mathrm{\Gamma },J^{}I_2}^+f_{I_1}^J^{}f_{I_2}^J^{}\left({\displaystyle \underset{\sigma I_1I_2}{}}{\displaystyle \frac{1}{1n_\sigma ^0}}\right)`$ (B6)
$`\times \left[\left(1\right)^{\left|J\overline{J}^{}\right|}{\displaystyle \underset{\sigma J^{}\mathrm{}J}{}}n_\sigma ^0{\displaystyle \underset{\sigma \overline{J}^{}\mathrm{}\{JI_1I_2\}}{}}(1n_\sigma ^0)\right].`$
In Appendix (C) we need an expression for vertices
$$\stackrel{~}{\xi }_𝔇^{𝔇^{}𝔇^{\prime \prime }}=\stackrel{~}{\xi }_{(\sigma _1\sigma _2)}^{(\sigma _3\sigma _4)(\sigma _5\sigma _6)}$$
(B7)
with three incoming and outgoing lines, which stem from a term like
$$y_{l;(\sigma _1,\sigma _3,\sigma _5)(\sigma _2,\sigma _4,\sigma _6)}\left\{\mathrm{}\left(\widehat{c}_{l,\sigma _1}^+\widehat{c}_{l,\sigma _2}\right)\left(\widehat{c}_{l,\sigma _3}^+\widehat{c}_{l,\sigma _4}\right)\left(\widehat{c}_{l,\sigma _5}^+\widehat{c}_{l,\sigma _6}\right)\mathrm{}\right\}_{\mathrm{\Phi }_0}^cf_{\sigma _1}^{\sigma _3}f_{\sigma _5}^{\sigma _1\sigma _3}f_{\sigma _2}^{\sigma _4}f_{\sigma _6}^{\sigma _2\sigma _4}.$$
(B8)
Hence, these vertices are given as
$$\stackrel{~}{\xi }_{(\sigma _1\sigma _2)}^{(\sigma _3\sigma _4)(\sigma _5\sigma _6)}=f_{\sigma _1}^{\sigma _3}f_{\sigma _5}^{\sigma _1\sigma _3}f_{\sigma _2}^{\sigma _4}f_{\sigma _6}^{\sigma _2\sigma _4}x_{(\sigma _1,\sigma _3,\sigma _5)(\sigma _2,\sigma _4,\sigma _6)}.$$
(B9)
### 2 Diagrams
In Appendix (C) and (D) we need some diagrams, which can be evaluated in momentum-space as follows:
1. $`W_{(\sigma _3\sigma _4)(\sigma _5\sigma _6)}^{(\sigma _1\sigma _2)}(\stackrel{}{q}){\displaystyle \underset{ijl}{}}e^{i\stackrel{}{q}(\stackrel{}{R}_j\stackrel{}{R}_l)}P_{i,j}^{\sigma _1\sigma _4}P_{j,l}^{\sigma _3\sigma _6}P_{l,i}^{\sigma _5\sigma _2}`$ (B11)
$`={\displaystyle \underset{\stackrel{}{k}}{}}n_\stackrel{}{k}^{\sigma _1\sigma _4}n_{\stackrel{}{k}+\stackrel{}{q}}^{\sigma _3\sigma _6}n_\stackrel{}{k}^{\sigma _5\sigma _2}\delta _{\sigma _3}^{\sigma _6}n_{\sigma _3}^0\stackrel{~}{P}_{(\sigma _5\sigma _4)}^{(\sigma _1\sigma _2)}(\stackrel{}{0})\delta _{\sigma _2}^{\sigma _5}n_{\sigma _2}^0\stackrel{~}{P}_{(\sigma _3\sigma _4)}^{(\sigma _1\sigma _6)}(\stackrel{}{q})`$ (B12)
$`\delta _{\sigma _1}^{\sigma _4}n_{\sigma _1}^0\stackrel{~}{P}_{(\sigma _5\sigma _6)}^{(\sigma _3\sigma _2)}(\stackrel{}{q})+2\delta _{\sigma _1}^{\sigma _4}\delta _{\sigma _3}^{\sigma _6}\delta _{\sigma _2}^{\sigma _5}n_{\sigma _1}^0n_{\sigma _2}^0n_{\sigma _3}^0`$ (B13)
2. $$E_{(\sigma _1\sigma _2)}^{(\sigma _3\sigma _4)}(\stackrel{}{q})\underset{i()j}{}e^{i\stackrel{}{q}(\stackrel{}{R}_i\stackrel{}{R}_j)}t_{i,j}^{\sigma _3\sigma _4}P_{i,j}^{\sigma _1\sigma _2}=\underset{\stackrel{}{k}}{}\epsilon _\stackrel{}{k}^{\sigma _3\sigma _4}n_{\stackrel{}{k}+\stackrel{}{q}}^{\sigma _1\sigma _2},$$
(B14)
with
$$\epsilon _\stackrel{}{k}^{\sigma _3\sigma _4}\underset{ij}{}e^{i\stackrel{}{k}(\stackrel{}{R}_i\stackrel{}{R}_j)}t_{i,j}^{\sigma _3\sigma _4}$$
(B15)
3. $`V_{(\sigma _3\sigma _4)(\sigma _5\sigma _6)}^{(\sigma _1\sigma _2)}(\stackrel{}{q}){\displaystyle \underset{ijl}{}}e^{i\stackrel{}{q}(\stackrel{}{R}_i\stackrel{}{R}_l)}t_{i,j}^{\sigma _5\sigma _6}P_{i,l}^{\sigma _4\sigma _2}P_{l,j}^{\sigma _1\sigma _3}`$ (B17)
$`={\displaystyle \underset{\stackrel{}{k}}{}}\epsilon _\stackrel{}{k}^{\sigma _5\sigma _6}n_{\stackrel{}{k}+\stackrel{}{q}}^{\sigma _4\sigma _2}n_\stackrel{}{k}^{\sigma _1\sigma _3}+\delta _{\sigma _2}^{\sigma _4}n_{\sigma _2}^0E_{(\sigma _1\sigma _3)}^{(\sigma _5\sigma _6)}(\stackrel{}{0})+\delta _{\sigma _1}^{\sigma _3}n_{\sigma _1}^0E_{(\sigma _4\sigma _2)}^{(\sigma _5\sigma _6)}(\stackrel{}{q})`$ (B18)
and
$`\overline{V}_{(\sigma _3\sigma _4)(\sigma _5\sigma _6)}^{(\sigma _1\sigma _2)}(\stackrel{}{q}){\displaystyle \underset{ijl}{}}e^{i\stackrel{}{q}(\stackrel{}{R}_l\stackrel{}{R}_j)}t_{i,j}^{\sigma _5\sigma _6}P_{i,l}^{\sigma _4\sigma _2}P_{l,j}^{\sigma _1\sigma _3}`$ (B20)
$`={\displaystyle \underset{\stackrel{}{k}}{}}\epsilon _\stackrel{}{k}^{\sigma _5\sigma _6}n_\stackrel{}{k}^{\sigma _4\sigma _2}n_{\stackrel{}{k}+\stackrel{}{q}}^{\sigma _1\sigma _3}+\delta _{\sigma _2}^{\sigma _4}n_{\sigma _2}^0E_{(\sigma _1\sigma _3)}^{(\sigma _5\sigma _6)}(\stackrel{}{q})+\delta _{\sigma _1}^{\sigma _3}n_{\sigma _1}^0E_{(\sigma _4\sigma _2)}^{(\sigma _5\sigma _6)}(\stackrel{}{0})`$ (B21)
4. $`U_{(\sigma _3\sigma _4)(\sigma _5\sigma _6)}^{(\sigma _1\sigma _2)}(\stackrel{}{q}){\displaystyle \underset{ijlm}{}}e^{i\stackrel{}{q}(\stackrel{}{R}_m\stackrel{}{R}_l)}t_{i,j}^{\sigma _1\sigma _2}P_{i,m}^{\sigma _1\sigma _4}P_{m,l}^{\sigma _3\sigma _6}P_{l,j}^{\sigma _5\sigma _2}`$ (B23)
$`={\displaystyle \underset{\stackrel{}{k}}{}}\epsilon _\stackrel{}{k}^{\sigma _1\sigma _2}n_\stackrel{}{k}^{\sigma _1\sigma _4}n_{\stackrel{}{k}+\stackrel{}{q}}^{\sigma _3\sigma _6}n_\stackrel{}{k}^{\sigma _5\sigma _2}\delta _{\sigma _3}^{\sigma _6}n_{\sigma _3}^0\overline{V}_{(\sigma _1\sigma _2)(\sigma _1\sigma _2)}^{(\sigma _5\sigma _4)}(\stackrel{}{0})\delta _{\sigma _1}^{\sigma _4}n_{\sigma _1}^0V_{(\sigma _3\sigma _2)(\sigma _1\sigma _2)}^{(\sigma _5\sigma _6)}(\stackrel{}{q})+`$ (B24)
$`\delta _{\sigma _2}^{\sigma _5}n_{\sigma _2}^0\overline{V}_{(\sigma _1\sigma _6)(\sigma _1\sigma _2)}^{(\sigma _3\sigma _4)}(\stackrel{}{q})\delta _{\sigma _1}^{\sigma _4}\delta _{\sigma _2}^{\sigma _5}n_{\sigma _1}^0n_{\sigma _2}^0E_{(\sigma _3\sigma _6)}^{(\sigma _1\sigma _2)}(\stackrel{}{q})\delta _{\sigma _1}^{\sigma _4}\delta _{\sigma _3}^{\sigma _6}n_{\sigma _1}^0n_{\sigma _3}^0E_{(\sigma _5\sigma _2)}^{(\sigma _1\sigma _2)}(\stackrel{}{0})+`$ (B25)
$`\delta _{\sigma _2}^{\sigma _5}\delta _{\sigma _3}^{\sigma _6}n_{\sigma _1}^0n_{\sigma _3}^0E_{(\sigma _1\sigma _4)}^{(\sigma _1\sigma _2)}(\stackrel{}{0})\stackrel{~}{P}_{(\sigma _3\sigma _4)}^{(\sigma _1\sigma _6)}(\stackrel{}{q})E_{(\sigma _5\sigma _2)}^{(\sigma _1\sigma _2)}(\stackrel{}{q})\stackrel{~}{P}_{(\sigma _3\sigma _2)}^{(\sigma _5\sigma _6)}(\stackrel{}{q})E_{(\sigma _1\sigma _4)}^{(\sigma _1\sigma _2)}(\stackrel{}{q})`$ (B26)
## C Evaluation of the atomic interaction
According to eq. (46) we have to analyze the expectation values
$`{\displaystyle \frac{1}{N_\stackrel{}{q}}}{\displaystyle \underset{k}{}}\mathrm{\Psi }_\stackrel{}{q}^\text{G}\left|\widehat{m}_{k;\mathrm{\Gamma }}\right|\mathrm{\Psi }_\stackrel{}{q}^\text{G}`$ $`=`$ $`{\displaystyle \frac{1}{N_\stackrel{}{q}}}{\displaystyle \underset{i,j,k}{}}e^{i\stackrel{}{q}(\stackrel{}{R}_i\stackrel{}{R}_j)}\mathrm{\Psi }_\text{G}\left|\widehat{S}_i^+\widehat{m}_{k;\mathrm{\Gamma }}\widehat{S}_j^{}\right|\mathrm{\Psi }_\text{G}`$ (C2)
$`=`$ $`{\displaystyle \frac{1}{N_\stackrel{}{q}}}{\displaystyle \underset{k}{}}[{\displaystyle \underset{i,j(k)}{}}e^{i\stackrel{}{q}(\stackrel{}{R}_i\stackrel{}{R}_j)}\mathrm{\Psi }_\text{G}\left|\widehat{S}_i^+\widehat{S}_j^{}\widehat{m}_{k;\mathrm{\Gamma }}\right|\mathrm{\Psi }_\text{G}`$ (C5)
$`+{\displaystyle \underset{j(k)}{}}(e^{i\stackrel{}{q}(\stackrel{}{R}_k\stackrel{}{R}_j)}\mathrm{\Psi }_\text{G}\left|\widehat{S}_k^+\widehat{m}_{k;\mathrm{\Gamma }}\widehat{S}_j^{}\right|\mathrm{\Psi }_\text{G}+c.c.)`$
$`+\mathrm{\Psi }_\text{G}\left|\widehat{S}_k^+\widehat{m}_{k;\mathrm{\Gamma }}\widehat{S}_k^{}\right|\mathrm{\Psi }_\text{G}]`$
This evaluation will be done separately for the three terms (C5)-(C5).
In the first term (C5), we have to distinguish connected and unconnected diagrams. Here, the term “connected” means, that the lattice site $`k`$ is connected to one of the lattice sites $`i,j`$. This condition necessarily requires that $`k`$ is connected to both lattice sites $`i`$ and $`j`$, because $`\widehat{S}_i^+`$ or $`\widehat{S}_j^{}`$ generate a spin-flip. Such a process cannot be compensated in a diagram with only one of these operators, neither by the external vertex-operator $`\widehat{m}_{k;\mathrm{\Gamma }}`$ nor by one of the internal vertex-operators $`\widehat{n}_{l;I,I^{}}^{\text{HF}}`$.
The unconnected terms can be written as
$`(\text{C5})^{\text{uc}}`$ $`=`$ $`{\displaystyle \frac{1}{N_\stackrel{}{q}}}{\displaystyle \underset{k}{}}{\displaystyle \frac{\mathrm{\Psi }_\text{G}\left|\widehat{m}_{k;\mathrm{\Gamma }}\right|\mathrm{\Psi }_\text{G}}{N_\text{G}}}`$ (C8)
$`\times [{\displaystyle \underset{ij(k)}{}}e^{i\stackrel{}{q}(\stackrel{}{R}_i\stackrel{}{R}_j)}{\displaystyle \underset{l(i,j,k)}{}}\left(\widehat{P}_{i;\text{G}}\widehat{S}_i^+\widehat{P}_{i;\text{G}}\right)\left(\widehat{P}_{j;\text{G}}\widehat{S}_j^{}\widehat{P}_{j;\text{G}}\right)\widehat{P}_{l;\text{G}}^2_{\mathrm{\Phi }_0}`$
$`+{\displaystyle \underset{i}{}}{\displaystyle \underset{l(i,k)}{}}\left(\widehat{P}_{i;\text{G}}\widehat{S}_i^+\widehat{S}_i^{}\widehat{P}_{i;\text{G}}\right)\widehat{P}_{l;\text{G}}^2_{\mathrm{\Phi }_0}].`$
When we ignore the restriction $`i,j,lk`$, the sum over $`i`$ and $`j`$ gives just $`N_\stackrel{}{q}`$. Thus, we find the correct result for $`(\text{C5})^{\text{uc}}`$ by substracting all diagrams with $`i=k`$ or $`j=k`$ for an external vertex, or $`l=k`$ for one of the internal vertices. These contributions are given as
$`i=j=k`$ $`:`$ $`N_\text{G}_\mathrm{\Gamma }^{}S_{}(\mathrm{\Gamma }^{})m_\mathrm{\Gamma }^{}`$ $`i=k\text{ or }j=k`$ $`:`$ $`2N_\text{G}\stackrel{}{S}^+\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\stackrel{}{S}^{}`$ $`l=k`$ $`:`$ $`N_\text{G}\stackrel{}{S}^+\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\stackrel{~}{x}\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\stackrel{}{S}^{}.`$ Altogether, we obtain the following expression for the contribution of the unconnected diagrams:
$`(\text{C5})^{\text{uc}}`$ $`=`$ $`\left[Lm_\mathrm{\Gamma }m_\mathrm{\Gamma }{\displaystyle \frac{LN_\text{G}}{N_\stackrel{}{q}}}\left(2\stackrel{}{S}^+\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\stackrel{}{S}^{}+{\displaystyle \underset{\mathrm{\Gamma }^{}}{}}S_{}(\mathrm{\Gamma }^{})m_\mathrm{\Gamma }^{}+\stackrel{}{S}^+\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\stackrel{~}{x}\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\stackrel{}{S}^{}\right)\right]`$ (C10)
$``$ $`Lm_\mathrm{\Gamma }+m_\mathrm{\Gamma }^1(\stackrel{}{q}),`$ (C11)
where $`\frac{LN_\text{G}}{N_\stackrel{}{q}}`$ is given in (93). Note that only the second term $`m_\mathrm{\Gamma }^1(\stackrel{}{q})`$ is relevant for our spin-wave-dispersion $`E_\stackrel{}{q}^{var}`$, since the first term $`Lm_\mathrm{\Gamma }`$ is canceled by the respective ground-state contribution in eq. (42).
For the connected diagrams in (C5) we may distinguish between those diagrams with two or four lines, which enter or leave the external vertex. The vertex-factors with two lines can be evaluated from (83), whereas the respective factor with four lines stems from the expansion
$`\left\{\mathrm{}\widehat{m}_{I_1,I_2}\mathrm{}\right\}_{\mathrm{\Phi }_0}^C`$ $``$ $`{\displaystyle \underset{\sigma _1,\sigma _2,\sigma _3,\sigma _4}{}}\stackrel{~}{𝔙}_{_{I_1,I_2}}^{(\sigma _1\sigma _2)(\sigma _3\sigma _4)}\left\{\mathrm{}\left(\widehat{c}_{\sigma _1}^+\widehat{c}_{\sigma _2}\right)\left(\widehat{c}_{\sigma _3}^+\widehat{c}_{\sigma _4}\right)\mathrm{}\right\}_{\mathrm{\Phi }_0}^c`$ (C12)
$`\stackrel{~}{𝔙}_{_{I_1,I_2}}^{(\sigma _1\sigma _2)(\sigma _3\sigma _4)}`$ $`=`$ $`{\displaystyle \underset{J(\sigma _1,\sigma _2,\sigma _3,\sigma _4I)}{}}f_{\sigma _1}^Jf_{\sigma _2}^Jf_{\sigma _3}^Jf_{\sigma _4}^Jm_J^0\left({\displaystyle \underset{\sigma (\sigma _1,\sigma _2,\sigma _3,\sigma _4)}{}}{\displaystyle \frac{1}{(1n_\sigma ^0)}}\right)`$ (C15)
$`\times [f_{\sigma _1}^{\sigma _3}f_{\sigma _2}^{\sigma _4}\delta _{J(\sigma _1,\sigma _3)}^{I_1}\delta _{J(\sigma _2,\sigma _4)}^{I_2}\delta _{\sigma _1}^{\sigma _2}\delta _{J\sigma _3}^{I_1}\delta _{J\sigma _4}^{I_2}\delta _{\sigma _3}^{\sigma _4}\delta _{J\sigma _1}^{I_1}\delta _{J\sigma _2}^{I_2}+`$
$`\delta _{\sigma _1}^{\sigma _4}\delta _{J\sigma _1}^{I_1}\delta _{J\sigma _2}^{I_2}\delta _{\sigma _2}^{\sigma _3}\delta _{J\sigma _1}^{I_1}\delta _{J\sigma _4}^{I_2}+(\delta _{\sigma _1}^{\sigma _2}\delta _{\sigma _3}^{\sigma _4}\delta _{\sigma _1}^{\sigma _4}\delta _{\sigma _2}^{\sigma _3})\delta _J^{I_1}\delta _J^{I_2}],`$
which is also derived for $`\left|I_1\right|=\left|I_2\right|`$. This expression together with eq. (79) leads to the following external vertex $`\underset{𝔇}{\overset{𝔇^{}}{\stackrel{+}{T}}}`$ for the operator $`\widehat{S}_i^+`$
$$\underset{(\sigma _1\sigma _2)}{\overset{(\sigma _3\sigma _4)}{\stackrel{+}{T}}}=\underset{\mathrm{\Gamma }}{}\lambda _{\mathrm{\Gamma }_+}\lambda _\mathrm{\Gamma }\sqrt{S_+(\mathrm{\Gamma })}\underset{I_1,I_2}{}T_{I_1,\mathrm{\Gamma }_+}T_{\mathrm{\Gamma },I_2}^+\stackrel{~}{𝔙}_{_{I_1,I_2}}^{(\sigma _1\sigma _2)(\sigma _3\sigma _4)},$$
(C16)
whereas the corresponding factor $`\underset{𝔇}{\overset{𝔇^{}}{\stackrel{}{T}}}`$ for the operator $`\widehat{S}_j^{}`$ is given as
$$\underset{(\sigma _1\sigma _2)}{\overset{(\sigma _3\sigma _4)}{\stackrel{}{T}}}=\left(\underset{(\sigma _2\sigma _1)}{\overset{(\sigma _4\sigma _3)}{\stackrel{+}{T}}}\right)^{}.$$
(C17)
Using eqs. (83), (C12), and
$$\widehat{P}_\text{G}\widehat{m}_\mathrm{\Gamma }\widehat{P}_\text{G}=\underset{I_1,I_2}{}\lambda _\mathrm{\Gamma }^2T_{\mathrm{\Gamma },I_1}T_{I_2,\mathrm{\Gamma }}^+\widehat{m}_{I_1,I_2}$$
(C18)
we obtain the following expression for the vertices of the external-operator in (C5) with one or two incoming and outgoing lines,
$`M_{(\sigma _1\sigma _2)}(\mathrm{\Gamma })`$ $`=`$ $`{\displaystyle \underset{I_1,I_2}{}}\lambda _\mathrm{\Gamma }^2T_{\mathrm{\Gamma },I_1}T_{I_2,\mathrm{\Gamma }}^+𝔙_{_{I_1,I_2}}^{(\sigma _1\sigma _2)},`$ (C20)
$`\stackrel{~}{M}_{(\sigma _1\sigma _2)}^{(\sigma _3\sigma _4)}(\mathrm{\Gamma })`$ $`=`$ $`{\displaystyle \underset{I_1,I_2}{}}\lambda _\mathrm{\Gamma }^2T_{\mathrm{\Gamma },I_1}T_{I_2,\mathrm{\Gamma }}^+\stackrel{~}{𝔙}_{_{I_1,I_2}}^{(\sigma _1\sigma _2)(\sigma _3\sigma _4)}.`$ (C21)
When we define the vector $`\stackrel{}{M}(\mathrm{\Gamma })`$ of components $`M_𝔇(\mathrm{\Gamma })`$, we may write the connected terms (C5) with $`ij`$ as
$`m_\mathrm{\Gamma }^2(\stackrel{}{q})`$ $`=`$ $`{\displaystyle \frac{LN_\text{G}}{N_\stackrel{}{q}}}{\displaystyle \underset{𝔇_1,𝔇_2,𝔇_3}{}}\left[\stackrel{}{M}(\mathrm{\Gamma })\left(\widehat{1}\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{0})\stackrel{~}{x}\right)\right]_{𝔇_1}\left(W_{𝔇_2,𝔇_3}^{𝔇_1}(\stackrel{}{q})+W_{𝔇_3,𝔇_2}^{𝔇_1}(\stackrel{}{q})\right)`$ (C24)
$`\times \left[\left(\widehat{1}+\stackrel{~}{x}\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\right)\stackrel{}{S}^+\right]_{𝔇_2}\left[\left(\widehat{1}+\stackrel{~}{x}\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\right)\stackrel{}{S}^{}\right]_{𝔇_3}`$
$`m_\mathrm{\Gamma }^3(\stackrel{}{q})`$ $`=`$ $`{\displaystyle \frac{LN_\text{G}}{N_\stackrel{}{q}}}{\displaystyle \underset{𝔇_1,𝔇_2,𝔇_3}{}}\left[\stackrel{}{M}(\mathrm{\Gamma })\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{0})\right]_{𝔇_1}\stackrel{~}{\xi }_{𝔇_1}^{𝔇_2,𝔇_3}\left[\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\stackrel{}{S}^+\right]_{𝔇_2}\left[\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\stackrel{}{S}^{}\right]_{𝔇_3}`$ (C25)
$`m_\mathrm{\Gamma }^4(\stackrel{}{q})`$ $`=`$ $`{\displaystyle \frac{LN_\text{G}}{N_\stackrel{}{q}}}\stackrel{}{M}(\mathrm{\Gamma })\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{0})(\stackrel{+}{T}\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\stackrel{}{S}^{}+\stackrel{}{T}\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\stackrel{}{S}^+)`$ (C26)
$`m_\mathrm{\Gamma }^5(\stackrel{}{q})`$ $`=`$ $`{\displaystyle \frac{LN_\text{G}}{N_\stackrel{}{q}}}\stackrel{}{S}^+\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\stackrel{~}{M}(\mathrm{\Gamma })\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\stackrel{}{S}^{}`$ (C27)
The tensors $`W_{𝔇_2,𝔇_3}^{𝔇_1}(\stackrel{}{q})`$, $`\stackrel{~}{\xi }_{𝔇_1}^{𝔇_2,𝔇_3}`$ are defined in Appendix (B). In fig. (4) all diagrams, which belong to the atomic interactions are presented.
The connected diagrams in (C5) with $`i=j`$ are determined by the external vertex-operator
$$\widehat{P}_\text{G}\widehat{S}^+\widehat{S}^{}\widehat{P}_\text{G}=\underset{\mathrm{\Gamma }^{}}{}\lambda _\mathrm{\Gamma }^{}^2S_{}(\mathrm{\Gamma }^{})\widehat{m}_\mathrm{\Gamma }^{}.$$
(C28)
This expression leads to the contribution
$$m_\mathrm{\Gamma }^6(\stackrel{}{q})=\frac{LN_\text{G}}{N_\stackrel{}{q}}\underset{\mathrm{\Gamma }^{}}{}S_{}(\mathrm{\Gamma }^{})\left[\stackrel{}{M}(\mathrm{\Gamma }^{})\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{0})\stackrel{}{M}(\mathrm{\Gamma })\right]=m_\mathrm{\Gamma }^6(\stackrel{}{0}).$$
(C29)
For the evaluation of (C5), we need the external vertex for the operator $`\widehat{P}_\text{G}\widehat{S}^+\widehat{m}_\mathrm{\Gamma }\widehat{P}_\text{G}`$, which is the same as the respective term in (86) for fixed $`\mathrm{\Gamma }`$,
$$A_{\left(\sigma _1\sigma _2\right)}^+\lambda _{\mathrm{\Gamma }_+}\lambda _\mathrm{\Gamma }\sqrt{S_+(\mathrm{\Gamma })}\underset{I_1,I_2}{}T_{I_1,\mathrm{\Gamma }_+}T_{\mathrm{\Gamma },I_2}^+𝔙_{_{I_1,I_2}}^{(\sigma _1\sigma _2)}.$$
(C30)
The only connected diagram in (C5) is therefore given as
$$m_\mathrm{\Gamma }^7(\stackrel{}{q})=\frac{LN_\text{G}}{N_\stackrel{}{q}}(\stackrel{}{A}^+\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\stackrel{}{S}^{}+c.c.).$$
(C31)
Finally, we determine the contribution (C5) as
$$m_\mathrm{\Gamma }^8(\stackrel{}{q})=\frac{LN_\text{G}}{N_\stackrel{}{q}}S_+(\mathrm{\Gamma })m_{\mathrm{\Gamma }_+}=m_\mathrm{\Gamma }^8(\stackrel{}{0}).$$
(C32)
Here, we used the relation
$$\widehat{S}^+\widehat{m}_\mathrm{\Gamma }\widehat{S}^{}=S_+(\mathrm{\Gamma })\widehat{m}_{\mathrm{\Gamma }_+}.$$
(C33)
To summarize, the expectation value (C5) for the atomic energy is given as
$$\frac{\mathrm{\Psi }_\stackrel{}{q}^\text{G}\left|\widehat{H}_{\text{at}}\right|\mathrm{\Psi }_\stackrel{}{q}^\text{G}}{N_\stackrel{}{q}}=L\underset{\mathrm{\Gamma }}{}E_\mathrm{\Gamma }m_\mathrm{\Gamma }+\underset{\mathrm{\Gamma }}{}E_\mathrm{\Gamma }\underset{c=1}{\overset{8}{}}m_\mathrm{\Gamma }^c(\stackrel{}{q}).$$
(C34)
## D Evaluation of the one-particle energy
For the diagrammatic evaluation of the one-particle-energy we write the expectation value (47) as
$`{\displaystyle \frac{\mathrm{\Psi }_\stackrel{}{q}^\text{G}\left|\widehat{H}_1\right|\mathrm{\Psi }_\stackrel{}{q}^\text{G}}{N_\stackrel{}{q}}}={\displaystyle \frac{1}{N_\stackrel{}{q}}}{\displaystyle \underset{k()l}{}}{\displaystyle \underset{\sigma _k,\sigma _l}{}}t_{k,l}^{\sigma _k,\sigma _l}\{{\displaystyle \underset{i,j(k,l)}{}}e^{i\stackrel{}{q}(\stackrel{}{R}_i\stackrel{}{R}_j)}\widehat{S}_i^+\widehat{c}_{k;\sigma _k}^+\widehat{c}_{l;\sigma _l}\widehat{S}_j^{}_{\mathrm{\Psi }_\text{G}}+`$ (D2)
$`+{\displaystyle \underset{i(k,l)}{}}\left(e^{i\stackrel{}{q}(\stackrel{}{R}_i\stackrel{}{R}_l)}\widehat{S}_i^+\widehat{c}_{k;\sigma _k}^+\widehat{c}_{l;\sigma _l}\widehat{S}_l^{}_{\mathrm{\Psi }_\text{G}}+e^{i\stackrel{}{q}(\stackrel{}{R}_i\stackrel{}{R}_k)}\widehat{S}_i^+\widehat{c}_{k;\sigma _k}^+\widehat{c}_{l;\sigma _l}\widehat{S}_k^{}_{\mathrm{\Psi }_\text{G}}\right)+c.c`$ (D3)
$`+\left(e^{i\stackrel{}{q}(\stackrel{}{R}_k\stackrel{}{R}_l)}\widehat{S}_k^+\widehat{c}_{k;\sigma _k}^+\widehat{c}_{l;\sigma _l}\widehat{S}_l^{}_{\mathrm{\Psi }_\text{G}}+e^{i\stackrel{}{q}(\stackrel{}{R}_l\stackrel{}{R}_k)}\widehat{c}_{k;\sigma _k}^+\widehat{S}_k^{}\widehat{S}_l^+\widehat{c}_{l;\sigma _l}_{\mathrm{\Psi }_\text{G}}\right)+c.c`$ (D4)
$`+\widehat{S}_k^+\widehat{c}_{k;\sigma _k}^+\widehat{S}_k^{}\widehat{c}_{l;\sigma _l}_{\mathrm{\Psi }_\text{G}}+c.c.\}`$ (D5)
Here, we already used the fact that the tight-binding parameters do not contain any local terms. In (D2) we have to distinguish connected and unconnected diagrams. The unconnected contributions may be written as
$`(\text{D2})^{\text{uc}}`$ $`=`$ $`{\displaystyle \frac{1}{N_\stackrel{}{q}}}{\displaystyle \underset{kl}{}}{\displaystyle \underset{\sigma _k,\sigma _l}{}}t_{k,l}^{\sigma _k,\sigma _l}\widehat{S}_i^+\widehat{c}_{k;\sigma _k}^+\widehat{c}_{l;\sigma _l}\widehat{S}_j^{}_{\mathrm{\Psi }_\text{G}}`$ (D8)
$`\times [{\displaystyle \underset{ij(k,l)}{}}e^{i\stackrel{}{q}(\stackrel{}{R}_i\stackrel{}{R}_j)}{\displaystyle \underset{m(i,j,k,l)}{}}\left(\widehat{P}_{i;\text{G}}\widehat{S}_i^+\widehat{P}_{i;\text{G}}\right)\left(\widehat{P}_{j;\text{G}}\widehat{S}_j^{}\widehat{P}_{j;\text{G}}\right)\widehat{P}_{m;\text{G}}^2_{\mathrm{\Phi }_0}`$
$`+{\displaystyle \underset{i}{}}{\displaystyle \underset{m(i,k,l)}{}}\left(\widehat{P}_{i;\text{G}}\widehat{S}_i^+\widehat{S}_i^{}\widehat{P}_{i;\text{G}}\right)\widehat{P}_{m;\text{G}}^2_{\mathrm{\Phi }_0}].`$
Eq. (D8) can be evaluated, using the same arguments as discussed in connection with eq. (C8). This evaluation leads to
$`(\text{D2})^{\text{uc}}`$ $`=`$ $`E_{\text{kin}}\epsilon ^1(\stackrel{}{q})`$ (D10)
$`\epsilon ^1(\stackrel{}{q})`$ $`=`$ $`{\displaystyle \frac{2E_{\text{kin}}}{L}}{\displaystyle \frac{LN_\text{G}}{N_\stackrel{}{q}}}\left(2\stackrel{}{S}^+\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\stackrel{}{S}^{}+{\displaystyle \underset{\mathrm{\Gamma }^{}}{}}S_{}(\mathrm{\Gamma }^{})m_\mathrm{\Gamma }^{}+\stackrel{}{S}^+\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\stackrel{~}{x}\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\stackrel{}{S}^{}\right)`$ (D11)
where
$$E_{\text{kin}}=\frac{1}{N_\text{G}}\underset{kl}{}\underset{\sigma _k,\sigma _l}{}t_{k,l}^{\sigma _k,\sigma _l}\widehat{c}_{k;\sigma _k}^+\widehat{c}_{l;\sigma _l}_{\mathrm{\Psi }_\text{G}}$$
(D12)
Before we start to discuss the connected diagrams, we should first consider the general structure of external vertices in (D2) at lattice sites $`k`$ and $`l`$, where electrons are created or annihilated. For example, the vertex-function of the operator $`\widehat{c}_\sigma ^+`$ leads to
$$\left\{\mathrm{}\widehat{P}_\text{G}\widehat{c}_\sigma ^+\widehat{P}_\text{G}\mathrm{}\right\}_{\mathrm{\Phi }_0}^C=\underset{\mathrm{\Gamma },\mathrm{\Gamma }^{}}{}\lambda _\mathrm{\Gamma }\lambda _\mathrm{\Gamma }^{}\underset{I(\sigma I)}{}f_\sigma ^IT_{\mathrm{\Gamma },I}^+T_{I\sigma ,\mathrm{\Gamma }^{}}\underset{I_1,I_2}{}T_{I_1,\mathrm{\Gamma }}T_{\mathrm{\Gamma }^{},I_4}^+\left\{\mathrm{}\widehat{m}_{I_1,I_2}\mathrm{}\right\}_{\mathrm{\Phi }_0}^C$$
(D13)
Hence, our general problem is the calculation of vertex-contributions for operators $`\widehat{m}_{I_1,I_2}`$ with $`\left|I_1\right|\left|I_2\right|=1`$. The first case is a vertex with only one incoming and no outgoing line,
$`\left\{\mathrm{}\widehat{m}_{I_1,I_2}\mathrm{}\right\}_{\mathrm{\Phi }_0}^C`$ $``$ $`{\displaystyle \underset{\sigma ^{}}{}}𝔥_{_{I_1,I_2}}(\sigma ^{})\left\{\mathrm{}\widehat{c}_{\sigma _1}^+\mathrm{}\right\}_{\mathrm{\Phi }_0}^c`$ (D15)
$`𝔥_{I_1,I_2}(\sigma ^{})`$ $`=`$ $`{\displaystyle \frac{1}{1n_\sigma ^{}^0}}{\displaystyle \underset{I^{}(\sigma ^{}I^{})}{}}f_\sigma ^{}^I^{}m_I^{}^0\delta _{I_1}^{I^{}\sigma ^{}}\delta _{I_2}^I^{}.`$ (D16)
Note that eq. (D13) together with (27) yields directly the $`q`$-factor (27), which occurs as the renormalization factor for hopping-processes.
Furthermore we have the cases with two (three) incoming and one (two) outgoing lines. This processes lead to the vertices
$`\left\{\mathrm{}\widehat{m}_{I_1,I_2}\mathrm{}\right\}_{\mathrm{\Phi }_0}^C`$ $``$ $`{\displaystyle \underset{\sigma ^{}}{}}{\displaystyle \underset{\sigma _1,\sigma _2}{}}_{I_1,I_2}^{(\sigma _1,\sigma _2)}(\sigma ^{})\left\{\mathrm{}\left(\widehat{c}_\sigma ^{}^+\right)\left(\widehat{c}_{\sigma _1}^+\widehat{c}_{\sigma _2}\right)\mathrm{}\right\}_{\mathrm{\Phi }_0}^c`$ (D18)
$`_{I_1,I_2}^{\left(\sigma _1\sigma _2\right)}(\sigma ^{})`$ $`=`$ $`{\displaystyle \underset{I^{}(\sigma ^{},\sigma _1,\sigma _2I^{})}{}}f_\sigma ^{}^I^{}f_{\sigma _1}^I^{}f_{\sigma _2}^I^{}m_^I^{}^0\left({\displaystyle \underset{\stackrel{~}{\sigma }(\sigma ^{},\sigma _1,\sigma _2)}{}}{\displaystyle \frac{1}{(1n_{\stackrel{~}{\sigma }}^0)}}\right)`$ (D20)
$`\times \left[f_\sigma ^{}^{\sigma _1}\delta _{I_1,}^{I^{}(\sigma ^{},\sigma _1)}\delta _{I_2}^{I^{}\sigma _2}\delta _{\sigma _1}^{\sigma _2}\delta _{I_1}^{I^{}\sigma ^{}}\delta _{I_2}^I^{}+\delta _\sigma ^{}^{\sigma _2}\delta _{I_1}^{I^{}\sigma _1}\delta _{I_2}^I^{}\right].`$
and
$`\left\{\mathrm{}\widehat{m}_{I_1,I_2}\mathrm{}\right\}_{\mathrm{\Phi }_0}^C{\displaystyle \underset{\sigma ^{}}{}}{\displaystyle \underset{\sigma _1,\sigma _2}{}}\stackrel{~}{}_{I_1,I_2}^{\left(\sigma _1\sigma _2\right)(\sigma _3\sigma _4)}(\sigma ^{})\left\{\mathrm{}\left(\widehat{c}_\sigma ^{}^+\right)\left(\widehat{c}_{\sigma _1}^+\widehat{c}_{\sigma _2}\right)\left(\widehat{c}_{\sigma _3}^+\widehat{c}_{\sigma _4}\right)\mathrm{}\right\}_{\mathrm{\Phi }_0}^c`$ (D22)
$`\stackrel{~}{}_{I_1,I_2}^{\left(\sigma _1\sigma _2\right)(\sigma _3\sigma _4)}(\sigma ^{})={\displaystyle \underset{I(\sigma _1,\sigma _2,\sigma _3,\sigma _4,\sigma ^{}I)}{}}f_{\sigma _1}^If_{\sigma _2}^If_{\sigma _3}^If_{\sigma _4}^If_\sigma ^{}^Im_I^0\left({\displaystyle \underset{\sigma (\sigma _1,\sigma _2,\sigma _3,\sigma _4,\sigma ^{})}{}}{\displaystyle \frac{1}{(1n_\sigma ^0)}}\right)f_\sigma ^{}^{\sigma _1}f_\sigma ^{}^{\sigma _3}`$ (D23)
$`\times [f_{\sigma _1}^{\sigma _3}f_{\sigma _2}^{\sigma _4}\delta _{I(\sigma _1,\sigma _3,\sigma ^{})}^{I_1}\delta _{I(\sigma _2,\sigma _4),}^{I_2}f_\sigma ^{}^{\sigma _1}(\delta _{\sigma _1}^{\sigma _2}\delta _{I(\sigma _3,\sigma ^{})}^{I_1}\delta _{I\sigma _4}^{I_2}\delta _{\sigma _1}^{\sigma _4}\delta _{I(\sigma _3,\sigma ^{})}^{I_1}\delta _{I\sigma _2}^{I_2})`$ (D24)
$`+f_\sigma ^{}^{\sigma _3}(\delta _{\sigma _2}^{\sigma _3}\delta _{I(\sigma _1,\sigma ^{})}^{I_1}\delta _{I\sigma _4}^{I_2}\delta _{\sigma _3}^{\sigma _4}\delta _{I(\sigma _1,\sigma ^{})}^{I_1}\delta _{I\sigma _2}^{I_2})+f_\sigma ^{}^{\sigma _1}f_\sigma ^{}^{\sigma _3}(\delta _{\sigma _1}^{\sigma _2}\delta _{\sigma _3}^{\sigma _4}\delta _{\sigma _2}^{\sigma _3}\delta _{\sigma _1}^{\sigma _4})\delta _{I\sigma ^{}}^{I_1}\delta _I^{I_2}]`$ (D25)
$`+\delta _\sigma ^{}^{\sigma _2}f_\sigma ^{}^{\sigma _1}f_\sigma ^{}^{\sigma _3}\left(f_{\sigma _1}^{\sigma _3}\delta _{I(\sigma _1,\sigma _3)}^{I_1}\delta _{I\sigma _4}^{I_2}+\delta _{\sigma _3}^{\sigma _4}\delta _{I\sigma _1}^{I_1}\delta _I^{I_2}\delta _{\sigma _1}^{\sigma _4}\delta _{I\sigma _3}^{I_1}\delta _I^{I_2}\right)+`$ (D26)
$`\delta _\sigma ^{}^{\sigma _4}f_\sigma ^{}^{\sigma _1}f_\sigma ^{}^{\sigma _3}(f_{\sigma _1}^{\sigma _3}\delta _{I(\sigma _1,\sigma _3)}^{I_1}\delta _{I\sigma _2}^{I_2}+\delta _{\sigma _3}^{\sigma _2}\delta _{I\sigma _1}^{I_1}\delta _I^{I_2}\delta _{\sigma _1}^{\sigma _2}\delta _{I\sigma _3}^{I_1}\delta _I^{I_2})].`$ (D27)
Using (D13) we may now define the following vertex-functions for a single creation-operator as it occurs in (D2)
$`Q_𝔇^+(\sigma ,\sigma ^{})`$ $`=`$ $`{\displaystyle \underset{\mathrm{\Gamma },\mathrm{\Gamma }^{}}{}}\lambda _\mathrm{\Gamma }\lambda _\mathrm{\Gamma }^{}{\displaystyle \underset{I(\sigma I)}{}}f_\sigma ^IT_{\mathrm{\Gamma },I}^+T_{I\sigma ,\mathrm{\Gamma }^{}}{\displaystyle \underset{I_1,I_2}{}}T_{I_1,\mathrm{\Gamma }}T_{\mathrm{\Gamma }^{},I_2}^+_{I_1,I_2}^𝔇(\sigma ^{}),`$ (D29)
$`\stackrel{~}{Q}_{𝔇𝔇^{}}^+(\sigma ,\sigma ^{})`$ $`=`$ $`{\displaystyle \underset{\mathrm{\Gamma },\mathrm{\Gamma }^{}}{}}\lambda _\mathrm{\Gamma }\lambda _\mathrm{\Gamma }^{}{\displaystyle \underset{I(\sigma I)}{}}f_\sigma ^IT_{\mathrm{\Gamma },I}^+T_{I\sigma ,\mathrm{\Gamma }^{}}{\displaystyle \underset{I_1,I_2}{}}T_{I_1,\mathrm{\Gamma }}T_{\mathrm{\Gamma }^{},I_2}^+\stackrel{~}{}_{I_1,I_2}^{𝔇𝔇^{}}(\sigma ^{})`$ (D30)
The respective vertices $`Q_𝔇(\sigma ,\sigma ^{})`$, $`\stackrel{~}{Q}_{𝔇𝔇^{}}(\sigma ,\sigma ^{})`$ for annihilation-operators are given as
$`Q_{\left(\sigma _1\sigma _2\right)}(\sigma ,\sigma ^{})`$ $`=`$ $`\left[Q_{\left(\sigma _2\sigma _1\right)}^+(\sigma ,\sigma ^{})\right]^{},`$ (D32)
$`\stackrel{~}{Q}_{\left(\sigma _1\sigma _2\right)(\sigma _3\sigma _4)}(\sigma ,\sigma ^{})`$ $`=`$ $`\left[\stackrel{~}{Q}_{\left(\sigma _2\sigma _1\right)(\sigma _4\sigma _3)}^+(\sigma ,\sigma ^{})\right]^{}.`$ (D33)
Now we can determine the connected diagrams in (D2). First, we have the following terms for $`ij`$
$`\epsilon ^2(\stackrel{}{q})`$ $`=`$ $`{\displaystyle \frac{N_\text{G}}{N_\stackrel{}{q}}}{\displaystyle \underset{\sigma _1,\sigma _2}{}}\sqrt{q_{\sigma _1}q_{\sigma _2}}{\displaystyle \underset{𝔇_1,𝔇_2}{}}\left(\stackrel{~}{U}_{𝔇_1,𝔇_2}^{\sigma _1,\sigma _2}(\stackrel{}{q})+\stackrel{~}{U}_{𝔇_2,𝔇_1}^{\sigma _1,\sigma _2}(\stackrel{}{q})\right)`$ (D36)
$`\times \left[\left(\widehat{1}+\stackrel{~}{x}\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\right)\stackrel{}{S}^+\right]_{𝔇_1}\left[\left(\widehat{1}+\stackrel{~}{x}\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\right)\stackrel{}{S}^{}\right]_{𝔇_2},`$
$`\epsilon ^3(\stackrel{}{q})`$ $`=`$ $`{\displaystyle \frac{N_\text{G}}{N_\stackrel{}{q}}}{\displaystyle \underset{\sigma _1,\sigma _2}{}}\sqrt{q_{\sigma _1}q_{\sigma _2}}{\displaystyle \underset{𝔇_1,𝔇_2,𝔇_3,𝔇_4}{}}V_{\left(\sigma _1\sigma _2\right)\left(\sigma _1\sigma _2\right)}^{𝔇_1}(\stackrel{}{0})\left[\stackrel{~}{x}\left(\widehat{1}+\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{0})\stackrel{~}{x}\right)\right]_{𝔇_1,𝔇_2}`$ (D39)
$`\times \left(\stackrel{~}{W}_{𝔇_3,𝔇_4}^{𝔇_2}(\stackrel{}{q})+\stackrel{~}{W}_{𝔇_4,𝔇_3}^{𝔇_2}(\stackrel{}{q})\right)`$
$`\times \left[\left(\widehat{1}+\stackrel{~}{x}\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\right)\stackrel{}{S}^+\right]_{𝔇_3}\left[\left(\widehat{1}+\stackrel{~}{x}\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\right)\stackrel{}{S}^{}\right]_{𝔇_4},`$
$`\epsilon ^4(\stackrel{}{q})`$ $`=`$ $`{\displaystyle \frac{N_\text{G}}{N_\stackrel{}{q}}}{\displaystyle \underset{\sigma _1,\sigma _2}{}}\sqrt{q_{\sigma _1}q_{\sigma _2}}{\displaystyle \underset{𝔇_1,𝔇_2,𝔇_3,𝔇_4}{}}V_{\left(\sigma _1\sigma _2\right)\left(\sigma _1\sigma _2\right)}^{𝔇_1}(\stackrel{}{0})\left[\widehat{1}+\stackrel{~}{x}\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{0})\right]_{𝔇_1,𝔇_2}`$ (D41)
$`\times \stackrel{~}{\xi }_{𝔇_2}^{𝔇_3,𝔇_4}\left[\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\stackrel{}{S}^+\right]_{𝔇_3}\left[\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\stackrel{}{S}^{}\right]_{𝔇_4},`$
$`\epsilon ^5(\stackrel{}{q})`$ $`=`$ $`{\displaystyle \frac{N_\text{G}}{N_\stackrel{}{q}}}{\displaystyle \underset{\sigma _1,\sigma _2,\sigma _1^{},\sigma _2^{}}{}}{\displaystyle \underset{𝔇_1,𝔇_2}{}}\left(Q_{\overline{𝔇}_1}^+(\sigma _1,\sigma _1^{})Q_{𝔇_2}(\sigma _2,\sigma _2^{})+Q_{\overline{𝔇}_2}^+(\sigma _1,\sigma _1^{})Q_{𝔇_1}(\sigma _2,\sigma _2^{})\right)`$ (D43)
$`\times E_{\sigma _1^{},\sigma _2^{}}^{\sigma _1,\sigma _2}(\stackrel{}{q})\left[\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\stackrel{}{S}^+\right]_{𝔇_1}\left[\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\stackrel{}{S}^{}\right]_{𝔇_2},`$
$`\epsilon ^6(\stackrel{}{q})`$ $`=`$ $`{\displaystyle \frac{N_\text{G}}{N_\stackrel{}{q}}}{\displaystyle \underset{\sigma _1,\sigma _2,\sigma _1^{}}{}}\sqrt{q_{\sigma _2}}{\displaystyle \underset{𝔇_1,𝔇_2}{}}\left(Q_{\overline{𝔇}_1}^+(\sigma _1,\sigma _1^{})\overline{V}_{\left(\sigma _1^{}\sigma _2\right)\left(\sigma _1\sigma _2\right)}^{𝔇_2}(\stackrel{}{q})+Q_{\overline{𝔇}_1}(\sigma _1,\sigma _1^{})V_{\left(\sigma _2\sigma _1^{}\right)\left(\sigma _2\sigma _1\right)}^{𝔇_2}(\stackrel{}{q})\right)`$ (D46)
$`\times \{[(\widehat{1}+\stackrel{~}{x}\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q}))\stackrel{}{S}^+]_{𝔇_2}[\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\stackrel{}{S}^{}]_{𝔇_1}`$
$`+[\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\stackrel{}{S}^+]_{𝔇_1}[(\widehat{1}+\stackrel{~}{x}\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q}))\stackrel{}{S}^{}]_{𝔇_2}\},`$
$`\epsilon ^7(\stackrel{}{q})`$ $`=`$ $`{\displaystyle \frac{N_\text{G}}{N_\stackrel{}{q}}}{\displaystyle \underset{\sigma _1,\sigma _2,\sigma _1^{}}{}}\sqrt{q_{\sigma _2}}{\displaystyle \underset{𝔇_1,𝔇_2,𝔇_3,𝔇_4}{}}\left(Q_{\overline{𝔇}_1}^+(\sigma _1,\sigma _1^{})E_{\left(\sigma _1^{}\sigma _2\right)}^{\left(\sigma _1\sigma _2\right)}(\stackrel{}{0})+Q_{\overline{𝔇}_1}(\sigma _1,\sigma _1^{})E_{\left(\sigma _2\sigma _1^{}\right)}^{\left(\sigma _2\sigma _1\right)}(\stackrel{}{0})\right)`$ (D49)
$`\times \left[\widehat{1}+\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{0})\stackrel{~}{x}\right]_{𝔇_1,𝔇_2}\left(W_{𝔇_3,𝔇_4}^{𝔇_2}(\stackrel{}{q})+W_{𝔇_4,𝔇_3}^{𝔇_2}(\stackrel{}{q})\right)`$
$`\times \left[\left(\widehat{1}+\stackrel{~}{x}\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\right)\stackrel{}{S}^+\right]_{𝔇_3}\left[\left(\widehat{1}+\stackrel{~}{x}\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\right)\stackrel{}{S}^{}\right]_{𝔇_4},`$
$`\epsilon ^8(\stackrel{}{q})`$ $`=`$ $`{\displaystyle \frac{N_\text{G}}{N_\stackrel{}{q}}}{\displaystyle \underset{\sigma _1,\sigma _2,\sigma _1^{}}{}}\sqrt{q_{\sigma _2}}{\displaystyle \underset{𝔇_1,𝔇_2,𝔇_3,𝔇_4}{}}\left(Q_{\overline{𝔇}_1}^+(\sigma _1,\sigma _1^{})E_{\left(\sigma _1^{}\sigma _2\right)}^{\left(\sigma _1\sigma _2\right)}(\stackrel{}{0})+Q_{\overline{𝔇}_1}(\sigma _1,\sigma _1^{})E_{\left(\sigma _2\sigma _1^{}\right)}^{\left(\sigma _2\sigma _1\right)}(\stackrel{}{0})\right)`$ (D51)
$`\times \stackrel{~}{\mathrm{\Omega }}_{\overline{𝔇}_1}^{𝔇_2}(\stackrel{}{q})\stackrel{~}{\xi }_{𝔇_2}^{𝔇_3,𝔇_4}\left[\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\stackrel{}{S}^+\right]_{𝔇_3}\left[\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\stackrel{}{S}^{}\right]_{𝔇_4},`$
$`\epsilon ^9(\stackrel{}{q})`$ $`=`$ $`{\displaystyle \frac{N_\text{G}}{N_\stackrel{}{q}}}{\displaystyle \underset{\sigma _1,\sigma _2,\sigma _1^{}}{}}\sqrt{q_{\sigma _2}}{\displaystyle \underset{𝔇_1,𝔇_2}{}}\left(\stackrel{~}{Q}_{\overline{𝔇}_1𝔇_2}^+(\sigma _1,\sigma _1^{})E_{\left(\sigma _1^{}\sigma _2\right)}^{\left(\sigma _1\sigma _2\right)}(\stackrel{}{0})+\stackrel{~}{Q}_{\overline{𝔇}_1𝔇_2}(\sigma _1,\sigma _1^{})E_{\left(\sigma _2\sigma _1^{}\right)}^{\left(\sigma _2\sigma _1\right)}(\stackrel{}{0})\right)`$ (D53)
$`\times \left[\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\stackrel{}{S}^+\right]_{𝔇_1}\left[\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\stackrel{}{S}^{}\right]_{𝔇_2},`$
$`\epsilon ^{10}(\stackrel{}{q})`$ $`=`$ $`{\displaystyle \frac{N_\text{G}}{N_\stackrel{}{q}}}{\displaystyle \underset{\sigma _1,\sigma _2,\sigma _1^{}}{}}\sqrt{q_{\sigma _2}}{\displaystyle \underset{𝔇_1}{}}\left(Q_{\overline{𝔇}_1}^+(\sigma _1,\sigma _1^{})E_{\left(\sigma _1^{}\sigma _2\right)}^{\left(\sigma _1\sigma _2\right)}(\stackrel{}{0})+Q_{\overline{𝔇}_1}(\sigma _1,\sigma _1^{})E_{\left(\sigma _2\sigma _1^{}\right)}^{\left(\sigma _2\sigma _1\right)}(\stackrel{}{0})\right)`$ (D55)
$`\times [\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})(\stackrel{+}{T}\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\stackrel{}{S}^{}+\stackrel{}{T}\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\stackrel{}{S}^+)]_{\overline{𝔇}_1},`$
$`\epsilon ^{11}(\stackrel{}{q})`$ $`=`$ $`{\displaystyle \frac{N_\text{G}}{N_\stackrel{}{q}}}{\displaystyle \underset{\sigma _1,\sigma _2}{}}\sqrt{q_{\sigma _1}q_{\sigma _2}}{\displaystyle \underset{𝔇_1}{}}V_{\left(\sigma _1\sigma _2\right)\left(\sigma _1\sigma _2\right)}^{𝔇_1}(\stackrel{}{0})`$ (D57)
$`\times [(\widehat{1}+\stackrel{~}{x}\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q}))(\stackrel{+}{T}\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\stackrel{}{S}^{}+\stackrel{}{T}\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\stackrel{}{S}^+)]_{𝔇_1}.`$
Expressions for the diagrams $`E_{𝔇_2}^{𝔇_1}`$, $`U_{𝔇_2𝔇_3}^{𝔇_1}`$, $`V_{𝔇_2𝔇_3}^{𝔇_1}`$, and $`\overline{V}_{𝔇_2𝔇_3}^{𝔇_1}`$ can be found in Appendix (B). Using eq. (C28) we derive the connected diagrams in (D2) with $`i=j`$ as
$`\epsilon ^{12}(\stackrel{}{q})`$ $`=`$ $`{\displaystyle \frac{N_\text{G}}{N_\stackrel{}{q}}}{\displaystyle \underset{\sigma _1,\sigma _2}{}}\sqrt{q_{\sigma _1}q_{\sigma _2}}{\displaystyle \underset{𝔇_1}{}}{\displaystyle \underset{\mathrm{\Gamma }^{}}{}}\lambda _\mathrm{\Gamma }^{}^2S_{}(\mathrm{\Gamma }^{})V_{\left(\sigma _1\sigma _2\right)\left(\sigma _1\sigma _2\right)}^{𝔇_1}(\stackrel{}{0})\left[\left(\widehat{1}+\stackrel{~}{x}\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{0})\right)\stackrel{}{M}(\mathrm{\Gamma }^{})\right],`$ (D59)
$`\epsilon ^{13}(\stackrel{}{q})`$ $`=`$ $`{\displaystyle \frac{N_\text{G}}{N_\stackrel{}{q}}}{\displaystyle \underset{\sigma _1,\sigma _2,\sigma _1^{}}{}}\sqrt{q_{\sigma _2}}{\displaystyle \underset{𝔇_1}{}}{\displaystyle \underset{\mathrm{\Gamma }^{}}{}}\lambda _\mathrm{\Gamma }^{}^2S_{}(\mathrm{\Gamma }^{})\left(Q_{\overline{𝔇}_1}^+(\sigma _1,\sigma _1^{})E_{\left(\sigma _1^{}\sigma _2\right)}^{\left(\sigma _1\sigma _2\right)}(\stackrel{}{0})+Q_{\overline{𝔇}_1}(\sigma _1,\sigma _1^{})E_{\left(\sigma _2\sigma _1^{}\right)}^{\left(\sigma _2\sigma _1\right)}(\stackrel{}{0})\right)`$ (D61)
$`\times \left[\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{0})\stackrel{}{M}(\mathrm{\Gamma }^{})\right]_{\overline{𝔇}_1}.`$
For the evaluation of (D3) we need the vertex-function for the operators
$`\widehat{P}_\text{G}\widehat{c}_\sigma \widehat{S}^{}\widehat{P}_\text{G}`$ $`=`$ $`{\displaystyle \underset{\mathrm{\Gamma },\mathrm{\Gamma }^{}}{}}\lambda _\mathrm{\Gamma }\lambda _\mathrm{\Gamma }^{}\sqrt{S_{}(\mathrm{\Gamma }^{})}{\displaystyle \underset{I(\sigma I)}{}}f_\sigma ^IT_{\mathrm{\Gamma },I}^+T_{I\sigma ,\mathrm{\Gamma }_{}^{}}{\displaystyle \underset{I_1,I_2}{}}T_{I_1,\mathrm{\Gamma }}T_{\mathrm{\Gamma }^{},I_2}^+\widehat{m}_{I_1,I_2},`$ (D63)
$`\widehat{P}_\text{G}\widehat{c}_\sigma ^+\widehat{S}^{}\widehat{P}_\text{G}`$ $`=`$ $`{\displaystyle \underset{\mathrm{\Gamma },\mathrm{\Gamma }^{}}{}}\lambda _\mathrm{\Gamma }\lambda _\mathrm{\Gamma }^{}\sqrt{S_{}(\mathrm{\Gamma }^{})}{\displaystyle \underset{I(\sigma I)}{}}f_\sigma ^IT_{\mathrm{\Gamma },I\sigma }^+T_{I,\mathrm{\Gamma }_{}^{}}{\displaystyle \underset{I_1,I_2}{}}T_{I_1,\mathrm{\Gamma }}T_{\mathrm{\Gamma }^{},I_2}^+\widehat{m}_{I_1,I_2}`$ (D64)
with one (or two) incoming or outgoing lines. These functions follow directly from eqs. (D18) and (D22)
$`r(\sigma ,\sigma ^{})`$ $`=`$ $`{\displaystyle \underset{\mathrm{\Gamma },\mathrm{\Gamma }^{}}{}}\lambda _\mathrm{\Gamma }\lambda _\mathrm{\Gamma }^{}\sqrt{S_{}(\mathrm{\Gamma }^{})}{\displaystyle \underset{I(\sigma I)}{}}f_\sigma ^IT_{\mathrm{\Gamma },I}^+T_{I\sigma ,\mathrm{\Gamma }_{}^{}}{\displaystyle \underset{I_1,I_2}{}}T_{I_1,\mathrm{\Gamma }}T_{\mathrm{\Gamma }^{},I_2}^+𝔥_{I_2,I_1}(\sigma ^{}),`$ (D66)
$`r^+(\sigma ,\sigma ^{})`$ $`=`$ $`{\displaystyle \underset{\mathrm{\Gamma },\mathrm{\Gamma }^{}}{}}\lambda _\mathrm{\Gamma }\lambda _\mathrm{\Gamma }^{}\sqrt{S_{}(\mathrm{\Gamma }^{})}{\displaystyle \underset{I(\sigma I)}{}}f_\sigma ^IT_{\mathrm{\Gamma },I\sigma }^+T_{I,\mathrm{\Gamma }_{}^{}}{\displaystyle \underset{I_1,I_2}{}}T_{I_1,\mathrm{\Gamma }}T_{\mathrm{\Gamma }^{},I_2}^+𝔥_{I_1,I_2}(\sigma ^{}),`$ (D67)
$`R_{\sigma _1,\sigma _2}(\sigma ,\sigma ^{})`$ $`=`$ $`{\displaystyle \underset{\mathrm{\Gamma },\mathrm{\Gamma }^{}}{}}\lambda _\mathrm{\Gamma }\lambda _\mathrm{\Gamma }^{}\sqrt{S_{}(\mathrm{\Gamma }^{})}{\displaystyle \underset{I(\sigma I)}{}}f_\sigma ^IT_{\mathrm{\Gamma },I}^+T_{I\sigma ,\mathrm{\Gamma }_{}^{}}{\displaystyle \underset{I_1,I_2}{}}T_{I_1,\mathrm{\Gamma }}T_{\mathrm{\Gamma }^{},I_2}^+_{I_2,I_1}^{\left(\sigma _2\sigma _1\right)}(\sigma ^{}),`$ (D68)
$`R_{\left(\sigma _1\sigma _2\right)}^+(\sigma ,\sigma ^{})`$ $`=`$ $`{\displaystyle \underset{\mathrm{\Gamma },\mathrm{\Gamma }^{}}{}}\lambda _\mathrm{\Gamma }\lambda _\mathrm{\Gamma }^{}\sqrt{S_{}(\mathrm{\Gamma }^{})}{\displaystyle \underset{I(\sigma I)}{}}f_\sigma ^IT_{\mathrm{\Gamma },I\sigma }^+T_{I,\mathrm{\Gamma }_{}^{}}{\displaystyle \underset{I_1,I_2}{}}T_{I_1,\mathrm{\Gamma }}T_{\mathrm{\Gamma }^{},I_2}^+_{I_1,I_2}^{\left(\sigma _1\sigma _2\right)}(\sigma ^{}).`$ (D69)
Note that for the vertices of the remaining operators $`\widehat{P}_\text{G}\widehat{S}^+\widehat{c}_\sigma \widehat{P}_\text{G}`$ and $`\widehat{P}_\text{G}\widehat{S}^+\widehat{c}_\sigma ^+\widehat{P}_\text{G}`$ the rules (D) apply. Now we are able to summarize the contributions (D3) as
$`\epsilon ^{14}(\stackrel{}{q})`$ $`=`$ $`{\displaystyle \frac{N_\text{G}}{N_\stackrel{}{q}}}{\displaystyle \underset{\sigma _1,\sigma _2,\sigma _1^{}}{}}\sqrt{q_{\sigma _2}}{\displaystyle \underset{𝔇_1}{}}\left(r^+(\sigma _1,\sigma _1^{})\overline{V}_{\left(\sigma _1^{}\sigma _2\right)\left(\sigma _1\sigma _2\right)}^{𝔇_1}(\stackrel{}{q})+r(\sigma _1,\sigma _1^{})V_{\left(\sigma _2\sigma _1^{}\right)\left(\sigma _2\sigma _1\right)}^{𝔇_1}(\stackrel{}{q})\right)`$ (D72)
$`\times [(\widehat{1}+\stackrel{~}{x}\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q}))\stackrel{}{S}^{}]_{𝔇_1}+c.c.,`$
$`\epsilon ^{15}(\stackrel{}{q})`$ $`=`$ $`{\displaystyle \frac{N_\text{G}}{N_\stackrel{}{q}}}{\displaystyle \underset{\sigma _1,\sigma _2,\sigma _1^{}}{}}\sqrt{q_{\sigma _2}}{\displaystyle \underset{𝔇_1}{}}\left(R_{\overline{𝔇}_1}^+(\sigma _1,\sigma _1^{})E_{\left(\sigma _1^{}\sigma _2\right)}^{\left(\sigma _1\sigma _2\right)}(\stackrel{}{0})+R_{\overline{𝔇}_1}(\sigma _1,\sigma _1^{})E_{\left(\sigma _2\sigma _1^{}\right)}^{\left(\sigma _2\sigma _1\right)}(\stackrel{}{0})\right)`$ (D74)
$`\times [\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\stackrel{}{S}^{}]_{𝔇_1}+c.c,`$
$`\epsilon ^{16}(\stackrel{}{q})`$ $`=`$ $`{\displaystyle \frac{N_\text{G}}{N_\stackrel{}{q}}}{\displaystyle \underset{\sigma _1,\sigma _2,\sigma _1^{},\sigma _2^{}}{}}{\displaystyle \underset{𝔇_1}{}}\left(r^+(\sigma _1,\sigma _1^{})Q_{\overline{𝔇}_1}(\sigma _2,\sigma _2^{})+r(\sigma _2,\sigma _2^{})Q_{\overline{𝔇}_1}^+(\sigma _1,\sigma _1^{})\right)E_{\left(\sigma _1^{}\sigma _2^{}\right)}^{\left(\sigma _1\sigma _2\right)}(\stackrel{}{q})`$ (D76)
$`\times [\stackrel{~}{\mathrm{\Omega }}(\stackrel{}{q})\stackrel{}{S}^{}]_{𝔇_1}+c.c..`$
The contributions from eq. (D4) are
$$\epsilon ^{17}(\stackrel{}{q})=\frac{N_\text{G}}{N_\stackrel{}{q}}\underset{\sigma _1,\sigma _2,\sigma _1^{},\sigma _2^{}}{}\left(r^+(\sigma _1,\sigma _1^{})\left(r^+(\sigma _2,\sigma _2^{})\right)^{}+\left(r(\sigma _1,\sigma _1^{})\right)^{}r(\sigma _2,\sigma _2^{})\right)E_{\left(\sigma _1^{}\sigma _2^{}\right)}^{\left(\sigma _1\sigma _2\right)}(\stackrel{}{q}).$$
(D77)
Finally, we need the vertex-function for the operator $`\widehat{P}_\text{G}\widehat{S}^+\widehat{c}_\sigma ^+\widehat{S}^{}\widehat{P}_\text{G}`$ in (D5) with one outgoing line
$$l(\sigma ,\sigma ^{})=\underset{\mathrm{\Gamma },\mathrm{\Gamma }^{}}{}\lambda _\mathrm{\Gamma }\lambda _\mathrm{\Gamma }^{}\sqrt{S_{}(\mathrm{\Gamma }^{})S_{}(\mathrm{\Gamma })}\underset{I(\sigma I)}{}f_\sigma ^IT_{\mathrm{\Gamma }_{},I}^+T_{I\sigma ,\mathrm{\Gamma }_{}^{}}\underset{I_1,I_2}{}T_{I_1,\mathrm{\Gamma }}T_{\mathrm{\Gamma }^{},I_2}^+𝔥_{I_1,I_2}(\sigma ^{}).$$
(D78)
in (D5). This gives us the contribution from eq. (D5),
$$\epsilon ^{18}(\stackrel{}{q})=\underset{\sigma _1,\sigma _2,\sigma _1^{}}{}\sqrt{q_{\sigma _2}}l(\sigma _1,\sigma _1^{})E_{\left(\sigma _1^{}\sigma _2\right)}^{\left(\sigma _1\sigma _2\right)}(\stackrel{}{0})+c.c.$$
(D79)
Altogether we may write the expectation value for the kinetic energy (D) as
$$\frac{\mathrm{\Psi }_\stackrel{}{q}^\text{G}\left|\widehat{H}_1\right|\mathrm{\Psi }_\stackrel{}{q}^\text{G}}{N_\stackrel{}{q}}=E_{\text{kin}}+\underset{c=1}{\overset{18}{}}\epsilon ^c(\stackrel{}{q}).$$
(D80)
|
warning/0005/astro-ph0005154.html
|
ar5iv
|
text
|
# Early Results from the Chandra X-ray Observatory
## 1 Introduction
The Chandra X-ray Observatory (CXO, or Chandra) was launched in July, 1999 into a high Earth orbit. Although there has been some degradation of the spectral resolution of the Advanced CCD Imaging Spectrometer (ACIS), all instruments are performing well. Here I will focus on data from ACIS and the High Energy Transmission Grating Spectrometer (HETGS), in order to demonstrate the overall performance of the telescope and some of the ways that Chandra will contribute to our understanding of active galactic nuclei (AGN).
## 2 Chandra Observations of AGN
The first target observed with Chandra was the quasar PKS 0637–752 (figure 1). It was rather surprising to find an X-ray jet extending 7-12” from the quasar core because the target had been chosen for focus measurements. The core point response function is small enough ($``$ 0.42”, half power radius), however, that the jet had no effect on determining the best detector focus position. New radio observations showed that the radio and X-ray jets were remarkably similar out to 12” where there is a significant bend in the radio emission (figure 1). The optical fluxes from the Hubble data are so low that models of the jet X-ray emission are difficult to constrain. SM00 .
The HETGS flight calibration program included observations of the late type star Capella in order to verify the spectral resolution by using emission lines, and an observation of the BL Lac object Mrk 421 was included in order to verify the effective area for point sources. An observation of 3C 273 was scheduled for January, 2000, which will be used for cross-calibration between Chandra spectrometers and other X-ray telescopes, including BeppoSAX and ASCA. The Capella observation C00 shows that the spectral resolution meets the pre-flight expectations. The calibration spectra of Mrk 421 are shown in Figure 2. The spectra from the MEG and the HEG are consistent with each other to within the 10-20% systematic uncertainties. The spectra are well fitted by a simple power law with $`\alpha =1.9`$ ($`f_\nu \nu ^\alpha `$), and the only deviation from this fit appears near the O-K edge, which will be removed as the detector calibration is refined.
The BL Lac object PKS 2155–304 was observed as part of the HETGS guaranteed time observation program. Spectra are shown in Figure 3. Again, the HEG and MEG spectra are consistent to within the 10-20% systematic unceratainties and are well fitted by a simple, pure power law model. The spectral index is $`1.70+/0.02`$ and there are no significant absorption features. A feature such as the one found previously CK84 would have been detected easily in the MEG portion of the Chandra HETGS spectrum, so we conclude that it must be variable.
Other AGN observed in the early phase of Chandra HETGS guaranteed time observations include NGC 1275 and PKS 2149-305. Preliminary analysis indicates that these HETGS spectra are all well fitted by simple power law models with absorption by neutral interstellar material. The spectral indices are somewhat smaller: 0.8 and 0.2, respectively. No Fe-K$`\alpha `$ lines are detected in any of their spectra.
## 3 Acknowledgements
I thank the Principal Investigator of the HETGS, Prof. Claude R. Canizares for his support and the rest of the HETGS team for their contributions to the development of the HETGS. This work was supported in part by the contract SAO SV1-61010.
|
warning/0005/hep-th0005087.html
|
ar5iv
|
text
|
# Marginal and Relevant Deformations of N=4 Field Theories and Non-Commutative Moduli Spaces of Vacua
## 1 Introduction
The study of the deformations of the $`N=4`$ $`U(M)`$ supersymmetric theory in four dimensions is of interest from several points of view. This theory is superconformally invariant, and it has been known for some time that exactly marginal deformations of this theory exist (Ref. and references therein) and should be described by interacting superconformal field theories (CFT). These CFT’s are largely unexplored.
In the large $`M`$-limit, the deformations of the $`N=4`$ theory have a nice description in terms of a supergravity dual and are obtained by the addition of operators which modify the boundary conditions at infinity. Each of the renormalizable deformations are reflected in the $`AdS/CFT`$ correspondence through backgrounds for massless and tachyonic excitations, including both $`RR`$ and $`NS`$ fields.
Among the marginal deformations, of particular interest is the $`q`$-deformation
$$W_q=\mathrm{tr}\left(\varphi _1\varphi _2\varphi _3q\varphi _2\varphi _1\varphi _3\right)$$
(1)
which is a deformation of the superpotential by the symmetric invariant preserving $`N=1`$ supersymmetry and a $`U(1)^3`$ global symmetry. For special values of $`q`$ these theories are described by the near-horizon geometries of orbifolds with discrete torsion. It was conjectured in Ref. that these orbifold theories are related by mirror symmetry to string theories on $`S^5`$-deformations of $`AdS_5\times S^5`$.
There are also relevant deformations which carry the theory away from the ultraviolet fixed point CFT. In some cases, the infrared theory is of interest—a prime example being the deformation by rank-one mass terms. From the supergravity point of view, the renormalization group flow is encoded as a dependence of the background on the radial scaling variable.
With a rank-three mass matrix, the field theory has been analyzed in many papers. More recently, the supergravity duals of these theories have been analyzed. There it was noticed that 5-brane sources resolve the would-be singularity in the dual supergravity background, an application of the dielectric effect.
In this paper, we will begin an exploration of these field theories obtained by marginal and relevant deformations of the $`N=4`$ theory. The analysis will concentrate on the classical vacua, particularly those aspects which depend upon holomorphic quantities. We introduce a new way of thinking about these moduli spaces that should be of quite general applicability.
Normally, the vacua of a supersymmetric gauge theory are parameterized using gauge invariant holomorphic polynomials in the fields. This is attractive because of the gauge invariance, but it is also unwieldy. As $`M`$ increases, the number of independent invariants increases dramatically. The $`F`$-term constraints on vacua are given, on the other hand, directly in terms of holomorphic matrix equations, and the proposal centres around using this description directly. Matrix variables have a number of technical advantages, principally that the analysis is independent of their dimension, $`M`$. The main problem with this approach is gauge invariance—the $`D`$-term constraints must be applied separately.
Given this, the $`F`$-terms can be thought of as a set of constraints on the algebra of $`M\times M`$ matrices. Generally, this is a non-commutative algebra. There is a technical simplicity to the choice of renormalizable superpotentials, namely that the constraints are quadratic, and this simplifies the algebraic analysis significantly. The constrained algebras that appear here in some cases bear some resemblance to algebras considered in the literature on quantum groups (see for example, Refs. ).
A related problem is the behavior of $`D`$-branes in dual descriptions of these field theories. For small deformations, these duals are close to $`AdS_5\times S^5`$, and therefore the moduli space of $`D`$-branes also has a description in this framework. The moduli space of probe $`D`$-branes is roughly a symmetric product space; indeed, classically, we can think of a single $`D`$-brane as moving on the moduli space of the corresponding field theory, and the moduli space for multiple branes can often be related to the direct product of this space, modded out by the permutation group. Realizing all aspects of the field theory analysis in these dual descriptions gives insight into many non-trivial aspects of $`D`$-brane geometry. In particular, a clear understanding of these issues reveals a T-duality transformation which realizes mirror symmetry between near-horizon geometries and orbifold theories.
In the present context, these remarks lead to the notions of non-commutative moduli spaces of vacua, and moduli spaces of $`D`$-brane configurations are symmetric products of a non-commutative space. We construct these notions algebraically; for example, points in the non-commutative space correspond to irreducible representations of the algebra or equivalently to maximal ideals with special properties. As we show later, these notions have some tremendous advantages over the standard points of view. In particular, it is often the case that we can think of a commutative subring (built out of the center of the algebra) as a sort of coarse view of the full moduli space. In fact, the phenomenon of $`D`$-brane fractionation at singularities follows precisely this rule: the fractionation is present in a commutative description, but from the full non-commutative point of view, the fractional nature is more readily understandable.
It is clear that from a string theory point of view, it is the open strings that see this non-commutative structure directly. Closed strings appear naturally within this framework as single trace operators, and thus should see only the commutative part of the space. The remnant of non-commutativity in the closed string sector is the presence of twisted states.
In general, when one studies more general configurations of $`D`$-branes, they should correspond to algebraic geometric objects and classes in K-theory. Because of our emphasis, one needs to develop a non-commutative version of algebraic geometry and K-theory. K-theory in this context is provided by the algebraic K-theory of the non-commutative ring. We give the rudimentary structure of such a definition of algebraic geometry; this definition apparently differs from others given in the mathematics literature, but we believe our proposal is more natural, as dictated by string theory.
Our version of non-commutative geometry is clearly different than that which has been recently studied extensively (see for example and citations thereof). In that case, the non-commutativity occurs in the base space of a super-Yang Mills theory, whereas here it is in the moduli space, namely the directions transverse to a brane. In many cases, the boundary state formalism (for a review, see Ref. ) is convenient to describe $`D`$-branes, but we do not use that technology here. It would be interesting to generalize our discussion to that formalism, although it is not clear to us how to solve for the boundary states in the absence of a spectrum-generating algebra.
The paper is organized as follows. In Section 2, we review the structure of marginal and relevant deformations of the $`N=4`$ theory, and discuss the interpretation of supersymmetric vacua in terms of non-commutative geometry. (We focus throughout on $`U(M)`$ gauge groups.) We also review the map between the superpotential deformations and supergravity backgrounds. In Section 3, we give our construction of non-commutative algebraic geometry and the resulting K-theory. This section is mathematically intensive; in order that the casual reader may skip this section if desired, we provide at the beginning, an overview of the key structures. In Section 4, we investigate the vacua of various field theories, using the non-commutative formalism. We begin with the $`q`$-deformed theory, and then investigate this theory with (a) a single mass term, (b) a mass term and a linear term, (c) three arbitrary linear terms, and (d) three mass terms. In each case, we work out the representation theory. We also consider the general case, and in particular consider the effects of the other independent marginal deformation. The general case is quite difficult, but we are able to identify a few interesting properties.
In Section 5, we turn our attention to string theory. For the $`q`$-deformed theory, the field theory predicts new branches in moduli space for arbitrarily small values of $`q1`$. To realize this branch in string theory, we need to consider BPS states corresponding to $`D5`$-branes with 3-brane charge in the deformed backgrounds; the physics here is reminiscent of the dielectric effect, but is more general. The new branch of moduli space is identified as a $`D5`$-brane wrapped on a degenerating 2-torus. One obtains a natural 2-torus fibration of the 5-sphere; T-duality on this torus leads to the mirror orbifold theory.
In Section 6, we consider the problem of identifying closed string physics directly from the field theory description. Closed string states are naturally identified with single trace operators. In this section, we also note several features of interest, including connections to quantum groups and the K-theory of the non-commutative geometry.
In Section 7, we make some final remarks and indicate avenues for further research.
## 2 Field theory deformations
Our first objective will be to analyze the marginal and relevant deformations of the $`N=4`$ super Yang-Mills (SYM) theory in four dimensions, with gauge group $`U(M)`$. As usual, we write this in terms of an $`N=1`$ SYM theory with three adjoint chiral superfields $`\varphi _i`$, $`i=1,2,3`$, coupled through the superpotential
$$W=g\mathrm{tr}\left([\varphi _1,\varphi _2]\varphi _3\right).$$
(2)
If we choose to preserve $`N=1`$ SCFT, there is a moduli space of marginal deformations, given by a general superpotential of the form
$$W_{marg}=a\mathrm{tr}\left(\varphi _1\varphi _2\varphi _3q\varphi _2\varphi _1\varphi _3+\frac{\lambda }{3}\left(\varphi _1^3+\varphi _2^3+\varphi _3^3\right)\right).$$
(3)
The Yang-Mills coupling $`g`$ measures how strongly interacting the theory is and is a function of $`a,q,\lambda `$ such that each of the $`\beta `$-functions are zero. The structure of the moduli space of vacua depends only on $`q,\lambda `$.
We will also consider relevant deformations of the form
$$W_{rel}=c_1\mathrm{tr}(\varphi _1^2)+c_2\mathrm{tr}(\varphi _2^2+\varphi _3^2)+\underset{j}{}\zeta _j\mathrm{tr}(\varphi _j).$$
(4)
For $`q1`$, general quadratic polynomials may always be brought to this form after a change of variables.
The vacua of the theory are found by solving the $`F`$-term constraints
$$\frac{W}{\varphi _j}=0.$$
(5)
In the present cases, these are quadratic matrix polynomial equations in the $`\varphi _j`$
$`\varphi _1\varphi _2q\varphi _2\varphi _1`$ $`=`$ $`\lambda \varphi _3^22c_2\varphi _3\zeta _3`$ (6)
$`\varphi _2\varphi _3q\varphi _3\varphi _2`$ $`=`$ $`\lambda \varphi _1^22c_1\varphi _1\zeta _1`$ (7)
$`\varphi _3\varphi _1q\varphi _1\varphi _3`$ $`=`$ $`\lambda \varphi _2^22c_2\varphi _2\zeta _2`$ (8)
These matrix equations are independent of $`M`$. In general, solutions will consist of a collection of points, but at special values of parameters, we get a full moduli space of vacua.
The equations (6)–(8) are a quite general class of relations. Note in particular that when $`q=1`$, we have a Poisson bracket structure, whereas if in addition $`\lambda =\zeta _i=0`$, we find $`SU(2)`$ commutation relations. When $`q1`$ and/or $`\lambda 0`$, with $`c_j=\zeta _j=0`$, the algebra is that of a quantum plane. When $`q1`$ and $`\lambda =c_j=0`$, these are $`q`$-deformations of Heisenberg algebras.
The moduli space of vacua is usually parameterized in terms of gauge-invariant polynomials in the fields $`\varphi _j`$. This has the feature that the non-holomorphic $`D`$-term constraints are automatically satisfied. The down side is that for large $`M`$, the number of polynomials required becomes very large, and when perturbations are present, the description of the space becomes quite complicated. Instead, we will choose to describe the moduli space of vacua directly in terms of matrix variables. This has the virtue that the equations are independent of $`M`$, as noted. Thus instead of considering the moduli space of vacua as an algebraic variety, for general values of parameters, we should think of this as a non-commutative algebraic variety.
Understanding the vacua of these theories then is equivalent to understanding the non-commutative geometry defined by the relations (6)–(8). The $`\varphi _j`$ can be thought of as the generators of the corresponding non-commutative algebra. $`M\times M`$ matrices which satisfy the relations are an $`M`$-dimensional representation of the abstract algebra. The general problem at hand then is to study the representation theory of the algebra. The basic representations of interest are those that are irreducible; given a finite set of such solutions $`(\varphi _1^i,\varphi _2^i,\varphi _3^i)`$ labeled by $`i`$, then
$$\stackrel{~}{\varphi }_k=_i\varphi _k^i$$
(9)
is also a solution of the matrix equations.
It is important to keep in mind however that we must also consider the $`D`$-term constraints. It is well-known that for every solution of the $`F`$-term constraints, there is a solution to the $`D`$-terms in the completion of the orbit of the complexified gauge group $`SL(M,)`$. If the solution occurs at a finite point in the orbit, then we get a true vacuum. If it occurs in the completion of the orbit (at infinity in the complexified gauge group), we need to check that the solution does not run away to infinity.
### 2.1 Relation to Deformations in $`AdS_5\times S^5`$ Geometry
Since the $`N=4`$ theory is related to superstring theory on $`AdS_5\times S^5`$, the marginal deformations correspond to deformations of $`S^5`$. In particular, these are related to massless states in the 5-dimensional supergravity. They transform in the $`\mathrm{𝟒𝟓}`$ of $`SU(4)`$ R-symmetry and are related to vevs for harmonics of $`RR`$ and $`NSNS`$ fields, $`F_{(3)}^{RR}`$ and $`H_{(3)}^{NS}`$, along the 5-sphere. Similarly, the relevant deformations correspond to tachyonic excitations of the 5-dimensional supergravity, and transform in the $`\mathrm{𝟏𝟎}`$ of $`SU(4)`$.
When $`q`$ is a root of unity, we will often for convenience say that $`q`$ is rational. In this case, the $`q`$-deformation is known to be dual to the near-horizon geometry of $`D`$-branes on an orbifold with discrete torsion, $`^3/(_n\times _n)`$.
The moduli space of vacua of the field theory is the moduli space of $`D`$-branes. Because of the $`RR`$ and $`NSNS`$ backgrounds, $`D`$-branes move on a non-commutative space. For small enough deformations, we expect the $`AdS_5\times S^5`$ geometry to be a close approximation and we can interpret the eigenvalues of matrices as the positions of $`D`$-branes, à la matrix theory.
Note that the superpotentials that we are considering are single trace operators. This suggests that these operators correspond to effects that may be seen in classical supergravity. This may be understood by looking at how background couplings to $`D`$-branes behave at weak coupling. The leading effect comes from a disk diagram as shown in Figure 1 where $`V`$ is the background vertex.
Multiple trace operators would then correspond to string loop diagrams, and are therefore suppressed by powers of $`g_{str}`$.<sup>1</sup><sup>1</sup>1We assume a large but finite number of branes, so that relations between traces of finite matrices appear only for very irrelevant perturbations. It is not clear that the string generates these effects perturbatively, but to avoid them we could work at weak string coupling. It is also possible that a non-perturbative non-renormalization theorem might keep multiple trace operators equal to zero, at least up to some number of derivatives.
## 3 Non-Commutative Algebraic Geometry
### 3.1 Overview
As discussed above, we usually think of the moduli spaces of vacua as varieties, namely commutative algebraic geometric objects. Because the $`F`$\- and $`D`$-term constraints may be recast in matrix form, it is more convenient to think of the same space as a non-commutative object. Although the physical problem to solve is the same and the space of solutions is the same, the non-commutative interpretation invokes extra structure (namely, the commutative algebra of individual matrix elements is organized into the non-commutative algebra of matrices). Because $`D`$-brane solutions are associated with algebraic geometric objects in general (and see for a review), we need a formulation of non-commutative algebraic geometry which captures $`D`$-brane physics correctly.
Several different versions of non-commutative algebraic geometry have been discussed in the mathematics literature, but none of these seem natural in the present context. In this section, our aim is to describe a definition of algebraic geometry appropriate to the moduli spaces of vacua of $`D`$-branes in the field theory limit.
We want to understand the holomorphic structure of the moduli space, so we will only concentrate on the $`F`$-term equations and will assume that given a solution to the $`F`$-terms, there is a solution to the $`D`$-term equations. From the physics point of view, we confine ourselves to those properties which are protected by supersymmetry; from a mathematical viewpoint, these are holomorphic structures and can be described in terms of algebraic geometry. Ordinarily, non-commutative geometries are related to $`C^{}`$-algebras, which include the adjoint operation; this is not a natural operation in a holomorphic framework, and we will discard it for the considerations of this paper. In supersymmetric theories, the holomorphic and anti-holomorphic features couple only through $`D`$-terms, and these effects are in general not protected by supersymmetry. In discarding the $`D`$-terms, we lose information about the metric in moduli space, but not topological features. Thus we need a framework where we can do non-commutative algebraic geometry without $`C^{}`$-algebras.
In the rest of this subsection, we give a brief outline of the mathematics involved, and its physical interpretations. For the reader who wishes to skip the details of the mathematics, this overview should suffice, and one may proceed to Section 4.
The building blocks for solutions are the finite dimensional irreducible representations of the algebra, as in eq. (9). In the non-commutative algebraic geometry that we will describe, these are defined to be points. In matrix theory, non-commuting matrices are interpreted as extended objects; here, these are considered point-like objects. In orbifold theories, D-branes in the bulk are also considered to be point-like even though they can be built out of fractional branes.
In our construction the center (the commutative sub-algebra of the Casimir operators) plays a pivotal role. In particular, one expects that for a given background in string theory, there are two descriptions, a commutative version relevant to closed strings and a non-commutative version for open strings. In our best-understood examples, the commutative space is the algebraic geometry associated to the center. By Schur’s lemma, on every irreducible representation the Casimirs are proportional to the identity. Because of this, one finds a map between non-commutative points and commutative points. For “good” algebras, the non-commutative geometry covers the commutative geometry; in this case, we will say that the non-commutative algebra is semi-classical. In this case we can think of the commutative space as a coarse-grained version of the full non-commutative space. For semi-classical algebras, our interpretation of irreducible representations as points is equivalent to the point-like properties of $`D`$-branes in orbifolds. In other cases, the non-commutative geometry may have little relation to the commutative geometry and it is not clear that one should interpret the $`D`$-brane states as pointlike.
A general solution of the $`F`$-term constraints is a direct sum of irreducible representations, and thus the natural non-commutative structure is an unordered finite collection of points. For commutative algebras, we would interpret this as a symmetric product space, and we will carry over this name in the non-commutative case. This symmetric product structure leads directly to an interpretation of $`D`$-brane fractionation at a singularity, whereby an irreducible representation can be continuously deformed and becomes reducible at a certain point. In this sense, in the non-commutative version, single-particle and multi-particle states are continuously connected. In commutative geometry on the other hand, this process would be singular.
The remainder of the formal discussion deals with an extension of this construction to subvarieties, sheaves, and the algebraic K-theory of the ring, relevant to an understanding of extended $`D`$-branes. The discussion presented here lays the outline for non-commutative algebraic geometry; a full account will appear elsewhere.
### 3.2 Preliminaries: Points and Topology
Consider an associative algebra $`𝒜`$ over the complex numbers $``$, generated by a set of operators subject to some relations. (In all the examples we consider, we will have three generators and a set of quadratic relations.) As is standard, the non-commutative algebra should be thought of as a ring of functions on some affine non-commutative space. Ring homomorphisms should correspond to holomorphic maps between affine non-commutative geometries. We will assume that all rings are Noetherian (so that any ideal always has a finite basis) and that the algebra is polynomial.
Given the matrix equations (6)–(8), we want to find solutions in terms of $`M\times M`$ matrices (with unspecified $`M`$). That is, we are interested in representations of the algebra, and we will assign a geometrical space to these solutions.
An element $`a𝒜`$ is central if it commutes with every other element in $`𝒜`$; that is, $`a`$ is a Casimir of the algebra. We usually think of the Casimir operators as sufficient to define a representation (e.g., as in the finite dimensional representations of $`SL(2,)`$) and thus will pay particular attention to the center of the algebra, denoted $`𝒵𝒜`$.
Since $`𝒵𝒜`$ is commutative, we can associate an ordinary commutative space to it, which is the general philosophy behind algebraic geometry (see for example, Ref. ). We interpret the center of the algebra as a coarse description of the full non-commutative geometry. The picture we have is that there is a map between the non-commutative geometry to the commutative one which forgets some of the structure (namely the functions that don’t commute). The natural inclusion $`i:𝒵𝒜𝒜`$ is to be thought of as the pullback of functions from the commutative space to the non-commutative space.
To describe the non-commutative space, we need to define the notions of points and open sets and to impose a topology. Loosely speaking, a point will be a solution of the constraint equations in finite matrices (just as points in varieties are solutions of the equations defining the variety). In commutative algebra, points are interpreted as maximal ideals of the algebra, and we want to incorporate both of these notions in our definition of a point.
Let us now be precise. A representation $`R`$ of dimension $`M`$ of the algebra $`𝒜`$ is an algebra homomorphism $`\mu `$ from $`𝒜`$ to the algebra of $`M\times M`$ matrices (i.e., $`\mu `$ respects the addition, product, and multiplication by scalars). For the map to be well defined, the relations ($`F`$-term constraints) must be satisfied in terms of the $`M\times M`$ matrices.
A representation is irreducible if there is no linear subspace of $`^M`$ which is invariant under multiplication by all the elements of the image of the algebra $`\mu (𝒜)`$. The representation is reducible otherwise. If a representation is irreducible, then the map is such that we have an exact sequence
$$𝒜\stackrel{\mu }{}M_M()0$$
(10)
with the map $`\mu `$ defined by the representation of the algebra. The kernel of this map is a double-sided ideal $``$ of $`𝒜`$ to which the representation is associated, and we have the isomorphism
$$𝒜/M_M().$$
(11)
This isomorphism is non-canonical (two representations are identified if they lie in the same orbit of the group $`GL(M,)`$ by similarity transformations).
The space associated to an algebra is constructed from the irreducible representations of $`𝒜`$ as follows. To each irreducible representation of finite dimension $`M`$, there is an associated ideal in the algebra $`𝒜`$, namely the ideal $`=\mathrm{ker}(\mu )`$. $``$ is a double-sided maximal ideal and is declared to be a point. This definition is borrowed from , but without the $`C^{}`$-algebra framework. In general, one would also allow infinite dimensional representations of the algebra; in that case, we would need some sense of convergence for sequences. For our physical problem, we are interested only in finite dimensional representations, and so we will simply discard this possibility. As a result, we have a space which is better behaved from an algebraic standpoint. Irreducible representations are considered equivalent if they are related by a change of basis (i.e., by orbits of the group $`GL(M,)`$). This equivalence is the fact that we have in supersymmetric field theories a complexified gauge group, and each point has an associated maximal double-sided ideal $``$ of the algebra $`𝒜`$ such that $`𝒜/`$ is non-canonically isomorphic to the algebra of $`M\times M`$ matrices. The variety associated to $`𝒜`$ will be labeled $`_𝒜`$.
A closed set is defined in terms of an arbitrary double-sided ideal $`^{}`$ in $`𝒜`$, as one does to define the Zariski topology of a space. A closed set is the collection of points (given by maximal ideals $``$) which contain $`^{}`$. By definition, points are closed sets, as one takes the maximal ideal $``$ associated to the point.
The union of two closed sets corresponds to the double-sided ideal $`=_1_2`$, and the intersection of two closed sets corresponds to the double-sided ideal $`_1+_2`$, which is the direct sum of the ideals. Indeed, direct sums may be extended to an infinite number of ideals, so arbitrary intersections and finite unions of closed sets are closed and define a topology on the set of points. This should be thought of as a model for the definition of the geometry, and the construction mimics the construction of algebraic varieties over $``$ as much as possible.
Note that the definition of a point has the following technical property. A point is Morita-equivalent to a point in a commutative algebraic variety. This is important for K-theory considerations, which we return to in a later subsection.
### 3.3 Naturalness of Symmetric Spaces
So far, we have defined a non-commutative space together with some topology. Given these definitions, additional structure naturally emerges, as we now discuss.
When one has the ring of functions of a variety, one can pull back functions between maps of varieties. Thus, a map between non-commutative spaces will correspond to ring homomorphisms. If we take two rings $`𝒜`$ and $``$ and consider a ring homomorphism $`\phi :𝒜`$, it will correspond to a continuous map from $`_{}_𝒜`$.
Now consider a point $`x_{}`$. By construction, it corresponds to an irreducible representation of $``$ in $`M\times M`$ matrices for some $`M`$. We label the corresponding representation $`r_x`$, i.e, a homomorphism $`\mu _x:M_M()`$. Thus we have a diagram of maps
$$𝒜\stackrel{\phi }{}\stackrel{\mu _x}{}M_M().$$
(12)
We wish to find the image of the point $`x`$ in $`_𝒜`$. By composing arrows, we get a natural representation of $`𝒜`$ in the ring of $`M\times M`$ matrices. The natural image of $`x`$ is the kernel of the composition map, but it is important to note that the corresponding $`M\times M`$ representation of $`𝒜`$ might be reducible. In general, to find irreducible representations associated to a given reducible one, we need to consider the composition series of the representation $`R`$.
That is, given a reducible representation $`R`$, there is an invariant linear subspace $`R^{}`$, and we have an exact sequence of vector spaces:
$$0R^{}RR/R^{}0.$$
(13)
By construction the dimensions of $`R^{}`$ and $`R/R^{}`$ are smaller than that of $`R`$. If any of $`R^{}`$, $`R/R^{}`$ is reducible, we repeat the procedure. Eventually, we obtain a collection of irreducible representations of $`𝒜`$, and any two such decompositions contain the same irreducibles.
The image of $`x`$ should be considered a positive sum of points in $`_𝒜`$ with multiplicities given by the number of times a particular irreducible representation of $``$ appears. Thus the map is multivalued, given the description used so far. If we wish to make such maps single-valued, we can either restrict the choice of maps, or modify the definition of a point. The latter possibility is most natural, and there is an obvious choice. We should consider instead a new space consisting of the free sums of points with coefficients in $`^+`$. These sums of points are generated by the points of $``$ such that the sums are finite. We should consider the maps between two such spaces as linear transformations between the two formal sums. Notice that the formal positive sums of $`n`$ points in a commutative variety $``$ corresponds to the symmetric product $`^n/S_n`$. Hence the natural object to understand in this version of non-commutative algebraic geometry is the symmetric product of the space $`_𝒜`$. This is also the framework necessary for matrix theory and matrix string theory, and this is why we find it a very appealing aspect of the construction. We will denote this symmetric product space by $`𝒮_𝒜`$. $`_𝒜`$ is a subset of its symmetric space, and it is the set which generates the formal sums.
There is a grading present here which gives a notion of the degree of a point, which we denote $`\mathrm{deg}(x)`$. There is a natural map from the formal sums of points to $``$; namely, for each irreducible representation $`x_𝒜`$, we consider the map that assigns to $`x`$ the dimension of the representation that $`x`$ is associated with (that is, the character $`\mathrm{tr}_{\mu _x}1`$). This extends by linearity to the symmetric product space, and the maps between the sums of points are such that they are degree-preserving. Indeed, given any function on the space (an element of the ring $`a𝒜`$), we consider the invariant of the function $`a`$ at the point $`x`$, $`\mathrm{tr}_{\mu _x}a`$. The trace is linear, and independent of the choice of basis for the local matrix ring and can therefore be extended to direct sums of representations (i.e., to the space $`𝒮_𝒜`$).
Each positive sum of points of $`_𝒜`$ is associated with a representation of the ring $`𝒜`$, which is the direct sum of irreducible representations. To each element in the ring, we associate a character in the representation. Consider the vector space associated to a representation on which the matrix ring acts as a left module of the ring of functions. Because the associated representations of points are only well defined up to conjugation, we should impose the same constraint on the modules; namely, we want isomorphism classes of modules for the ring $`𝒜`$. This is very reminiscent of algebraic K-theory, and we will expand on this idea later, making the connection precise.
Now, although we have found it natural to extend $`_𝒜`$ to $`𝒮_𝒜`$, it is not the case that $`𝒮_𝒜`$ automatically inherits a topology from that of $`_𝒜`$. Rather, we should repeat the construction of a Zariski topology, by giving a definition of closed sets. For any function $`a𝒜`$ and for every complex number $`z`$, we define the following set
$$𝒵=\{p𝒮_𝒜|\mathrm{tr}_{R_p}(a)=z\}$$
(14)
to be closed. These sets form a basis of closed sets for the topology of $`𝒮_𝒜`$. This topology coincides with the natural topology in the commutative algebra of functions generated by traces of operators, which are the polynomials in the gauge invariant superfields, and thus gives us the same topological information that we would want for the moduli space if we just consider the ring of holomorphic functions on the moduli space with values in $``$. For later use, we define the support of a character as
$$Supp(\mathrm{tr}a)=\overline{\left\{p𝒮_𝒜\right|\mathrm{tr}_{R_p}a0\}}$$
(15)
where the overline denotes the closure operation.
We would like to be able to say that the space is foliated by sets of degree $`m`$ (which will count the $`D`$-brane charge of the point). Because the Zariski topology is coarse, we must then additionally declare that the sets defined by $`\mathrm{deg}(x)=m`$ are both open and closed.
Now recall that for two different points in an algebraic variety $`V`$, there is some function on the variety which distinguishes them. Only a finite number of these functions is needed to determine a point exactly; the ring of polynomials (with relations) is finitely generated. This construction is also sufficient to determine a collection of $`n`$ unordered points of the variety. By examining $`n`$th order polynomials in a function $`f`$, we can determine the values that $`f`$ takes at the $`n`$ points. If the values are different for all the points for some function $`f`$, then one can use $`f`$ as a coordinate, and one has a collection of $`n`$ non-overlapping algebraic subsets of $`V`$, with one point chosen from each one. Thus we can construct a function which vanishes at all but one of the subsets, which we call $`f_1`$ and by multiplying $`f_1`$ by all the basis functions of the ring associated to $`V`$, we can identify one of the points. The procedure can be repeated if no one function is able to tell them all apart, and then we get the multiplicities of the points.
In the non-commutative case, we say that we can always distinguish two irreducible representations by some collection of characters (traces) of the ring $`𝒜`$. Thus there is a given finite number of functions with which we can distinguish $`n`$ points.
Recall that we wish to think of the non-commutative symmetric space as a refined version of the commutative space. Topologically, the two spaces are the same. It is clear that, at least locally, given the characters of enough elements of the ring, we can fully reconstruct a representation by holomorphic matrices on the commuting variables. This endows the symmetric space locally with a holomorphic vector bundle structure.
### 3.4 The Role of the Center
Now let us apply the above construction to maps between the spaces $`𝒮_𝒜`$ and $`𝒮_{𝒵𝒜}`$. Consider in particular the inclusion map $`𝒵𝒜𝒜`$, which is the pullback of functions on $`_{𝒵𝒜}`$ to $`_𝒜`$.
We want to know the image of a point $`p`$ in $`_𝒜`$. Given the point $`p`$, there is an associated $`M`$-dimensional irreducible representation $`\mu _p`$. Consider composing the maps $`𝒵𝒜\stackrel{i}{}𝒜\stackrel{\mu _p}{}M_M()`$. As the last map is onto, if $`a𝒵𝒜`$, it commutes in the image of the composition of maps and by Schur’s lemma is proportional to the identity. The representation associated to $`p`$ splits into $`M`$ identical copies of a single representation of $`𝒵𝒜`$, namely, into $`M`$ copies of a single point. We write this as
$$pM\overline{p}$$
(16)
where $`\overline{p}`$ is the associated maximal ideal of $`𝒵𝒜`$, which is the kernel of the inclusion map.
We will call this the natural map, as it respects degree. Notice that we can also define a map between the symmetric spaces $`p\overline{p}`$ which forgets the degree. In this case, the image of a point is a point, so we can restrict the maps to $`_𝒜`$ and $`_{𝒵𝒜}`$. We will call this the forgetful map.
The center of the algebra will play an important role in the physics. We wish to restrict the maps between rings $`𝒜`$ and $``$ in the following way. Note that we have the following diagrams
$$\begin{array}{ccc}𝒜& & \\ & & \\ 𝒵𝒜& & 𝒵\end{array}\begin{array}{ccc}𝒮_{}& & 𝒮_𝒜\\ & & \\ 𝒮_𝒵& & 𝒮_{𝒵𝒜}\end{array}$$
(17)
We require that these diagrams be commutative; namely, the ring homomorphism $`𝒜`$ induces a map $`𝒵𝒜𝒵`$, and consequently $`𝒮_𝒵𝒮_{𝒵𝒜}`$. Thus, we want the map to be such that central elements are central not just in the subalgebra of the image of $`𝒜`$ but in the algebra $``$ itself.
Now we are at a point where we can build a category for the non-commutative algebraic geometry. The objects will be rings $`(𝒜)`$ with the center identified and the inclusion map singled out.<sup>2</sup><sup>2</sup>2It is appropriate then to use the larger notation $`(𝒜)(𝒜,𝒵𝒜,i:𝒵𝒜𝒜)`$. The allowed maps between rings are such that they produce commuting squares
$$\begin{array}{ccc}𝒜& & \\ & & \\ 𝒵𝒜& & 𝒵\end{array}$$
(18)
with the upwards arrows the natural inclusion maps.
The non-commutative space is a contravariant functor from this category of rings to a category of ‘symmetric spaces’ as we have defined previously (including the degree map and the degree-preserving property). Thus we have the diagram
$$\begin{array}{ccc}𝒮_{}& & 𝒮_𝒜\\ & & \\ 𝒮_𝒵& & 𝒮_{𝒵𝒜}\end{array}$$
(19)
Together with the center preserving property, the forgetful map induces
$$_𝒵_{𝒵𝒜}.$$
(20)
This is a map of commutative algebraic varieties, to which we can apply intuition. It is also clear that this map between varieties has all the data required to specify the map between their symmetric spaces. With the topology of these spaces, all the arrows are continuous maps, and composition of maps is a map that respects the properties of the category.
To make a full connection with algebraic geometry, we want to be able to glue rings on open sets. This should be done by a process of localization. These details will be left for a future publication.
It is useful to notice that all the irreducible representations of the center may not appear when we consider the projection map, $`𝒮_𝒜𝒮_{𝒵𝒜}`$. If most<sup>3</sup><sup>3</sup>3That is, an open set of $`_{𝒵𝒜}`$ in the Zariski topology. do appear, then we will call the algebra semi-classical, because to the points in the variety associated to the center, we can lift to points in the non-commutative variety. The non-commutative variety covers the commutative one and this notion will be important from several perspectives below. In particular, there are applications involving $`D`$-branes in which phenomena on orbifold spaces are more precisely described by non-commutative geometry.
### 3.5 $`D`$-brane Fractionation
The technology developed so far contains some interesting aspects of $`D`$-brane physics. In particular, we wish to show that a $`D`$-brane fractionates as we move to a singular point of a non-commutative algebraic variety. In fact we define singular points via this process of fractionation. We will consider in this subsection $`D`$-branes which correspond to points in $`_𝒜`$, and the degree of the point is identified with the $`D`$-brane charge. The moduli space of supersymmetric configurations of $`D`$-branes is identified with $`𝒮_𝒜`$.
Let $`R`$ be an irreducible representation of dimension $`M`$ in $`𝒜`$. Consider its image in $`𝒮_𝒜`$ as a single point of degree $`M`$. Because $`𝒮_𝒜`$ is an algebraic variety, it will consist of several components or branches. The branch of $`𝒮_𝒜`$ where $`R`$ is located is a closed set of some complex dimension $`d`$ which is not a closed subset of any set with larger local dimension.
On this branch, we can define a local function which is the dimension of the commutant $`𝒵R`$ of the representation $`R`$. As we move along the branch, this function is semi-continuous—it may jump in value on closed sets.
Clearly, the sets with $`dim(𝒵R)>1`$ are closed. In this case, we have at least two linearly independent matrices which commute with everything in the image of $`𝒜`$, and thus the representation cannot be irreducible. For irreducible representations, $`dim(𝒵R)`$ must be unity. Thus, if we start at a point on a branch of the variety that is irreducible, as we continuously deform along it, we can reach a special point as a limit point, where the representation becomes reducible.
Parametrize this deformation by $`z`$; on the symmetric product space we have the process
$$\underset{zz_0}{lim}x(z)=x_1+\mathrm{}+x_n$$
(21)
if $`z_0`$ is such a limit point, and where $`n`$ is the number of irreducible representations that $`R(z)`$ splits into. Then, we say that the $`D`$-brane has fractionated, and there may be additional branches that intersect that point, corresponding to separating the fractional branes.
From the point of view of the center of the algebra, each element is proportional to the identity throughout the branch of the symmetric space, and thus there is no splitting seen in $`𝒮_{𝒵𝒜}`$. In this sense, the non-commutative geometry is a finer description of the $`D`$-brane moduli space than the associated commutative geometry.
In the cases where there are branches corresponding to separating the branes at $`z_0`$, if we think in terms of the forgetful map, we would have a single point splitting into $`n`$ points. From the point of view of the commutative algebraic variety, there is a jump in dimension as we go from one branch to the next; this is naturally associated with a singularity. We will see explicit examples of this in Section 4.
### 3.6 Higher dimensional branes
So far, we have considered $`D`$-branes that are point-like on the moduli space. We would also like to identify more general brane configurations, such as those wrapped on holomorphic subspaces; hence we need to construct such objects algebraically. It is natural to consider coherent sheaves: these are the modules over the ring $`𝒜`$ which locally have a finite presentation and are well-behaved when considered from the commutative standpoint. Extended BPS brane solutions usually correspond to stable sheaves, given some appropriate notion of stability. Moreover, they are also well-behaved as far as K-theory is concerned. However, in order to define these structures, it is most convenient to have a semi-classical ring. This does not mean that there is no useful way to define these objects for more general rings, but on some rings where the points are discrete there is no obvious notion of an extended object. For the rest of this section, we will assume that we are indeed working in a semi-classical ring.
As the ring is semi-classical, we can try to construct extended objects by first building them over a holomorphic subspace of the commutative structure, and then try to lift them up to the non-commutative geometry. In the commutative case, a $`D`$-brane corresponds to a coherent sheaf with support on a commutative subvariety. For any notion of non-commutative sheaf, it must be the case that it is also a coherent sheaf over the commutative ring. In the commutative case, the $`D`$-brane is a module over the ring $`𝒵𝒜`$, such that if $`𝒵`$ is the ideal corresponding to the support of the sheaf, the module action of $`𝒵𝒜`$ factors through $`𝒵𝒜/𝒵`$, which is considered to be the coordinate ring of the closed set associated to the ideal $`𝒵`$.
On ‘good’ varieties we always have a presentation of a sheaf $`𝒮`$ as the right-hand term of some exact sequence
$$𝒵𝒜^m𝒵𝒜^n𝒮0$$
(22)
that is, as a module with $`n`$ generators with relations induced by the images of $`𝒵𝒜^m`$.
We want to mimic this construction for the non-commutative version of the $`D`$-brane. Note that in the non-commutative case, we have a choice of left-, right- or bi-modules of the algebra $`𝒜`$. However, physically, we need to consider only bi-modules, as both ends of open strings end on a $`D`$-brane. That is, gauge transformations (which act locally) act both on the left and right, and therefore the algebra has to be able to accommodate both types of actions on the modules. Referring to the bi-module as $``$, we want them to arise from exact sequences
$$𝒜^m𝒜^n0$$
(23)
in analogy to eq. (22). This defines locally<sup>4</sup><sup>4</sup>4Recall that we are not concerned with gluing. the coherent sheaves over $`𝒜`$.
The annihilator of a bi-module $``$ is defined as the largest ideal $``$ of $`𝒜`$ such that $`==0`$. This is a double-sided ideal, and thus defines a closed set, in the topology of $`_𝒜`$. For any point $`p`$ in $`_𝒜`$ which does not belong to $``$, $`+_p`$ is equal to $`𝒜`$ (because $`_p`$ is maximal). We will refer to this ideal $``$ as $`\text{Ann}()`$.
We can find how a bi-module restricts to a closed subset (described by an ideal $``$) by noticing that if $``$ is a bi-module over $`𝒜`$, then $`/(+)`$ is a bi-module over $`𝒜/`$. For maximal $`_p`$ the restriction is zero if $`\text{Ann}()/_p`$. We can define the support of a sheaf to be the set of points such that
$$\text{Ann}()_p.$$
(24)
We study the sheaves locally by restricting to a point. We will look at two different notions of the rank of a sheaf. Each non-commutative point is Morita equivalent to a commutative point. This tells us that modules over the functions restricted to a point $`p`$ (the ring of $`n\times n`$ matrices) behave just as vector spaces over $``$. Thus the sheaf restricted to a point is a bi-module over the ring of $`\mathrm{deg}(p)\times \mathrm{deg}(p)`$ matrices, and thus $`|_p`$ is isomorphic to $`\left(𝒜|_p\right)^k`$, for some $`k`$. One possible definition of non-commutative rank of a sheaf at the point $`p`$ is just $`k`$. However, as seen from the commutative standpoint, the dimension of the representation associated to $`\left(𝒜|_p\right)^k`$ is $`k\mathrm{deg}p`$, and this then serves as another definition of rank, which we will refer to as the commutative rank.
As usual, rank is upper semi-continuous (the points where $`rank()|_p>m`$ form a closed set). The rank can jump in value on some closed subset, and this is interpreted in terms of an additional $`D`$-brane of smaller dimension stuck to the brane, as follows from the anomalous couplings of $`D`$-branes.
For the non-commutative points, we have to take into account that a limit set of a collection of points might be a sum of points. Consider the trivial bi-module of $`𝒜`$, namely $`𝒜`$. The non-commutative rank (equal to 1) does not jump under the fractionation process. The commutative rank on the other hand, does jump at this singularity. The non-commutative rank then is the natural definition of rank for non-commutative algebras.
Note, however, that if we look just at the center $`𝒵𝒜`$, the commutative rank is the natural definition, as it does not jump in a splitting process; we have
$$\mathrm{deg}(p)=\underset{i}{}\mathrm{deg}p_i.$$
(25)
With these definitions, a $`D`$-brane is a coherent sheaf over both the non-commutative ring and the commutative sub-ring.
### 3.7 K-theory interpretation
We have seen that our approach to non-commutative geometry has led us to some definitions of $`D`$-brane states. We now want to add $`K`$-theory to the discussion. Because we have a ring and we have bi-modules, we get automatically a K-theory associated to this structure, namely, the algebraic K-theory of the ring $`𝒜`$ (see for example). From a mathematical standpoint this is review material, and we will just glimpse at the dynamics in terms of brane-antibrane systems.
Indeed, let us start with the construction of the symmetric space. We had formal sums of points which makes the non-commutative geometry a semigroup. We can make it into a group by adding minus signs, and a rule for cancellation. This group is the equivalent of zero-chains of points.
Now, the idea of the group structure is to understand that $`p+(p)=0`$. So if $`p`$ is a point, we interpret it as a point like $`D`$-brane, and $`p`$ is an anti-$`D`$-brane. The cancellation law of the group is the statement that a $`D`$-brane anti-$`D`$-brane pair can be created from the vacuum.
Given these minus signs, the degree function now maps to the integers and gives us an invariant, which is a group homomorphism of Abelian groups, $`\mathrm{deg}:p`$. This number is the total $`D`$-brane charge of a configuration. It is possible to give a topology to finite formal sums of points by the same construction we used based on characters. The extra ingredient to make $`pp=0`$ is to add minus signs for the anti-$`D`$-brane.
Thus dynamically $`pp=0`$ is the statement that any character of $`pp`$ is the same as a character of zero, and thus the configurations are connected. When we create a $`D`$-brane anti-D brane pair we can separate them if there are moduli available, and thus the process is continuous in this topology. Dynamical information would include the energy required for this process. (A generic point is such that this energy is much less than the energy required to move off of the moduli space of sums of points. A non-generic point is where the mass matrix has zero eigenvalues.)
The idea now is to define the K-theory of points as a homotopy invariant which respects the additivity of branes. On one hand, we have the mathematical definition of the $`K_0^p`$-theory of points as the formal abelian group of homotopy classes of finite dimensional representations of the algebra $`𝒜`$, such that if $`a,b`$ are such representations, then the K-theory classes associated to $`ab,a,b`$ satisfy
$$K(ab)=K(a)+K(b)$$
(26)
and if one has a homotopy between the two representations $`ab`$, then $`K(a)=K(b)`$. Indeed, to the point $`p`$ we associate the bi-module $`𝒜|_p`$, and this is thought of as the skyscraper sheaf over $`p`$. With this extra relation this is part of the K-theory of bi-modules of the algebra.
The other way to define K-theory is to say $`K(a)=K(b)`$ if there is $`c`$ such that $`acbc`$. Both of these definitions agree.
There are a few possible choices of K-theory depending on the type of modules one chooses. As we have stated above, in this paper we are interested in finitely presented bi-modules over the ring $`𝒜`$. Generically, one defines the K-theory associated to projective bi-modules of the algebra. The K-theory of projective bi-modules is the same as the K-theory of finitely presented bi-modules as long as every bi-module admits a projective resolution (this is true, for example, in smooth manifolds where every vector bundle is projective over the coordinate ring of the manifold). Thus, as long as there is a long exact sequence
$$0P_1\mathrm{}P_kM0$$
(27)
with the $`P_i`$ projective, then the K-theory class of $`M`$ is defined. This requires the ring to be regular. Whether or not we will always get regular rings in string theory in this framework is not clear. (Singular varieties are not regular as commutative algebras, yet they do appear in string theory.)
The $`K_0`$-theory is defined as the set of formal sums of bi-modules modulo homotopy, and modulo the relations
$$K(b)=K(a)+K(c)$$
(28)
whenever there is a short exact sequence
$$0abc0$$
(29)
of the bi-modules we have described. We think of this as the statement that $`bac=0`$.
If a module $`M`$ admits a projective resolution as in (27) then it is a simple exercise to show that
$$K(P_1)K(P_2)+\mathrm{}(1)^kK(P_k)+(1)^kK(M)=0$$
(30)
Physically, we say the dynamics of brane-antibrane configurations is such that given a short exact sequence (29), the process
$$XX\pm ab\pm c$$
(31)
is allowed, namely $`ab+c`$ carries no $`D`$-brane charge. In particular, the exact sequence
$$0aa00$$
(32)
will allow any of the two processes
$$XX+a+(a)X$$
(33)
which correspond to the creation of brane-antibrane pairs. Of course, the real dynamics of these processes is not available to us, but the topology of allowed transitions is correctly reproduced.
Note that taking tensor products of bi-modules is locally a good operation (at each point we are taking a tensor product of finite dimensional spaces, and we get a finite dimensional space). Thus the K-theory is not just an additive group but we have a multiplication as well, and this permits us to do intersection theory (that is, we can count strings when branes intersect). This type of information can often be enough to calculate topological quantities in string theory.
As a final comment, we have to give some warnings to the reader. The K-theory we have constructed here is the one associated to the holomorphic algebra, and thus is a version which is relevant for the algebraic geometry. This K-theory is an invariant of a holomorphic space which is much finer than the topological K-theory, and thus contains a lot more information. The K-theory which is relevant for $`D`$-brane charge is the one associated to both the holomorphic and anti-holomorphic structures, namely, the K-theory of a $`^{}`$ algebra of which $`𝒜`$ is a subalgebra. This $`^{}`$ algebra includes the D-term constraints, and by theorems on existence of solutions the geometric space associated to the $`^{}`$ algebra has just as many non-commutative points as the one associated to $`𝒜`$. So just as in commutative cases, the non-commutative holomorphic data parametrize the full variety. The K-theory of the two algebras does differ. The holomorphic K-theory is therefore more appropriate to count BPS states, rather than just account for the $`D`$-brane charge.
This is the end of the mathematical preliminaries. We believe that we have presented a fairly general account of how applications of these techniques might be pursued. We will see that this approach is not just a big machine which describes things we already knew in a complicated manner. Indeed, once we have examined the examples in the next sections, it will be clear that the formulation brings sound intuition and gives a very nice picture of how string geometry behaves.
## 4 Examples
In this section, we consider a variety of examples in order to build a picture of the generic behavior of the geometry, which is not present in the simplest case. The presentation is given in terms of the language of Section 3; the reader will find it necessary to read, at least, the overview in Section 3.1. In the first few examples, we first calculate the commutative algebra of the center which reproduces the string geometry associated with the field theory (e.g., the orbifold). We then attempt to build the irreducible representations of the full non-commutative algebra by exploiting knowledge of the center. A posteriori, the structures that we find here and the relations to the physics of $`D`$-branes in these geometries discussed in later sections, motivates the formal constructions of Section 3. In more general examples, the calculation of the center is difficult, and we present only partial results.
### 4.1 Orbifolds with discrete torsion: the $`q`$-deformation
Our first example to study will be orbifolds with discrete torsion. In particular, we consider the orbifold $`^3/_n\times _n`$ with maximal discrete torsion. To construct the low energy effective field theory of a point-like brane one can use a quiver construction with projective representations of the orbifold group. The use of projective representations was justified in Ref. . The algebraic variety associated to the orbifold singularity is given by the solutions of one complex equation in four variables,
$$xyz=w^n.$$
(34)
As Douglas showed, the theory has $`N=1`$ supersymmetry in four dimensions and consists of a quiver with one node, gauge group $`U(M)`$, three adjoint superfields and a superpotential
$$W_q=\mathrm{tr}(\varphi _1\varphi _2\varphi _3)q\mathrm{tr}(\varphi _2\varphi _1\varphi _3)$$
(35)
with $`q`$ a primitive $`n`$-th root of unity. This theory can be obtained by a marginal deformation of the $`N=4`$ supersymmetric field theory as shown in , and as such, when studied under the AdS/CFT correspondence, displays a duality between two totally different near-horizon geometries, describing the same field theory.
The $`F`$-term constraints are given by
$`\varphi _1\varphi _2q\varphi _2\varphi _1`$ $`=`$ $`0`$ (36)
$`\varphi _2\varphi _3q\varphi _3\varphi _2`$ $`=`$ $`0`$ (37)
$`\varphi _3\varphi _1q\varphi _1\varphi _3`$ $`=`$ $`0`$ (38)
We will often write these using the $`q`$-commutator notation $`[\varphi _1,\varphi _2]_q=0`$, etc. These equations are exactly the type of relations seen in the algebras related to quantum planes, and have been very well studied. Let us analyze the algebra using the tools described in Section 3.
Because of the $`F`$-term constraints, we can always write any monomial in ‘standard order’
$$\varphi _1^{k_1}\varphi _2^{k_2}\varphi _3^{k_3}.$$
(39)
We associate to this monomial the vector $`(k_1,k_2,k_3)`$.
Note that if an element commutes with $`\varphi _1,\varphi _2,\varphi _3`$, then it commutes with any of the monomials, and thus is an element of the center of the algebra. Monomials may be multiplied, and up to phases, we have
$$(k_1,k_2,k_3).(s_1,s_2,s_3)(k_1+s_1,k_2+s_2,k_3+s_3).$$
(40)
Because of the phases, generators of the center are monomials.
It is easy to see that $`(k_1,k_2,k_3).\varphi _1=\varphi _1.(k_1,k_2,k_3)q^{k_3k_2}`$, so that $`k_3=k_2modn`$ for $`(k_1,k_2,k_3)`$ to be in the center. Similarly one proves $`k_2=k_1modn`$ and thus the center is given by the condition
$$𝒵𝒜=\left\{(k_1,k_2,k_3)\right|k_1=k_2=k_3modn\}$$
(41)
This is a sub-lattice of the lattice of monomials, and it is generated by the vectors $`(1,1,1),(n,0,0),(0,n,0),(0,0,n)`$. Call $`w=(1,1,1)`$, $`x=(n,0,0)`$, $`y=(0,n,0)`$ and $`z=(0,0,n)`$. Clearly we have the relation
$$(w)^n+xyz=0$$
(42)
so we see the orbifold space is described by the center of the algebra. The singularities occur along branches where two of $`x,y,z`$ are zero.
Now that we have the commutative points, let us consider the non-commutative points of the geometry. We should consider the irreducible finite dimensional representations of the algebra. Because $`x,y,z`$ are central, on an irreducible representation of the algebra they act by multiples of the identity.
Suppose at least two of $`x,y,z`$ are non-zero (say $`x,y`$). In this case $`(1,0,0)`$ and $`(0,1,0)`$ are invertible matrices. By a linear transformation, we can diagonalize $`(1,0,0)`$. Consider an eigenvector $`|a`$ of $`(1,0,0)`$ with eigenvalue $`a`$. We see that $`|qa(0,1,0)|a`$ is an eigenvector of $`(1,0,0)`$ with eigenvalue $`qa`$. Thus we get a collection of states $`|a,|qa,\mathrm{},|q^{n1}a`$ constructed as $`|a,(0,1,0)|a,\mathrm{},(0,n1,0)|a`$. This sequence terminates (and thus the representation is of dimension $`n`$) because $`(0,n,0)`$ is central, and $`q^n=1`$. A set of matrices which satisfies these conditions is
$`(1,0,0)=aP`$ (43)
$`(0,1,0)=bQ`$ (44)
with $`P`$ and $`Q`$ defined by<sup>5</sup><sup>5</sup>5We have changed basis compared to Ref. .
$$P=\left(\begin{array}{ccccc}1& 0& 0& \mathrm{}& 0\\ 0& q& 0& \mathrm{}& 0\\ 0& 0& q^2& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& \mathrm{}& q^{n1}\end{array}\right),Q=\left(\begin{array}{ccccc}0& 0& \mathrm{}& 0& 1\\ 1& 0& \mathrm{}& 0& 0\\ 0& 1& \mathrm{}& 0& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& \mathrm{}& 1& 0\end{array}\right)$$
(45)
As $`w=(1,1,1)`$ is central, it is proportional to the identity. It trivially follows that $`(0,0,1)=cQ^1P^1`$.
Notice that our solutions are parameterized by three complex numbers, namely $`a,b,c`$. It is easy to see that $`x=a^nI`$, $`y=b^nI`$ and $`z=(c)^nI`$, $`w=abcI`$, and that one can cover the full orbifold with these solutions, except for the singularities (where two out of the three $`x,y,z`$ are zero). Notice also that the covering is done smoothly, so any two points can be connected by a path which does not touch the singularities.
If we label the representation by $`R(a,b,c)`$, it is easy to see that $`R(a,b,c)`$ is equivalent under a similarity transformation to $`R(qa,q^1b,c)`$ and $`R(qa,b,q^1c)`$. Thus the eigenvalues of the center completely describe the representation. That is, for any commutative point which is non-singular, we have a unique non-commutative point of degree $`n`$ sitting over it.
Let us now analyze the case where two of the three $`x,y,z`$ are zero. Then $`w=0`$ as well, and we are along one of the singular branches of the orbifold. Let us assume that $`x0`$; then $`(1,0,0)`$ is invertible, and can be diagonalized. On the other hand $`(0,1,0).(0,0,1)=(0,0,1).(0,1,0)=(0,n,0)=(0,0,n)=0`$ in the representation. Given any vector $`v`$ in the representation, $`v^{}(0,n1,0)v`$ is annihilated by $`(0,0,1)`$ and $`(0,1,0)`$, and any other vector obtained by multiplying with $`(1,0,0)`$ enjoys this same property. Thus given a representation, we find a sub-representation where both $`(0,0,1)`$ and $`(0,1,0)`$ act by zero. As $`(1,0,0)`$ is invertible it can be diagonalized in this subrepresentation. Clearly the representation is irreducible only if it is one dimensional, and determined by the eigenvalue of $`(1,0,0)`$, which is a free parameter that we call $`a`$. The value of $`x`$ is $`a^n`$, and for each point in the singular complex line $`y=z=0`$ we find $`n`$ irreducible representations of the algebra, except at the origin. The same result holds when we go to any of the other complex lines of singularities. We label these representations by $`R(a,0,0)`$, etc.
Here $`R(a,0,0)`$ is not equivalent to $`R(qa,0,0)`$. They are equivalent as far as the commutative points are concerned, because both of these representations have the same characters over the center of the algebra. But as far as the non-commutative points are concerned, the characters of the non-central element $`(1,0,0)`$ differ. That is
$$\mathrm{tr}_{R(a,0,0)}(1,0,0)=a$$
(46)
Thus we have two distinct points. It is also clear that any one of these representations can be continuously connected to any other.
These smaller representations are not regular for the $`^3/_n\times _n`$ orbifold and may be identified with the fractional branes. Notice also that
$$\mathrm{tr}_{R(a,b,c)}(1,0,0)=a\mathrm{tr}P=0$$
(47)
so that this character is different from zero only at the classical singularity. This is the primary reason for adopting the convention for the support of a character in eq. (15).
To summarize, for each point in the classical moduli space we have at least one point in the non-commutative space which sits over it. This is an example of a semi-classical geometry (see Section 3). The commutative singular lines are covered by an $`n`$-fold non-commutative complex plane branched at the origin.
Now consider what happens when we bring a point from the regular part of the orbifold towards the singularity. The representation behaves in this limit as
$$\underset{b,c0}{lim}R(a,b,c)=R(a,0,0)R(qa,0,0)\mathrm{}R(q^{n1}a,0,0)$$
(48)
In our description of moduli space, this corresponds to the branes becoming fractional at the orbifold fixed lines, as we have discussed previously. Indeed, once we reach this point we can separate the fractional branes, and the non-commutative symmetric product is the right tool for describing the moduli space in full.
We can also see the quiver of the singularity type by consideration of this same limit. Indeed, we assign a node to each irreducible representation in the right-hand side of eq. (48). We draw an arrow between any two nodes appropriate to the non-zero entries in eqs. (45) and we obtain Figure 2 which is indeed the quiver diagram of the orbifold in the neighborhood of a point in the singular complex line.
Thus the singularities can be said to be locally quiver. From the point of view of the center of the algebra, the nodes of the quiver are at the same point, but they are distinct in the non-commutative algebra. The behavior of the field theory near the singularities is precisely what we would get from the orbifold analysis.
Recall that the commutative singular lines are covered by $`n`$ non-commutative branches. The monodromies of the quiver diagram are encoded in this structure, and thus their calculation is geometrically obvious. This compares quite favorably to the rather cumbersome procedure employed in . Indeed, we can change $`a\omega a`$ for $`\omega =e^{2\pi i/n}`$. This results in a permutation of the factors appearing on the right-hand side of eq. (48). This permutation is the monodromy.
### 4.2 Adding one mass term
Next, we consider a relevant deformation of the last theory, obtained by the addition of a single mass term. This theory is a $`q`$-deformed version of the theory which flows in the infrared to an $`N=1`$ conformal field theory. The superpotential is
$$W=\mathrm{tr}\left(\varphi _1\varphi _2\varphi _3q\varphi _2\varphi _1\varphi _3+\frac{m}{2}\varphi _3^2\right)$$
(49)
Again, we assume that $`q`$ is an $`n`$-th root of unity. The $`F`$-term constraints are given by
$`[\varphi _1,\varphi _2]_q`$ $`=`$ $`m\varphi _3`$ (50)
$`[\varphi _2,\varphi _3]_q`$ $`=`$ $`0`$ (51)
$`[\varphi _3,\varphi _1]_q`$ $`=`$ $`0`$ (52)
As in the previous case, we look for the center of the algebra to obtain the commutative manifold. It is easy to see that $`z=\varphi _3^n`$ is still in the center. We can also show that
$`[\varphi _1^n,\varphi _2]`$ $`=`$ $`\varphi _1^n\varphi _2q\varphi _1^{n1}\varphi _2\varphi _1+q\varphi _1^{n1}\varphi _2\varphi _1q^2\varphi _1^{n2}\varphi _2\varphi _1^2+\mathrm{}`$
$`=`$ $`\varphi _1^{n1}(m\varphi _3)+q\varphi _1^{n2}(m\varphi _3)\varphi _1+\mathrm{}`$
$`=`$ $`m\varphi _1^{n1}\varphi _3{\displaystyle \underset{r=0}{\overset{n1}{}}}q^{2r}`$
which vanishes, apart from at the special values $`q=\pm 1`$. Similarly one proves that $`\varphi _2^n`$ is central, away from $`q=\pm 1`$. For now, we will assume that $`q^21`$, and return to these cases later. Thus we have at least three central variables $`x=\varphi _1^n`$, $`y=\varphi _2^n`$, $`z=\varphi _3^n`$.
The variable $`w`$ is modified by the presence of the mass term. Consider the commutator
$`[\varphi _1\varphi _2\varphi _3,\varphi _1]`$ $`=`$ $`\varphi _1\varphi _2\varphi _3\varphi _1q\varphi _1\varphi _2\varphi _1\varphi _3+q\varphi _1\varphi _2\varphi _1\varphi _3\varphi _1\varphi _1\varphi _2\varphi _3`$
$`=`$ $`m\varphi _1\varphi _3^2`$
This result may be rewritten as a commutator for $`q\pm 1`$, and thus we see that
$$w=\varphi _1\varphi _2\varphi _3+\frac{m}{1q^2}\varphi _3^2$$
(54)
is central.
The four variables $`x,y,z,w`$ are related by
$$xyz=(w)^n\left(\frac{m}{1q^2}\right)^nz^2$$
(55)
This is a deformation of the complex structure of (42). It is easy to see that we now have singularities at $`w=xz=yz=xy+2tz=0`$ with $`t=(m/(1q^2))^n`$. Thus, the singularities are at $`xy=0,w=0,z=0`$, and so we have two lines of singularities $`x=0`$ and $`y=0`$. The mass term has resolved one of the three complex lines of singularities (for $`q^21`$).
It is easy to check that a general solution is of the form
$`\varphi _1`$ $`=`$ $`aP`$ (56)
$`\varphi _2`$ $`=`$ $`bP^1QcP^1Q^1`$ (57)
$`\varphi _3`$ $`=`$ $`dQ^1`$ (58)
with $`P,Q`$ defined as in (45), and where $`a,b,c,d`$ are numbers satisfying
$$ac(1q^2)=md$$
(59)
One then gets $`x=a^nI`$, $`y=(b^n+c^n)I`$, $`z=d^nI`$ and $`w=abdI`$. Note that this representation has been chosen such that $`\varphi _1`$ is diagonal at the singularity $`y=0,z=0`$.
Because we have a three complex parameter solution of the equations, we at least cover an open patch of the commutative variety, and we are again in a semi-classical ring. Indeed, we cover everything by finite matrices except $`x=0`$, as then $`c`$ is infinite.
A patch which does cover $`x=0`$ is given by
$`\varphi _1`$ $`=`$ $`aPQ^1cPQ`$ (60)
$`\varphi _2`$ $`=`$ $`bP^1`$ (61)
$`\varphi _3`$ $`=`$ $`dQ^1`$ (62)
with $`ab(1q^2)=mqd`$. This will be a good description for $`y0`$. The two patches cover the two lines of singularities. There is still the closed set $`x=y=0`$ which is not covered by either patch. We can find solutions for this set by taking $`\varphi _3`$ diagonal and making an ansatz for $`\varphi _1`$ which is upper triangular with entries just off-diagonal and $`\varphi _2`$ a similar lower triangular matrix. The dimension of this representation is also $`n`$ and depends on one complex parameter, namely the eigenvalues of $`\varphi _3`$.
On approaching the singularity $`y=z=0`$ from the bulk, we again get a split set of irreducible representations as follows:
$$\underset{b,d0}{lim}R(a,b,d)=R(a,0,0)R(qa,0,0)\mathrm{}$$
(63)
#### 4.2.1 Comments on the Infrared CFT
This case is also very interesting from the field theory perspective because by adding one mass term to a theory with three adjoints, we obtain a nontrivial conformal field theory in the infrared.
On the moduli space, we have the following $`U(1)`$ symmetries
$`a,b,c`$ $``$ $`\lambda \gamma a,\lambda \gamma ^1b,\lambda \gamma ^1c`$ (64)
$`d`$ $``$ $`\lambda ^2d`$ (65)
where we refer to the parameterization for $`x0`$. The transformation given by $`\gamma `$ is an ordinary $`U(1)`$, while $`\lambda `$ is the $`U(1)_R`$ symmetry, that of the superpotential in the infrared. The $`R`$ charges can be chosen in such a way that the superpotential is invariant at the infrared fixed point. Indeed, one can integrate out $`\varphi _3`$ and one finds a theory in the infrared with a quartic superpotential
$$\frac{1}{2m}\mathrm{tr}(\varphi _1\varphi _2q\varphi _2\varphi _1)^2$$
(66)
This superpotential is a marginal deformation of the infrared theory. Note that we also have, in the infrared, a $`_2`$ symmetry $`\varphi _1\varphi _2`$ which ensures that the anomalous dimensions of $`\varphi _1`$ and $`\varphi _2`$ are equal. (In the ultraviolet, this $`_2`$ symmetry is absent, as we would also have to simultaneously exchange $`qq^1`$ and rescale $`m`$.) This symmetry exchanges the two singular complex lines of the commutative moduli space.
#### 4.2.2 Special Cases: $`q=\pm 1`$
Let us return to discuss the moduli space for the cases $`q=\pm 1`$ from the algebraic point of view.
For $`q=1`$, which is the $`N=4`$ theory with one mass term, the moduli space is the set of solutions to
$$[\varphi _1,\varphi _2]=m\varphi _3$$
(67)
with all other commutators vanishing. For an irreducible representation, $`\varphi _3`$ is central and thus a constant. Because the commutator of $`\varphi _1,\varphi _2`$ is a constant, we get the Heisenberg algebra, and the only finite dimensional representations are those with $`\varphi _3=0`$. Thus the moduli space is a commutative space consisting of the symmetric product of the complex plane, $`^2`$. Notice that this space is of complex dimension two and not complex dimension three as in the generic case studied above. Indeed, in this case the center of the algebra is generated by $`\varphi _3`$. Because $`\varphi _3=0`$ on the moduli space, we can actually relax the condition for an element being central: we can take $`\varphi _1,\varphi _2`$ as central elements, which makes the moduli space commutative.
Indeed, this is a case where the algebra is not semi-classical. The variety associated to the center is the algebra of $``$. The non-commutative space is $`^2`$, which projects to the origin of $``$. The two have almost nothing in common.
As far as the commutative variety is concerned, the moduli space is a point. Notice that in this case when we integrate out the field $`\varphi _3`$ we get the correct dimension of the moduli space by counting fields. This does not happen for generic $`q`$.
For $`q=1`$, we can find the two dimensional solution
$$\varphi _1=a\sigma _1,\varphi _2=b\sigma _2,\varphi _3=0$$
(68)
plus two one-dimensional branches where either $`\varphi _1`$ or $`\varphi _2`$ is zero.
Here, the center is generated by $`\varphi _3^2`$. Indeed, it can be shown that these solutions exhaust the list of irreducible representations of the $`q=1`$ algebra. This result follows from the fact that $`[zx,y]z^2`$, so a finite dimensional representation must have $`z^2=0`$.<sup>6</sup><sup>6</sup>6Nilpotent possibilities, such as $`\varphi _3\left(\begin{array}{cc}0& a\\ 0& 0\end{array}\right)`$ are ruled out by $`D`$-terms.
The lesson to be learned from these special examples is that the commutative and non-commutative spaces may contain little or no information about each other when the center of the associated algebra is small. In this case, the full algebra is an infinite dimensional vector space over the center, and by considering only finite dimensional representations, we miss a lot of information.
### 4.3 One mass term and a linear term
We can easily modify the previously studied cases by adding a linear term to the superpotential
$$W=\mathrm{tr}\left(\varphi _1\varphi _2\varphi _3q\varphi _2\varphi _1\varphi _3+\frac{m}{2}\varphi _3^2+\zeta _3\varphi _3\right)$$
(69)
Note that by a field redefinition of $`\varphi _3`$, this is equivalent to adding a mass term $`\frac{\zeta _3}{m}(q1)\mathrm{tr}\varphi _1\varphi _2`$. We will see that the usual intuition for mass terms fails in this case, namely, that the moduli space is not destroyed by the quadratic terms. On the other hand, if we had added $`\mathrm{tr}\varphi _1\varphi _2`$ for $`q=1`$, we would indeed expect the space of vacua to be reduced to a set of points.
It is straightforward to show that $`x=\varphi _1^n,y=\varphi _2^n`$, and $`z=\varphi _3^n`$ are central, and that
$$w=\varphi _1\varphi _2\varphi _3+\frac{m}{1q^2}\varphi _3^2+\frac{\zeta _3}{1q}\varphi _3$$
(70)
is also central, provided that $`q\pm 1`$.
The relation between the central elements is
$$xyz=(w)^n\left(\frac{m}{1q^2}\right)^nz^2+\left(\frac{\zeta _3}{q1}\right)^nz$$
(71)
and a generic solution of the equations is provided by
$`\varphi _1`$ $`=`$ $`aP`$ (72)
$`\varphi _2`$ $`=`$ $`bP^1Q+cP^1dP^1Q^1`$ (73)
$`\varphi _3`$ $`=`$ $`eQ^1`$ (74)
with $`ac(1q)=\zeta _3`$, $`ad(1q^2)=me`$, and $`x=a^n`$, $`y=b^n+c^nd^n`$, $`z=e^n`$, $`w=abe`$. The singularities now occur at $`z=w=0`$ and
$$xy=\left(\frac{\zeta _3}{q1}\right)^n$$
(75)
The two complex lines of singularities that met at the origin when $`\zeta _3=0`$ are now replaced by a single $`^{}`$, a cylinder. In the parameterization above, this corresponds to $`b,d,e=0`$. We see that the non-commutative $`^{}`$ is an $`n`$-fold cover of the cylinder without branch points, and again the monodromies of the cover are manifest, since we chose $`\varphi _1`$ diagonal. This is again a semi-classical ring.
In addition, there are finite dimensional representations which may be thought of as deformations of $`SU(2)`$ representations. These occur for $`x=y=0`$ and cover regions not captured by the parameterization above. Some solutions give rise to isolated fractional branes at $`x=y=0`$. A similar effect occurs in Section 4.5 and we will return to a full discussion there.
The values $`q=\pm 1`$ are special, as in previous cases, in the sense that singularities occur, and the non-commutative algebra is not semi-classical.
### 4.4 Three linear terms
Consider the superpotential
$$W=\mathrm{tr}(\varphi _1\varphi _2\varphi _3)q\mathrm{tr}(\varphi _2\varphi _1\varphi _3)+\underset{i}{}(q1)\zeta _i\mathrm{tr}\varphi _i.$$
(76)
This case was studied in Ref. using gauge invariant variables. Our conclusions will be consistent with that analysis.
For convenience, we have rescaled the $`\zeta `$ parameters by a factor of $`(q1)`$. The $`F`$-terms give
$$[\varphi _1,\varphi _2]_q=(1q)\zeta _3,[\varphi _2,\varphi _3]_q=(1q)\zeta _1,[\varphi _3,\varphi _1]_q=(1q)\zeta _2$$
(77)
A possible parameterization is
$$\varphi _1=aP\frac{\zeta _3}{b}Q^1P,\varphi _2=bP^1Q+\frac{\zeta _1}{c}Q,\varphi _3=cQ^1+\frac{\zeta _2}{a}P^1$$
(78)
Note that $`x_1=\varphi _1^n`$, $`x_2=\varphi _2^n`$ and $`x_3=\varphi _3^n`$ are central, while the fourth central variable takes the form
$$w=\varphi _1\varphi _2\varphi _3\zeta _1\varphi _1q\zeta _2\varphi _2\zeta _3\varphi _3$$
(79)
In the given basis, we find $`x_1=a^n(\zeta _3/b)^n`$, $`x_2=b^n+(\zeta _1/c)^n`$, $`x_3=c^n+(\zeta _2/a)^n`$ and $`w=abc+q\frac{\zeta _1\zeta _2\zeta _3}{abc}`$.
These four variables are related on the moduli space by
$$x_1x_2x_3\underset{i}{}\zeta _i^nx_i+2\beta ^nT_n\left(\frac{w}{2\beta }\right)=0,$$
(80)
where $`\beta (q\zeta _1\zeta _2\zeta _3)^{1/2}`$ and $`T_n(x)=\mathrm{cos}(n\mathrm{cos}^1x)`$ is the $`n`$-th Chebyshev polynomial of the first kind.
### 4.5 Three mass terms
Next, we consider a rank 3 mass term of the form
$$W=\mathrm{tr}(\varphi _1\varphi _2\varphi _3)q\mathrm{tr}(\varphi _2\varphi _1\varphi _3)+\frac{1}{2}m\underset{i}{}\mathrm{tr}\varphi _i^2.$$
(81)
This superpotential has a $`_2\times _2`$ symmetry that changes two of the $`\varphi _i\varphi _i`$, and a $`_3`$ cyclic symmetry that permutes the $`\varphi _i`$. This is the remnant of the $`SU(4)_R`$ symmetry group of the $`N=4`$ SYM theory. The group generators do not commute with each other, and this symmetry is enhanced to $`SU(2)`$ when $`q=1`$. Thus the symmetry is a subgroup of $`SU(2)`$ which contains a $`_2\times _2`$ and a $`_3`$ subgroup. These are the symmetries of the tetrahedron, $`\widehat{E}_6`$, and since they arise from the $`SU(4)`$ R-symmetry they are chiral.
This superpotential yields the $`F`$-flatness conditions (cyclic on $`j`$, mod 3)
$$[\varphi _j,\varphi _{j+1}]_q=\varphi _{j+2},$$
(82)
where we have rescaled the fields in order to eliminate a factor of $`m`$.
We wish to find representations of this algebra; we will not immediately assume that $`q`$ is a root of unity. There is a certain class of solutions which may be thought of as deformations of representations of $`SL(2,)`$.
Note that (for $`q1`$) there is a one-dimensional representation
$$\varphi _j=\frac{1}{1q}$$
(83)
Higher dimensional representations may always be constructed as $`\varphi _i=\frac{1}{1q}I`$, but this is clearly reducible. An irreducible 2-dimensional representation (for $`q1`$) is given by
$$\varphi _j=\frac{i}{q+1}\sigma _j,$$
(84)
where the $`\sigma _j`$ are the Pauli matrices. We can construct higher dimensional irreducible representations by making the following ansatz: we suppose that one of the fields, $`\varphi _3`$, is diagonal and traceless, and that the other two fields only have non-zero elements just off the diagonals. (For $`q=1`$, these reduce to standard $`M`$-dimensional $`SL(2,)`$ generators). We have not been able to construct a proof that all such irreps may be obtained this way. These are the representations which respect the discrete chiral symmetry of the system, and are all obtained from the deformation of the representations of $`SL(2,)`$. The eigenvalues will thus be paired $`\pm \alpha _k`$ and will be the same for all three matrices because the symmetries are respected.
The explicit forms for the representation matrices fall into two classes, with dimensions $`M=2p`$ and $`M=2p+1`$, the analogues of half-integer and integer spins.
For $`M=2p`$, one finds
$`(\varphi _1)_k\mathrm{}`$ $`=`$ $`\delta _{k+1,\mathrm{}}a_k+\delta _{k1,\mathrm{}}{\displaystyle \frac{b_{\mathrm{}}}{a_{\mathrm{}}}},`$ (85)
$`(\varphi _2)_k\mathrm{}`$ $`=`$ $`iq^{kp}\delta _{k+1,\mathrm{}}a_kiq^{pk+1}\delta _{k1,\mathrm{}}{\displaystyle \frac{b_{\mathrm{}}}{a_{\mathrm{}}}},`$ (86)
$`(\varphi _3)_k\mathrm{}`$ $`=`$ $`i\alpha _k\delta _k\mathrm{}`$ (87)
and we have $`b_{p+j}=b_{pj}`$ for $`j=1,2,\mathrm{},p1`$, and $`\alpha _{p+n}=\alpha _{pn+1}`$ for $`n=1,2,\mathrm{},p`$.
The $`a_k`$’s may all be set to, say, unity, by $`SL(M)`$ transformations. The $`b_j`$’s are determined recursively by the formula<sup>7</sup><sup>7</sup>7We’ve defined $`\sigma _x[q]=1+q+q^2+\mathrm{}+q^x`$.
$$b_j=\frac{\frac{q}{1+q}\sigma _{2(pj)}[q]+b_{j1}\left(1+q^{2(pj)+3}\right)}{q^2\left(1+q^{2(pj)1}\right)};b_0=0,$$
(88)
for $`j=1,2,\mathrm{},p1`$. The recursion relation is solved by
$$b_{k,p}=\frac{q(q^{4p}q^{2k})(q^{2k}1)}{(q^21)^2(q^{2p}+q^{2k1})(q^{2p}+q^{2k+1})}$$
(89)
and notice that all singularities (poles and zeroes) happen for $`q`$ a root of unity.
All three matrices have the eigenvalues
$$\pm \alpha _n=\pm \frac{1}{q^{pn}(1+q)}\sigma _{2(pn)}[q].$$
(90)
for $`n=1,2,\mathrm{},p`$.
When $`M=2p+1`$, we have instead
$`(\varphi _1)_k\mathrm{}`$ $`=`$ $`\delta _{k+1,\mathrm{}}a_k+\delta _{k1,\mathrm{}}{\displaystyle \frac{b_{\mathrm{}}}{a_{\mathrm{}}}},`$ (91)
$`(\varphi _2)_k\mathrm{}`$ $`=`$ $`iq^{kp1/2}\delta _{k+1,\mathrm{}}a_kiq^{p\mathrm{}+1/2}\delta _{k1,\mathrm{}}{\displaystyle \frac{b_{\mathrm{}}}{a_{\mathrm{}}}},`$ (92)
$`(\varphi _3)_k\mathrm{}`$ $`=`$ $`i\alpha _k\delta _k\mathrm{}`$ (93)
where $`b_{p+n}=b_{pn+1}`$ and
$$b_n=\frac{q\sigma _{pn}[q^2]+b_{n1}\left(1+q^{2(pn+2)}\right)}{q^2\left(1+q^{2(pn)}\right)};b_0=0,$$
(94)
for $`n=1,2,\mathrm{},p`$. The recursion relation is solved by
$$b_{k,p}=\frac{q(q^{2k}1)(q^{4p}q^{2k2})}{(q^21)^2(q^{2p}+q^{2k2})(q^{2p}+q^{2k})}$$
(95)
and again we see that all singularities happen for roots of unity. We also have $`\alpha _{p+r+1}=\alpha _{pr+1}`$ for $`r=0,1,\mathrm{},p`$ and the eigenvalues of each matrix are in this case
$$0,\pm \alpha _n=\pm \frac{\sigma _{pn}[q^2]}{q^{(M2n)/2}}$$
(96)
for $`n=1,2,\mathrm{},p`$.
Note that the solutions that we have written here are not $`D`$-flat. However, by standard theorems, there exists such a solution, which is an $`SL(M)`$ transformation of the stated solutions. Still, we must be careful in drawing conclusions based on these solutions. In particular, there are apparent singularities at special values of $`q`$. We will analyze this point further in Section 4.5.2.
#### 4.5.1 Finding more solutions
So far, we have found representations of the algebra which in the limit $`q1`$ reduce to finite dimensional representations of the $`SL(2,)`$ algebra. We also noted an additional one-dimensional representation which becomes singular in this limit, and therefore corresponds to a vacuum of the theory, which goes to infinity in the limit. This additional solution is characterized by the property $`\mathrm{tr}\varphi _10`$, whereas for all the other solutions $`\mathrm{tr}\varphi _1=0`$.
We should ask if there are more irreducible representations of this algebra, that we have not found above. The answer must be yes, because for $`q1`$ many of the solutions which correspond to irreducible representations of $`SL(2,)`$ go away to infinity (the eigenvalues of the matrices are rational functions of $`q`$ with finite numerator and $`q+1`$ in the denominator. Thus they are infinitely far away in field space, and do not describe vacua of the theory.)
We can construct additional irreps that do not disappear in the $`q1`$ limit as follows. The discrete subgroup of $`SU(2)`$ has a three dimensional representation in terms of Pauli matrices, which suggests the following Ansatz for the representations.
The following satisfy the algebra (82)
$`\varphi _1`$ $`=`$ $`\varphi _1^{}(i\sigma _1)`$ (97)
$`\varphi _2`$ $`=`$ $`\varphi _2^{}(i\sigma _2)`$ (98)
$`\varphi _3`$ $`=`$ $`\varphi _3^{}(i\sigma _3)`$ (99)
if we have
$$[\varphi _j^{},\varphi _{j+1}^{}]_q=\varphi _{j+2}^{}.$$
(100)
Thus, if we know solutions for a given $`q`$, we generate solutions for $`q`$ in this way. These representations are reducible. We will refer to the the irreducible representations obtained in this way as twisted. There are two cases to consider, ‘half integer’ spin and ‘integer spin’ representations.
The integer spin representations have each eigenvalue repeated twice, including zero and are split into two irreducible representations with eigenvalues for $`\varphi _3`$ in the succession
$$\pm i\alpha _1i\alpha _2\pm i\alpha _3\mathrm{}0\mathrm{}i\alpha _2\pm i\alpha _1$$
(101)
These satisfy $`\mathrm{tr}(\varphi _3)0`$, and $`\mathrm{tr}\varphi _{1,2}=0`$, as these are off-diagonal. The broken $`_2`$ exchanges these two representations. By acting with the $`_3`$ symmetry we get a total of six new representations for each even-spin irreducible representation of $`SU(2)`$.
The ‘half-integer’ cases satisfy $`\mathrm{tr}(\varphi _{1,2,3})0`$. One can clearly see a splitting into two irreducible representations, but because there is no eigenvalue $`0`$, this splitting into two is reducible and in total we get four new representations of the algebra. One of these is a $`_3`$ singlet, and the other three form a triplet.
#### 4.5.2 Interpreting the singularities
In this section, we will study properties of representations. In general there are two classes of representations, irreducible and reducible. In the reducible case, there is no mass gap (classically) as some part of the gauge group is unbroken (apart from the decoupled $`U(1)`$). The case of irreducible representations are potentially more interesting as they confine magnetic degrees of freedom. We will exploit $`S`$-duality to find dual configurations that are electrically confining. Note that as we have not been able to prove that all irreducible representations are accounted for, we cannot be sure that we see all of the vacua. For the sake of the present argument, we will assume that the classification is complete and try to extract conclusions about the non-perturbative behavior of the theory.
The representations we have found are all matrices which are rational functions of $`q`$. From the solutions (89),(95), we see that there are poles at roots of unity, $`q^n=1`$.
In the case where we have zeroes and not poles, one observes that as we take the limit to an appropriate root of unity, the matrix decomposes in block-diagonal form. Thus the representation becomes reducible in the limit, and we get various copies of the same type representations (of lower dimension). These singularities are interpreted in the field theory as having enhanced gauge symmetry, because the commutant of the representation is larger. If one pictures the vacua of fixed rank as a covering of the $`q`$-plane, we have branch points at some roots of unity.
There are other singularities at roots of unity in the denominators of the fields $`\varphi _{2,3}`$. As these are not singularities in the eigenvalues of the matrices, it is not clear that these are singular solutions. This may correspond to an unfortunate choice of basis for the representation.
Considering that the roots of unity are special, in the sense that they are related to orbifolds with discrete torsion which have a very nice semi-classical geometry associated to them, and also considering that in the limit $`q\pm 1`$ an infinite family of solutions to the vacua disappear (in this case there are singularities in the eigenvalues of the matrices), it is plausible that these are actually bona-fide singularities and the vacua go to infinity. As we will see, at these values of $`q`$, there are moduli spaces of vacua and this is how we interpret the singularities.
Let us begin with a discussion of $`q=\pm 1`$. First, we know that at $`q=1`$ all of the states which break the chiral symmetry disappear. Thus we get a jump in the Witten index at this special value. It is also the case that here for some representations one sees no signal of the eigenvalues of the matrices being badly behaved, but it is true that we get poles in the off-diagonal elements.
Let us now discuss $`q=1`$, paying particular attention to discrete chiral symmetry breaking. For $`U(M)`$, $`M`$ even, the $`q`$-deformed $`SU(2)`$ representations move off to infinity at $`q=1`$, and thus all the irreducible representations come from the ‘half integer’ twisted case. Thus the Higgs vacua break the $`_2\times _2`$ subgroup completely, and the vacuum has an unbroken $`_3`$ subgroup. Each of the four vacua have the $`_3`$ embedded differently.
For $`U(M)`$, $`M`$ odd, there are irreducible representations of either integer or half-integer twisted type. Thus, some of the Higgs vacua break the group to an unbroken $`_3`$ as in the previous case, and some leave an unbroken $`_2`$ if they are constructed from the ‘integer spin’ type representations. In addition, the $`q`$-deformed representations survive for $`q1`$ but are reducible (the matrix elements $`b_{k,p}0`$).
Notice that in the previous arguments we have used only the perturbative symmetries of the theory. We believe that the $`SL(2,)`$ S-duality of $`N=4`$ SYM is realized and perhaps enlarged in the present case in some way. We will not address that interesting question here; instead, we confine ourselves to a few remarks based on $`SL(2,)`$ alone.
Because of the $`SL(2,)`$ symmetry, at the $`N=4`$ point we can make a map of gauge invariant operators between the different dual theories. Thus we can follow the deformations of the theory for any S-dual configuration of the $`N=4`$ theory we start with.
Because of the symmetries preserved by the superpotential, changing from one dual picture to another keeps the general form of the Lagrangian invariant. Thus we have a map between couplings $`(g,q)(g^{},q^{})`$, and $`mm^{}(g,q)`$. Because at the roots of unity the theory is special (many vacua collide), the roots of unity must be preserved by the S-duality action on the space of field theories, thus the most general holomorphic transformation that keeps $`q=1`$ fixed and the structure of the singularities is of the form $`qq^{\pm 1}`$.
Given a vacuum that disappears at a root of unity, let’s say a Higgs vacuum, any of the vacua related to it by S-duality also disappear. For $`q=1`$ and $`M`$ even, the trivial vacuum is $`S`$-dual to the $`q`$-deformed Higgs vacuum which moves off to infinity as $`q1`$, and thus the trivial vacuum is also removed. For $`M`$ odd, again the trivial vacuum is $`S`$-dual to the $`q`$-deformed Higgs vacuum. The latter is reducible, and thus does not appear to have a mass gap; we conclude that the trivial vacuum is not confining. There are still the twisted representations, and thus at $`q=1`$, confinement implies (discrete) chiral symmetry breaking.
If $`q`$ is a more general $`n`$th root of unity, even though we get poles in the $`b_{k,p}`$, we have not been able to find any gauge invariant chiral quantity which becomes singular. This suggests that the poles are obtained from a coordinate singularity. In any case, there seems to be an upper bound on the number and dimension of irreducible representations, as each of these general representations seems to decompose into irreducibles of smaller rank. The bound is given in terms of $`n`$.
Thus as $`q`$ goes to a root of unity, we can obtain enhanced gauge symmetry. If we do an S-duality transformation and use some more general combination of electric and magnetic condensates, there will still be an upper bound on the dimension of irreducibles and thus no mass gap.
The upper bound on the irreducibles also suggests that one can construct a large center for the algebra. Indeed, one can take the direct sum of all the irreducible representations of the algebra we have constructed. If there are no more irreducible representations, this is a finite dimensional reducible representation, and the subalgebra of the $`\varphi _i`$ which is the inverse image of the center of the representation is a large center for the full algebra.
Experience with the example in Section 4.3 suggests that in this case one might actually get a moduli space of vacua. As we have argued that we get a finite number of discrete vacua, let us now show that there is a moduli space for roots of unity $`q^n=1`$ with $`n>2`$.
We would want the moduli space to be built out of the $`P,Q`$ matrices in some simple fashion. Let us choose $`\varphi _1`$ to be diagonal. Without the mass deformation, the $`\varphi _i`$ contain $`P,Q,P^1Q^1`$. Indeed, one can see that only with the powers $`P^{\pm 1}`$, $`Q^{\pm 1}`$ can one get a single factor of $`q`$ in the commutation relations, and there is a potential to get a cancellation of terms. Thus we take
$$\varphi _1=a_1P+a_2P^1$$
(102)
Because of the symmetry between $`P,Q`$, we also take
$$\varphi _2=a_3Q+a_4Q^1$$
(103)
and the $`q`$ commutation relations are as follows
$$\varphi _1\varphi _2q\varphi _2\varphi _1PQ^1+QP^1$$
(104)
so we take
$$\varphi _3=a_5PQ^1+a_6QP^1$$
(105)
The parameters are related by
$`a_1a_4(1q^2)`$ $`=`$ $`ma_5`$ (106)
$`q^1a_5a_2(1q^2)`$ $`=`$ $`ma_4`$ (107)
and thus it follows that
$$a_5a_6=a_3a_4=a_1a_2=\frac{qm^2}{(1q^2)^2}$$
(108)
so apart from factors depending on $`m,q`$, $`a_{2i+1}a_{2i}1`$. This cuts the number of variables from six down to three, and (106) provides one more constraint. Thus we are left with a two parameter solution of the $`F`$-term constraints. More surprisingly, these also solve the $`D`$-term constraints. These representations are inequivalent as one can show that the gauge invariant vacuum expectation value $`\mathrm{tr}(\varphi _1^n)`$ is not independent of the $`a_i`$. For $`q=\pm 1`$, eq. (108) shows that the $`a_i`$ are singular, and $`\mathrm{tr}(\varphi _1^2)`$ is singular for $`q=1`$, thus this branch of moduli space does not appear at these roots of unity, and one only has isolated vacua.
One can also explicitly show that for $`q^3=1`$, the element $`x_i=\varphi _i^3+\frac{m^2}{q}\varphi _i`$ is central. Thus here one gets a large center, as we have four Casimir operators and one relation. The fourth Casimir is of the form
$$w=A\varphi _1\varphi _2\varphi _3+\alpha _1\varphi _1^2+\alpha _2\varphi _2^2+\alpha _3\varphi _3^2$$
(109)
and it is invariant under the full discrete group of symmetries of the potential. A Casimir of this form exists for all $`q`$, and when $`m=0`$ it is the familiar $`\varphi _1\varphi _2\varphi _3`$; it also reduces to the quadratic Casimir of the $`SU(2)`$ algebra when $`q1`$.
The commutative space associated to the algebra is again a deformation of the $`^3/_3\times _3`$ orbifold, and it is three complex dimensional. We have only found a two parameter solution of the equations; we believe that this is because we chose a very special form for the solutions, and not necessarily because the ring fails to be semi-classical.
### 4.6 The General Superpotential
In this section we will try to make progress towards understanding the general deformation, eqs. (3,4). Solving for the center of the general algebra and also finding the most general finite dimensional irreducible representations of the algebra can be quite difficult. There are some cases which are worth singling out among these, because at least we can find some partial solutions to the moduli space problem. We have also seen that semi-classical rings are better behaved than others, as they lead to nice commutative geometries. Finding all possible semi-classical geometries from our sets of constraints is very important as they might correspond to the behavior of $`D`$-branes at new dual singularities (not necessarily orbifolds with discrete torsion) which can be connected to $`AdS_5\times S^5`$. Of particular importance are configurations with conformal invariance, as they might provide new non-spherical horizons. Our analysis is quite incomplete due to the difficulties of the algebraic program involved, but some general comments will be made here.
As the deformations are taken to zero, the algebra looks like a Poisson algebra if we interpret commutators as Poisson brackets. Because we have three variables, and Poisson manifolds are foliated by symplectic manifolds (which are of even dimension), the symplectic form in the full algebra is degenerate and therefore there is at least one constant of motion. This suggests that there is at least one element of the center which can be easily computed. For the $`q`$-deformations, the element of the center $`w=\varphi _1\varphi _2\varphi _3`$ exists for arbitrary values of $`q`$, which suggests that the element of the center is a polynomial of degree less than or equal to three, depending on the chosen perturbation. For marginal deformations, it is indeed of degree three, as will be shown later; for the deformation by three mass terms it is quadratic (the Casimir of the $`SU(2)`$ algebra).
Because we have commutators we can think of the algebra as deformation-quantization of the Poisson structure. This suggests that we can standard order operators and establish a correspondence between the full algebra, and the algebra of three commuting variables. In standard constructions, this is given by formal power series expansions in a small parameter $`\mathrm{}`$. As we have argued before, we want to avoid infinite power series, and rather give an explicit solution which shows that the constraints can be standard ordered in some open set. In order to do this, we separate at each order the polynomials in $`\varphi _1,\varphi _2,\varphi _3`$ which can be considered as standard ordered.
A choice of standard ordering is important. If we want to find elements of the center, we need to check that their commutators are zero for all the generators of the algebra. Without standard ordering a given expression, it is very hard to decide if it is zero or not in the algebra.
By using the constraints, an arbitrary polynomial operator $`𝒪`$ can be re-ordered into standard ordered form up to small corrections. We write this as
$$𝒪=𝒪_{so}+\mathrm{}𝒪^{}$$
(110)
$`𝒪_{so}`$ is a linear combination of standard ordered monomials, and it is polynomial in $`\mathrm{}`$. Similarly, we can expand $`𝒪^{}a_iM^i`$, where the $`a_i`$ are polynomial in $`\mathrm{}`$ and the $`M^i`$ are a collection of non-standard ordered monomials. Because of the form of the algebras, the degree of $`𝒪^{}`$ as a polynomial in the variables of the algebra is smaller than or equal to the degree of $`𝒪`$. Taking all the possible non-standard ordered monomials of degree less than or equal to some fixed number $`g`$, we obtain a matrix equation
$$M^i=M_{so}^i+\mathrm{}a_j^iM^j$$
(111)
Now $`\mathrm{}`$ is a small parameter, so the matrix
$$A_i^j=\delta _i^j\mathrm{}a_i^j$$
(112)
is finite dimensional and invertible. Hence, any non-standard ordered operators may be written as linear combinations of the standard ordered operators, where the coefficients are rational functions in the deformation parameters, the denominators coming from $`A^1`$.
Since the parameters are complex, more generally we need only worry about the possibility of poles in this construction. At such poles, one of two things can happen. Either the basis for standard ordered polynomials is badly chosen, (e.g., the elements become linearly dependent), or there is a true obstruction to standard ordering independent of the basis. This second possibility can happen, if we take $`q=0`$ for example.
Thus, in principle we can proceed order-by-order in the degree of polynomials to find central elements. Every element of the algebra can be written in standard ordered form, and as the degree of the element is preserved or lowered by the commutation relations, it is a matter of linear algebra to calculate the elements of a given order which are in the center.
Although the procedure is well-defined, it is not efficient, as we need to calculate the matrix $`A`$ at each order to resolve this problem. Thus a general solution of how the center depends on the parameters is at best difficult to calculate. Also notice that in the examples we have studied, there is no upper bound in degree for elements of the center.
In some cases, we may find a large center; that is, the center is generated by more than one element of the algebra. If the center is large enough, then we may obtain a semi-classical algebra.
Let us consider the case where the algebra $`𝒜`$ is a finitely generated module over its center, with generators $`e_i`$. We can choose one of the generators to be the identity in the ring, and the others will satisfy a multiplication rule of the type
$$e_ie_j=f_{ijk}e_k$$
(113)
with $`f_{ijk}𝒵𝒜`$. On a given irreducible representation of the algebra, the elements of the center can be treated as numbers, and thus we can argue that we have a family of algebras parametrized by the algebraic variety corresponding to the center of the algebra.
Because of the form of eq. (113), we can see that given a vector in the representation of the algebra, its orbit under the action of the $`e_i`$ is finite dimensional. Thus there is an upper bound on the dimensions of the irreducible representations. We can imagine that this upper bound is realized by the branes living in the bulk, and that any other representation with smaller dimension is a fractional brane of some sort. The finite dimensionality of the irreps suggests that the ring is semi-classical, although we have no proof of this assertion. The semi-classical rings that we have studied all have this property, and this suggests that the two conditions might be equivalent.
Let us now consider a few more examples.
#### 4.6.1 General marginal deformations
As an example of the general difficulties that one faces, let us consider a general marginal deformation of the $`N=4`$ theory. The superpotential is given by
$$W=\mathrm{tr}\left(\varphi _1\varphi _2\varphi _3q\varphi _2\varphi _1\varphi _3+\frac{\lambda }{3}\left(\varphi _1^3+\varphi _2^3+\varphi _3^3\right)\right)$$
(114)
and the equations we need to solve for the moduli space are (cyclic)
$$[\varphi _j,\varphi _{j+1}]_q=\lambda \varphi _{j+2}^2$$
(115)
This algebra is homogeneous, and thus if we were able to find a non-trivial irrep, we could scale it to zero: this implies that the moduli space is connected.
For $`\lambda =0`$ the element of the center that is always present is $`w=\varphi _1\varphi _2\varphi _3`$, and this suggests that the element of the center is cubic in general. Indeed, a direct calculation shows that
$$(1q)\varphi _1\varphi _2\varphi _3\lambda \varphi _1^3+q\lambda \varphi _2^3\lambda \varphi _3^3$$
(116)
is central.
Let us first consider one-dimensional irreps. These will satisfy (cyclic)
$$(q1)\varphi _j\varphi _{j+1}=\lambda \varphi _{j+2}^2$$
(117)
A non-trivial solution will have the $`\varphi _j`$ all non-zero complex numbers. We can easily see that this is only solvable provided that
$$(q1)^3=\lambda ^3$$
(118)
so $`(q1)/\lambda `$ is a cube root of unity. Given $`\lambda `$ and $`q`$ satisfying these constraints, one can find solutions to the equations where $`\varphi _1`$ and $`\varphi _2`$ are equal up to cube roots of unity, and then $`\varphi _3`$ is determined from the other two. Thus we get three complex lines meeting at the origin, reminiscent of the moduli space for $`q`$-deformations. Indeed when $`\lambda `$ and $`q`$ are related in this way, there is a linear change of basis of the fields which returns the superpotential to a $`q`$-deformation. Therefore we have new semi-classical rings, but they are related by a change of basis to the ones we already know.
Another thing that we can do is exploit the $`_3`$ symmetry which permutes $`\varphi _1,\varphi _2,\varphi _3`$. Set $`\varphi _1=aP`$, $`\varphi _2=bQ`$, $`\varphi _3=cP^1Q^1`$; for the $`P,Q`$ matrices associated to the cube roots of unity, we also have $`\varphi _3^2\varphi _1\varphi _2`$, so three-dimensional irreps of the algebra may exist.
In this case we want to find solutions to (cyclic on $`a,b,c`$)
$$ab(q\omega 1)=\omega \lambda c^2$$
(119)
with $`\omega `$ a cube root of unity. One can see that this gives us the constraint
$$(q\omega 1)^3=\lambda ^3$$
(120)
For $`\lambda 0`$ we associate the geometry to $`^3/(_3\times _3)`$ which happens to be one of the orbifolds one can realize globally on a three-dimensional complex torus.
If for $`\lambda 0`$ we can find a large center, we might be able to compute the full geometry of moduli space, and treat this solution as a fractional brane. Notice that if this is the case, it does not correspond to a $`^3/_n\times _n`$ singularity, as the fractional branes in that case behave differently. Further exploration of this model will be left for future work.
## 5 $`D`$-branes in near-horizon geometries
So far, we have mainly discussed moduli spaces of vacua and how to include extended objects in the discussion. The analysis has been done directly in the field theory. Now we will try to understand the background and the moduli space that the $`D`$-branes realize from the AdS/CFT perspective.
The field theory is understood as being dual to the near-horizon geometry of a brane configuration. Because the moduli space is of the form of a symmetric product (built out of smaller components), one can think of adding these small component $`D`$-branes as probes in the near horizon geometry and testing how the field theory moduli space is realized on these probes.
We will carefully compare the field theory marginal and relevant deformations to the corresponding deformations of $`AdS_5\times S^5`$ geometry. In doing so, we uncover and solve several puzzles. In particular, there are new branches of moduli space in the field theory which open up for arbitrarily small values of $`q1`$, as discussed in Section 4.1. Uncovering this structure in the string theory will have several bonuses. This nongeneric branch is realized by wrapping a 5-brane on a 2-torus and using this information, we will argue that the mirror symmetry between deformed 5-spheres and orbifolds can be understood as a standard T-duality operation. The two supergravity descriptions are valid in different areas of parameter space. It also becomes clear in this analysis that there is no sense in which the field theories are dual to supergravity on a space; rather, string theory is absolutely necessary for a consistent duality.
We have seen that the moduli space of vacua has very non-trivial behavior in response to the deformations. In particular, the somewhat artificial separation between the center of the algebra and other elements of the algebra is very subtle in field theory. This will be addressed later and we will find a satisfactory solution. If we look at the same construction from the $`AdS_5\times S^5`$, each of these perturbations is in the bulk of the $`S^5`$ geometry, and there is no reason to single out any special elements of the algebra.
### 5.1 Effects of the Background on $`D`$-branes
It is important to notice that as seen from the $`AdS_5\times S^5`$ perspective when one deforms the theory, the $`D`$-brane moduli space changes drastically. To first approximation, this is because the added potential localizes the $`D`$-branes to the ‘fixed planes’ of the deformation. But even for very small deformations $`q1`$, we can find rational solutions of $`q^n=1`$ for large $`n`$, and thus the moduli space has non-generic behavior for a large enough number of branes; indeed, we need $`n`$ such branes to find extra components of moduli space. These new branches can be seen from (48), and predict that the $`D`$-branes are going to be uniformly distributed on a circle. We take the eigenvalues of matrices to determine the coordinates of the $`D`$-branes, as in matrix theory.
If the branes are point-like then the open string states stretching between them would be massive and one would not find the new branch in moduli space. However, this is clearly inconsistent with our field theory results, and thus we are motivated to find a satisfactory solution within string theory. Note that these extra components of moduli space do not just appear in the vicinity of the origin; rather, they extend to infinity with the rest of $`D`$-brane moduli space.
The resolution of these issues bears close resemblance to recent results of Myers concerning the dielectric properties of branes in background fields. Since the deformations of the field theory superpotential correspond to non-zero vevs of fields on the 5-sphere, we do indeed expect these phenomena to occur. Roughly speaking, the $`D3`$-branes should be thought of as $`D5`$-branes on $`^4\times S^2`$, where, as we show below, the $`S^2`$ is contained in $`S^5`$. In order to find new branches of the moduli space, we want to argue that there are configurations which support massless open string modes, and topologically this will happen when different spheres intersect each other. Thus their centers can be separated, and we can still have massless string states stretching between them.
Now let us begin by analyzing in some detail the map between superpotential deformations and vevs. This material is of course not new, but is included here for completeness.
The $`q1`$ and $`m`$ deformations correspond to background values for magnetic potentials $`F_{(3)}^{RR}`$ and $`H_{(3)}^{NS}`$. The mass deformation is not marginal, and will therefore depend on the radial direction of $`AdS_5`$. The field $`\tau =C+ie^\varphi `$ gives the gauge coupling, and will be kept constant. The field $`G_{(3)}=F_{(3)}\tau H_{(3)}`$ is related directly to the superpotential deformations. The harmonic in the $`\mathrm{𝟏𝟎}`$ of $`SU(4)`$ is a tachyon state in the $`AdS`$, thus this perturbation blows-up in the infrared. The marginal cubic operators correspond to a harmonic of $`G_{(3)}`$ in the $`\mathrm{𝟒𝟓}`$ of $`SU(4)`$. In this case, there will be no dependence on the radial direction of $`AdS`$ as the associated scalar is massless in five dimensions; this fact guarantees that we preserve the conformal group to leading order.
Let us now specialize to the marginal deformations. As explained in Ref. , $`D3`$-branes in the presence of $`RR`$ background fields pick up a dipole moment for higher brane charge, and become extended in two additional dimensions. The simplest topological shape, and the one with the lowest energy, is a 2-sphere centered at the position of the $`D3`$-brane. Since we are considering a weakly coupled string theory regime, we should take these to be $`D5`$-branes. More precisely, the $`F_{(3)}`$ background is dual to $`F_{(7)}`$ which couples to a $`D5`$-brane. $`F_{(7)}`$ has support on $`^4\times D^3`$, where $`D^3`$ is the 3-disk with $`S^2`$ boundary. The $`D3`$-branes are extended in the $`^4`$, which in near-horizon geometry is contained in $`AdS_5`$. We thus write $`F_{(7)}=\stackrel{~}{F}_{(3)}dVol_4`$, and integrating, we can normalize it such that
$$_{^4\times D^3}F_{(7)}=_{D^3}\stackrel{~}{F}_{(3)}$$
(121)
The 3-disk extends along the radial direction of $`AdS_5`$ plus two directions along the the $`S^5`$. As a result, we can write
$$\stackrel{~}{F}_{(3)}=d\rho \stackrel{~}{C}_{(2)}$$
(122)
As such, if the effect were solely due to the dielectric effect it is hard to understand how the $`D`$-branes can have massless states at different angles along the $`S^5`$, as the stretching happens mostly in the radial direction. The $`D3`$-brane charge of this 5-brane is obtained from a flux through the 2-sphere, $`\frac{1}{2\pi }_{S^2}F=n`$.
As follows from Ref. , there will also be a background $`H_{(3)}^{NS}`$ turned on in the presence of the superpotential deformations. If we expect some energy contribution from the integral of $`H_{(3)}^{NS}`$ over the disk, then the 2-sphere prefers to be stretched along the 5-sphere, because $`H_{(3)}^{NS}`$ does not have any component along the $`AdS`$ directions. In general, then, the radius of the disk $`D^3`$ is oriented partially in the radial direction of $`AdS_5`$ and partially in $`S^5`$, as there are two competing effects deforming the branes.
We want to look for configurations where $`D3`$-branes are intersecting in the sense of intersections of their $`S^2`$’s. This is where we can expect massless string states, at least topologically. The $`H^{NS}`$ deformation is the one that gives us the deformation of the $`D`$-branes in the appropriate direction. We will assume that these configurations are supersymmetric and that probes do not affect the background.
For rational $`q`$, the moduli space has a scaling direction, which follows from the conformal invariance: in the language of Section 4.1, we have
$$a,b,cta,tb,tc$$
(123)
This is reflected in the near-horizon geometry by the fact that if we move $`D3`$-branes along the radial direction of $`AdS_5`$, they simply rescale–in particular, if we have intersecting branes, they remain intersecting as we perform this motion.
For relevant deformations, such as a mass term, the $`RR`$ and $`NS`$ backgrounds grow as we move in along the $`AdS_5`$, and thus we expect the 2-spheres to grow in size along the 5-sphere as we go to the infrared. As in this case the $`H_{(3)}`$ fields will also have a radial component, then both types of fields $`H_{(3)}`$ and $`\stackrel{~}{F}_{(3)}`$ help the 2-sphere to grow along the $`S^5`$ and the radial direction. Eventually, the 2-spheres will be of comparable size to the 5-sphere, and at this point, the notion of point-like $`D`$-branes loses any meaning. To avoid these issues and for ease of calculation, we will treat only marginal deformations in the following sections.
### 5.2 Size and Configurations of 5-branes
First, let us find the expected size $`r`$ of an $`S^2`$ associated with a $`D3`$-brane. We will assume that this $`S^2`$ is small compared to $`R_5`$, the radius of $`S^5`$, but that it is large enough that we can neglect its self-interaction (from opposite sides of the $`S^2`$). We will show that for small deformations, the size grows linearly with the potential. To this effect, we do a probe calculation. We have a geometry which is almost $`AdS_5\times S^5`$ generated by some $`D`$-branes which are at the origin, and we have a small extra $`D`$-brane which turns into a sphere on which we are going to do our analysis. Because conformal invariance is preserved by the marginal deformations, there can be no dependence on the $`AdS`$ radial direction in the physical quantities of interest, apart from setting the scale of the physics. We can therefore work in a local frame and ignore redshift factors, etc.
The DBI action for a $`D5`$-brane determines the energy
$$E_{DBI}=\frac{\mu _5}{g}d^2\mathrm{\Omega }\sqrt{det(GB+2\pi \alpha ^{}F)}\mu _5_{D^3}\stackrel{~}{F}_{(3)}\mu _5(2\pi \alpha ^{}FB)C_4$$
(124)
The metric $`G`$ scales as $`r^2`$, whereas $`F`$ behaves as $`r^0`$ (by the flux quantization). By expansion of the DBI part for small $`r`$, we find an energy of the form
$$E=E(D3)+\frac{\alpha }{g}r^4\beta r^3+o(r^5)$$
(125)
We write
$$_{D_3}H_{(3)}^{NS}=c_{NS}r^3_{D_3}\stackrel{~}{F}_{(3)}=c_Rr^3$$
(126)
The field strengths are constant over the disk. We will do the analysis ignoring the five-form field strength. At the end, we will compensate for this omission. The general features of the result should not depend on how far we are from the origin. This is how we can justify this omission.
From the energy (124), we see that the $`D3`$-brane charge is given by the coupling to $`C_4`$
$`Q_3`$ $`=`$ $`n{\displaystyle \frac{1}{4\pi ^2\alpha ^{}}}{\displaystyle _{S^2}}B`$ (127)
$`=`$ $`n{\displaystyle \frac{c_{NS}}{4\pi ^2\alpha ^{}}}r^3`$ (128)
The expansion of the energy in powers of $`r`$ now reads
$$E=\frac{\mu _3}{g}\left(Q_3\frac{c_R}{4\pi ^2\alpha ^{}}r^3+\frac{1}{(4\pi ^2\alpha ^{})^22n}r^4+\mathrm{}\right)$$
(129)
where $`Q_3`$ is a constant plus small coorrections in $`r^3`$. The result is minimized at
$$r\frac{3}{2}(c_{NS}+gc_R)(4\pi ^2\alpha ^{})^2n$$
(130)
The energy at this radius satisfies
$$E=\frac{\mu _3}{g}Q_3\left(1+\frac{1}{4Q_3}r^3(3c_{NS}gc_R)+\mathrm{}\right)$$
(131)
This result is puzzling, since it suggests that $`n`$ $`D3`$-branes extend to a single $`S^2`$ of radius proportional to $`n`$, as opposed to a sphere wrapped $`n`$ times around the solution for a single brane. This result is wrong from several points of view. First, this solution cannot give enhanced $`U(n)`$ gauge symmetry, as there are no massless states apparent, and suggests a totally different picture of the moduli space, very different for each value of $`n`$. We must be more careful in interpreting eq. (131).
We interpret the branes as a black hole in the supergravity which is almost pointlike as far as the $`S^5`$ is concerned. One minimizes the energy (129) and then compares the ratios of energy to $`D3`$-brane charge of two configurations to determine which may be BPS. In fact, $`n`$ $`D5`$-branes wrapping an $`S^2`$ of radius $`r`$ have lower $`E/Q_3`$ than a single $`D5`$-brane wrapping an $`S^2`$ of radius $`nr`$. This indicates that the former configuration has the better chance of being BPS.
The stabilization mechanism which impedes the spheres from shrinking further is that the flux of $`F`$ is quantized. This mechanism has been found when studying $`D`$-branes from the boundary state formalism for group manifolds, but it is clear that it should be a general phenomenon for $`D`$-branes in non-trivial $`H_{(3)}^{NS}`$ backgrounds.
Here we see also that the $`RR`$ charge for the large sphere is not quantized in general as it gets an anomalous defect proportional to $`H^4n^3`$. These can be meta-stable boundary states, and in group manifolds these can be calculated exactly, where a similar defect in the brane charge quantization condition occurs. When $`H_{NS}=0`$ the $`D`$-brane charge is related to K-theory, and then we expect a quantization condition. This puzzle was recently solved in Ref. where there is a back reaction from the bulk which contributes to the 3-brane charge. Thus, eq. (131) is incomplete as it does not take into account the energy associated to this back reaction. However, the ratio $`E/Q_3`$ on the horizon of the brane probe is exact, being the local tension divided by charge. Here the BPS $`D`$-branes behave better, as we get an anomalous charge which is proportional to the $`D`$-brane number.
### 5.3 Large $`n`$ branches
We now want to find the new branches of moduli space for $`q^n=1`$, by finding the geometric configuration into which it can be deformed. Because the marginal deformation preserves the conformal group, in the near-horizon geometry, the $`H_{(3)}^{NS}`$ lies entirely along the 5-sphere, and thus the $`D`$-brane becomes spherical along 5-sphere directions.
For the $`q`$-deformation, this means that the $`D`$-branes grow in size linearly with respect to $`q1`$, which is the small parameter. This is the important point of the calculation in the previous section. Consider a configuration of $`n`$ of these 2-spheres, distributed around a circle in $`S^5`$ so that they touch each other. The value of $`n`$ is proportional to $`(q1)^1`$, in accordance with $`q^n=1`$.
In order for the $`D`$-branes to touch, we need to know the shape of the 2-spheres well. For $`|q|=1`$, one finds massless states between 2-spheres which are at the same distance from the origin in $`AdS`$ space. To see this, we can calculate the masses of the off-diagonal states from the superpotential
$$\mathrm{tr}\varphi _1\varphi _2\varphi _3q\mathrm{tr}\varphi _2\varphi _1\varphi _3$$
(132)
with
$$\varphi _1=\left(\begin{array}{cc}a& 0\\ 0& b\end{array}\right)$$
(133)
These masses are proportional to $`aqb`$ and $`bqa`$, and thus in order to have massless states for $`|b|=|a|`$ we need $`|q|=1`$.
In order to get the 2-spheres to touch when they are at different radii, the dielectric effect on the $`D`$-brane has to be included, as it is responsible for extending the $`D`$-brane in the radial direction. We now want to argue that the $`D5`$-branes laying flat on the $`S^5`$ actually do touch at another point. The reason why this is important is that moving apart a pair of 2-spheres on $`S^5`$ might make it impossible for them to touch again. Because of the geometric setup, if we consider two $`D3`$-branes at the same location and we move one with respect to the other in moduli space (of one real dimension on the $`S^5`$), we will get two $`2`$-spheres. Because the solution of the linearized supergravity equations rotates $`H_{(3)}^{NS}`$ as we move along this one parameter, the spheres become linked on the $`S^5`$. This is explicitly shown in Figure 3.
As the spheres are unlinked when they are very far from each other, they necessarily pass through a point where they touch. This is, topologically, the place where the extra states become massless. With the dielectric effect turned on the spheres are tilted with respect to the $`S^5`$ and that is why they touch at different values of their radial position. The tuning required to make the $`D`$-branes lie flat on the $`S^5`$ is precisely the action of removing the dielectric effect on the $`D`$-branes, and corresponds to one real condition on a one complex parameter deformation of the theory.
Thus, we arrive at a configuration of spherical $`D5`$-branes which touch at points. Now there are configurations, for rational $`q`$, with $`n`$ spheres where each touches the next one and they stack on a circle. This is the configuration where the new branch of moduli space opens up, as in eq. (48). This structure should be thought of as a 2-torus with $`n`$ pinches. Indeed the massless states at the intersection of the $`D`$-branes are such that they resolve the pinching points into tubes, as in Figure 4.
This resolution of these configurations is equivalent to moving onto the new branches of moduli space.
A pinching torus with $`n`$ nodes is also exactly the degeneration which produces fractional branes in an ALE singularity or on an elliptically fibered Calabi-Yau in F-theory. Thus this configuration of branes seems to be the right one to deform into the extra branches of moduli space for the rational values of $`q`$.
Notice also that this semi-classical torus is reflected also in equations (45), where we see a realization of a non-commutative torus algebra via clock and shift operators. Thus the non-commutative geometry description of the moduli space knows that the $`D`$-brane in $`AdS_5\times S^5`$ is shaped like a torus, and that when we deform to the degeneration, we split the torus into $`n`$ spheres (fractional branes), as required by the ALE singularity type of the orbifold in moduli space.
The picture presented above is meant as a topological argument for the branches of moduli space in string theory. These arguments rely upon a few technical assumptions, which we think are reasonable. We have assumed that the different $`D`$-branes do not affect each other and that they intersect at supersymmetric angles. Although it would be nice to assert this, as it would make our whole construction a purely topological argument, there is no natural complex structure for the spheres which would guarantee this property, and we have to rely on a dynamical mechanism instead. For completeness we should study the possibility that the 2-spheres might interact strongly with each other near the intersection point. In that case, the 2-spheres would develop a throat between them; so, topologically, we have a sphere, with a line bundle of degree two to count the number of D3 branes. When we move in moduli space we deform the line bundle and the metric. For a non-generic bundle, one can get extra states which are massless, and these would be the extra massless modes one needs. Of course, because the field theory tells us that the massless states are there, we believe that these constructions are sensible.
A second point which needs to be made is that although we made an argument based on $`D5`$-branes, by the $`SL(2,)`$ duality we can make an argument with any $`(p,q)`$-5 brane. Thus the fact that the $`RR`$ and $`NS`$ mix in the near-horizon geometry is necessary to implement the $`SL(2,)`$ duality on the field theory space of deformations as we change the string coupling $`g`$ and make different $`(p,q)`$-strings light. The reason we get a description purely in terms of $`D`$-branes is that we are using weakly coupled string theory, and for any other brane with $`NS5`$-brane charge the fundamental strings cannot end on it. This ambiguity in the description has also been found in Ref. . In their case, only one configuration would be such that the supergravity degrees of freedom were weakly coupled through most of the geometry.
### 5.4 Mirror Symmetry
We have seen that the construction of moduli space suggests a two torus fibration of the five sphere. This torus can be made explicit by using the following invariant coordinates
$$r_1^2=|\varphi _1|^2,r_2^2=|\varphi _2|^2,r_3^2=|\varphi _3|^2,w=\varphi _1\varphi _2\varphi _3$$
(134)
indeed, $`\rho ^2=r_i^2`$ is the radial direction in $`AdS_5`$ and $`w`$ is equal to $`r_1r_2r_3`$ except for a phase. We get a total of three real coordinates on the $`S^5`$, and we are left with two phases to determine, which are the arguments of $`\varphi _1/\varphi _2`$ and $`\varphi _1/\varphi _3`$. These two phases determine the two-torus fibration on the $`S^5`$, and the fibration is independent of how many branes are stacked together to get the new branches of moduli space.
Note that the $`T^2`$ so obtained may have $`n`$ nodes (related to the $`D3`$-brane charge) but the $`T^2`$ may wrap $`m`$ times around the $`S^5`$ before closing. The latter clearly corresponds to 5-brane charge. In Ref. a mirror symmetry was noted between string theory on a deformed 5-sphere and an orbifold theory. We are now in position to demonstrate that this mirror symmetry may be obtained by T-duality<sup>8</sup><sup>8</sup>8This is expected from the work of Strominger, Yau and Zaslow. on the near-horizon geometry, where the T-duality is taken fiberwise on the $`T^2`$ fibration.
The T-duality acts on the 2-torus that we have described above. Explicitly, the charges $`(m,n)`$ transform as a doublet under the $`SL(2,)`$ T-duality. Choose a mapping that takes $`(m,n)`$, with $`m,n`$ relatively prime, to $`(0,1)`$; this mapping will single out a point-like $`D3`$-brane on the mirror. This is achieved by the matrix
$$M=\left(\begin{array}{cc}a& m\\ b& n\end{array}\right)$$
(135)
where $`a,b`$ are fixed numbers, modulo $`m,n`$ respectively. The torus with complexified Kähler form $`K=B+iA`$ is taken to a torus with a different value of $`K`$. Explicitly, we have
$$KK^{}=\frac{aK+m}{bK+n}$$
(136)
The area of the torus goes to zero and $`B_{NS}`$ is smooth at the singularities (where only one phase remains). We can examine the effect of the transformation on $`K`$ near this limit. Indeed, we get that in the dual torus
$$K^{}\frac{m}{n}$$
(137)
which signals a constant $`B`$-field of strength $`m/n`$. This value is quantized and its fractional part corresponds to the discrete torsion phase.
As the area of the two torus is not constant, if we start without $`H_{NS}`$ then upon the T-duality, we will get a varying $`B_{NS}`$ flux through the dual torus, and thus we have generated an $`H_{NS}`$ in the T-dual geometry. If we want to cancel this quantity, there is a choice of $`B_{NS}`$ which makes $`\mathrm{Re}(K^{})=\frac{m}{n}`$ constant over the dual fibration. This determines explicitly the $`H_{NS}`$ field needed to perform the marginal deformation on the field theory, from $`q=1`$ to a given value of $`q`$.
Notice that at the singularities we have the allowed degeneration into fractional branes from the splitting of $`(m,n)n(m/n,1)`$. Thus the T-dual fibration has singularities of the $`A_{n1}`$ type. As the fractional branes can be connected to each other in moduli space, we get a circle of such singularities and the monodromies around that circle are exactly the ones associated to orbifolds with discrete torsion.
Thus we have both the fractional $`B`$-field on the T-dual torus, and the monodromies of the singularities so that we can identify the T-dual geometry as the orbifold with discrete torsion.
As we have a T-duality description of the relation between the two compactifications, if the Kähler form is generically large in one setup, it is small in the other. It is therefore necessary to understand which description can be accounted for by supergravity calculations at a given point.
This question can be answered in $`AdS_5\times S^5`$. If we want $`n`$ $`D`$-branes to become one of these 2-tori, then $`q^n1`$, and as we saw before $`n1/H`$. The calculation of the $`D`$-brane action was done in string units, thus the $`D`$-branes are generically of a size commensurate with the string scale.
When we go to the supergravity regime, the string length is related to the supergravity background by the relation
$$l_s\frac{1}{\sqrt[4]{gN}}l_p.$$
(138)
In the configurations that we have discussed, we have $`nrR_5\sqrt[4]{gN}`$ in string units. Now, in order for $`\alpha ^{}`$-corrections to be small, we must have $`rH\mathrm{}_s`$, which implies
$$n\sqrt[4]{gN}$$
(139)
Thus if we want $`n`$ relatively small, $`AdS_5\times S^5`$ is a poor description of the geometry unless the total number of branes $`N`$ is such that $`\sqrt[4]{gN}<<n`$.
Similarly, for the $`S^5/\mathrm{\Gamma }`$ to be large, we need a very large number $`N`$ of $`D`$ branes. Indeed, the size of $`S^5/\mathrm{\Gamma }`$ is of order
$$rl_p/n$$
(140)
For the supergravity approximation to be valid here, we should require that twisted sector states are massive; this is the condition $`l_s<<l_p/n`$. As a result, the crossover region is at the same place, $`n\sqrt[4]{gN}`$. Thus, in the orbifold frame, we need to have $`N`$ large enough so that $`l_p>nl_s`$, whereas for the deformed 5-sphere, we needed $`N`$ small enough so that $`nl_s>l_p`$.
For a general $`q`$-deformation which is not a root of unity, no supergravity description will be good, and one is forced to take into account all of the stringy corrections to the supegravity equations of motion in order to determine the background.
It is also clear that the strength of the perturbation in string units needed to change the value of $`q`$ at the boundary is related to the number of branes in the configuration. Thus the limit is not uniform in supergravity. In this sense, it is hard to separate vevs from expectation values, as the supergravity boundary conditions are changed drastically when we change the number of branes.
## 6 Closed Strings and K-theory
So far we have described features of the moduli space of vacua for point-like (in the sense of non-commutative geometry) $`D`$-branes in deformed geometries. We want now to present a more complete picture of the field theory. This will have two aspects. First, we discuss the chiral ring of the field theory which has a clear interpretation in terms of the supergravity background and we give an interpretation in terms of the algebra itself. The second point that we wish to address is some of the features of extended branes which are accessible by topological considerations. In particular, this involves a somewhat more detailed understanding of K-theory and of discrete anomalies.
### 6.1 Closed Strings for Near-Horizon Geometry
Next, we will use ideas from the geometry/field theory correspondence to describe the physics of closed strings from the field theory point of view. This closed string theory is to be thought of as the dual string theory to the field theory of some $`D`$-branes near a singularity. Our aim is to understand the open string – closed string duality a little better, and how one might expect to realize it in the field theory. We have dealt with four-dimensional field theories so far in the classical regime. Our purpose is now to extract a closed string theory out of the quantum dynamics of the field theory.
The near-horizon geometry will have certain boundary conditions which control the superpotential, and some additional set of boundary conditions which specify the vacuum. That is, there are two contributions to the boundary conditions: those that decay sufficiently fast are related to the moduli of the branes, and those that decay more slowly are related to changes in the superpotential. To fully specify the field theory, we need in addition the correlation functions of operators. First, though, we need an identification of those operators.
We take the closed string states to be single trace operators in the field theory. This is in accordance with the AdS/CFT correspondence in that closed string states are gauge invariant operators in the field theory. The idea is to restrict ourselves now to the chiral ring of the field theory for simplicity, and because in all of our analysis we have kept only the parts which are protected by supersymmetry.
Let us assume first that we have a conformal field theory, and that its associated algebra is semi-classical (e.g., orbifolds with discrete torsion). We will exploit the following idea: the vevs of the closed string states (corresponding to states that decay quickly enough at the $`AdS`$ boundary) are generated by the stack of $`D`$-branes being at different locations in the moduli space. With the asymptotic values one reconstructs the near-horizon geometry of a set of parallel $`D`$-branes by summing over holes with given boundary conditions. Thus we can identify different tadpoles of the string states by motion in the moduli space of vacua. The right question to ask is what region of moduli space gives a vev to an operator.
We will combine this knowledge with the identification of the chiral ring for some geometries. Let us review a few results from Ref. . In that paper it was noticed that for orbifolds with discrete torsion, one could see the twisted and untwisted string states in the near-horizon geometry as coming from traces of different chiral operators. We will review the case of the orbifold $`^3/_n\times _n`$ with maximal discrete torsion.
Chiral operators come in two types
$$𝒪(k_1,k_2,k_3)=\mathrm{tr}(\varphi _1^{k_1}\varphi _2^{k_2}\varphi _3^{k_3})$$
(141)
with $`k_1=k_2=k_3mod(n)`$, which are untwisted states, and
$$𝒪_j(k)=\mathrm{tr}\varphi _j^k$$
(142)
which are twisted states so long as $`k0mod(n)`$.
The constraint on the $`k_j`$ for untwisted states is familiar from eq. (41). That is, the center of the algebra is associated with the untwisted states. This shows why the center of the algebra is so important to understand the geometry. Namely, the algebraic geometry of the center of the algebra is the geometry that the closed string sector sees. Here again we see that the geometry of the closed strings is commutative, as in Ref. . The non-commutativity of the moduli space appears from the closed string theory point of view because we have twisted sectors.
Notice that in (46) it is clear that it is the fractional branes which give vevs to the twisted sector strings. This is just as it should be, as we always think of coupling twisted sector strings to fractional branes living at the singularities of the classical space. Although we have discussed chiral operators here, it is more generally possible to distinguish twisted and untwisted states. As well, the same statement may be made if we do not have a conformal field theory.
In the case of $`AdS_5\times S^5`$, the $`F`$-term and $`D`$-term constraints give us a commutative geometry. Thus the center is the whole algebra, and every closed string state is untwisted and lives in the bulk.
We saw in Section 4 that the behavior under mass deformations was special for $`q=\pm 1`$. From our analysis, we can now see why this is the case. Namely, for $`q=\pm 1`$, the mass perturbation is untwisted, and therefore affects the bulk of moduli space. For any other rational $`q`$, the mass perturbation is twisted, and we expect that it will only affect the vicinity of a singularity.
### 6.2 Chiral ring revisited and Quantum Groups
Let us analyze the chiral ring in more detail. We have already learned that twisted and untwisted states are associated with traces of central (non-central) elements of the algebra, respectively.
States in the chiral ring are made by taking traces of holomorphic elements of the algebra. There are two steps for this construction. First we need a description of the elements of the algebra, and then we need to interpret the properties of the trace.
Any operator (for the deformations we have studied) can always be written in monomial ordered form for a small enough deformation, as we shown in Section 4.6. The difference between two possible orderings is given by $`F`$-terms and therefore they correspond to derivatives of other fields. In a conformal theory, these would be descendants and not primaries. For the topological chiral ring, we set all $`F`$-terms to zero, so the operators are identified as traces of elements of the algebra.
Let us consider the case where we have a conformal field theory in the ultraviolet. Because we have an algebra described by quadratic relations, we have a quantum hyperplane geometry. The operators with the same conformal dimension are homogeneous. On every quadratic algebra of the type described, there is an associated quantum group acting on the algebra. The states of the same degree are associated to the representations of this quantum group. This suggests that there might be a relation between operators in the closed string theory and representations of the quantum group. If this is indeed the case, then the fusion rules of the closed string operators will be related to the fusion rules of the representations of the quantum group algebra. This relation would give testable predictions for $`3`$-point functions in the deformed $`AdS_5\times S^5`$ supergravity. Quantum groups have also made an appearance in near-horizon geometry in the work of Ref. in connection with the stringy exclusion principle.
We do have to remember that we associate an operator to an element of the algebra, and that it is not the element of the algebra itself which is the gauge invariant operator. The association is by taking
$$𝒪(a)=\mathrm{tr}(a)$$
(143)
Because of the cyclic property of the trace, we need the following rule
$$𝒪(ab)=\mathrm{tr}(ab)=\mathrm{tr}(ba)=𝒪(ba)$$
(144)
thus the map from the algebra to the operators factors through
$$𝒜𝒜/[𝒜,𝒜]$$
(145)
as a vector space. It is the class $`[a]`$ in $`𝒜/[𝒜,𝒜]`$ that matters, and not $`a`$ itself.
The space
$$𝒜/[𝒜,𝒜]=HH_0(𝒜)$$
(146)
is actually a homology group of the Hochschild complex and suggests that the chiral ring is in general a cohomology group of the non-commutative space (so long as we have some sort of Poincaré duality). Because of our knowledge of Calabi-Yau manifolds, we can think of the chiral ring as a ring of deformations of a non-commutative complex structure, because we have found a relation with homology. Indeed, for a non-compact orbifold space the ring of deformations of the complex structure is infinite-dimensional because of the non-compactness, and it is associated to a cohomology group of the manifold $`H^{2,1}(M)`$. This suggests that orbifolds with discrete torsion may be better understood as a non-commutative Calabi-Yau space.
### 6.3 K-theory
Let us now make a few remarks about K-theory. To this effect we will review some of the results of Section 4.
Let us analyze the results of the $`q`$-deformation for rational $`q`$. There we found two types of finite dimensional representations: the representation of a non-commutative point associated to the bulk and some other representations which correspond to fractional branes at a singularity.
The set of non-singular points are all connected, and thus each point defines the same K-theory class. On going to the singularity, the points would split as
$$limR_{reg}=_iR_{sing}^i$$
(147)
where the subscript indicates that the point belongs to the regular part of the variety, or the singular part.
It so happens that the $`R^i`$ are homotopic to each other. That is, they can be deformed continuously into each other. If $`q^n=1`$, then there are $`n`$ representations on the right hand side of (147). In K-theory we thus have
$$K(R_{reg})=nK(R_{sing})$$
(148)
and the K-theory of points is generated by the K-theory class of a single singular point. Thus $`K_0^p(𝒜)=`$.
If we add one mass deformation, we find two coordinate patches that cover all of the variety except for a single complex line. This complex line is the complex line of singularities that was resolved by the deformation. One can also find solutions that cover this line of singularities. One still has two complex lines of singularities meeting at the origin, and the K-theory of points is still $``$. In both of these cases the degree of a point is enough to determine its K-theory class.
For the other rings, we find different phenomena. There are isolated points which correspond to fractional branes which cannot be connected to other singular points. For a rank three mass deformation and $`q=\pm 1`$ or $`q`$ not a root of unity, the moduli space is completely destroyed and the number of finite dimensional irreducible representations of the algbera is infinite. These are examples of rings which are not semi-classical, and in these cases the K-theory of points consists of an infinite number of copies of $``$, one for each irreducible representation. In the other cases, for $`q`$ a root of unity, the number of isolated points is finite.
The reason why the K-theory of points is not preserved under the deformations of the algebra relies on the fact that this is the K-theory appropriate to algebraic geometry, and not real geometry. This stems from the fact that we are restricting ourselves to the moduli spce of vacua, and we are forbidding transitions that go between the different components in moduli space. This is only appropriate if we are studying BPS objects, so this K-theory would serve to count BPS states, and not brane charges.
The full K-theory that we would need to understand brane-charge properly requires the inclusion of anti-holomorphic data and is less refined. This new K-theory would be the algebraic K-theory of the $`^{}`$ algebra associated to the string compactification. That is, the holomorphic K-theory construction gives too many K-theory classes, and does not give classes for the objects which cannot be represented in the holomorphic setting (e.g., odd dimensional $`D`$-branes).
The second statement that we want to make in K-theory has to do with extended classes. Indeed, based on discrete anomalies, the orbifolds with discrete torsion have a different K-theory than the commutative one associated to the ordinary orbifold. Our K-theory of points reproduces this result. We can also see the anomaly for extended objects.
Consider the orbifold with discrete torsion $`^3/_n\times _n`$, and consider trying to wrap a brane along the singular complex line $`x=y=0`$. This is a holomorphic subspace of the manifold. For the brane to cover the complex line, we need to have a lift to the non-commutative geometry, but the non-commutative geometry covers the singular complex line by an n-fold cover. Thus if we write a brane solution which would cover the singular complex line only once (which corresponds to a sheaf of rank $`1`$), this solution would correspond to a fractional brane. The lifting of this solution is obstructed, because if one lifts a point and does an analytic continuation, the brane would be broken in the non-commutative space. Indeed, we need a sheaf of rank $`1`$ in the non-commutative sense, and this is a sheaf of commutative rank $`n`$. Thus the brane charge is quantized in units of $`n`$ larger than in the standard orbifold, just as expected from the discrete anomaly. We believe that the non-commutative analysis makes the calculation of the anomaly more transparent.
As a final point, note that in principle, we have defined a K-theory that is capable of extending to $`B^{NS}0`$. As seen from the AdS/CFT, the deformations corresponding to the superpotentials are obtained by addition of antisymmetric tensors to the background. Our results suggest that the K-theory necessary to study these background is the algebraic K-theory of a non-commutative algebra.
## 7 Conclusions
In this paper, we have studied relevant and marginal deformations of the $`N=4`$ SYM theory from a non-commutative algebraic point of view. The moduli space looks like a symmetric product of a non-commutative geometry. This is interesting because it implies that $`D`$-branes may be considered as independent to a certain extent in the weakly coupled regime. This symmetric space captures well the phenomena of $`D`$-brane fractionation at singularities. Our approach has led us to the beginnings of a new definition of non-commutative algebraic geometry, which is still under investigation. The center of the algebra plays an important role in this construction, and, indeed, in a string theory picture the commutative subalgebra is related to closed strings, while the full non-commutative algebra is needed for open strings. When studied from the AdS/CFT point of view, the field theories that we studied present new dualities between distinct near-horizon geometries. These dualities are realized by T-duality of a 2-torus fibration of the 5-sphere. Different choices of T-duality lead to different dual near-horizon geometries. These results imply that AdS/CFT is inherently a stringy phenomenon, as they exhibit T-dualities which are not symmetries of classical supergravities. In order to understand this duality, we have constructed the $`D`$-brane configurations which realize the moduli space. We have found that the point-like $`D3`$-branes of the $`AdS_5\times S^5`$ become non-commutative 5-branes wrapping the torus fibration.
The non-commutative geometric framework suggests a natural formulation of K-theory appropriate to holomorphic data, and this is successful in reproducing the physics of discrete anomalies. This suggests that in general backgrounds the K-theory appropriate to $`D`$-brane charge is that derived from non-commutative algebra.
Our work suggests several avenues for future research. In particular, it would be of interest to understand the general problem of classifying what we have termed semi-classical algebras. A thorough understanding of this problem should provide new backgrounds in which $`D`$-branes can propagate, and would shed light on the existence of other dualities in near-horizon geometries.
In more generality, one should study the full problem of non-commutative algebraic geometry, including global questions. With a precise notion of gluing and compactness for example, we could entertain the idea of non-commutative Calabi-Yaus and their stringy geometry.
There are also interesting questions concerning non-perturbative effects, which we would need to understand $`S`$-dualities for example. We also must be concerned about the possibility of non-perturbative effects modifying our results, through, for example, the appearance of multi-trace operators in the superpotential.
For relevant deformations, there will be renormalization group flows which are reflected in near-horizon geometry. It would be of interest to construct these flows for the examples that we have studied, particularly since one expects that stringy corrections become important in the infrared. The study of correlation functions should also be of interest, with a possible connection to the representation theory of quantum groups.
Generalizations of our work to more complicated quivers is possible and will be explored elsewhere.
Acknowledgments: We wish to thank M. Ando, D. Grayson, A. Hashimoto, A. Jevicki, D. Kutasov, J. Maldacena, A. Strominger and C. Vafa for discussions. We are especially indebted to M. Strassler for a critical reading of this paper. DB thanks Harvard University for hospitality. Work supported in part by U.S. Department of Energy, grant DE-FG02-91ER40677 and an Outstanding Junior Investigator Award.
|
warning/0005/astro-ph0005551.html
|
ar5iv
|
text
|
# Towards global phase-delay VLBI astrometry: Observations of QSO 1150+812 and BL 1803+784
## 1 Introduction
The establishment of a radio reference frame with submilliarcsecond accuracy has been a major goal of astrometrists for the last several decades. Centimeter-wavelength very-long-baseline interferometry (VLBI) group-delay astrometry of extragalactic radio sources routinely provides precisions at the milliarcsecond (mas) level, thus allowing a celestial reference frame to be built with corresponding accuracy (e.g., Ma et al. ma98 (1998)). The use of phase-delay difference astrometry (see, e.g., Shapiro et al. shapiro (1979)) allows us to determine angular separations between pairs of radio sources at the submilliarcsecond level. One important advantage of phase-delay astrometry is that we can identify with sufficient accuracy suitable reference points within the structures of the radio sources whose relative positions we wish to determine, whereas in group-delay astrometry suitable reference points are not so easily identified from epoch to epoch with the desired accuracy.
With the phase-delay technique, the differences between phase delays for two radio sources are used to determine their relative positions, because these differences have reduced sensitivity to unmodeled effects. For sources nearby to one another on the sky, this technique yields statistical standard errors of only a few microarcseconds ($`\mu `$as), as in the case of quasars 1038+528 A and B (Marcaide & Shapiro jmm83 (1983); Marcaide et al. jmm94 (1994)), whose components are separated by only 33$`\stackrel{}{.}`$ However, the overall standard errors are dominated by inaccuracies in the reference-point identification (e.g., Rioja et al. rioja97 (1997)). For sources with larger separations, the main contributions to the astrometric standard error in the relative position of two sources comes from uncertainties in the coordinates of the source chosen as the reference, in the value used for UT1–UTC, and in the effects of the propagation medium. For radio sources with separations ranging from $`0.5\mathrm{°}`$ to $`7\mathrm{°}`$ on the sky, standard errors in relative positions of about 0.1–0.3 mas have been obtained (Bartel et al. bartel86 (1986); Guirado et al. 1995a ,b, gui98 (1998); Lara et al. lara96 (1996); Ros et al. ros99 (1999)).
A limiting factor in centimeter-wavelength VLBI astrometry is the uncertainty in the contribution of the ionosphere to the astrometric observables, even though its dispersive character makes this contribution scale as $`\nu ^2`$ and, in principle, allows it to be determined accurately (see, e.g., Thompson et al. tms (1986)). One strategy to take advantage of this scaling is to make simultaneous VLBI measurements in two frequency bands. The main disadvantage of this option is that the (fixed) bandwidth of the recording equipment has to be split between two frequency bands, decreasing the signal-to-noise-ratio (SNR) for each. Alternatively, one may compute corrections based on estimates of the ionosphere total electron content (TEC) obtained independently from Faraday-rotation measurements and, more recently, from the Global Positioning System (GPS). The advantage of the latter approach (e.g., Guirado et al. 1995b , Ros et al. ros2000 (2000)), is that only single-band VLBI observations are needed, avoiding the loss of sensitivity mentioned above.
In this paper, we study the applicability of the phase-difference technique to the strong radio sources 1150+812 and 1803+784, separated on the sky by almost 15°. We show that reliable phase connection is feasible at such a large angular separation, and estimate the relative position of the two sources with submilliarcsecond accuracy. We compare the estimates of the relative angular separation that result from use of two different methods of removing the ionosphere contribution based on two different types of data, namely, phase delays from dual-frequency-band VLBI measurements and TEC values from GPS measurements. The estimates from the two methods agree to within the standard errors from each method, showing that single-frequency astrometric VLBI experiments can be confidently carried out.
In Sect. 2, we briefly describe the observations; in Sect. 3 we describe the radio sources, and in Sect. 4 the data reduction process. In Sect. 5 we discuss our estimate of the relative position of the source 1803+784 with respect to 1150+812, and in Sect. 6 we carry out a sensitivity analysis of the estimated position of 1803+784 to errors in other model parameters. Finally, in Sect. 7 we summarize our main results and discuss their implications.
## 2 Observations
We observed the radio sources 1150+812 and 1803+784 on 18 November 1993 for 12 hours, in right circular polarization, simultaneously recording at two frequency bands (X-band $``$ 8.4 GHz and S-band $``$ 2.3 GHz). We used the following radio telescopes (in parenthesis we give the site symbol used in this paper, the diameter, and the location of the telecope): Effelsberg (B, 100 m, Germany); Medicina (L, 32 m, Italy); Onsala (S, 20 m, Sweden); Fort Davis (F, 25 m, Texas); Hancock (H, 25 m, New Hampshire); North Liberty (I, 25 m, Iowa); Owens Valley (O, 25 m, California); Los Alamos (X, 25 m, New Mexico); and the phased VLA (Y, 130 m equivalent, New Mexico). The European antennas (B, L, and S) recorded in Mark III mode A, covering a total bandwidth of 56 MHz, via seven adjacent channels spanning 8,403 to 8,431 MHz and seven such channels spanning 2,273 to 2,301 MHz. The American antennas used the Mark III recording system in mode B, covering a total bandwidth of 28 MHz, with four contiguous channels spanning 8,403 to 8,419 MHz and three such channels spanning 2,289 to 2,301 MHz. We used an observing cycle consisting of 2 min observing 1150+812, 1.5 min antenna slewing, 2 min observing 1803+784, and 1.5 min antenna slewing back to 1150+812, making a total cycle duration of 7 min.
The data were correlated at the MPIfR correlator in Bonn, Germany. Since both sources were strong ($`>`$ 1 Jy) at the epoch of observation, we detected each of them with high SNR within the 2 min integration time, at both frequencies and for each baseline, except those baselines involving the Onsala antenna. After correlation, we exported the VLBI observables (fringe amplitude and phase, group delay, and phase-delay rate) for the reference frequencies $`\nu _\mathrm{X}`$ = 8413 MHz and $`\nu _\mathrm{S}`$ = 2295 MHz, and calibrated the fringe amplitudes using the information provided by the staffs at the observing antennas. In our astrometric analysis, we discarded the data from Onsala due to the low SNR of the detections with this antenna; and from the phased VLA, since this antenna recorded only at 8.4 GHz and was used solely to increase the SNR of our hybrid maps. In addition, we discarded data from Effelsberg for the interval 18:30 until 23:28, since this antenna, too, only observed at 8.4 GHz during that interval. To obtain our hybrid maps, we self-calibrated, Fourier inverted, and CLEANed the visibility data using the Caltech program DIFMAP (Shepherd et al. shepherd (1995)).
## 3 The Radio Sources
The radio source 1803+784 has been optically identified with a BL Lac object at z=0.684 (Stickel et al. stickel91 (1991)). This source has a nearly flat spectrum from millimeter to decimeter wavelengths (Strom & Biermann strom91 (1991)) and displays many interesting properties (Biermann et al. biermann81 (1981)), including rapid variability, compactness, and strong X-ray emission.
Our 8.4 and 2.3 GHz maps of 1803+784 (Fig. 1) show a jet pointing westward. The hybrid map at 8.4 GHz shows the jet extending about 3 mas, with three components, labeled X1, X2, and X3. Most of the emission comes from components X1, identified as the core, and X2, about 1.5 mas west of X1. Component X3 lies about 2.8 mas west of X1 and is much fainter than X1 and X2. Our map at 2.3 GHz also shows the jet, which, as in the maps of Ros et al. (ros99 (1999)) and Eckart et al. (eckart87 (1987)), extends up to about 25 mas and is slightly bent towards the southwest. At 2.3 GHz, most of the emission is concentrated within 3 mas of the peak, but emission is also discernible at points about 5, 8, and 10 mas west of S1 (Fig. 1, middle). These latter bright regions, especially the innermost and the outermost, may correspond to real features of the jet, but like the questionable features farther out, they are not relevant for our astrometric purposes and, therefore, are not labeled.
The radio source 1150+812 is a QSO with redshift $`z=1.25`$ (Hewitt & Burbidge hb93 (1993)). For this source, our 8.4 and 2.3 GHz maps (Fig. 2) indicate a curved jet directed to the south and west. The 8.4 GHz map contains at least four components, labeled X1 through X4, from north to south. Most of the emission comes from the components X1 and X2, which are separated by $`1.5`$ mas. Extended emission is still discernible at a distance of $`4`$ mas from X1 (components X3 and X4). At 2.3 GHz, the emission is concentrated within the innermost 5 mas, but emission is also evident farther to the southeast.
Our maps of 1803+784 and 1150+812 are in good agreement with maps of similar resolution from 1995 published by Fey & Charlot (fc97 (1997)).
## 4 Astrometric Data Reduction
Our strategy in the astrometric data analysis was to obtain a set of “connected” phase delays that could be analyzed via weighted-least-squares to estimate the relative position of the two sources.
The observables in our experiment, after correlating and fringe-fitting the data, were the interferometric phase, $`\varphi `$, the group delay, $`\tau _G`$, and the phase-delay rate, $`\dot{\tau }_\varphi `$. We based our analysis primarily on $`\varphi `$, converted first to the phase delay, $`\tau _\varphi `$, via removal of the “$`2\pi n`$” ambiguities (where $`n`$ is an integer). The delay equivalent of $`2\pi `$ in phase is $`120`$ ps at 8.4 GHz and $`440`$ ps at 2.3 GHz. To remove these ambiguities we used a “phase connection” technique (see, e.g., Shapiro et al. shapiro (1979)) that can be outlined as follows: We constructed a model of the phase delays, based on a model of the geometry of our interferometric array and the radio sources (Table 1); the propagation medium; and the clock behavior at each station relative to a reference. We estimated the parameters of this model via weighted-least-squares analysis of the observed phase-delay rates and group delays. We used an improved version of the program VLBI3 (Robertson robertson (1975)) to carry out this weighted-least-squares analysis. The resultant model of the phase delays was accurate enough to allow phase connection, i.e., the elimination of the $`2\pi n`$ ambiguities for almost all of the observations (see below), via a suitable iterative scheme that took advantage of the closure relations over triangles of baselines. Since we phase-connected the data independently for each source, we used weighted-least-squares to verify – successfully – the consistency of the “overall resolution of ambiguities” of the delays for the two sources (see Ros et al. ros99 (1999) for a detailed discussion of this step). During our analysis, we discarded all data from the Hancock antenna, since our phase connection failed at both frequencies. We also discarded the data from the Medicina antenna from 18:30 to 23:48, since in this time period our phase-connection appeared to be unreliable. (These failures are likely the result of some combination of tropospheric, ionospheric, and instrumental effects.)
### 4.1 Source-Structure Contribution
The extended structure of even a “compact” radio source often makes a significant contribution to the phase delays. Such a contribution depends on the point chosen as a reference on the map. Identifying a reference point in a reliably epoch-independent manner is crucial for the use of our method to compare positions obtained from different epochs. For each source and frequency band, our procedure was to choose the peak of brightness (POB) as the reference point. To find the POB, we constructed fine-grained maps of the two radio sources using pixel size 0.01 mas at 8.4 GHz and 0.03 mas at 2.3 GHz. We identified the brightest pixel in each such map as the POB and defined a new coordinate system with that reference point as its origin. We then computed from these maps the structure-phase contribution to each phase delay. These contributions – up to $``$ 25 ps at 8.4 GHz and $``$ 55 ps at 2.3 GHz – were removed from the phase delays, to effect a point-like source at each POB. The standard deviations assigned to each POB location included uncertainties due to (a) the pixel size used for each map, which was about 0.01 (0.03) mas at 8.4 (2.3) GHz, for each source; and (b) SNR for the peak, which was about 0.04 (0.27) mas at 8.4 (2.3) GHz, for each source. Note, however, that when the delays at the two frequencies are combined to form “plasma-free” delays, errors in the phase delays at 2.3 GHz are scaled down by a factor $`[(\nu _X/\nu _S)^21]^10.08`$ (see Sect. 4.4), and therefore the uncertainty of the POB at 2.3 GHz is less significant than that of the POB at 8.4 GHz. The root-sum-squares of these contributions, at each frequency, are indicated as the first (8.4 GHz) and second (2.3 GHz) entries in Table 2, and are dominated by the relatively low SNR of the hybrid maps.
### 4.2 Neutral Atmosphere Contribution
The neutral atmosphere (primarily the troposphere) adds an extra delay to the incoming radio waves, the equivalent of up to a few meters in pathlength. We monitored the pressure, relative humidity, and temperature at each observing antenna to track the atmosphere behavior during the observations. We used the model by Saastamoinen saasta (1993), in which the atmosphere is separated into two components: a dry component and a wet component (due to the water vapor in the atmosphere). Based on this model, we calculated a priori values for the delays for the wet and dry atmosphere components in the zenith direction for each antenna site (Table 1), and then used the Chao (chao (1974)) mapping function to determine delays at other elevations. This mapping function agrees to within about 1 cm with ray-tracing computations (Davis et al. davis85 (1985)) for antenna elevations larger than 20; all of our observations had elevation angles above this limit. To specify adjustments to the combined (dry and wet) atmosphere delays during the observations, we used a piecewise-linear function characterized by zenith delay values at epochs about three hours apart. Errors in the combined neutral-atmosphere delays – mainly due to the wet troposphere – are likely to be up to 0.1 ns/hr (D. Lebach, priv. comm.). Since the wet atmosphere zenith delay fluctuates approximately as a random walk (Treuhaft & Lanyi wet (1987)), an error of 0.1 ns/hr transforms into a standard error of $`\sqrt{3}\times 0.10.17`$ ns every three hours. Therefore, in our sensitivity analysis, we allowed the atmosphere parameters at each antenna location to vary with this standard error.
### 4.3 Ionosphere Contribution
As explained in Sect. 1, we removed most of the ionosphere contribution in two alternative ways, first by using our dual-frequency-band observations and, second, by using the total electron content (TEC) values deduced from GPS measurements (see, e.g., Sardón et al. sardon (1994)). For the latter removal, we followed the same procedure as described in Ros et al. (2000), modeling the ionosphere as a thin layer located at an altitude of 350 km. We used GPS data from Wettzell (Germany) to obtain TEC values for Effelsberg and Medicina, and from Goldstone (California) to obtain TEC values for Owens Valley, Los Alamos, and Fort Davis; we were fortunate to have a GPS antenna collocated with our VLBI antenna at North Liberty to obtain TEC values for this site. In Fig. 3 we compare the estimated ionosphere delays as deduced from both dual-frequency-band observations (8.4 GHz/2.3 GHz delays) and GPS-based measurements (TEC delays), for a representative subset of baselines, intracontinental (upper two plots) and intercontinental (lower two plots). Dual-frequency corrections refer only to baselines, not individual antennas. On the other hand, TEC ionosphere corrections are calculated for individual antennas. Therefore, for each antenna there can be a constant offset between the TEC ionosphere correction and the dual-frequency one as a result of instrumental effects, but, for clarity, we subtracted the mean difference between each of the two corrections (Effelsberg, 1.35 ns; Medicina, 0.55 ns; North Liberty, 7.80 ns; Fort Davis and Los Alamos had no offset). We assumed a statistical standard error for the vertical TEC at each GPS antenna of $`2\times 10^{16}\mathrm{m}^2`$ (Sardón, priv. comm.), corresponding to a standard error of $`40`$ ps for the (vertical) phase delays at 8.4 GHz. For each antenna, this error is multiplied by the value from the mapping function for the elevation angle of each radio source (here, the secant of the zenith angle at the ionospheric point) of each radio source. The resultant mapped TEC-based ionosphere corrections had standard errors ranging from $`70`$ ps to $`120`$ ps, highly correlated from point to point (Fig. 3). The standard errors for our dual-frequency-band corrections are the appropriate combination of the statistical errors of the phase delays at each frequency, and ranged from $`5`$ ps for the Effelsberg–Medicina baseline up to $`30`$ ps for the Fort Davis–Owens Valley baseline. The root-mean-square (rms) of the differences between the corrections obtained by the two methods ranged from $`20`$ ps up to $`90`$ ps, depending on baseline. The maximum difference was for intercontinental baselines and was $`110`$ ps, when it was night in Europe and noon in North America. The level of agreement in the results from the two methods is gratifying, and foreshadows the agreement in the astrometric results obtained below with the two methods.
### 4.4 Opacity Contribution
When we use phase delays at both 8.4 and 2.3 GHz (see Sect. 4.3), the reference points on the 2.3 GHz maps should be chosen to correspond to the same points in the sky as those on the 8.4 GHz maps. Unfortunately, opacity effects may introduce an offset between the POB for one frequency band with respect to the POB for the other frequency band (see, e.g., Marcaide & Shapiro jmm84 (1984)). Since jet components are less self-absorbed than core components, their associated peaks should be less affected by opacity effects and should therefore be better suited for use as guides for the alignment of the maps at 8.4 and 2.3 GHz. We registered the peaks of the less self-absorbed components in our 8.4 GHz maps with their corresponding peaks in (twofold overresolved) 2.3 GHz maps.
For 1803+784, we identified component X3 at 8.4 GHz with component S3 at 2.3 GHz (Fig. 1) on a twofold overresolved map (not shown). Based on this map, we concluded that blended components S1, S2, and S3 at 2.3 GHz correspond to components X1, X2, and X3, respectively, in the 8.4 GHz map. This identification implied that in our maps the POB of 1803+784 at 2.3 GHz (S1) is offset by $``$ 0.6 mas westward from the POB at 8.4 GHz (X1). For 1150+812, we also used a twofold overresolved map at 2.3 GHz to make a plausible registration of the maps for the two bands. We identified components X3 and X4 at 8.4 GHz with components S3 and S4 at 2.3 GHz (Fig. 2), implying a shift of the POB at 2.3 GHz (S1) of $``$ 0.5 mas northward from the peak of the core in the 8.4 GHz map (X1). To estimate the standard error in our determination of those shifts we followed a similar procedure to that described for the determination of the standard error associated with the reference-point identification (POB), but applied to components X3 and S3, for each source. This procedure yielded a standard error of 0.4 mas for each source. Finally, to take into account the possible additional shifts due to opacity effects of the POB of component S3 with respect to X3, for each source, we increased the estimated standard error in each coordinate of the registration estimate from 0.4 mas to 0.7 mas. Estimates of the shift between the POBs at 8.4 and 2.3 GHz (where opacity effects are expected to be larger than in the jet components) of other radio sources – 1038+528 A (Marcaide & Shapiro jmm84 (1984)), 1226+023 (Charlot charlot93 (1993)), 1901+319 (Lara et al. lara94 (1994)) – are also no greater than 0.7 mas. Thus, our estimates of these shifts for 1803+784 and 1150+812 seem reasonable.
In this dual-frequency-band method to reduce the ionosphere contribution, the resultant “ionosphere-free” phase delays have the form:
$$\tau _{\varphi ,free}=\frac{R\tau _{\varphi ,1}\tau _{\varphi ,2}}{R1}$$
where the subscripts 1 and 2 refer to the two frequency bands, with $`\nu _1>\nu _2`$ and $`R=(\nu _1/\nu _2)^2`$. In our case, $`R(8.4/2.3)^213`$; hence standard errors in the phase delays at 8.4 GHz are scaled up by a factor $`R/(R1)1.08`$, whereas corresponding errors in the phase delays at 2.3 GHz are scaled down by $`1/(R1)0.08`$ in their effect on $`\tau _{\varphi ,\mathrm{free}}`$. Likewise, uncertainties accounting for registration errors due to opacity effects have to be scaled by $`1/(R1)`$. Thus, the contribution to the total error budget due to the use of dual-frequency-band measurements, in particular the seemingly large error at 2.3 GHz discussed above, is not significant (see entry “Opacity effects” in Table 2).
## 5 Relative Position of the Two Sources
We obtained a final set of phase delays for each source by correcting for source-structure, opacity-effect, and ionosphere contributions as described above. We then formed a set of differenced phase delays by subtracting the delay for each observation of 1803+784 from the corresponding delay for the previous observation of 1150+812. The use of differenced phase delays is in general more effective the closer together the sources are in the sky, since differencing results from neighboring observations tends to cancel effects which for each source alone cannot be accurately enough described by theoretical models. The best such pair of sources for such cancellation so far studied is 1038+528 A and B (Marcaide & Shapiro jmm83 (1983)), because this pair simultaneously lies inside the beam of each antenna, yielding almost complete cancellation of several sources of error. For sources separated by increasingly larger angular distances, the cancellation of unmodeled effects in general lessens, due directly to the increase in angular separation of the sources and indirectly to the increase of the cycle length of the observations.
From a weighted-least-squares analysis of the differenced phase delays, we estimated the coordinates of 1803+784 relative to those of 1150+812 (see Table 3). In this analysis, we also included the (undifferenced) phase delays of 1150+812, suitably weighted, to estimate the relative behavior of the site clocks. (The opposite scheme, i.e., using the differenced phase delays of 1803+784 minus those for 1150+812, and the undifferenced phase delays of 1803+784, did not alter our estimate significantly.) The total number of parameters whose values were estimated was 26, and include those of a third-degree polynomial to model the clock behavior at each station (except for North Liberty, which was taken as our reference station), the coordinates of 1803+784, and four atmosphere parameters (see Sect. 4.2) for the Effelsberg antenna, since the information on the atmosphere conditions provided for this site was very sparse, and thus of doubtful utility. Nonetheless, the estimated atmosphere parameters for Effelsberg agreed with the a priori values within 0.11 ns.
Our choice of 1150+812 as the reference source was motivated by its being one of the defining sources of the quasi-inertial International Celestial Reference Frame. As a test of the effect of this choice, we solved for the coordinates of 1150+812–keeping those of 1803+784 fixed at their IERS values–and obtained different results for $`\mathrm{\Delta }\alpha `$ and $`\mathrm{\Delta }\delta `$. As a simple check on the significance of those differences – one that avoids the effects of the reference system in which the coordinates are expressed – we calculated the arclength between the two sources for each of the two solutions. We expected to obtain nearly the same result for the two cases and, indeed, the difference was a mere 5 $`\mu \mathrm{as}`$, negligible compared to the overall standard errors of our results.
For our nominal, dual-frequency-band solution described above, we scaled separately the standard deviations of the differenced phase delays and of the phase delays of 1803+784, so that for the data for each type, and for each baseline, the weighted-mean-square of the postfit residuals was nearly unity. As a result, the error bars assigned to the differenced phase delays were, in general, smaller by a factor $``$ 1.2 with respect to the error bars assigned to the phase delays of 1150+812 (see Figs. 4 and 5). We studied the effect of varying the relative weights of our two data types and verified that such variations did not change our result significantly. In fact, any ratio of the weights of the differenced to the undifferenced data between 0.1 and 10 with respect to our “nominal” solution did not alter our estimate of the angular separation significantly. Figs. 4 and 5 also illustrate a virtue of the difference phase-delay technique mentioned above: although the separation of the sources is $`14^{}\mathrm{\hspace{0.17em}50}^{}\mathrm{\hspace{0.17em}21}\stackrel{}{.}147559`$ and the cycle length of the observations is correspondingly large, the differenced phase delays are partially free from some incompletely modeled effects as attested by the larger residuals seen in Fig. 4 than in Fig. 5, in particular for the baselines involving the Medicina antenna.
The standard errors in Table 3 are larger than the differences between our estimates of $`\mathrm{\Delta }\alpha `$ and $`\mathrm{\Delta }\delta `$, and those provided by IERS, indicating that the a priori IERS values are quite accurate. This accuracy is not surprising, since both radio sources are dominated by emission from the core, and so the IERS values – mainly based on averages of results from observations at many epochs using group-delay astrometry – should be accurate, except for a small offset that may be contributed by persistent opacity effects in the core components. For sources with strong emitting features well separated from their cores, larger differences from the a priori IERS values can be expected, given the greater sensitivity of the “structure group delay” to components farther from the cores (Thomas thomas80 (1980)). Those differences can be largely removed by using the more precise, structure-free phase delays.
## 6 Sensitivity Analysis
We carried out a sensitivity analysis to estimate the contributions of individual effects to the standard errors in the relative-position determination; see Table LABEL:tab:sens for a list of these effects. We found that the primary source of error is due to the uncertainty in our knowledge of the neutral atmosphere, which by itself contributes a standard error of about 0.4 mas in right ascension and about 0.6 mas in declination to the estimate of the position of 1803+784. To be conservative in our estimate of the error contributed by the neutral atmosphere, we took for each site the larger of the following two values: the root-sum-square of the delay contribution from all atmosphere parameters for that site (one every three hours), or the algebraic sum of the signed partials. In this way, we made some allowance for the potential effects of possible correlations among the estimates of the atmosphere parameters for each site. The second largest contribution to the total error budget of each coordinate comes from the corresponding standard error in the a priori position used for the reference source (1150+812). This large contribution is due to the geometry of the sources on the sky. For sources with smaller angular separations, the contribution of the a priori uncertainties of the coordinates of the reference source to the standard error in relative position is usually a small fraction of the a priori uncertainty of the reference coordinates (e.g., Guirado et al. 1995b ), because there is a positive correlation between the right ascensions, and similarly the declinations, of the two sources. However, because 1150+812 and 1803+784 are within $`12\mathrm{°}`$ of the north celestial pole and lie almost exactly six hours apart in right ascension, an error in the position of one of the sources in right ascension (or declination) translates into a similar error for the other source in declination (or right ascension). (See Table LABEL:tab:sens.) A covariance analysis confirms this “tradeoff”, showing a fairly strong correlation between the estimates of the right ascension of one source and the declination of the other source ($`\rho (\alpha _{1150+812},\delta _{1803+784})0.7;\rho (\delta _{1150+812},\alpha _{1803+784})0.7`$), but little correlation between the estimates of the two right ascensions (and between the two declinations). Moreover, the nearly complete cancellation of the geometric errors usually obtained in astrometric studies of radio sources with smaller angular separations is not so effective in the present case. For example, the uncertainties in UT1–UTC values do contribute significantly to the standard error in the estimate of the angular separation. The remaining significant contributions to the standard error are attributable to the uncertainties in the estimates of the antenna locations.
Software limitations prevented us from taking into account in our analysis such effects as ocean and atmosphere loading, tidal terms in polar motion and UT1, and arbitrary variations in the atmosphere zenith delays. Based upon a surrogate test of the sensitivity of the relative position of a different pair of radio sources (observed in an unrelated experiment) to these effects, performed using the CALC (Ma et al. ma86 (1986)) and SOLVK (Herring et al. herring90 (1990)) packages, we estimated that the known limitations of our software contribute an uncertainty of about 0.4 mas to each coordinate of our estimate of the sources’ relative position.
## 7 Summary
We observed the strong radio sources 1150+812 and 1803+784 with a VLBI array on 18 November 1993, an epoch of mild solar activity. The antennas recorded data simultaneously at 8.4 and 2.3 GHz, which allowed us to estimate the ionosphere contributions to the phase delays; we also used TEC values from GPS measurements to estimate such contributions. We estimated, and thereby partly removed from the phase-delay data, contributions due to the dry and wet components of the atmosphere, and the brightness structure of the sources. The phase-connection process, required to extract the precision inherent in the phase delays, did not pose special difficulties for our 7 min cycle time, indicating that spatial and temporal fluctuations due to the atmosphere and the ionosphere were not large enough to prevent reliable phase connection in either frequency band.
We then estimated the relative position of the sources (Table 3) via a weighted-least-squares analysis of a combined data set of undifferenced and differenced phase delays. The estimates resulting from use of GPS-based TEC values to account for ionosphere effects agree with those obtained from our dual-frequency-band VLBI measurements to within the standard errors for each method. Ros et al. (ros99 (1999)ros2000 (2000)) also successfully used GPS-based TEC values to remove the ionosphere’s contribution to the phase delays at 8.4 GHz. The checking of their removal in that experiment was limited to the precision of the dual-frequency-band (8.4 and 2.3 GHz) group delays, since Ros et al. could not reliably connect their phase delays at 2.3 GHz. In our case, we successfully connected the phase delays at both frequencies, and thus for the first time compared TEC ionosphere delays against phase-delay-based ionosphere corrections (Fig. 3). This agreement supports the use of single-band ($`8.4`$ GHz) VLBI observations for astrometry purposes, an option of particular interest for experiments to be carried out at epochs of strong solar activity, when phase connection for a frequency band $`2.3`$ GHz is likely to fail. Single-frequency VLBI observations along with GPS-derived TEC ionosphere corrections can be confidently used instead.
Our standard errors shown in Table 3 are larger than the difference between our estimates of relative source position and those provided by IERS. The high level of agreement is likely due to the fact that both radio sources are strongly core-dominated. However, were a bright emitting feature present several milliarcseconds from one of the cores, large offsets could exist between the true core position and the IERS position. With VLBI phase-delay astrometry, such errors can be effectively reduced using structure-phase corrections computed from self-calibrated maps derived from the same VLBI observations. Moreover, phase-delay astrometry has the advantage of allowing us to know very accurately to which feature inside each source we are referring when calculating relative positions.
In summary, we have determined with submilliarcsecond accuracy the angular separation of two radio sources separated by almost $`15\mathrm{°}`$, using phase delays from dual-frequency-band VLBI measurements. Since within $`15\mathrm{°}`$ of any given source there are almost always two or more reasonably strong radio sources, we have thus demonstrated that phase delays should be usable for full-sky astrometry of radio sources. We have also shown that GPS-based measurements can be used to obtain reliable ionosphere corrections to the phase delays, thus demonstrating the feasibility of conducting VLBI solely observations at frequencies $`8.4`$ GHz for astrometric purposes. Future improvements in the modeling of the atmospheric delay contribution, and in the knowledge of Earth rotation and pole position, as well as antenna location, should result in increasingly accurate estimates of the relative positions of sources far apart on the sky.
We stress that phase delays are more reliably corrected for structure phase than are group delays. Therefore, differenced phase-delay astrometry is better suited than group-delay astrometry for carrying out astrometric studies of extended radio sources. To date, however, group-delay astrometry has been used to establish a quasi-inertial celestial reference frame based on estimates of the positions of a number of relatively compact extragalactic radio sources from many years of regular observations. Our results open the avenue to an alternative, potentially more accurate, approach, namely that of carrying out a suitable series of observing sessions and using difference phase-delay astrometry to obtain submilliarcsecond positions for the cores of these sources.
###### Acknowledgements.
We are grateful to the referee, Jim Ulvestad, for valuable comments and suggestions. We thank the staffs of all the observatories for their contribution to the observations, and in particular the MPIfR staff for their efforts during the correlation. We also thank Dan Lebach for his work with the CALC/SOLVK package to test the quantitative significance of the limitations of our VLBI3 program. This work has been partially supported by the Spanish DGICYT grants PB93-0030 and PB96-0782, and by the European Comission’s TMR-LSF program under contract No. ERBFMGECT950012. The VLA is a facility of the National Radio Astronomy Observatory. The NRAO is a facility of the National Science Foundation operated under cooperative agreement by Associated Universities, Inc.
|
warning/0005/hep-th0005203.html
|
ar5iv
|
text
|
# Note on the dual BRST Symmetry in U(1) Gauge Theory
## I INTRODUCTION
Nowadays the concept of BRST symmetry plays an essential role in the quantization of gauge theories. As is well known the BRST formalism has been very useful in the framework of path integral quantization, where the BRST generator (charge) is a key ingredient of the effective action, as well as it finds interesting applications in the operator formulation of the theory. An illustrative example on this subject arises when one considers string theory. In fact, following the BRST inspired approach it was possible to derive the string critical dimensions in a straightforward and economical way. In this context it may be recalled that there exists two approaches to the BRST formulation of gauge theories. One is based on the Hamiltonian formulation, where the BRST charge is constructed in terms of the constraints and the higher order structure functions in a gauge independent way. The other approach is based on the Lagrangian formulation, in such a case the BRST charge is computed from a gauge-fixed Lagrangian by using Noether’s prescription. In passing we also recall that in the path integral quantization formalism (both Lagrangian and Hamiltonian) the original gauge invariance is incorporated by means of the extension of the phase space including ghost fields. Thus the main idea is to substitute the local gauge invariance by a rigid Grassmannian symmetry (or global supersymmetry) known as the BRST symmetry. In this way one assigns a global nilpotent charge to this symmetry whose cohomology produces the physical states.
On the other hand, recently a great deal of attention has been devoted to the study of new symmetries in gauge theories. For instance, Lavelle and McMullan found that QED displays a new nonlocal and noncovariant symmetry. In such a case the symmetry transformations are compatible with the gauge-fixing conditions. At the same time Tang and Finkelstein constructed a nonlocal but covariant symmetry for QED. Let us also mention here that Yang and Lee derived a noncovariant but nonlocal symmetry of QED. More recently, Malik showed that in two dimensions of spacetime there exists a local, covariant and nilpotent BRST symmetry, the so-called dual symmetry, under which the gauge-fixing term remains invariant for a free U(1) gauge theory and QED. Furthermore, this author claimed that this symmetry transformation is not the generalization of the above symmetries in two dimensions of spacetime. It is worth stressing at this stage that despite their relevance these studies have been, however, carried out in the gauge fixed scheme only.
Meanwhile, in a previous paper we have discussed the relation between the Lagrangian and Hamiltonian symmetry generators for the Lavelle and McMullan’s symmetry in a gauge invariant way using the Batalin-Fradkin-Vilkovisky formalism. In particular, we have showed that the Lavelle and McMullan’s symmetry may be derived from a canonical transformation in the ghost sector. We also recall that there are definite advantages of the Hamiltonian approach over the conventional gauge fixed analysis. Ambiguities related to gauge fixing conditions are avoided and it does not need an auxiliary field to construct an off-shell nilpotent symmetry transformation. Let us also mention here that a similar analysis has been made, independently, by Rivelles. We are thus motivated to investigate in this paper whether the so-called dual symmetry is a new symmetry or it is merely an artifact of the canonical transformation in the ghost sector.
The outline of this paper is as follows. In Sect.2 we briefly recap the BFV-BRST formalism for a free U(1) gauge theory in four dimensions of spacetime. This will form the basis of our subsequent considerations. In Sect.3 we will focus our attention to the two-dimensional case. Particular care is paid to establish a direct connection between the Lagrangian and Hamiltonian BRST symmetry generators for the so-called dual symmetry.
## II General Considerations on the BFV-BRST Formalism
Let us commence our considerations with a short presentation of the BFV-BRST formalism for a free U(1) gauge theory in four dimensions. It should be noted that this method is a general procedure for quantizing systems with first class constraints. A detailed discussion of the formalism can be found in . We summarize the essence of this formalism in terms of a finite number of phase-space variables, this makes the discussion simpler. In such a case the action for the theory under consideration is taken to be
$$S=𝑑t\left(p^\mu \stackrel{}{q_\mu }H_0\lambda ^a\phi _a\right),$$
(1)
where the coordinates $`(q_\mu ,p^\mu )`$ are the canonical variables describing the theory. The canonical Hamiltonian is $`H_0`$, and $`\lambda ^a`$ are the Lagrange multipliers associated with the first class constraints $`\phi _a`$. As prescribed by the general theory the Lagrange multipliers are treated in the same foot as the canonical variables, thus we introduce conjugate canonical momenta to $`\lambda ^a`$, say $`p_a`$. Evidently, the $`p_a`$’s must be imposed as new constraints in order that the dynamics of the theory does not change. Now, the BFV approach introduces a pair of canonically conjugate ghosts $`(C^a(x),𝒫^a(x))`$ for each constraints. The Poisson algebra of these ghosts is
$$[C(𝐱,t),𝒫(𝐲,t)]=\delta \left(𝐱𝐲\right),$$
(2)
where $`C`$ and $`𝒫`$ has ghost number $`1`$ and $`1`$, respectively. These considerations naturally lead to an extended phase space, where we have substituted the local gauge invariance by a global supersymmetry invariance (BRST invariance). In this extended phase space the generator of the BRST symmetry for a theory with first class constraints has the form
$$\mathrm{\Omega }=C_a\phi ^a+\frac{1}{2}P^af_a^{bc}C_bC_c+\mathrm{},$$
(3)
where $`f_a^{bc}`$ are the structure functions, and $`\mathrm{\Omega }`$ is by construction nilpotent ($`[\mathrm{\Omega },\mathrm{\Omega }]=0`$). We also recall that, at the quantum level, in the extended phase space there exists the Fradkin-Vilkovisky theorem . This theorem states that the functional integral
$$𝒵_\mathrm{\Psi }=𝒟\mu \mathrm{exp}\left(iS_{eff}\right),$$
(4)
where the effective action $`S_{eff}`$ is given by
$$S_{eff}=𝑑t\left(p^\mu \stackrel{}{q_\mu }+C^a\underset{a}{\overset{}{𝒫}}+p^a\underset{a}{\overset{}{\lambda }}H_0[\mathrm{\Omega },\mathrm{\Psi }]\right),$$
(5)
being independent of the choice of $`\mathrm{\Psi }`$. Here $`\mathrm{\Psi }`$ is an arbitrary fermionic gauge-fixing function, and $`𝒟\mu `$ is the Liouville measure on the phase space. This concludes our brief review of the BFV formalism.
Let us now proceed to apply the above procedure for a free U(1) gauge theory, in other words,
$$=\frac{1}{4}F_{\mu \nu }F^{\mu \nu }.$$
(6)
Then, from (1) the canonical action takes the form
$$S=𝑑x\left(\stackrel{}{A_i}\mathrm{\Pi }^iH_0\lambda \phi \right),$$
(7)
where $`\mathrm{\Pi }^i`$ is the momenta conjugate to $`A_i`$. $`H_0`$ is the canonical Hamiltonian, that is,
$$H_0=d^3x\left(\frac{1}{2}\mathrm{\Pi }_i\mathrm{\Pi }^i+\frac{1}{4}F_{ij}F^{ij}+\mathrm{\Pi }_i^iA_0\right),$$
(8)
and it is straightforward to see that the preservation in time of the constraint primary $`(\mathrm{\Pi }^0=0)`$ leads to the secondary constraint
$$\phi =_i\mathrm{\Pi }^i=0.$$
(9)
We mention in passing that in the action (7) the canonical variables $`A_0`$ and $`\mathrm{\Pi }^0`$ have been omitted because $`\mathrm{\Pi }^0=0`$, which does not represent a true dynamical degree of freedom of the theory. Thus, $`A_0`$ can be absorbed by redefining the multiplier $`\lambda `$, i. e., $`\lambda `$ and $`A_0`$ do not need to be treated as independent variables. With this at hand, the effective action then reads:
$$S_{eff}=d^4x\left(\mathrm{\Pi }^i\stackrel{}{A_i}+𝒫\stackrel{}{C}+\overline{𝒫}\stackrel{}{\overline{C}}+\mathrm{\Pi }^0\stackrel{}{A_0}H_0[\mathrm{\Omega },\mathrm{\Psi }]\right),$$
(10)
where we have introduced the antighost pair $`(\overline{C}(x),\overline{𝒫}(x))`$ with respective ghost numbers $`1`$ and $`+1`$, and satisfying the Poisson algebra (2). The BRST charge $`\mathrm{\Omega }`$ can be easily given as
$$\mathrm{\Omega }=d^3x\left(C_i\mathrm{\Pi }^ii\overline{𝒫}\mathrm{\Pi }_0\right).$$
(11)
We can now write the corresponding transformations generated by $`\mathrm{\Omega }`$, that is,
$$\delta A_i=\epsilon _iC,$$
(12)
$$\delta A_0=i\epsilon \overline{𝒫},$$
(13)
$$\delta \mathrm{\Pi }_i=0,$$
(14)
$$\delta \mathrm{\Pi }_0=0,$$
(15)
$$\delta C=0,$$
(16)
$$\delta \overline{C}=i\epsilon \mathrm{\Pi }_0,$$
(17)
$$\delta \overline{𝒫}=0,$$
(18)
$$\delta 𝒫=\epsilon _i\mathrm{\Pi }^i,$$
(19)
where $`\epsilon `$ is an anticommuting spacetime independent infinitesimal parameter. In order to compute the effective action (10), we have to select the gauge fixing function $`\mathrm{\Psi }`$. There are a variety of these which have been found useful and convenient in different calculational context. We can choose, for example, $`\mathrm{\Psi }`$ in the form
$$\mathrm{\Psi }=d^3x(𝒫A_0i\overline{C}(\frac{x_iA^i}{x^2}\frac{\xi }{2}\mathrm{\Pi }^0\underset{0}{\overset{}{A}})),$$
(20)
which leads to the modified Fock-Schwinger gauge . However, of this turn, we take $`\mathrm{\Psi }`$ as
$$\mathrm{\Psi }=d^3x\left(𝒫A_0i\overline{C}\left(_i\mathrm{\Pi }^i\frac{\xi }{2}\mathrm{\Pi }_0\right)\right),$$
(21)
where $`\xi `$ is a real parameter that describes a set of gauges. Explicitly, for $`\xi =0,1`$ and infinity we obtain the Landau, Feynman and unitary gauges, respectively. Plugging this expression into (10), we find that the resulting effective action is given by
$$S_{eff}=d^4x\left(\frac{1}{4}F^{\mu \nu }F_{\mu \nu }+i\overline{C}_\mu ^\mu C+\frac{1}{2\xi }\left(_\mu A^\mu \right)^2\right).$$
(22)
We immediately recognize the above to be the same as the Lagrangian effective action.
Before concluding this section we call attention to the fact that in contrast to the gauge-fixing term, the gauge field $`F_{\mu \nu }`$ remains invariant under the transformation generated by $`\mathrm{\Omega }`$ (11), that is, $`\delta F_{\mu \nu }=0`$ and $`\delta \left(_\mu A^\mu \right)=\epsilon \left(i\stackrel{}{\overline{𝒫}}^2C\right)`$. However, it is possible to recast the BRST charge (11) which corresponds to a nilpotent symmetry transformation under which the gauge-fixing term remains invariant. This can be done by a suitable canonical transformation in the BFV phase space, in such a way that any two BRST generators are related by such transformations . In the work of Ref., we had showed that by performing the following canonical transformation in the ghost sector:
$$C^{}=\frac{1}{^2}𝒫,$$
(23)
$$𝒫^{}=^2C,$$
(24)
$$\overline{C}^{}=\overline{𝒫},$$
(25)
$$\overline{𝒫}^{}=\overline{C},$$
(26)
the gauge-fixing term remains invariant. In effect, as a consequence of this canonical transformation, the new BRST charge $`\mathrm{\Omega }^{}`$ then becomes
$$\mathrm{\Omega }^{}=d^3x\left(\frac{1}{^2}𝒫_i\mathrm{\Pi }^i+i\overline{C}\mathrm{\Pi }_0\right).$$
(27)
Hence we see that the corresponding transformations generated by $`\mathrm{\Omega }^{}`$ are:
$$\delta ^{}A_i=\epsilon _i\frac{1}{^2}𝒫,$$
(28)
$$\delta ^{}A_0=i\epsilon \overline{C},$$
(29)
$$\delta ^{}\mathrm{\Pi }_\mu =0,$$
(30)
$$\delta ^{}C=\epsilon \frac{1}{^2}_i\mathrm{\Pi }^i,$$
(31)
$$\delta ^{}\overline{C}=0,$$
(32)
$$\delta ^{}𝒫=0,$$
(33)
$$\delta ^{}\overline{𝒫}=i\epsilon \mathrm{\Pi }_0.$$
(34)
Thus it follows that on integration over the momenta the gauge-fixing term remains invariant under the transformation generated by $`\mathrm{\Omega }^{}`$, that is, $`\delta ^{}(_\mu A^\mu )=0`$. The above expressions coincide with the Lavelle and McMullan’s result . However, these symmetry transformations turn out to be nonlocal. The preceding analysis opens up the way to a stimulating discussion of how the so-called dual BRST symmetry appears. This is precisely the task that we shall carry out in the next section.
## III Dual BRST symmetry
As already mentioned, our immediate objective is to implement the above general considerations to the two-dimensional case. With this in mind, we start by considering
$$=\frac{1}{4}F_{\mu \nu }F^{\mu \nu },$$
(35)
in two dimensions of spacetime. Just as for the four-dimensional case, the canonical action is
$$S=𝑑x\left(\stackrel{}{A_1}\mathrm{\Pi }^1H_0\lambda \phi \right),$$
(36)
where $`\mathrm{\Pi }^1`$ is the momenta conjugate to $`A_1`$. Here it is important to realize that the corresponding canonical Hamiltonian is now
$$H_0=𝑑x\left(\frac{1}{2}\mathrm{\Pi }_1\mathrm{\Pi }^1+\mathrm{\Pi }_1^1A_0\right),$$
(37)
The constraint structure for the gauge field naturally remains identical to the previous case ( See Eq. (9) ). Thus it follows that
$$\phi =_1\mathrm{\Pi }^1=0.$$
(38)
Again we find that the effective action can be written as
$$S_{eff}=𝑑x\left(\mathrm{\Pi }^1\stackrel{}{A_1}+𝒫\stackrel{}{C}+\overline{𝒫}\stackrel{}{\overline{C}}+\mathrm{\Pi }^0\stackrel{}{A_0}H_0[\mathrm{\Omega },\mathrm{\Psi }]\right).$$
(39)
As in the preceding section, the BRST generator reduces to
$$\mathrm{\Omega }=𝑑x\left(C_1\mathrm{\Pi }^1i\overline{𝒫}\mathrm{\Pi }_0\right).$$
(40)
We can now write the corresponding transformations generated by $`\mathrm{\Omega }`$, that is,
$$\delta A_1=\epsilon _1C,$$
(41)
$$\delta A_0=i\epsilon \overline{𝒫},$$
(42)
$$\delta \mathrm{\Pi }_1=0,$$
(43)
$$\delta \mathrm{\Pi }_0=0,$$
(44)
$$\delta C=0,$$
(45)
$$\delta \overline{C}=i\epsilon \mathrm{\Pi }_0,$$
(46)
$$\delta \overline{𝒫}=0,$$
(47)
$$\delta 𝒫=\epsilon _1\mathrm{\Pi }^1.$$
(48)
Following our procedure we now calculate the effective action (39). As in the previous section, we choose the gauge-fixing function in the form
$$\mathrm{\Psi }=𝑑x\left(𝒫A_0i\overline{C}\left(_1A^1\frac{\xi }{2}\mathrm{\Pi }_0\right)\right).$$
(49)
In the present case, the effective action is found to be
$$S_{eff}=d^2x\left(\frac{1}{4}F_{\mu \nu }F^{\mu \nu }\frac{1}{2\xi }\left(_\mu A^\mu \right)^2+i\overline{C}_\mu ^\mu C\right).$$
(50)
It has been recently claimed that the effective action is also invariant under local variations:
$$\delta _DA_\mu =\eta \epsilon _{\mu \nu }^\nu \overline{C},$$
(51)
$$\delta _D\overline{C}=0,$$
(52)
$$\delta _DC=i\eta ,$$
(53)
where $``$ is an auxiliary field. Accordingly, we have that
$$\delta _D\left(_\mu A^\mu \right)=0.$$
(54)
At this point it is reasonable to ask how the transformations (51\- 53) are related with the ones (41\- 48). In view of this situation and on the basis of the discussion in the previous section, we now proceed to perform a canonical transformation in the ghost sector. In that case, we propose the following canonical transformation
$$C^{}=i\frac{𝒫}{_1},$$
(55)
$$P^{}=i_1C,$$
(56)
$$\overline{C}^{}=i\frac{\overline{𝒫}}{^1},$$
(57)
$$\overline{P}^\iota =i^1\overline{C}.$$
(58)
As before, we keep the notation $`\mathrm{\Omega }^{}`$ for the charge which results from a canonical transformation. Thus, the new charge may be rewritten as
$$\mathrm{\Omega }^{}=𝑑x\left(i\frac{𝒫}{_1}\left(_1\mathrm{\Pi }^1\right)^1\overline{C}\mathrm{\Pi }_0\right).$$
(59)
It is now once again straightforward to work out the transformations generated by (59). They are
$$\delta ^{}A_1=i\epsilon 𝒫,$$
(60)
$$\delta ^{}A_0=\epsilon ^1\overline{C},$$
(61)
$$\delta ^{}\mathrm{\Pi }_\mu =0,$$
(62)
$$\delta ^{}\overline{C}=0,$$
(63)
$$\delta ^{}C=i\epsilon \mathrm{\Pi }^1,$$
(64)
$$\delta ^{}𝒫=0,$$
(65)
$$\delta ^{}\overline{𝒫}=\epsilon ^1\mathrm{\Pi }_0.$$
(66)
One immediately sees that, on integration over the momenta, the above transformations (60-66) reduce to the ones found in . It is important to realize that, after integration over the momenta, the new transformations yield $`\delta ^{}(_\mu A^\mu )=0`$ off shell. It is worthwhile sketching at this point our procedure. As mentioned before, in the extended phase space we have $`\delta \left(_\mu A^\mu \right)=\epsilon \left(i\stackrel{}{\overline{𝒫}}+_1^1C\right)`$, which in the configuration space reads $`\delta \left(_\mu A^\mu \right)=\epsilon _\mu ^\mu C`$, but this is just the classical equation of motion of $`C`$. From our above analysis, we see that the proposed canonical transformation makes a change of the ghost equations, that is, $`\delta ^{}\left(_\mu A^\mu \right)=\epsilon ^1\left(\stackrel{}{\overline{C}}i𝒫\right)`$ which after integration over the momenta gives zero, turning the variation of the gauge-fixing term null on shell to null off shell. Since the canonical transformation has been carried out in the ghost sector, all the basics processes that can be explained by the old effective action, should likewise be obtained from the new effective action. It is satisfying to notice the simplicity and directness of this derivation, which is manifestly gauge-independent.
## IV ACKNOWLEDGMENTS
The author would like to thank I. Schmidt for his support, and J.J. Yang for reading the manuscript. Work supported in part by Fondecyt (Chile) grant 1000710, and by a Cátedra Presidencial (Chile).
|
warning/0005/quant-ph0005101.html
|
ar5iv
|
text
|
# Optimal local implementation of non-local quantum gates
## I Introduction
A quantum computer allows, in principle, for the efficient solution of some problems that are intractable on a classical computer, the most striking example being the factorization of large numbers . However, the practical problems involved in the actual construction of a quantum computer of an interesting size (certainly more than 50 qubits) that is capable of performing a sufficiently large number of logical gates (a few hundred appear as a lower limit for an interesting problem involving 50 qubits) are daunting. Problems range from fundamental effects such as decoherence and dissipation, experimental imperfections for example in the timing, length and intensity of the laser pulses to the non-trivial task of storing and isolating reliably a large number of qubits . In fact, in proposals such as ion trap or the cavity QED implementations it seems problematic to store and process very large numbers of qubits in a single ‘processor’. A possible way out would be the construction of a quantum computer not as a local device that contains all qubits in a single processor, but to build it from the outset as a multi-processor device where each processor contains only a small number of qubits. Such a ’distributed quantum computer’ can be viewed as a generalization of a quantum communication network in which each node can act as a sender or receiver and contains only a small number of qubits. Distributed quantum computation has been considered previously by Grover , and he demonstrated that the solution of a phase estimation problem can be obtained efficiently with such a device assuming ideal conditions. It was later shown, that even under non-ideal conditions, i.e., in the presence of decoherence, a distributed quantum computer can be superior to a classical computer in terms of the resources that are required for the solution of the phase estimation problem . However, these investigations considered the specific problem of phase estimation and did not address the question of universal quantum computation. Before one is able to consider the physical resource efficiency of a distributed quantum computer in general, it is necessary to establish first optimal implementations of quantum gates between qubits that are located in different nodes of the distributed quantum computer. This problem is addressed in this paper. We present optimal protocols implementing gates that affect qubits in different nodes (here dubbed non-local gates) only using local operations and classical communication (LOCC) and previously shared entanglement. Optimality is measured in terms of the consumption of the basic experimental resources of entanglement and classical communication between nodes. We present general theorems that give lower bounds on the resources required for the implementation of quantum gates and for several universal quantum gates we present optimal implementations. We also discuss the general structure of the classical communication transfer in these implementations.
It should be noted that the issue addressed in the present paper is different from the question as to whether (and how) a particular entanglement transformation is possible under local quantum operations and classical communication in that in the course of the non-local implementation of a quantum gate the initial state is not known in advance. Instead, with the use of shared entanglement particular joint unitary operations between several parties are simulated.
In Section II we begin with an investigation of two-qubit gates. We establish some lower bounds on the resources that are required to implement two-qubit gates and present optimal implementations for a number of important gates. In particular we present a protocol that implements a CNOT gate consuming one ebit of entanglement and using only one classical bit of communication between the two parties. We then proceed in Section III to study multi-party gates such as Toffoli gates and other more general multi-party quantum gates again presenting bounds on the required physical resources and optimal protocols for some important classes of gates.
## II Non-local two-qubit gates
General single-bit rotations together with a CNOT gate are sufficient to implement any multi-qubit unitary transformation. This implies that the resource requirements for the implementation of a CNOT gate are a limiting factor in the construction of general unitary transformations in a distributed quantum computer. For this reason we investigate first the CNOT gate.
###### Theorem 1
One bit of classical communication in each direction and one shared ebit is necessary and sufficient for the non-local implementation of a quantum CNOT gate.
Proof: (i) Necessity: To demonstrate that one bit of communication in each direction is necessary we first note that the procedure consists of local operations and classical communication. As local operations cannot transmit information from Alice to Bob, or vice versa, all information which has been sent at the end of the operation must have been sent classically. Consider now the CNOT quantum gate. If the target qubit is initialised in the state $`|0`$, then its final state will be $`|0`$ or $`|1`$ depending on the initial state of the control qubit being $`|0`$ or $`|1`$ respectively. Therefore, the final result of the gate in this case is the communication of one bit of information from Alice (holding the control qubit) to Bob (holding the target qubit). Consequently, in the non-local implementation, one bit of classical information must have been sent classically from Alice to Bob. The reason for this can be seen from an elegant argument presented in the figure caption of the last figure in (see for more details). In short, assume that Alice needs to send less than one bit. In that case she could omit sending the bit and force Bob to make a guess. As he would guess the correct answer with a probability larger than $`1/2`$, Alice and Bob could then use error correction codes to establish a perfect channel and would end up with a superluminal communication channel. To see that one bit must also have been sent from Bob to Alice, we need merely note that in the basis $`|\pm =(|0\pm |1)/\sqrt{2}`$ the role of control and target in a CNOT gate are reversed. Consequently, if Alice’s particle is prepared in the standard state $`|+`$ and Bob chooses to prepare his particle either in state $`|+`$ or $`|`$, Alice will, after the application of the CNOT gate, hold a particle which is either in state $`|+`$ or $`|`$ depending on the state Bob’s particle has been prepared in. Therefore one bit of information has been transmitted from Bob to Alice. As the implementation of the CNOT must be independent of the initial state, the procedure must allow for one bit of communication in each direction, and as a consequence the non-local implementation must involve, as a minimum, one bit of communication in both directions.
That one ebit is required can be seen from the fact that a CNOT gate acting on the initial state $`(|0_A+|1_A)|0_B`$ leads to a maximally entangled state $`(|00_{AB}+|11_{AB}`$. As the amount of entanglement cannot be increased by local operations, this implies that the non-local implementation of a CNOT gate must consume at least one ebit.
(ii) Sufficiency: In the following we construct a quantum circuit which performs the CNOT non-locally using one e-bit and the transmission of one classical bit in each direction. This quantum circuit is given in figure 1. The CNOT is performed between the qubits $`A`$ and $`B`$. Alice holds the qubits $`A`$ and $`A_1`$, and Bob holds the qubits $`B`$ and $`B_1`$. The wavy line connecting $`A_1`$ and $`B_1`$ signifies that they are entangled. In particular we will choose their initial state to be $`(|00+|11)/\sqrt{2}`$. The initial state of $`A`$ is necessarily arbitrary, and so is given by $`\alpha |0_A+\beta |1_A`$. The initial state of $`B`$ is also arbitrary, and is given by $`\gamma |0_B+\delta |1_B`$. Time now flows from left to right in figure 1. First a local CNOT is performed with $`A`$ as the control and $`A_1`$ as the target. After this the combined state of $`A`$, $`A_1`$ and $`B_1`$ is
$$\frac{1}{\sqrt{2}}(\alpha |000+\alpha |011+\beta |110+\beta |101)_{AA_1B_1}.$$
(1)
Alice then performs a measurement on $`A_1`$ in the computational basis, and the line corresponding to this qubit terminates. The result of the measurement is one bit of information, which is communicated to Bob, and this communication is denoted by the dashed line. If the result is $`|0`$ Bob does nothing, and if the result is $`|1`$ Bob performs the not operation. At this point the combined state of $`A`$ and $`B_1`$ is $`\alpha |00_{AB_1}+\beta |11_{AB_1}`$. That is, we have now effectively performed a CNOT between $`A`$ and $`B_1`$, in which the initial state of $`2`$ was $`|0`$. Now particle $`B_1`$ contains the necessary information about the state of $`A`$. We can now perform a CNOT between $`B_1`$ and $`B`$. The combined state of $`A`$, $`B_1`$ and $`B`$ is now
$$\frac{1}{\sqrt{2}}(\alpha \gamma |000+\alpha \delta |001+\beta \delta |110+\beta \gamma |111)_{AB_1B}.$$
(2)
All we have to do is to remove $`B_1`$ from the state. This is done by performing a Hadamard transformation on $`B_1`$, and then measuring $`B_1`$ in the computational basis, at which point the line denoting $`B_1`$ terminates. The result of the measurement (one bit) is communicated to Alice. If the result is ’0’ Alice does nothing, and if the result is ’1’ she performs a (state-independent) $`\sigma _z`$ operation on particle $`A`$. This completes the non-local CNOT.
###### Theorem 2
A control-U gate can be implemented using one shared ebit and one bit of classical communication in each direction.
Proof: A control-U gate is defined as a gate that applies the identity on the target qubit if the control bit is in state $`|0`$ and it applies the unitary operator $`U`$ to the target if the control qubit is in state $`|1`$. The same quantum circuit as in Fig. 1 can be used except that the CNOT gate on Bobs side is replaced by a control-U gate.
In general a single application of a control-U gate cannot be employed to create one e-bit from an initial product state of two qubits. Furthermore, the amount of classical information that can be sent from Alice to Bob via a general control-U gate is less than one bit. This raises the question as to whether such a control-U gate can be implemented with less resources than a full ebit and one classical bit of communication in each direction. Clearly this will not be possible when we only wish to implement a single instance of a control-U gate. However, it may be conceivable that one has a situation in which one needs to carry out a large number of control-U gates simultaneously. In that case it is conceivable that this could be done with less than 1 ebit of entanglement per gate and less than one bit of classical communication in each direction. However, this turns out to be a difficult question and we have been unable to find such a scheme.
Let us now move on to investigate general two-qubit quantum gates to establish the minimum resource requirements for their implementation.
###### Theorem 3
Two bits of classical communication in both directions and two shared ebits is sufficient for the non-local implementation of a general two-bit gate.
Proof: To demonstrate that this amount of communication is sufficient to implement all quantum operations we need merely invoke quantum teleportation. Any operation may be performed by teleporting Alice’s state to Bob, at which point Bob may locally perform the operation, and then teleport the resulting state back to Alice. This procedure requires two bits of communication in each direction and $`2`$ shared ebits
Moreover, there are two-qubit gates that require two bits of classical communication in each direction and consumes $`2`$ bits. An example is the state-swapper, which may be written as three CNOT gates, one after the other, with Alice as the control, target, and then control, in that order (see Fig. 2) . To show that two bits of classical communication are required (each way) in the non-local implementation of this gate, we need to show that this amount of information may be communicated from Alice to Bob (and vice versa) when the gate is performed. To do this we merely have to note that at the completion of the gate Alice has sent her state to Bob. Now, this state could have been initially in a maximally entangled state with a qubit that Bob possesses. Dense coding tells us that this enables Alice to send two bits of information to Bob . Naturally, Bob can use the same procedure to send two bits of information to Alice. Therefore, in a non-local implementation, the state swapper requires at least two bits of communication in each direction. An analogous argument shows that the state swapper would also require two shared ebits, as a state swapper can be used to establish two ebits from a product state. To achieve this one simply applies the state swapper to particles $`A_2`$ and $`B_2`$ of the state $`(|00_{A_1A_2}+|11_{A_1A_2})(|00_{B_1B_2}+|11_{B_1B_2})`$.
It is remarkable that the swap gate requires only two shared ebits as it can be shown that three CNOT gates are necessary to implement it when one employs the ordinary gate array picture using a universal set of quantum gates that is made up of CNOT gates and local unitary operations . This observation may be useful, as it demonstrates that in some cases the use of entanglement can be replaced partially by local measurements and classical communication.
Before we move on to investigate the implementation of non-local multi-party gates we would like to analyze the structure of the classical information transfer involved in the gate implementation somewhat further. In both examples discussed above it turned out that the classical information transfer between the two parties is symmetric, i.e., the same number of bits need to be sent from Alice to Bob and vice versa. Likewise, the amount of classical information that can be sent using these two-qubit gates is also the same in each direction. It is therefore quite natural to ask whether this is the case in general. Indeed we have not been able to find a counter-example and we therefore make the following two closely related conjectures.
###### Conjecture 4
The minimal amount of classical communication required to implement any two-party quantum gate with one qubit associated with each party and shared $`M`$ ebits, $`M=1,2`$, is always the same in each direction.
###### Conjecture 5
The amount of classical information that can be sent via any two-qubit gate is the same in each direction.
While these conjectures appear natural, we have not been able to find general proofs for them. However, we have been able to confirm both of them for a number of classes of two-qubit quantum gates. An example of a gate which has the same classical information capacity in both directions is the CNOT gate whose optimal implementation has been described above. How can we see that a quantum gate is symmetric with respect to its capability for classical information transfer? Before we move on to the most general case, let us consider the CNOT gate. Imagine we have the ability to perform a CNOT gate with Alice as the control and Bob as the target. Using this gate and local operations only, we can then also implement a CNOT with Alice as a target and Bob as a control, simply by applying a Hadamard gate to each qubit both before and after the CNOT, see Fig. 3.
The two versions of the CNOT gate are also related via the (nonlocal) state swapper.
$`U_{\text{CNOT}}^{BA}`$ $`=`$ $`U_{ss}U_{\text{CNOT}}^{AB}U_{ss}^{}`$ (3)
$`=`$ $`(HH)U_{\text{CNOT}}^{AB}(HH).`$ (4)
where $`U_{\text{CNOT}}^{AB}`$ represents the CNOT gate with $`A`$ as a control and $`B`$ as a target and $`U_{ss}`$ denotes the state swapper. In general if we can achieve the transformation $`U_{BA}U_{ss}U_{AB}U_{ss}^{}`$ from $`U_{AB}`$ and purely local operations, i.e., if there exist local one-qubit unitary operators $`U_1`$, $`U_2`$, $`U_3`$ and $`U_4`$ for which we have
$`U_{BA}`$ $`=`$ $`U_{ss}U_{AB}U_{ss}^{}`$ (5)
$`=`$ $`(U_1U_2)U_{AB}(U_3U_4)`$ (6)
then Eq. (5) is a sufficient condition for the classical information transmission capacities in each direction to be equal. In the following we will determine some sets of quantum gates $`U_{AB}`$ for which Eq. (5) holds.
Let us begin with a slightly simpler problem. Suppose that we have a two-qubit quantum gate $`V_1U(4)`$. $`V_1`$ can be expressed in terms of its generator as $`V_1=\mathrm{exp}(iH_1)`$, where the generator $`H_1`$ is a Hermitean operator. We now define another quantum gate $`V_2`$ as
$`V_2U_{ss}V_1U_{ss}^{}`$ $`=`$ $`U_{ss}e^{iH_1}U_{ss}^{}`$ (7)
$`=`$ $`e^{iU_{ss}H_1U_{ss}^{}}e^{iH_2}`$ (8)
where the generator $`H_2`$ of $`V_2`$ is clearly a Hermitean operator. Our goal can therefore be reformulated as: For which unitary operators $`V_1`$ can we write $`V_2`$ as $`V_2=(U_1U_2)V_1(U_1^{}U_2^{})`$, or equivalently for which generators $`H_1`$ of $`V_1`$ can we write
$$H_2U_{ss}H_1U_{ss}^{}=(U_1U_2)H_1(U_1^{}U_2^{}).$$
(9)
Note that this is less general than the transformation in Eq. (5). It is useful to realize that both the unitary operator $`V_1`$ and its generator $`H_1`$ are diagonal in the same basis, say $`\{|\varphi _i,i=1,\mathrm{\hspace{0.33em}2},\mathrm{\hspace{0.33em}3},\mathrm{\hspace{0.33em}4}\}`$. Furthermore, we can decompose $`H_1`$ with respect to its eigenvectors as $`H_1=_i\lambda _i|\varphi _i\varphi _i|_i\lambda _i\rho _i`$, where $`\lambda _i`$ is the eigenvalue of $`H_1`$ corresponding to the eigenvector $`|\varphi _i`$. Consequently, Eq. (9) becomes
$$\underset{i}{}\lambda _iU_{ss}\rho _iU_{ss}^{}=\underset{i}{}\lambda _i(U_1U_2)\rho _i(U_1^{}U_2^{})$$
(10)
We can now prove a number of lemmas. We begin with
###### Lemma 6
Any two-qubit quantum gate that has a generator with a single non-vanishing eigenvalue is symmetric with respect to its classical information transfer capacity.
Proof: Suppose that the only non-vanishing eigenvalue of the generator $`H_1`$ is $`\lambda _1`$ . In that case we can always find one-qubit unitary operators $`U_1`$ and $`U_2`$ such that Eq. (10) holds. To see this, note that the eigenstate $`|\varphi _i`$ is actually a pure state describing a system composed by two qubits. Therefore, it has the Schmidt decomposition $`|\varphi _1=_k\sqrt{p_k}|k_A\stackrel{~}{|k}_B_k\sqrt{p_k}|k\stackrel{~}{|k}`$. Furthermore, in this case we have
$`{\displaystyle \underset{i}{}}`$ $`\lambda _i`$ $`U_{ss}|\varphi _i\varphi _i\left|U_{ss}^{}=\lambda _1{\displaystyle \underset{k,l}{}}\sqrt{p_kp_l}\stackrel{~}{|k}\right|k\stackrel{~}{l|}l|`$ (11)
$`=(\stackrel{~}{U}`$ $``$ $`U)({\displaystyle \underset{k,l}{}}\lambda _1\sqrt{p_kp_l}|k\stackrel{~}{|k}l|\stackrel{~}{l|})(\stackrel{~}{U}^{}U^{})`$ (12)
$`=(\stackrel{~}{U}`$ $``$ $`U)({\displaystyle \underset{i}{}}\lambda _i|\varphi _i\varphi _i|)(\stackrel{~}{U}^{}U^{}),`$ (13)
where $`U`$ is defined to be the unitary operator which maps each basis vector $`|i`$ to its corresponding $`\stackrel{~}{|i}`$. Similarly, the unitary operator $`\stackrel{~}{U}`$ maps each basis vector $`\stackrel{~}{|i}`$ to its corresponding $`|i`$, i.e., $`\stackrel{~}{U}=U^{}`$.
Another non-trivial class of quantum gates $`U_{bd}`$ for which condition (9) holds, is the one whose generator is Bell diagonal, i.e., we have
###### Lemma 7
Any two-qubit quantum gate that has a generator which is Bell-diagonal is symmetric with respect to its classical information transfer capacity.
Proof: If $`|\mathrm{\Psi }`$ is any of the Bell states, the reader can easily verify that
$$|\mathrm{\Psi }\mathrm{\Psi }|=U_{ss}|\mathrm{\Psi }\mathrm{\Psi }\left|U_{ss}^{}=(\sigma _z\sigma _z)\right|\mathrm{\Psi }\mathrm{\Psi }|(\sigma _z\sigma _z)$$
Therefore, for the quantum gate $`U_{bd}`$, condition (9) is satisfied by either choosing $`U_1=U_2=\mathrm{𝟙}`$ or $`U_1=U_2=\sigma _z`$. Recall that $`\sigma _z`$ is the Pauli matrix corresponding to the arbitrarily chosen z direction.
Note however, that condition (9) is not satisfied for all quantum gates $`U_{AB}`$. A counterexample is the gate
$`U_{AB}`$ $`=`$ $`e^{i\lambda _1}|\mathrm{\hspace{0.17em}0}+\mathrm{\hspace{0.17em}0}+\left|+e^{i\lambda _2}\right|\mathrm{\hspace{0.17em}0}\mathrm{\hspace{0.17em}0}|`$ (14)
$`+`$ $`e^{i\lambda _3}|\mathrm{\hspace{0.17em}10}\mathrm{\hspace{0.17em}10}\left|+e^{i\lambda _4}\right|\mathrm{\hspace{0.17em}11}\mathrm{\hspace{0.17em}11}|.`$ (15)
For $`\lambda _1=\lambda _2=0`$ and non-trivial choice of $`\lambda _3`$ and $`\lambda _4`$ it is not possible to find local unitary operators $`U_1`$ and $`U_2`$ such that Eq. (9) is satisfied. Nevertheless, it is possible to find local unitary operators $`U_1`$, $`U_2`$, $`U_3`$ and $`U_4`$ which satisfy the more general condition (5). The local unitary operators will be of the form :
$`U_1=e^{i\lambda _4}|\mathrm{\hspace{0.17em}1}\mathrm{\hspace{0.17em}1}\left|+e^{i(\lambda _3\lambda _4)}\right|\mathrm{\hspace{0.17em}0}\mathrm{\hspace{0.17em}0}|,`$ (16)
$`U_2=|\mathrm{\hspace{0.17em}1}\mathrm{\hspace{0.17em}1}\left|+e^{i(\lambda _3\lambda _4)}\right|\mathrm{\hspace{0.17em}0}\mathrm{\hspace{0.17em}0}|,`$ (17)
$`U_3=\mathrm{𝟙},`$ (18)
$`U_4=e^{i\lambda _4}|\mathrm{\hspace{0.17em}1}\mathrm{\hspace{0.17em}1}|+|\mathrm{\hspace{0.17em}0}\mathrm{\hspace{0.17em}0}|.`$ (19)
We can then conclude to the following lemma:
###### Lemma 8
The amount of classical information that can be sent via any control-U gate of the form
$$U=|\mathrm{\hspace{0.17em}0}\mathrm{\hspace{0.17em}0}|\mathrm{𝟙}+|\mathrm{\hspace{0.17em}1}\mathrm{\hspace{0.17em}1}|\left(e^{i\lambda _3}|\mathrm{\hspace{0.17em}0}\mathrm{\hspace{0.17em}0}\left|+e^{i\lambda _4}\right|\mathrm{\hspace{0.17em}1}\mathrm{\hspace{0.17em}1}|\right)$$
is the same in each direction.
It should be noted that this does not mean that the amount of information transferred in any particular operation of the gate will be the same in both directions, as this will depend upon the choice of initial states. However, an implementation of the gate must work for all possible initial states, (in particular it must work for the case where both qubits are pure and therefore contain their maximum capacity), and this is what puts the limit on the minimal communication requirement.
It is clear that we may now put 2-bit quantum gates into two classes. Those which require no more than one bit of two-way communication, and those that require more than one bit (but no more than two bits). The CNOT falls into the first category, and the state-swapper falls into the second. Two other standard gates which fall into the first category are the c-U (which performs a unitary transformation on one system depending on the state of the other), and the state-preparer.
## III Non-local multi-party gates
In the previous section we have presented a number of results concerning the implementation of non-local two-qubit quantum gates in a distributed quantum computer. In the following we will generalize these ideas to local implementation of multi-qubit gates, i.e., gates where more than two parties are involved. To illuminate the system behind the construction, we explain the implementation of the Toffoli gate from which the generalization to other multi-party gates will be evident.
###### Theorem 9
Two shared ebits and a total of four bits of classical communication are necessary and sufficient for the local implementation of a non-local three-party quantum Toffoli gate.
Proof: (i) Necessity: A Toffoli gate can be reduced to an ordinary CNOT gate when one fixes the state of one of the control qubits to be $`|1`$. Chose the state of party $`A`$ to be $`|1`$. Then the initial state is
$$|\psi _{\text{ini}}=|1_A(\alpha |0+\beta |1)(\gamma |0+\delta |1),$$
(20)
and after the application of the Toffoli gate we find
$$|\psi _{\text{ini}}=|1_A(\alpha \gamma |00+\alpha \delta |01+\beta \gamma |11+\beta \delta |10)_{BC},$$
(21)
which shows that we have implemented a CNOT between parties $`B`$ and $`C`$. Therefore, Theorem 1 implies that one classical bit has to be exchanged in both directions between $`A`$ and the target party $`C`$ and one ebit has to be shared between them. The same argument applies when we fix the state of qubit $`B`$ to be $`|1`$.
Sufficiency. The implementation of the Toffoli gate with these minimal resources is presented in Fig. 4. Assume that Alice and Clare share a pair $`A_1`$, $`C_1`$ of qubits in a maximally entangled state $`|\varphi ^+=(|00+|11)/\sqrt{2}`$, and that Bob and Clare share another pair of particles $`B_1`$ and $`C_2`$ in the same state. Then the initial state of the whole system consisting of particles $`A`$, $`B`$, $`C`$, $`A_1`$, $`B_1`$, $`C_1`$ and $`C_2`$ is of the form
$$|\psi =|\psi _A|\psi _B|\psi _C|\varphi ^+_{A_1C_1}|\varphi ^+_{B_1C_2},$$
(22)
where
$`|\psi _A`$ $`=`$ $`\alpha |0+\beta |1,`$ (23)
$`|\psi _B`$ $`=`$ $`\gamma |0+\delta |1,`$ (24)
$`|\psi _C`$ $`=`$ $`\eta |0+\xi |1.`$ (25)
The first step is a local quantum CNOT gate on $`A`$ and $`A_1`$ with $`A`$ as control. Then Alice measures particle $`A_1`$ and Clare performs a NOT operation on her particle $`C_1`$ if Alice finds $`|1`$ and the identity if Alice finds $`|0`$. Qubit $`A_1`$ is subsequently discarded. Now Bob applies a local CNOT with $`B`$ being the control and $`B_1`$ being the target. Then Bob measures particle $`B_1`$ and Clare performs a NOT operation on her particle $`C_2`$ if Bob finds $`|1`$ and the identity if Bob finds $`|0`$. Qubit $`B_1`$ is subsequently discarded. Now the state of the remaining qubits $`A`$, $`B`$, $`C`$, $`C_1`$ and $`C_2`$ is given by
$`(\alpha |00+\beta |11)_{AC_1}(\gamma |00+\delta |11)_{BC_2}|\psi _C.`$ (26)
In a further step Clare applies locally a Toffoli with $`C_1`$ and $`C_2`$ being the control qubits. Subsequently Clare applies Hadamard gates to the qubits $`C_1`$ and $`C_2`$. Then she measures $`C_2`$ and applies $`\sigma _z`$ or the identity $`\mathrm{𝟏}`$ to $`B`$ if her result is $`|1`$ or $`|0`$ respectively. Finally she measures $`C_1`$ and applies $`\sigma _z`$ or the identity to $`A`$ if her result is $`|1`$ or $`|0`$ respectively. This completes the Toffoli gate.
The total number of classical bits which have to be communicated is four, and only two shared ebits of entanglement are consumed. Again, these results can be generalized to three-party control-U operations that can be represented in matrix form with respect to the computational basis as
$$\mathrm{𝟙}_6\left(\begin{array}{cc}u_{00}& u_{01}\\ u_{10}& u_{11}\end{array}\right),$$
(27)
where
$$\left(\begin{array}{cc}u_{00}& u_{01}\\ u_{10}& u_{11}\end{array}\right)$$
(28)
is the matrix representation of a unitary operator $`U`$. We only need to replace the local Toffoli gate by a local three-party control-U. This gives rise to
###### Lemma 10
A three party control-U gate can be implemented using four bits of classical communication and two shared ebits.
Using Theorem 9 and Lemma 10 we are now in a position to construct every possible quantum gate array using only ebits, classical communication and local operations. In particular one could use the results in to construct $`N`$-party controlled gates from CNOTs and single bit rotations. This, however, is not optimal in terms of physical resources. While it will be difficult to construct the optimal procedure for general quantum gates, for some gates we are able to find these procedures. We find for example
###### Theorem 11
An $`N`$ party control-U gate can be implemented using $`2(N1)`$ bits of classical communication and $`N1`$ shared ebits.
Proof: The control parties are enumerated from $`P_1`$ to $`P_{N1}`$ and each of them is carrying one ancilla numerated by $`P_1^{}`$ to $`P_{N1}^{}`$. The target qubit is denoted by $`T`$ and the target party possesses $`N1`$ further ancillary qubits. The first $`N1`$ steps of the protocol are essentially analogous. In the k-th step a local quantum CNOT gate is applied on $`P_k`$ and $`P_k^{}`$ with $`P_k`$ as control. Then this party measures particle $`P_k^{}`$ and the target party performs a NOT operation on her ancillary qubit $`T_k`$ if Alice finds $`|1`$ and the identity if Alice finds $`|0`$. Qubit $`P_k^{}`$ is subsequently discarded. Now we apply an $`N`$-party controlled U gate on Clares particles, with the ancillas $`C_1,\mathrm{},C_{N1}`$ being the control qubits and $`T`$ the target. Subsequently the target party performs Hadamard gates on each of its ancillas.
This is then followed by $`N1`$ steps involving measurements. In the k-th step qubit $`T_k`$ is measured in the $`|0,|1`$ basis. If the outcome is $`|1`$, then $`\sigma _z`$ is applied to the qubit $`P_k`$; if the outcome is $`|0`$ then no action is taken on qubit $`P_k`$. Qubit $`T_k`$ is subsequently discarded. Hence, the total required resources are $`2(N1)\text{ bits of classical information}`$ and $`N1`$ initially shared ebits.
The amount of consumed resources in the latter protocol is rather surprizing. In an inefficient non-local implementation of the above $`N`$-party gate one could employ the simulation of the gate with the use of two-party control-U gates and CNOT gates as in Ref. , but such that each step is realized non-locally. In such a procedure a supply of $`3\times 2^{N1}4`$ ebits would be necessary.
A more efficient teleportation-based protocol in which the respective states of the qubits at different nodes are twice teleported would still use $`2(N1)`$ ebits and $`4(N1)`$ bits of classical information.
## IV Conclusions
In this work we have addressed the problem of the local implementation of non-local gates in a distributed quantum computer, i.e. a computer which is composed of many subunits (local processors). Such a configuration may be useful, as it requires only a small number of qubits (e.g. ions) to be stored at each site which may be experimentally more feasible than storing a large number of qubits in a single site. However, this raised the issue of the non-local implementation of quantum gates. We have addressed this question and have shown what the minimal resources for the implementation of two-qubit quantum gates are. We have presented explicit optimal constructions for the local implementation of non-local control-U gates. We have generalized these results to multi-party gates such as for example the Toffoli gate. We have also adressed some issues concerning the structure of the information exchange that is required in these implementations. We hope that this work will be useful for the assessment of the viability of distributed quantum computation.
We acknowledge useful discussion with Daniel Jonathan and John Vaccaro. This work was supported by the Deutsche Forschungsgemeinschaft (DFG), the UK engineering and physical sciences research council (EPSRC), The Leverhulme Trust, the European Science Foundation (ESF) programme on quantum information processing, the EQUIP programme of the European Union and the State Scholarships Foundation of Greece.
Endnote: During completion of this work we became aware of the closely related work by D. Collins, N. Linden, and S. Popescu, Phys. Rev. A 64, 032302 (2001), quant-ph/0005102.
|
warning/0005/hep-ph0005307.html
|
ar5iv
|
text
|
# Loops and legs beyond perturbation theoryContributed to the Loops and Legs in Quantum Field Theory conference, April 9–14, 2000, Bastei-Königstein, Germany.
## 1 INTRODUCTION
There has been a lot of work lately in extending the calculation of radiative corrections to higher loop-orders and to multi-leg processes. The proceedings of this Workshop are suggestive for the complexity this field has reached since the pioneering work of Martinus Veltman and Gerardus ’t Hooft. Higher order QCD radiative corrections are necessary for interpreting the experimental data which the LHC experiments will provide. Higher order electroweak corrections are needed for interpreting precision electroweak measurements. Due to the time frame for the construction of future colliders, precision measurements will continue to be a main tool for new physics searches.
In this contribution we would like to discuss the use of multi-loop Feynman graphs for recovering results which are normally not within the scope of perturbation theory, namely deriving scattering amplitudes in the case of strongly coupled field theories.
A central question in particle physics is how is the electroweak symmetry broken. The experimental data currently available is compatible with a standard model scalar sector, and favors a light Higgs boson. However, additional degrees of freedom beyond the standard model have the potential to shift these fits. It is not at all excluded that the LHC experiments will see a scalar resonance at a considerably higher energy.
At strong coupling, the scalar sector was insufficiently explored. Qualitatively, one may expect new phenomena to appear: a Higgs particle strongly coupled to the vector bosons and to itself, anomalous vector boson self-interactions, and possibly the appearance of a spectrum of additional resonances in the scalar sector. Perturbative calculations loose their predictive power when radiative corrections blow up in higher loop orders, and the dependency on the renormalization scheme becomes substantial .
Ideally, one needs a solution for the scalar sector which is valid at strong coupling as well as at weak, and which is free of renormalization scheme uncertainty. This can be accomplished by explicitly summing up all loop orders within a non-perturbative $`1/N`$ expansion. If a solution of sufficient accuracy is desired, such that it can be used in phenomenological calculations to compete with ordinary perturbation theory, then the $`1/N`$ expansion must be carried out at higher order.
## 2 THE $`1/N`$ SOLUTION
We are interested in a solution of the scalar sector of the standard model when the self-coupling of the Higgs field becomes strong. The non-perturbative effects are thus given by the scalar field self-interaction, while the gauge coupling remains perturbative. It is then natural to resort to the equivalence theorem to relate amplitudes involving longitudinal electroweak gauge bosons to the corresponding amplitudes involving would-be Goldstone bosons. The problem is being reduced to solving a linear sigma model non-perturbatively.
### 2.1 The auxiliary field formalism
It is sufficient to consider the scalar sector alone, which is an $`SU(2)`$-symmetric linear sigma model. In the following, the scalar sector is extended to an $`O(N)`$-symmetric sigma model — with the standard model case recovered for $`N=4`$ — such that an expansion in powers of $`1/N`$ can be performed:
$`_1`$ $`=`$ $`{\displaystyle \frac{1}{2}}_\nu \mathrm{\Phi }_0^\nu \mathrm{\Phi }_0{\displaystyle \frac{\mu _0^2}{2}}\mathrm{\Phi }_0^2{\displaystyle \frac{\lambda _0}{4!N}}\mathrm{\Phi }_0^4,`$
$`\mathrm{\Phi }_0`$ $``$ $`(\varphi _0^1,\varphi _0^2,\mathrm{},\varphi _0^N)`$ (1)
From the Lagrangian above, it is easy to derive scattering amplitudes at leading order in $`1/N`$. Typically, this involves summing up a geometric series of one-loop bubble diagrams which have Goldstone bosons in the loop.
However, the auxiliary field formalism which was proposed in ref. proves to be most useful in organizing sub-leading orders in the $`1/N`$ expansion in a diagrammatically manageable way. It consists in introducing an additional unphysical field $`\chi `$ in the Lagrangian:
$`_2`$ $`=`$ $`_1+{\displaystyle \frac{3N}{2\lambda _0}}(\chi _0{\displaystyle \frac{\lambda _0}{6N}}\mathrm{\Phi }_0^2\mu _0^2)^2`$ (2)
$`=`$ $`{\displaystyle \frac{1}{2}}_\nu \mathrm{\Phi }_0^\nu \mathrm{\Phi }_0{\displaystyle \frac{1}{2}}\chi _0\mathrm{\Phi }_0^2+{\displaystyle \frac{3N}{2\lambda _0}}\chi _0^2`$
$`{\displaystyle \frac{3\mu _0^2N}{\lambda _0}}\chi _0+const.`$
Because the equation of motion for the auxiliary field $`\chi `$ is a constant, this addition does not change the dynamics. Green’s functions having Higgs and Goldstone bosons on the legs are the same, whether calculated with the Feynman rules given by $`_1`$ or by $`_2`$.
The Feynman rules, on the other hand, are changed. $`_1`$ contains the trilinear and quartic vertices $`\sigma \pi \pi `$, $`\sigma \sigma \sigma `$, $`\sigma \sigma \sigma \sigma `$, $`\sigma \sigma \pi \pi `$, and $`\pi \pi \pi \pi `$. Here $`\sigma `$ and $`\pi `$ are the massive and massless modes stemming from Lagrangian $`_1`$ after spontaneous symmetry breaking, respectively. $`_2`$ contains only the trilinear couplings $`\chi \sigma \sigma `$ and $`\chi \pi \pi `$.
This simplifies enormously the topological classification of Feynman diagrams according to their power of $`1/N`$ beyond leading order. The reader can easily convince himself of the utility of the auxiliary field formalism by writing down the diagrams which contribute to Goldstone-Goldstone scattering at NLO in $`1/N`$ in both formalisms.
### 2.2 Tachyonic regularization
By summing up the chains of one-loop bubble self-energy insertions, one encounters an ultraviolet renormalon. This gives rise to an additional, tachyonic pole in the propagators, apart from the expected physical spectrum containing one Higgs boson and $`N1`$ Goldstone modes.
By direct evaluation of the leading order contribution in $`1/N`$ to the two-point functions, one obtains the following propagators :
$`D_{\sigma \sigma }(s)`$ $`=`$ $`{\displaystyle \frac{i}{sm^2(s)}}`$
$`D_{\chi \chi }(s)`$ $`=`$ $`{\displaystyle \frac{1}{Nv^2}}{\displaystyle \frac{ism^2(s)}{sm^2(s)}}`$
$`D_{\chi \sigma }(s)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{N}v}}{\displaystyle \frac{im^2(s)}{sm^2(s)}}`$
$`D_{\pi _i\pi _j}(s)`$ $`=`$ $`{\displaystyle \frac{i}{s}}\delta _{ij},`$ (3)
where
$`m^2(s)`$ $`=`$ $`{\displaystyle \frac{v^2}{\frac{3}{\lambda }+\widehat{\alpha }^{(0)}(s)}}`$ (4)
$``$ $`{\displaystyle \frac{v^2}{\frac{3}{\lambda }\frac{1}{32\pi ^2}\mathrm{log}\left(\frac{s+i\eta }{\mu ^2}\right)}}.`$
Here $`\widehat{\alpha }^{(0)}(s)`$ is the ultraviolet finite part of the one-loop self-energy bubble diagram, with a Goldstone boson in the loop. $`\mu `$ is the ultraviolet subtraction scale.
The propagators contain an Euclidian pole in the ultraviolet region at an energy $`s=\mathrm{\Lambda }_t^2`$ given by the following transcendental equation:
$$\frac{v^2}{\mathrm{\Lambda }_t^2}\frac{1}{32\pi ^2}\mathrm{log}\left(\frac{\mathrm{\Lambda }_t^2}{\mu ^2}\right)+\frac{3}{\lambda }=0.$$
(5)
The tachyon scale is in the ultraviolet region for low values of the coupling, and tends to move towards low energy when the coupling is increased.
In higher order $`1/N`$ corrections, the Euclidian pole appears in the loop momentum integration. This induces causality violating contributions, even though these effects are numerically small as long as the tachyon scale is high enough. For this reason, it is necessary to investigate the origin of this singularity and find a way for dealing with it.
At any finite order of perturbation theory, the tachyon pole does not exist. It is an artifact of the bubble diagram summation. In the process of summing up all loop orders, an ambiguity is present, however. The summation of the perturbative series determines the result only up to functions which vanish in perturbation theory, such as $`e^{1/\lambda }`$. Because the residuum of the tachyonic pole is precisely such a function which vanishes in perturbation theory, its presence or absence is completely arbitrary and cannot be determined within perturbation theory. For this reason, it is justified to restore causality by minimally subtracting the tachyon at its pole since the original information stemming from Feynman diagrams remains unchanged :
$`D_{\sigma \sigma }(s)`$ $`=`$ $`i\left[{\displaystyle \frac{1}{sm^2(s)}}{\displaystyle \frac{\kappa }{s+\mathrm{\Lambda }_t^2}}\right]`$
$`D_{\chi \chi }(s)`$ $`=`$ $`{\displaystyle \frac{is}{Nv^2}}\left[{\displaystyle \frac{m^2(s)}{sm^2(s)}}+{\displaystyle \frac{\kappa \mathrm{\Lambda }_t^2}{s+\mathrm{\Lambda }_t^2}}\right]`$
$`D_{\chi \sigma }(s)`$ $`=`$ $`{\displaystyle \frac{i}{\sqrt{N}v}}\left[{\displaystyle \frac{m^2(s)}{sm^2(s)}}+{\displaystyle \frac{\kappa \mathrm{\Lambda }_t^2}{s+\mathrm{\Lambda }_t^2}}\right],`$ (6)
were $`\kappa =[1+\mathrm{\Lambda }_t^2/(32\pi ^2v^2)]^1`$ is the residuum of the tachyonic pole.
The tachyonic regularization can be seen as a prescription for summing up the perturbative series in a way which preserves causality.
### 2.3 $`1/N`$ renormalization
Performing renormalization at NLO in the $`1/N`$ expansion involves the treatment of ultraviolet divergences of diagrams with various numbers of loops. One way of keeping track of various loop order counterterms is to group them into $`1/N`$ counterterms. Each of the $`1/N`$ order counterterms $`\delta \lambda `$, $`\delta t`$, $`\delta t_\chi `$, $`\delta Z_{\pi ,\sigma ,\chi }`$ is a power series in the coupling constant $`\lambda `$ :
$`{\displaystyle \frac{3}{\lambda _0}}`$ $`=`$ $`{\displaystyle \frac{3}{\lambda }}+\mathrm{\Delta }\lambda `$
$``$ $`{\displaystyle \frac{3}{\lambda }}+\delta \lambda ^{(0)}+{\displaystyle \frac{1}{N}}\delta \lambda +𝒪\left({\displaystyle \frac{1}{N^2}}\right)`$
$`{\displaystyle \frac{3\mu _0}{\lambda _0}}`$ $`=`$ $`{\displaystyle \frac{v^2}{2}}(1+\mathrm{\Delta }t)`$
$``$ $`{\displaystyle \frac{v^2}{2}}\left[1+{\displaystyle \frac{1}{N}}\delta t+𝒪\left({\displaystyle \frac{1}{N^2}}\right)\right]`$
$`\varphi _0^i`$ $`=`$ $`\pi _iZ_\pi ,i=1,\mathrm{},N1`$
$``$ $`\pi _i\left[1+{\displaystyle \frac{1}{N}}\delta Z_\pi +𝒪\left({\displaystyle \frac{1}{N^2}}\right)\right]`$
$`\varphi _0^N`$ $`=`$ $`\sigma Z_\sigma +\sqrt{N}v`$
$``$ $`\sigma \left[1+{\displaystyle \frac{1}{N}}\delta Z_\sigma +𝒪\left({\displaystyle \frac{1}{N^2}}\right)\right]+\sqrt{N}v`$
$`\chi _0`$ $`=`$ $`\chi Z_\chi +\widehat{\chi }+\mathrm{\Delta }t_\chi `$ (7)
$``$ $`\chi \left(1+{\displaystyle \frac{1}{N}}\delta Z_\chi \right)+{\displaystyle \frac{v^2}{N}}\delta t_\chi +𝒪\left({\displaystyle \frac{1}{N^2}}\right)`$
In principle, these counterterms are sufficient to absorb the UV divergences of all diagrams of NLO in $`1/N`$.
When calculating $`1/N`$ diagrams at NLO, we resort to numerical integration. The ultraviolet divergence of the diagram must be removed before the numerical integration can be carried out, but the explicit $`ϵ`$ expansion and isolation of UV poles is cumbersome because of the complexity of the diagrams. For this reason we perform a BPHZ-type renormalization . We subtract the divergences of the diagrams according to the forest formula, as shown in figure 1. The subtracted expressions, being finite in the ultraviolet, can be integrated numerically. For a given physical process, the renormalization program then means to combine all UV subtraction terms from various diagrams with each other. Most subtractions cancel with each other, and the small remaining set of subtractions is trivially absorbed into the $`1/N`$ local counterterms above. We have checked that all UV subtractions are polynomial, as they ought to be in order to be absorbed into local counterterms.
### 2.4 Numerical solution
The finite expressions depicted in figure 1 are calculated by numerical integration because so far no analytical approach exists for dealing with this type of diagrams.
For a numerical solution, it is advantageous to identify first all loop integrations which can be performed analytically. These are associated with closed Goldstone loops. For the diagrams shown in figure 1, it is possible to perform analytically all integrations except for one final loop integration. The final integration involves the resummed and tachyonically subtracted propagators of eqs. 6 and form factors from one-loop triangle or box diagrams involving massless Goldstone bosons in the loop. The final loop integration needs to be performed numerically. It can be reduced to a two-fold integral with the methods of refs. , which do not resort to Feynman parameters.
### 2.5 Scheme independence of physical predictions
For a given order in the expansion parameter $`1/N`$, the Feynman diagrams of all loop orders are explicitly summed up. For this reason, the final result for a physical scattering amplitude is free of renormalization scheme dependence. In usual perturbative calculations, a renormalization scheme ambiguity is present because of the truncation of the perturbative expansion; this ambiguity is of higher order in the coupling constant.
We note that subtracted diagrams of the type shown in figure 1 are not individually free of renormalization scheme dependency. The internal subtractions which are performed for making the integrals finite in the ultraviolet, are calculated at a given subtraction point. This defines an intermediary renormalization scheme.
This renormalization scheme dependence cancels out only in physical results such as scattering amplitudes. In the following section we give the results for two scattering processes. We have checked explicitly that the two functions involved, $`f_1`$ and $`f_2`$ of eqns. 9, are independent of the subtraction point.
## 3 $`\mu ^+\mu ^{}`$ COLLIDERS AND THE HIGGS MASS SATURATION EFFECT
Recently, feasibility studies for muon colliders attracted quite some attention. A muon collider would be an ideal $`s`$-channel Higgs factory. For a heavy Higgs boson, two processes dominate: $`\mu ^+\mu ^{}Ht\overline{t}`$ and $`\mu ^+\mu ^{}HW_L^+W_L^{},Z_LZ_L`$.
Within the $`1/N`$ expansion, the amplitudes for these processes at NLO are given by the following expressions :
$`_{f\overline{f}}`$ $`=`$ $`{\displaystyle \frac{1}{sm^2(s)\left[1\frac{1}{N}f_1(s)\right]}}`$
$`_{ZZ}`$ $`=`$ $`{\displaystyle \frac{m^2(s)}{\sqrt{N}v}}{\displaystyle \frac{1\frac{1}{N}f_2(s)}{sm^2(s)\left[1\frac{1}{N}f_1(s)\right]}}.`$ (8)
Here, the correction functions $`f_1`$ and $`f_2`$ are given by a combination of the two- and three-point functions defined in figure 1 (see ):
$`f_1(s)`$ $`=`$ $`{\displaystyle \frac{m^2(s)}{v^2}}\widehat{\alpha }(s)+2\widehat{\gamma }(s)+{\displaystyle \frac{v^2}{m^2(s)}}[\widehat{\beta }(s)`$
$`2{\displaystyle \frac{sm^2(s)}{v^2}}(\delta Z_\sigma \delta Z_\pi )]`$
$`f_2(s)`$ $`=`$ $`{\displaystyle \frac{m^2(s)}{v^2}}\widehat{\alpha }(s)+\widehat{\gamma }(s)`$ (9)
$`\widehat{\varphi }(s){\displaystyle \frac{v^2}{m^2(s)}}\widehat{\eta }(s).`$
In figure 2 we give numerical results for these expressions. We plot the shape of the Higgs resonance for various strengths of the coupling. From the positions of the resonance maxima, it can be seen clearly that when the coupling is increased, a mass saturation effect appears. The position of the resonance’s peak does not increase above a saturation value just under 1 TeV. The precise resonance shape and implicitly the saturation value are process dependent.
For comparison, we indicate in figure 2 the corresponding peak maxima which would be obtained by using usual perturbation theory at NNLO. A saturation is present qualitatively in this curve, too. However, it should be noticed that the perturbative curve is affected by large radiative corrections and large scheme uncertainties at such large coupling, and therefore is not reliable in the saturation region.
## 4 SATURATION EFFECT AT THE LHC
The main Higgs production mechanism at the LHC is gluon fusion which proceeds through a top loop . The existing perturbative analyses indicate that the vector boson fusion becomes competitive at an energy of the order of 1 TeV.
In the presence of the saturation effect, the observation of a heavy Higgs resonance at the LHC will be different. In the saturation zone, the mass of the resonance remains more or less the same while the width increases with the coupling. The resonance becomes flatter and more difficult to detect.
In figure 3 we show a plot of the Higgs width as a function of the Higgs mass. The width and mass used in these plots, $`M_{PEAK}`$ and $`\mathrm{\Gamma }_{PEAK}`$, are defined from the line shape of the resonance as seen in fermion-fermion scattering. They are derived from the position and height of the resonance as if it were of Breit-Wigner type. In this picture we show the calculations available so far in usual perturbation theory (LO, NLO, NNLO) , and in the $`1/N`$ expansion (LO and NLO) . It can be seen that the two expansions converge nicely towards each other and display a mass saturation.
We studied the discovery potential of the LHC in the presence of the saturation effect. In ref. we performed a Monte-Carlo simulation for the LHC. Only leptonic channels were considered, namely the “golden plated” channel $`(l^+l^{})(l^+l^{})`$, and the $`l^+l^{}\nu \overline{\nu }`$ channel. When neutrinos in the final state are involved, the Higgs resonance appears as a Jacobian peak in the distribution of missing $`p_T`$. We included the gluon fusion process together with the relevant background. The strong interacting Higgs correction was included by using the NNLO perturbative calculation - it results into a faster Monte-Carlo while being a fair approximation of the non-perturbative $`1/N`$ result.
We used usual assumptions about the LHC energy and luminosity, and asked for a $`5\sigma `$ effect with respect to the background. The discovery potential estimated for the “golden plated” channel corresponds then to an on-shell Higgs mass of 830 GeV. The missing $`p_T`$ channel reaches up to an on-shell Higgs mass of 1030 GeV. Note that these values for the on-shell mass deviate considerably from the actual position of the resonance. The actual peak can be read from figure 3, where the on-shell mass is mapped onto the saturation curve.
## 5 CONCLUSIONS
The combinatorial structure of the sigma model makes possible an explicit calculation of the non-perturbative $`1/N`$ expansion at NLO.
The $`1/N`$ solution is valid at strong coupling as well as at weak coupling. It is also independent of the intermediate renormalization scheme which is being used. At NLO it also provides for ultraviolet finite wave function renormalization constants.
The non-perturbative solution of the sigma model is relevant for the standard model via the equivalence theorem. It implies that a mass saturation effect is present in the scalar sector. When the Higgs self-interaction is increased, the mass of the resonance increases only up to a maximum value under 1 TeV, while the width increases continuously.
An interesting question is how cutoff effects stemming from new physics at higher energies can modify the $`1/N`$ solution. This clearly deserves further investigation.
Acknowledgement The authors are grateful to the Organizers of the Loops and Legs in Quantum Field Theory 2000 conference for having organized an exciting and stimulating meeting. The work of T. B. is supported by the EU Fourth Training Programme ”Training and Mobility of Researchers”, Network ”Quantum Chromodynamics and the Deep Structure of Elementary Particles”, contract FMRX–CT98–0194 (DG 12 - MIHT).
|
warning/0005/astro-ph0005080.html
|
ar5iv
|
text
|
# 1 Introduction
## 1 Introduction
The existence of an isotropic, diffuse gamma background radiation (GBR) was first suggested by data from the SAS 2 satellite (Thompson & Fichtel 1982). The EGRET instrument on the Compton Gamma Ray Observatory confirmed this finding: by removal of point sources and of the galactic-disk and galactic-centre emission, and after an extrapolation to zero local column density, a uniformly distributed GBR was found, of alleged extragalactic origin (Sreekumar et al. 1998). Above an energy of $`10`$ MeV, this radiation –to which we shall refer throughout simply as “the GBR”– has a featureless spectrum, shown in Fig. 1, which is very well described by a simple power-law form, $`\mathrm{dF}/\mathrm{dE}\mathrm{E}^\beta `$, with $`\beta 2.10\pm 0.03`$ (Sreekumar et al. 1998).
The origin of the GBR is still unknown. The published candidate sources range from the quite conventional to the decisively speculative. Perhaps the most conservative hypothesis for the origin of an isotropic GBR is that it is extragalactic, and originates from active galaxies (Bignami et al. 1979; Kazanas & Protheroe 1983; Stecker & Salamon 1996). The fact that blazars have a $`\gamma `$-ray spectrum with an average index $`2.15\pm 0.04`$, compatible with that of the GBR, supports this hypothesis (Chiang & Mukerjee 1998). The possibility has also been discussed that Geminga-type pulsars, expelled into the galactic halo by asymmetric supernova explosions, be abundant enough to explain the GBR (Dixon et al. 1998; Hartmann 1995). More exotic hypotheses include a baryon-symmetric universe (Stecker et al. 1971), now excluded (Cohen et al. 1998), primordial black hole evaporation (Page & Hawking 1976; Hawking 1977), supermassive black holes formed at very high redshift (Gnedin & Ostriker 1992), annihilation of weakly interactive big-bang remnants (Silk & Srednicki 1984; Rudaz & Stecker 1991), and a long etc.
However, the EGRET GBR data in directions above the galactic disk and centre show a significant deviation from isotropy, correlated with the structure of our galaxy and our position relative to its centre (Dar et al. 1999). This advocates a local (as opposed to cosmological) origin for the GBR. Indications of a large galactic contribution to the GBR at large latitudes were independently found by Dixon et al. (1998) by means of a wavelet-based “non-parametric” approach that makes no reference to a particular model. Strong & Moskalenko (1998) and Moskalenko & Strong (2000) also found that the contribution of inverse Compton scattering of galactic cosmic ray electrons to the diffuse $`\gamma `$-ray background is presumably much larger than previously thought. In this paper we go one step further and explore in detail the possibility (Dar et al. 1999) that the diffuse gamma-ray background radiation at high galactic latitudes could be dominated by inverse Compton scattering of cosmic ray (CR) electrons on the cosmic microwave background radiation and on starlight from our own galaxy. In Section 2 we briefly review the GBR data and the evidence for its correlation with our position in the Galaxy.
The CR-proton and CR-electron spectra are briefly reviewed in Section 3. The origin, spectrum and composition of non-solar cosmic ray protons and nuclei have been debated for almost a century. The measurements now extend over some 30 orders of magnitude in flux and some 15 orders of magnitude in energy, up to an astonishing $`\mathrm{E}3\times 10^{11}`$ GeV (Bird et al. 1995, Takeda et al. 1998, Berezinskii et al. 1990 and references therein). Above $`5`$ GeV, this spectrum has also a power-law form $`\mathrm{E}^\beta `$, with two small variations in the “index” $`\beta `$ at the so-called “CR knee” and “CR ankle”. The local spectrum of CR electrons, shown in Fig. 2, is much harder to measure; it is only known up to $`10^3`$ GeV and, above $`5`$ GeV, it is also well described by a simple power law.
In Sections 4 and 5 we discuss relations between the indices of the GBR and the CR electron and proton spectra. In so doing, we make few and very simple assumptions: that the mechanism accelerating CR hadrons and CR electrons is the same (a moving magnetic “mirror”), that the locally-measured electron spectrum is representative of its average form throughout the Galaxy, that above a certain energy, inevitably, the electron spectrum is modulated by inverse Compton scattering on starlight and on the microwave background radiation, and that the GBR is dominated by the resulting Compton up-scattered photons. This allows one to derive, successfully, the GBR index from the electron index and the electron index from the proton index. The GBR index, as observed by EGRET, is uncannily directionally uniform. We interpret this fact as strong support for our simple assumptions.
In Section 6 we tackle a more difficult and potentially controversial subject: the origin and magnitude of the GBR. In a sense, our proposed explanation –that the GBR originates from inverse Compton scattering in our own galaxy (Dar et al. 1999) – is more conservative than any of the previously suggested origins.
The non-conventional aspect of our hypothesis is that, in order to reproduce the observed intensity of the GBR, we must assume the scale height of our galaxy’s CR-electron distribution to be almost twice the traditionally-accepted upper limit. Because of this, in Section 6, we briefly review the basis of the conventional wisdom and our critical view of it, whose main points are the following. Moskalenko, Strong and their collaborators have developed a very detailed understanding of the CR, radio and $`\gamma `$ observations of our galaxy. To fit the data, their models require a freely parametrized reacceleration of electrons, presumably by the motion of turbulent magnetic fields (e.g., Seo & Ptuskin, 1994). Strong & Moskalenko (1998) introduce a cutoff $`\mathrm{z}_\mathrm{h}`$ for the height above the galactic plane above which cosmic rays freely escape. They find an upper limit $`\mathrm{z}_\mathrm{h}<12`$ kpc, on the basis of a fit to the $`{}_{}{}^{10}\mathrm{Be}/^9\mathrm{Be}`$ ratio observed by Ulysses (Connell 1998). This result is “soft”: twice the upper limit would still be compatible with the ensemble of data (Lukasiak et al. 1994). Moreover, the galactic CR proton distribution extracted from a fit to EGRET $`\gamma `$-ray data, actually favours (Strong & Moskalenko 1998) an ad hoc distribution of CR sources that is not as well localized in the disk as the conventional supernova-remnant sources are (Webber 1997), even if $`\mathrm{z}_\mathrm{h}=20`$ kpc or more. This point, and the necessity to invoke CR reacceleration, indicate that scale heights of the CR electron distribution in excess of the 12 kpc “upper limit” may not be out of the question. Our results are optimized by a scale height of roughly 20 kpc. Such a large scale height is not in contradiction with radio synchrotron-emission from our galaxy if the galactic disk and its magnetic field are embedded in a larger magnetic halo with a much weaker field.
In studying the possibility that the diffuse GBR is not extragalactic, one has two choices. The first is to extend to high galactic latitudes the elaborate models (with many parameters, reacceleration, and ad hoc modifications of the CR-proton and CR-electron energy and source distributions) that have been developed to describe the intricate nature of the observations at low galactic latitudes (Strong & Moskalenko 1998; Moskalenko & Strong 2000). The second is to adopt our very naive set of hypotheses and employ a simple cosmic-ray model with, by conventional standards, a large scale height for CR-electrons. Models of this type (Dar & Plaga 1999), wherein cosmic ray sources are directly injected at high galactic latitudes, have actually been proposed<sup>1</sup><sup>1</sup>1The injector agents would be highly relativistic jets from the birth of compact objects in supernova explosions, leading to a CR population permeating a magnetized region of galactic-halo proportions and constituting a putative solution to the problem of the origin of the highest-energy cosmic rays, a qualitative description of the nuclear CR spectrum, and a possible explanation of jetted gamma-ray bursts..
In Section 7 we discuss the magnitude and angular-dependence of the two dominant contributions to the GBR within our model: inverse Compton scattering of galactic CR-electrons off the cosmic background radiation and starlight. In Section 8 we compute the small additive effect of sunlight, and in Section 9 we estimate the contribution from external galaxies, which is also sub-dominant. In Section 10 we compare our predictions with the data on the intensity and the angular dependence of the GBR. The results are very satisfactory and, within our model, lead to the conclusion that the GBR can be dominated by the emission from our own galaxy. We summarize our conclusions and predictions in Section 11.
## 2 The GBR data
We call “the GBR” the diffuse emission observed by EGRET by masking the galactic plane at latitudes $`|\mathrm{b}|10^\mathrm{o}`$, as well as the galactic centre at $`|\mathrm{b}|30^\mathrm{o}`$ for longitudes $`|\mathrm{l}|40^\mathrm{o}`$, and by extrapolating to zero column density, to eliminate the $`\pi ^0`$ and bremsstrahlung contributions to the observed radiation and to tame the model-dependence of the results. Outside the mask, the GBR flux integrated over all directions in the observed energy range of $`30`$ MeV to $`120\mathrm{GeV}`$, shown in Fig. 1, is well described by a power law:
$$\frac{\mathrm{dF}_\gamma }{\mathrm{dE}}(2.74\pm 0.11)\times 10^3\left[\frac{\mathrm{E}}{\mathrm{MeV}}\right]^{2.10\pm 0.03}\mathrm{cm}^2\mathrm{s}^1\mathrm{sr}^1\mathrm{MeV}^1.$$
(1)
The overall magnitude in Eq. (1) is sensitive to the model used to subtract the foreground (Sreekumar et al. 1998; Strong et al. 1998), but the spectral index is not. The EGRET data are given in Sreekumar et al. (1998) for 36 $`(\mathrm{b},\mathrm{l})`$ domains, 9 values for each half-hemisphere. The spectral index is, within errors, extremely directionally uniform, as shown in Fig. 3, where we have plotted the EGRET results as functions of $`\theta `$, the observation angle relative to the direction to the galactic centre ($`\mathrm{cos}\theta =\mathrm{cos}[\mathrm{b}]\mathrm{cos}[\mathrm{l}]`$). The normalization is less homogeneous, but in directions well above the galactic disk and away from the galactic-centre region it has been found to be consistent with a normal distribution around the mean value: thus the claim of a possible extragalactic origin (Sreekumar et al. 1998).
In Fig. 4 we have plotted, as a function of $`\theta `$, the EGRET GBR counting-rate above 100 MeV. This figure clearly shows, in three out of the four quarters of the celestial sphere, an increase of the counting rate towards the galactic centre. How significant is this effect? Let $`\overline{\chi }^2\chi ^2/\mathrm{d}.\mathrm{o}.\mathrm{f}.`$ be the “reduced” $`\chi ^2`$ per degree of freedom. The $`\overline{\chi }^2`$ value for constant flux is 2.6: very unsatisfactory. A best fit of the form $`\mathrm{F}=\mathrm{F}_0+\mathrm{F}_1(1\mathrm{cos}\theta )`$ yields $`\overline{\chi }^2=1.3`$, a very large amelioration (for higher polynomials in $`\mathrm{cos}\theta `$ the higher-order coefficients are compatible with zero: the fit does not significantly improve). Note also that at angles with $`\mathrm{cos}\theta `$ larger than its mean value $`\mathrm{cos}\theta =0.0246`$ ($`\theta <88.6^\mathrm{o}`$), 10 out of the 12 data points are above the average flux, while at angles with $`\theta >88.6^\mathrm{o}`$, 18 out of the 24 data points are below the average. The probability for a uniform distribution to produce this large or larger a fluctuation is $`1.5\times 10^4`$.
Even in directions pointing to the galactic disk and the galactic centre, EGRET data on $`\gamma `$-rays above 1 GeV show an excess over the expectation from galactic cosmic-ray production of $`\pi ^0`$’s (Pohl & Esposito 1998). Electron bremsstrahlung in gas is not the source of the 1–30 MeV inner-Galaxy $`\gamma `$-rays observed by COMPTEL (Strong et al. 1997), since their galactic latitude distribution is broader than that of the gas. These findings also imply that inverse Compton scattering may be much more important than previously believed (Strong & Moskalenko 1998; Moskalenko and Strong, 2000; Dar et al. 1999).
## 3 The CR data
The cosmic ray nuclei have a power-law spectral flux $`\mathrm{dF}/\mathrm{dE}\mathrm{E}^\beta `$ with an index $`\beta `$ that changes at two break-point energies. In the interval $`10^{10}\mathrm{eV}<\mathrm{E}<\mathrm{E}_{\mathrm{knee}}`$ $`3\times 10^{15}`$ eV, protons constitute $`96\%`$ of the CRs at fixed energy per nucleon, and their flux is (Berezinskii et al. 1990, and references therein):
$$\frac{\mathrm{dF}_\mathrm{p}}{\mathrm{dE}}1.8\left[\frac{\mathrm{E}}{\mathrm{GeV}}\right]^{2.70\pm 0.05}\mathrm{cm}^2\mathrm{s}^1\mathrm{sr}^1\mathrm{GeV}^1.$$
(2)
In the interval $`\mathrm{E}_{\mathrm{knee}}<\mathrm{E}<\mathrm{E}_{\mathrm{ankle}}`$ $`3\times 10^{18}`$ eV, the spectrum steepens from $`\beta _12.7`$ to $`\beta _23.0`$, flattening again to $`\beta _32.5`$ above $`\mathrm{E}_{\mathrm{ankle}}`$.
The CR flux of electrons (Prince 1979; Nishimura et al. 1980; Tang 1984; Golden et al. 1984; Evenson & Meyers 1984; Golden et al. 1994; Ferrando et al. 1996; Barwick et al. 1998; Wiebel-Sooth & Biermann 1998), shown in Fig. 2, is well fitted, from $`\mathrm{E}10\mathrm{GeV}`$ to $`2`$ TeV by:
$$\frac{\mathrm{dF}_\mathrm{e}}{\mathrm{dE}}(2.5\pm 0.5)\times 10^5\left[\frac{\mathrm{E}}{\mathrm{MeV}}\right]^{3.2\pm 0.10}\mathrm{cm}^2\mathrm{s}^1\mathrm{sr}^1\mathrm{MeV}^1.$$
(3)
The terrestrial and solar magnetic fields and the solar wind modify the electron spectrum below $`\mathrm{E}10`$ GeV, so that the direct observations at those energies may deviate from the local interstellar spectral shape.
Cosmic ray electrons undergo inverse Compton scattering (ICS) off the ambient photon baths: starlight and the cosmic background radiation. The spectral indices of the GBR and electron spectra can be very simply and successfully related (Dar et al. 1999), if the GBR dominantly consists of photons whose energy has been uplifted by ICS, as we proceed to show.
## 4 The index of the GBR spectrum
The current temperature, number density and mean energy of the CMB are $`\mathrm{T}_0=2.728`$ K, $`\mathrm{n}_0411\mathrm{cm}^3`$, and $`ϵ_0`$ $`2.7\mathrm{kT}_0`$ $`6.36\times 10^{10}\mathrm{MeV}`$ (Mather et al. 1993; Fixsen et al. 1996). The galactic starlight (SL) distribution is highly non-uniform, its average energy is $`ϵ_{}1`$ eV. Consider the ICS of high energy electrons on these radiations. Assume the shape of the electron flux, Eq. (3), observed at $`\mathrm{E}>10`$ GeV, to be representative of the average galactic spectrum. For the energy range of EGRET the Thomson limit is accurate even for ICS on SL, and the $`\mathrm{e}\gamma `$ cross section is $`\sigma __\mathrm{T}0.65\times 10^{24}\mathrm{cm}^2`$. The mean energy $`\mathrm{E}_\gamma `$ of the upscattered photons, –or $`\mathrm{\Delta }\mathrm{E}_\mathrm{e}`$, the mean energy loss per collision– is:
$$\mathrm{E}_\gamma (ϵ_\mathrm{i})\mathrm{\Delta }\mathrm{E}_\mathrm{e}(ϵ_\mathrm{i})\frac{4}{3}\left(\frac{\mathrm{E}_\mathrm{e}}{\mathrm{m}_\mathrm{e}\mathrm{c}^2}\right)^2ϵ_\mathrm{i},$$
(4)
with $`ϵ_\mathrm{i}=ϵ_0`$ or $`ϵ_{}`$.
The ICS photon spectrum originating in our galaxy is the sum of CMB and SL contributions:
$$\frac{\mathrm{dF}_\gamma }{\mathrm{dE}}=\frac{\mathrm{dF}_\gamma ^0}{\mathrm{dE}}+\frac{\mathrm{dF}_\gamma ^{}}{\mathrm{dE}},$$
(5)
and is a function of the galactic latitude (b) and longitude (l) coordinates. The ICS final-photon spectrum –a cumbersome convolution (Felten & Morrison 1966) of a CR power spectrum with a photon thermal distribution– can be approximated very simply. Using again the index “i” to label the CMB and SL fluxes:
$$\frac{\mathrm{dF}_\gamma ^\mathrm{i}}{\mathrm{dE}_\gamma }\mathrm{N}_\mathrm{i}(\mathrm{b},\mathrm{l})\sigma __\mathrm{T}\frac{\mathrm{dE}_\mathrm{e}^\mathrm{i}}{\mathrm{dE}_\gamma }\left[\frac{\mathrm{dF}_\mathrm{e}}{\mathrm{dE}_\mathrm{e}}\right]_{\mathrm{E}_\mathrm{e}=\mathrm{E}_\mathrm{e}^\mathrm{i}};\mathrm{E}_\mathrm{e}^\mathrm{i}\mathrm{m}_\mathrm{e}\mathrm{c}^2\sqrt{\frac{3\mathrm{E}_\gamma }{4ϵ_\mathrm{i}}},$$
(6)
where $`\mathrm{E}_\mathrm{e}^\mathrm{i}`$ is obtained from Eqs. (4) by inverting $`\mathrm{E}_\gamma (ϵ_\mathrm{i})`$. We postpone to Section 6 the discussion of the model-dependent normalization factors $`\mathrm{N}_{}(\mathrm{b},\mathrm{l})`$ and $`\mathrm{N}_0(\mathrm{b},\mathrm{l})`$: effective column densities resulting from the convolution of the space distribution of CR electrons with those of starlight and of the CMB. Introducing the CR-electron flux of Eq. (3), of the form $`\mathrm{dF}_\mathrm{e}/\mathrm{dE}=\mathrm{A}[\mathrm{E}/\mathrm{MeV}]^{\beta _\mathrm{e}}`$, into Eqs. (6), we obtain:
$$\frac{\mathrm{dF}_\gamma ^\mathrm{i}}{\mathrm{dE}}=\frac{\mathrm{N}_\mathrm{i}(\mathrm{b},\mathrm{l})\sigma __\mathrm{T}\mathrm{A}}{2}\left[\frac{4ϵ_\mathrm{i}\mathrm{MeV}}{3\mathrm{m}_\mathrm{e}^2\mathrm{c}^4}\right]^{\frac{\beta _\mathrm{e}1}{2}}\left[\frac{\mathrm{E}}{\mathrm{MeV}}\right]^{\frac{\beta _\mathrm{e}+1}{2}}[\mathrm{E}]^{2.10\pm 0.05}.$$
(7)
In the energy-range of EGRET, the CMB and SL contributions have the same spectral index, as do the small sunlight and external-galaxy contributions discussed in Sections 8 and 9.
The photon spectral index of Eqs. (7), which is related to that of the CR-electrons through $`\beta _\gamma =(\beta _e+1)/2`$, coincides with the measured one, Eq. (1). The electron spectrum of Eq. (3) describes the data in the range $`\mathrm{E}_\mathrm{e}>5`$ GeV, so that Eq. (7) should be valid above $`\mathrm{E}_\gamma 100`$ keV, the typical energy of photons up-scattered from the CMB. At $`\mathrm{E}_\gamma >50`$ GeV, at the upper end of the EGRET data, $`\sigma _\mathrm{T}`$ in the SL contribution should be replaced by the complete Klein–Nishina cross section, implying a steepening of the spectrum. The corresponding effect for the CMB contribution is at energy above the EGRET energy range.
In deriving Eqs. (7), we have assumed that the locally-measured slope of Eq. (3) is representative of the index of the spectrum of the electrons suffering ICS to produce the GBR, wherever they may be. The spectral index of the diffuse GBR observed by EGRET is independent of direction, as shown in Fig. 3. The statistical test for a flat distribution is surprisingly good: $`\overline{\chi }^20.5`$. This is encouraging support for our working hypothesis of an electron spectrum with a universal shape, and of a simple and dominant mechanism –ICS– to generate the GBR.
## 5 The index of the electron spectrum
To relate the spectra of CR electrons and protons, we need an estimate of the protons’ spectrum at their source. A source spectrum $`\mathrm{dF}^\mathrm{s}/\mathrm{dE}`$ with index $`\beta _\mathrm{s}2.2`$ is obtained from collisionless shock simulations (Bednarz & Ostrowski 1998) or analytical estimates of acceleration by relativistic jets (Dar 1998). The CR spectrum of nuclei is modulated by their residence time in the Galaxy, $`\tau _{\mathrm{gal}}(\mathrm{E})`$. For a steady source of CRs the energy dependence of the observed flux is roughly that of $`\tau _{\mathrm{gal}}\mathrm{dF}^\mathrm{s}/\mathrm{dE}`$. Observations of astrophysical and solar plasmas and of nuclear abundances as functions of energy (e.g. Swordy et al. 1990) indicate that $`\tau _{\mathrm{gal}}(\mathrm{E})\mathrm{E}^{0.5\pm 0.1}`$, explaining $`\beta _1\beta _\mathrm{s}+0.52.7`$, as in Eq. (2).
Practically all CR acceleration mechanisms invoke an ionized medium that is swept by a moving magnetic field, such as would be carried by the rarefied plasma in a supernova shell (Bhattacharjee & Sigl 2000) or by a ‘plasmoid’ of jetted ejecta (Dar & Plaga 1999). The magnetic field acts as a moving ‘mirror’ that imparts the same distribution in velocity, or Lorentz factor $`\gamma =\mathrm{E}/\mathrm{m}\mathrm{c}^2`$, to all charged particles. To the extent that particle-specific losses (such as synchrotron radiation) can be neglected at the acceleration stage, all source fluxes have the same energy-dependence. For electrons below the anticipated ‘electron’s knee’ at $`\mathrm{E}_\mathrm{e}=(\mathrm{m}_\mathrm{e}/\mathrm{m}_\mathrm{p})\mathrm{E}_{\mathrm{knee}}`$$`2`$ TeV, we expect $`\mathrm{dF}_\mathrm{e}^\mathrm{s}/\mathrm{dE}\mathrm{E}^{\beta _\mathrm{s}}`$, with $`\beta _s2.2`$. Confinement effects preserve this equality for ultrarelativistic electrons and protons: their behaviour in a magnetic maze is the same. But, unlike for hadrons, the ‘cooling’ time of electrons –that are significantly affected by the ambient radiation and magnetic fields– is shorter than their galactic confinement time, $`\tau _{\mathrm{gal}}(\mathrm{E})`$, above a relatively low energy. This implies that the CR electron spectrum is modulated mainly by the ICS, and not by the confinement time.
Electrons lose energy not only by ICS on starlight and the CMB, but also by synchrotron radiation on magnetic fields. All of these processes are essentially the same: scattering off photons, either real or virtual. The energy loss is governed by the rate at which a single electron interacts with the ambient electromagnetic fields, weighted by the corresponding average energy density: $`\mathrm{P}=\sigma __\mathrm{T}\mathrm{c}[\mathrm{n}_{}ϵ_{}+\mathrm{n}_0ϵ_0+\mathrm{B}^2/(8\pi )]`$. Let $`\mathrm{R}_\mathrm{p}`$ (an inverse time) be the production rate of CR electrons, assumed to be constant (Berezinskii et al. 1990), and let $`\mathrm{dn}_\mathrm{e}^\mathrm{s}/\mathrm{dE}`$ be their source number-density spectrum. The actual density $`\mathrm{dn}_\mathrm{e}/\mathrm{dE}`$ in an interval $`\mathrm{dE}`$ about $`\mathrm{E}`$ is continuously replenished and depleted by electrons whose energy is being degraded by interactions. This leads to a steady-state situation in which production and losses are in balance. Using Eq. (4) we obtain:
$$\frac{4}{3}\frac{\mathrm{P}}{(\mathrm{m}_\mathrm{e}\mathrm{c}^2)^2}\frac{\mathrm{d}}{\mathrm{dE}}\left(\mathrm{E}^2\frac{\mathrm{dn}_\mathrm{e}}{\mathrm{dE}}\right)=\mathrm{R}_\mathrm{p}\frac{\mathrm{dn}_\mathrm{e}^\mathrm{s}}{\mathrm{dE}}.$$
(8)
For a relatively uniform galactic CR electron density, Eq. (8) also applies to the local electron flux $`\mathrm{dF}_\mathrm{e}(\mathrm{c}/4\pi )\mathrm{dn}_\mathrm{e}`$. Substitute the spectrum $`\mathrm{dn}_\mathrm{e}^\mathrm{s}/\mathrm{dE}\mathrm{E}^{\beta _\mathrm{s}}`$ into the flux version of Eq. (8) to obtain:
$$\frac{\mathrm{dF}_\mathrm{e}}{\mathrm{dE}}=\frac{3\mathrm{m}_\mathrm{e}^2\mathrm{c}^4\mathrm{R}}{4(\beta _\mathrm{s}1)\mathrm{P}}\frac{\mathrm{dF}_\mathrm{e}^\mathrm{s}}{\mathrm{E}\mathrm{dE}}\mathrm{E}^{(\beta _\mathrm{s}+1)}.$$
(9)
For electrons with $`\mathrm{E}_\mathrm{e}<(\mathrm{m}_\mathrm{e}/\mathrm{m}_\mathrm{p})\mathrm{E}_{\mathrm{knee}}`$ we deduced that $`\beta _s2.2.`$ Thus, $`\beta _s+1=3.2`$, in agreement with the data: Eq. (3) and Fig. 2. Above the ‘electron’s knee’ at $`\mathrm{E}_\mathrm{e}2`$ TeV the spectrum should steepen up by $`\mathrm{\Delta }\beta 0.25`$, like that of CR hadrons (Dar 1998). The available spectral measurements extend only to $`\mathrm{E}_\mathrm{e}1.5`$ TeV.
The energy density in the CMB is $`\mathrm{n}_0ϵ_0=0.24`$ eV cm<sup>-3</sup>, coincidentally similar to that in starlight at our location: $`\mathrm{n}_{}ϵ_{}0.22`$ eV cm<sup>-3</sup>. If the local CR and magnetic energy densities are in equipartition, $`\mathrm{B}^2/(8\pi )1`$ eV cm<sup>-3</sup>, again in the same ballpark. The cooling time of electrons in the ensemble of these fields is:
$$\tau _{_{\mathrm{cool}}}(\mathrm{E})\frac{3\mathrm{m}_\mathrm{e}^2\mathrm{c}^4}{4\mathrm{P}\mathrm{E}}0.22\times \left[\frac{\mathrm{E}}{\mathrm{GeV}}\right]^1\mathrm{Gy}.$$
(10)
The galactic escape time of GeV electrons, which should be similar to that of CR protons $`\tau _{\mathrm{gal}}(\mathrm{E})\mathrm{E}^{0.5\pm 0.1}`$ (Swordy et al. 1990), has a weaker energy dependence than that of $`\tau _{_{\mathrm{cool}}}`$. At sufficiently low energy, then, $`\tau _{\mathrm{gal}}<\tau _{_{\mathrm{cool}}}`$, and processes other than Compton- or synchrotron cooling (such as Coulomb scattering, ionization losses and bremsstrahlung) become relevant. The slope of Eq. (9) should change as the energy is lowered. The spectrum of Fig. 2 shows such a change, but it occurs at $`\mathrm{E}<10\mathrm{GeV}`$, a range in which local modulations would mask the effect.
## 6 The scale height of CR electrons
The radio emission of galaxies seen edge-on –interpreted as synchrotron radiation by electrons on their local magnetic field– offers direct observational evidence for CR electrons well above galactic disks (e.g. Duric et al. 1998). For the particularly well observed case of NGC 5755, the exponential scale height of the synchrotron radiation is $`𝒪(4)`$ kpc. If the CRs and the magnetic field energy are in equilibrium, they should have similar distributions, and the exponential scale height $`\mathrm{h}_\mathrm{e}`$ of the electrons ought to be roughly twice that of the synchrotron intensity, which reflects the convolution of the electron- and magnetic-field distributions. The inferred value $`\mathrm{h}_\mathrm{e}8`$ kpc for NGC 5755 may not be universal for spirals, since $`\mathrm{h}_\mathrm{e}`$ is very sensitive to the density and distribution of CR sources, gas and plasma in each particular galaxy. Moreover, the magnetic field may be in equipartition with cosmic rays only where the interstellar plasma is dense enough. It is quite possible for the CR electrons to be confined in a large magnetic halo with a field much smaller than that in the disk. For these reasons we must discuss the observations of our own particular galaxy.
Traditionally CR electrons and nuclei were assumed to have a distribution that snugly fit that of the visible part of the Galaxy –where their conventional sources lie– implying a scale height above the plane of the disk of $`𝒪(1)`$ kpc (Broadbend et al. 1989). As the data and their analysis became more elaborate, scale heights more than one order of magnitude larger were discussed (e.g. Strong et al. 1998). Since electrons lose energy to the ambient radiation close to their sources, which have traditionally been located in the disk, not very well understood CR-reacceleration phenomena have had to be invoked (e.g. Seo & Ptuskin 1994). Even with reacceleration, a conventional distribution of cosmic-ray sources fails to describe the observed GBR (Strong & Moskalenko 1998).
Over the years, Moskalenko, Strong and their collaborators have developed what is presumably the most elaborate and detailed understanding of the CR, radio and $`\gamma `$ observations of our galaxy (Moskalenko et al. 1998; Moskalenko and Strong, 2000; Strong and Moskalenko, 1998; Strong et al. 1997; Strong et al. 1998). A crucial parameter in their models is the scale $`\mathrm{z}_\mathrm{h}`$ of the CR distribution orthogonal to the galactic plane, defined as the height above which CRs freely escape, as in a leaky-box model. Strong & Moskalenko (1998) conclude that $`\mathrm{z}_\mathrm{h}`$ lies between 4 and 12 kpc. The limits are based on the comparison of the $`{}_{}{}^{10}\mathrm{Be}/^9\mathrm{Be}`$ ratio observed by Ulysses (Connell 1998) with model predictions as a function of $`\mathrm{z}_\mathrm{h}`$, being all other parameters fixed at their adopted values. The dependence of the $`{}_{}{}^{10}\mathrm{Be}/^9\mathrm{Be}`$ ratio on $`\mathrm{z}_\mathrm{h}`$, shown in Fig. 9 of Strong & Moskalenko (1998) and reproduced here as Fig. 5, is very weak for $`\mathrm{z}_\mathrm{h}>10`$ kpc. At $`\mathrm{z}_\mathrm{h}=20`$ kpc, the prediction would be only some 1.3 standard deviations below the Ulysses central value, and even $`\mathrm{z}_\mathrm{h}=40`$ would be viable: the average of all previous and somewhat less precise observations, compiled in Lukasiak et al. (1994) and shown in Fig. 5a, would be in agreement with $`\mathrm{z}_\mathrm{h}=`$ 20 or 40 kpc. For all these reasons and the ones stated in the introduction, we shall not refrain from considering scale heights above the 12 kpc upper limit quoted by Strong & Moskalenko (1998).
## 7 The CBB and SL contributions to the GBR
The spectral index of the GBR, derived in Section 4, is independent of the details of the spatial distribution of starlight. We have argued that the EGRET GBR data support the simple hypothesis of an electron spectral index that is independent of location. The predicted GBR index is then also independent of the magnitude of the electron spectrum as a function of position. In this section we use a simplified model of the electron and starlight distributions to compute the magnitude and angular dependence of the CMB and SL contributions to the GBR.
We adopt $`\mathrm{h}_\mathrm{e}=20`$ kpc (a value obtained from a rough fit of our results to the angularly-averaged fluence of the GBR) for the Gaussian scale height of the CR electron distribution of our galaxy in the direction perpendicular to the galactic plane. For the distribution in $`\rho `$ –the radial coordinate orthogonal to the galactic axis– we adopt a Gaussian scale height $`\rho _\mathrm{e}=35`$ kpc; the results are quite insensitive to this parameter. The EGRET GBR data are not precise enough to be “invertible”, that is, for the actual high-latitude CR-electron distribution (Gaussian, exponential or otherwise) to be disentangled; a fact to be rediscussed anon, in view of our results. The distance of the solar system to the galactic centre is $`\mathrm{d}_{_{}}8.5`$ kpc. The factor $`\mathrm{N}_0(\theta ,\varphi )`$ in Eq. (6), which describes the angular dependence of the GBR photons due to ICS on the (uniformly distributed) CMB, is:
$`\mathrm{N}_0(\mathrm{b},\mathrm{l})={\displaystyle _0^{\mathrm{}}}\mathrm{dr}\mathrm{n}_0\mathrm{Exp}\left[+\left({\displaystyle \frac{\mathrm{d}_{_{}}}{\rho _\mathrm{e}}}\right)^2\right]\mathrm{Exp}\left[\left({\displaystyle \frac{\mathrm{h}(\mathrm{r},\mathrm{b})}{\mathrm{h}_\mathrm{e}}}\right)^2\left({\displaystyle \frac{\rho (\mathrm{r},\mathrm{b},\mathrm{l})}{\rho _\mathrm{e}}}\right)^2\right],`$
$`\mathrm{h}(\mathrm{r},\mathrm{b})\mathrm{r}\mathrm{sin}\mathrm{b},`$
$`\rho (\mathrm{r},\mathrm{b},\mathrm{l})\left([\mathrm{r}\mathrm{cos}(\mathrm{b})\mathrm{cos}(\mathrm{l})\mathrm{d}_{_{}}]^2+[\mathrm{r}\mathrm{cos}(\mathrm{b})\mathrm{sin}(\mathrm{l})]^2\right)^{1/2},`$ (11)
where r is the distance in the direction along the line of sight.
It is difficult to model in detail the contributionn from ICS on starlight (Hunter et al. 1997, Sreekumar et al. 1998). But we are only concerned with this light at high galactic latitudes, since the diffuse GBR of interest to us is that measured by EGRET by masking the galactic plane and centre. We make a coarse estimate by approximating the Galaxy’s starlight as that produced by a source at its centre with the galactic luminosity $`\mathrm{L}_{}=2.3\times 10^{10}`$ $`\mathrm{L}_{_{}}`$ $`5.510^{55}`$ eV s<sup>-1</sup> (Pritchet & van den Bergh 1999). The starlight contribution in Eq. (5) is then of the same form as Eq. (11), with $`\mathrm{N}_0`$ traded for $`\mathrm{N}_{}`$ by the substitution:
$`\mathrm{n}_0{\displaystyle \frac{\mathrm{L}_{}}{4\pi \mathrm{c}ϵ_{}}}{\displaystyle \frac{1}{(\mathrm{r}^22\mathrm{r}\mathrm{d}_{_{}}\mathrm{cos}(\mathrm{b})\mathrm{cos}(\mathrm{l})+\mathrm{d}_{_{}}^2)}}.`$ (12)
For the CMB and starlight contributions to the GBR, averaged over the EGRET unmasked domain, we obtain, by integration of Eqs. (5), (6), (11), (12):
$$\frac{\mathrm{dF}_\gamma }{\mathrm{dE}}(2.41\pm 0.55)\times 10^3\left[\frac{\mathrm{E}}{\mathrm{MeV}}\right]^{2.10\pm 0.05}\mathrm{cm}^2\mathrm{s}^1\mathrm{sr}^1\mathrm{MeV}^1.$$
(13)
For scale heights $`\mathrm{h}_\mathrm{e}`$ and $`\rho _\mathrm{e}`$ similar to the ones adopted (20 and 35 kpc, respectively), the CMB and SL contributions are comparable in magnitude, the first scales approximately linearly with $`\mathrm{h}_\mathrm{e}`$ while the second is rather insensitive to this parameter. The contribution to the CMB from sunlight and external galaxies, discussed in Section 8 and 9, adds corrections of 6% and $`10\%`$ (respectively) to Eq. (13), the total result is shown in Fig. 2. The fitted value of $`\mathrm{h}_\mathrm{e}`$ is imprecise: the starlight to CMB ratio is proportional to $`ϵ_{}/ϵ_0`$ raised to a very poorly determined power, $`0.10\pm 0.05`$.
We can use our assumed Gaussian distribution of electrons in a halo, with vertical and radial scale heights $`\mathrm{h}_\mathrm{e}`$ and $`\rho _\mathrm{e}`$, to compute the diffuse $`\gamma `$-ray luminosity of our galaxy, which in our model is dominated by ICS on CMB and SL photons. Using Eqs. (5), (6), (11), (12) we obtain, for the luminosity in $`\gamma `$-rays of energy above E:
$`\mathrm{L}_\gamma (>\mathrm{E})\mathrm{L}_\gamma ^0(>\mathrm{E})+\mathrm{L}_\gamma ^{}(>\mathrm{E}),`$
$`\mathrm{L}_\gamma ^0(>\mathrm{E})=1.31\times 10^{40}\left[{\displaystyle \frac{\rho _\mathrm{e}}{35\mathrm{kpc}}}\right]^2\left[{\displaystyle \frac{\mathrm{h}_\mathrm{e}}{20\mathrm{kpc}}}\right]\left[{\displaystyle \frac{\mathrm{E}}{\mathrm{MeV}}}\right]^{0.10\pm 0.05}\mathrm{erg}/\mathrm{s},`$
$`\mathrm{L}_\gamma ^{}(>\mathrm{E})=3.56\times 10^{39}\left[{\displaystyle \frac{\mathrm{h}_\mathrm{e}}{20\mathrm{kpc}}}\right]\left[{\displaystyle \frac{1}{2\mathrm{u}}}\mathrm{ln}{\displaystyle \frac{1+\mathrm{u}}{1\mathrm{u}}}\right]\left[{\displaystyle \frac{\mathrm{E}}{\mathrm{MeV}}}\right]^{0.10\pm 0.05}\mathrm{erg}/\mathrm{s},`$ (14)
where $`\mathrm{u}\sqrt{1\mathrm{h}_\mathrm{e}^2/\rho _\mathrm{e}^2}`$. A future $`\gamma `$-ray telescope, such as GLAST, could possibly see the corresponding glow of Andromeda’s halo.
## 8 Sunlight contribution to the local GBR
We are only at a distance $`\mathrm{l}_{_{}}=1.5\times 10^{13}\mathrm{cm}`$ from the sun. This entails a small but non-negligible contribution to the locally-observed GBR, resulting from ICS off photons in the heliosphere. The corresponding photon flux is described by Eq. (7), with the substitution of $`ϵ_i`$ by the mean energy $`ϵ_{_{}}1.35\mathrm{eV}`$ of solar photons, and of $`\mathrm{N}_\mathrm{i}`$ by $`\mathrm{N}_{_{}}`$, the solar-photon column density along the line of sight. Let $`\theta _{_{}}`$ be the angle between the line of sight and the direction to the sun. Then:
$$\mathrm{N}_{_{}}(\mathrm{cos}\theta _{_{}})=\frac{\mathrm{L}_{_{}}}{4\pi \mathrm{c}\mathrm{l}_{_{}}ϵ_{_{}}}\left(\frac{\pi \theta _{_{}}}{\mathrm{sin}\theta _{_{}}}\right).$$
(15)
For a uniform $`\mathrm{cos}\theta _{_{}}`$ distribution during the EGRET data taking, the average column density is $`\overline{\mathrm{N}}_{_{}}=\pi \mathrm{L}_{_{}}/(16\mathrm{c}\mathrm{l}_{_{}}ϵ_{_{}})`$, resulting in a sunlight-induced GBR flux:
$$\frac{\mathrm{dF}_\gamma ^{}}{\mathrm{dE}}1.32\times 10^4\left[\frac{\mathrm{E}}{\mathrm{MeV}}\right]\mathrm{cm}^2\mathrm{s}^1\mathrm{sr}^1\mathrm{MeV}^1.$$
(16)
This contribution is roughly 6% of our galaxy’s result, Eq. (13). At $`\mathrm{E}_\gamma >75`$ GeV, the spectrum of Eq. (16) should steepen, since ICS should then be described by the Klein–Nishina cross section, and not by its low energy Thomson limit.
## 9 Extragalactic contribution to the GBR
To estimate this contribution, some concepts and numbers need to be recalled. Hubble’s constan’ is $`\mathrm{H}_0=100\mathrm{h}\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$, with $`\mathrm{h}0.65`$; $`\mathrm{\Omega }_\mathrm{m}`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$ are matter and vacuum cosmic densities in critical units: $`\mathrm{\Omega }\mathrm{\Omega }_\mathrm{m}+\mathrm{\Omega }_\mathrm{\Lambda }`$; $`y1+z`$ is the redshift factor. In a Friedman model, the time to redshift relation is $`\mathrm{dy}/\mathrm{dt}=\mathrm{H}_0\mathrm{f}(\mathrm{y})\mathrm{y}`$, with $`\mathrm{f}(\mathrm{y})[(1\mathrm{\Omega })\mathrm{y}^2+\mathrm{\Omega }_\mathrm{m}\mathrm{y}^3+\mathrm{\Omega }_\mathrm{\Lambda }]^{1/2}`$. The luminosity density of the local universe (Ellis 1997) is $`\rho __\mathrm{L}=(2.0\pm 0.4)\times 10^8\mathrm{h}\mathrm{L}_{_{}}\mathrm{Mpc}^3`$. The combination $`\rho __\mathrm{L}/\mathrm{L}_{}`$ provides an estimate of the average number density of ‘Milky-Way-equivalent’ galaxies. If the main sources of CRs are young supernova remnants or gamma-ray bursts, the CR production rate ought to be proportional (e.g. Wijers et al. 1997) to the star formation rate $`\mathrm{R}_{\mathrm{SFR}}[\mathrm{y}]`$, recently measured up to redshift $`\mathrm{z}4.5`$ (Steidel et al. 1998).
The energy of CMB photons up-scattered by electrons at ‘epoch y’ is proportional to $`\mathrm{T}(\mathrm{y})=\mathrm{y}\mathrm{T}_0`$ and it is subsequently redshifted by the same factor; hence the spectra from distant galaxies should have the same energy dependence as from our galaxy. The situation for SL photons is more complicated. Young galaxies are bluer than older ones, but this effect is overcompensated by the expansion redshift from a relatively low y, onwards. Yet, at the energies observed by EGRET, and for the redshift values of $`𝒪`$(1) that dominate the extragalactic contribution, all these blue- and red-shifts simply relocate the photon energy, while roughly maintaining the slope of the spectrum. For the sum of all galaxies, we estimate:
$$\frac{\mathrm{dF}_\gamma ^{^{\mathrm{EG}}}}{\mathrm{dE}}\frac{1}{4\pi }\frac{\mathrm{dL}_\gamma }{\mathrm{E}\mathrm{dE}}\frac{\rho __\mathrm{L}}{\mathrm{L}_{}}\frac{\mathrm{c}}{\mathrm{H}_0}_1\frac{\mathrm{R}_{\mathrm{SFR}}(\mathrm{y})}{\mathrm{R}_{\mathrm{SFR}}(0)}\frac{\mathrm{y}}{\mathrm{f}(\mathrm{y})}\frac{\mathrm{dy}}{\mathrm{y}^3},$$
(17)
where $`\mathrm{dL}_\gamma /\mathrm{dE}`$ is to be obtained from the luminosity of a Milky-Way-like galaxy, Eq. (14). For $`\mathrm{R}_{\mathrm{SFR}}[\mathrm{y}]`$ we interpolate the summary values of Steidel et al. (1998). In writing Eq. (17) we have ignored the fact that, above $`\mathrm{E}10`$ GeV, absorption by $`e^+e^{}`$ production on the IR-to-UV background becomes relevant (Salamon & Stecker 1998), so that the extragalactic contribution should be quenched.
For $`\mathrm{\Omega }=\mathrm{\Omega }_\mathrm{m}=1`$ the value of the integral in Eq. (17) is $`0.82`$; it increases to $`1.08`$ for a currently more fashionable universe with $`\mathrm{\Omega }=1`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$, $`\mathrm{\Omega }_\mathrm{m}=0.3`$. For the latter case, the result is:
$$\frac{\mathrm{dF}_\gamma ^{^{\mathrm{EG}}}}{\mathrm{dE}}=2.48\times 10^4\left[\frac{\mathrm{E}}{\mathrm{MeV}}\right]^{2.10\pm 0.05}\mathrm{cm}^2\mathrm{s}^1\mathrm{sr}^1\mathrm{MeV}^1,$$
(18)
roughly $`10\%`$ of our galaxy’s angularly-averaged result, Eq. (13).
## 10 Detailed comparison with the EGRET data
Our predictions for the magnitude of the GBR and its directional dependence on b and l are shown in Figs. 5 and 6. In Fig. 5 we display separately the contributions from ICS off CMB and SL photons in our galaxy, as well as the uniformly distributed sunlight and extragalactic components. In Fig. 7 we compare the total GBR flux:
$$\frac{\mathrm{dF}_\gamma }{\mathrm{dE}}=\frac{\mathrm{dF}_\gamma ^0}{\mathrm{dE}}+\frac{\mathrm{dF}_\gamma ^{}}{\mathrm{dE}}+\frac{\mathrm{dF}_\gamma ^{}}{\mathrm{dE}}+\frac{\mathrm{dF}_\gamma ^{^{\mathrm{EG}}}}{\mathrm{dE}},$$
(19)
obtained by summing Eqs. (7), Eq. (16) and Eq. (13), with the EGRET data. Our result is a satisfactory fit to the observed magnitude and angular trend of the GBR ($`\overline{\chi }^2=0.98`$), a vast improvement over the result for a constant (extragalactic) ansatz, for which $`\overline{\chi }^2=2.6`$. Although this agreement would be more meaningful, had we used a more realistic model of starlight, a more careful treatment may be premature, for the EGRET error bars are large enough to accommodate considerable variations in the input modelling. In a previous analysis (Dar et al. 1999), for instance, we obtained a similarly good fit with an assumed constant-density, spherical CR-electron halo of radius 25 kpc, for which the results have the advantage of being simple analytical functions.
We have neglected various putative extragalactic contributions to the GBR. Blazars, because of their beamed emission, may not be very relevant. But CR electrons injected directly into intergalactic space by active galactic nuclei, radio galaxies or gamma ray bursters, may give rise to a contribution of comparable magnitude and shape to that of the CR electrons in external galaxies. These or other potential sources of GBR photons may imply that our parameters $`\mathrm{h}_\mathrm{e}`$ and $`\rho _\mathrm{e}`$ have been overestimated. But this effect cannot be very large, given our success at describing the non-trivial angular dependence of the EGRET data.
## 11 Conclusions and predictions
We have presented a simple understanding of the relation between the spectral indices of cosmic-ray protons, electrons and the GBR. Accepting the possibility that the CR-electron distribution in our galaxy may have a scale height larger than conventionally believed, we have also argued that the bulk of the GBR could originate in our own galaxy. Our modelling is extremely simplistic, but quite successful.
The predictions specific to our scenario are:
* The GBR should reflect the asymmetry of our off-centre position in the Galaxy.
* The halo of Andromeda should shine in gamma rays above a few MeV, with a luminosity comparable to that in Eq. (14). Likewise, very nearby star-burst Galaxies, such as M82, and radio galaxies with large CR production rates, such as Cygnus A, may be visible in gamma rays.
* If the CR-proton and electron acceleration mechanisms are the same, the existence of a knee in the observed proton spectrum translates into a related result for the power index $`\beta _\mathrm{e}`$ of the electron spectrum, which should steepen above $`\mathrm{E}1.6`$ TeV by $`\mathrm{\Delta }\beta 1/4`$.
* The GBR spectrum should not have the sharp cutoff, above $`\mathrm{E}100`$ GeV, expected (Salamon & Stecker 1998) for cosmological sources. But it should nonetheless steepen around 10–100 GeV, because of the anticipated “knee” in the electron spectrum and of the energy-dependence of the Klein-Nishina cross section.
These features of our scenario should be testable when the next generation of cosmic-ray and $`\gamma `$-ray satellites (AMS-02 and GLAST) are operational, hopefully by 2005. In spite of their maturity, cosmic-ray physics and $`\gamma `$-ray astrophysics are still young, and thriving.
ACKNOWLEDGEMENTS
We are indebted to G. Bignami, S. Dado and G. Raffelt for discussions and to I. Moskalenko and A. Strong for permission to reproduce their results in our Fig. 5.
REFERENCES
Barwick S. W. et al., 1998, ApJ, 498, 779
Bednarz J., Ostrowski M., 1998, PRL, 80, 3911
Berezinskii V. S. et al., 1990, Astrophysics of cosmic rays (North Holland, Amsterdam, 1990)
Bhattacharjee P., Sigl G., 2000, Phys. Rep., 327, 109
Bignami G. et al., 1979, ApJ, 232, 649
Bird D. J. et al., 1995, ApJ, 441, 144.
Broadbend A., Haslam C. G. T., Osborne, J. L., 1989, MNRAS, 237, 381
Chiang J., Mukerjee R., 1998, ApJ, 496, 772
Cohen A., De Rújula A., Glashow S. L., 1998, ApJ, 495, 539
Connell J. J., 1998, ApJ, 501, L59
Dar A., 1998, astro-ph/9809163, in Proceedings of the Rencontres de la Vallée d’Aoste, 1998 (ed. M. Greco), Frascati Physics Series, INFN Pubs, page 23
Dar A., De Rújula A., Antoniou N., 1999, astro-ph/9901004
Dar A., Plaga R., 1999, A&A, 349, 259
Dixon D. D. et al., 1998, New Astron. 3, 539
Duric N., Irwin J., Bloemen H., 1998, A& A, 331, 428
Ellis R. S., 1997, ARA&A, 35, 389
Evenson P., Meyers P., 1984, J. Geophys. Res., 89 A5, 2647
Felten J. E., Morrison P., 1966, ApJ, 146, 686
Ferrando P. et al., 1996, A&A 316, 528
Fixsen D.J. et al., 1996, ApJ, 473, 576
Gnedin N. Y., Ostriker J. P., 1992, ApJ, 400, 1
Golden R. L. et al., 1984, ApJ, 287, 622
Golden R. L., et al., 1994, ApJ, 436, 739
Hartmann D. H., 1995, ApJ, 447, 646
Hawking S. W., 1977, Scientific American, 236, 34
Hunter S. D. et al., 1997, ApJ, 481, 205
Kazanas D., Protheroe J. P., 1983, Nature, 302, 228
Lukasiak A. et al., 1994, ApJ, 423, 426
Mather J. C. et al., 1993, ApJ, 432, L15
Moskalenko I. V., Strong A. W., Reimer O., 1998, astro-ph/9811221
Moskalenko I. V., Strong A. W., 2000, ApJ, 528, 357
Nishimura J. et al., 1980, ApJ, 238, 394
Page D. N., Hawking S. W., 1976, ApJ, 206, 1
Pohl M., Esposito J. A., 1998, ApJ, 507, 327
Prince T. A., 1979, ApJ, 227, 676
Pritchet C. J., van den Bergh S., 1999, AJ, 118, 833
Rudaz S., Stecker F. W., 1991, ApJ, 368, 40
Salamon M. H., Stecker F. W., 1998, ApJ, 493, 547
Seo E. S., Ptuskin V. S., 1994, ApJ, 431, 705
Silk J., Srednicki M., 1984, PRL, 53, 264
Sreekumar P. et al., 1998, ApJ, 494, 523
Stecker F. W., Salamon M. H., 1996, ApJ, 464, 600
Stecker F. W., Morgan D. L., Bredekamp J., 1971, PRL, 27, 1469
Steidel C.C. et al., 1999, ApJ, 519, 1
Strong A., Moskalenko I. V., 1998, ApJ, 509, 212
Strong A. W. et al., 1997, in Proceedings of the 4th Compton Symposium, AIP, 410, 1198
Strong A. W., Moskalenko I. V., Reimer O., 1998, astro-ph/9811296
Swordy S. P. et al., 1990, ApJ, 330, 625
Takeda M. et al., 1998, PRL, 81, 1163
Tang K. K., 1984, ApJ, 278, 881
Thompson D. J., Fichtel C. E., 1982, A&A, 109, 352
Webber W. R., 1997, Sp. Sci. Rev., 81, 107
Wiebel-Sooth B., Biermann, P. L., 1998, Landolt-Börnstein, (Springer Verlag, Heidelberg 1998, in press)
Wijers R. A. M. J. et al., 1997, MNRAS, 294, L13
|
warning/0005/math0005231.html
|
ar5iv
|
text
|
# Untitled Document
Occupation Time Fluctuations in Branching Systems<sup>1</sup>
| D.A. Dawson<sup>2</sup> | L.G. Gorostiza<sup>3</sup> | and | A. Wakolbinger<sup>4</sup> |
| --- | --- | --- | --- |
## Abstract
We consider particle systems in locally compact Abelian groups with particles moving according to a process with symmetric stationary independent increments and undergoing one and two levels of critical branching. We obtain long time fluctuation limits for the occupation time process of the one–and two–level systems. We give complete results for the case of finite variance branching, where the fluctuation limits are Gaussian random fields, and partial results for an example of infinite variance branching, where the fluctuation limits are stable random fields. The asymptotics of the occupation time fluctuations are determined by the Green potential operator $`G`$ of the individual particle motion and its powers $`G^2,G^3`$, and by the growth as $`t\mathrm{}`$ of the operator $`G_t=_0^tT_s𝑑s`$ and its powers, where $`T_t`$ is the semigroup of the motion. The results are illustrated with two examples of motions: the symmetric $`\alpha `$–stable Lévy process in $`\text{}^d`$ $`(0<\alpha 2)`$, and the so called $`c`$–hierarchical random walk in the hierarchical group of order $`N`$ ($`0<c<N`$). We show that the two motions have analogous asymptotics of $`G_t`$ and its powers that depend on an order parameter $`\gamma `$ for their transience/recurrence behavior. This parameter is $`\gamma =d/\alpha 1`$ for the $`\alpha `$–stable motion, and $`\gamma =\mathrm{log}c/\mathrm{log}(N/c)`$ for the $`c`$–hierarchical random walk. As a consequence of these analogies, the asymptotics of the occupation time fluctuations of the corresponding branching particle systems are also analogous. In the case of the $`c`$–hierarchical random walk, however, the growth of $`G_t`$ and its powers is modulated by oscillations on a logarithmic time scale.
Key words: multilevel branching particle system, occupation time, fluctuation, Green potential, weak and strong transience, stable Lévy process, hierarchical random walk, critical dimensions. <sup>1</sup> Research partially supported by a Max Planck award to D.A. Dawson, CONACyT grant 27932–E (Mexico), and a DFG
grant (Germany).
<sup>2</sup> The Fields Institute, Toronto, Canada, e.mail: ddawson@fields.utoronto.ca
<sup>3</sup> Centro de Investigación y de Estudios Avanzados, A.P. 14–740, Mexico 07000 D.F., Mexico,
e.mail: gortega@servidor.unam.mx
<sup>4</sup> Johann Wolfgang Goethe–Universität, Frankfurt am Main, Germany, e.mail: wakolbin@math.uni-frankfurt.de
Table of contents
1. Intoduction
2. General Notions, Results and Comments
2.1. The individual motion: powers of the Green potential, strong and weak transience
2.2. Main results: Occupation time fluctuations for $`k`$-level critical binary branching particle systems
($`k=0,1,2`$)
2.3. Comments on the assumptions and the results
2.4. Order of transience and recurrence, asymptotics of powers of $`G_t`$, and some special growth functions
2.5. Infinite variance branching
2.6. Occupation time fluctuations of superprocesses
3. Two examples of individual motions: Symmetric $`\alpha `$-stable process and $`c`$-hierarchical random walk
3.1. Transience/recurrence properties of the motions
3.2. Occupation time fluctuations limits
3.3. Comments on the results
4. Definitions of constants and functions for the examples
4.1. Notation for $`\alpha `$-stable motion 4.2. Notation for $`c`$-hierarchical random walk
5. Proofs
5.1. Asymptotics of the powers of $`G_t`$
5.2. Main results 5.3. Examples
5.4. Conditions for the results on infinite variance branching
Appendix
A.1. Background on 1- and 2-level branching systems
A.2. The Palm formula
A.3. Tree representation of the Palm measures of $`R_t^1`$ and $`R_{\mathrm{}}^1`$
A.4. Second and third moments of $`R_{\mathrm{}}^1`$
Acknowledgements
References
1. INTRODUCTION Consider a particle system described by a random counting measure $`X_t`$ on a space of sites $`S`$, with the same intensity measure $`EX_t`$ for all $`t`$ which is denoted by $`\rho `$, and such that $`X_t`$ converges in distribution as $`t\mathrm{}`$ towards an equilibrium state which also has the same intensity $`\rho `$, i.e., the system is persistent. Then under mild conditions the occupation time fluctuation $`Y_t=_0^t(X_s\rho )𝑑s`$ obeys a law of large numbers, i.e. $`{\displaystyle \frac{1}{t}}Y_t0`$ as $`t\mathrm{}`$ (see e.g. Méléard and Roelly<sup>(32)</sup> for a branching particle system in $`\text{}^d`$). The question for which norming $`a_t`$ does a non–trivial limit of $`{\displaystyle \frac{1}{a_t}}Y_t`$ exist in distribution as $`t\mathrm{}`$ and what is the limiting random field depends on more specific properties of the system. In this paper we investigate this question for branching particle systems in locally compact Abelian groups, where the individual particle motion is a process with symmetric stationary independent increments, and for the so called “2–level” branching systems in which not only the individual particles but also whole families of particles undergo critical branching. Multilevel branching systems were introduced by Dawson and Hochberg<sup>(9)</sup>, and they have been studied by several authors: Dawson et al<sup>(10)</sup>, Gorostiza<sup>(16)</sup>, Gorostiza et al<sup>(17)</sup>, Greven and Hochberg<sup>(24)</sup>, Hochberg<sup>(25)</sup>, Hochberg and Wakolbinger<sup>(26)</sup>, Wu<sup>(40)</sup>.
For a transient motion and finite variance branching, the simple branching particle system converges as $`t\mathrm{}`$ to a Poisson system of independently evolving “clans”, each of which contributes to the occupation time. The asymptotics of the occupation time fluctuations should be determined by the space–time correlations within single clans. Also, the growth of the occupation time fluctuations as $`t\mathrm{}`$ should depend on whether there are long time dependencies caused by recurrent visits of single clans to bounded sets.
Let us first recall some known results. For a critical finite variance branching Brownian system $`X_t`$ in $`\text{}^d`$, started off from a Poisson system with Lebesgue intensity $`\lambda `$, the right norming $`a_t`$ for the occupation time fluctuation is $`t^{3/4}`$ for $`d=3`$, $`(t\mathrm{log}t)^{1/2}`$ for $`d=4`$, and $`t^{1/2}`$ for $`d>4`$ (see Cox and Griffeath<sup>(5)</sup>, and also Iscoe<sup>(27)</sup>, where the corresponding superprocess scenario is treated). We refer to this as the 1–level branching case.
The same normings appear for the occupation time fluctuations of Poisson systems of Brownian particles without branching, which we call 0–level systems, but two dimensions lower (Cox and Griffeath<sup>(4)</sup>, Deuschel and Wang<sup>(14)</sup>), and we shall see that they also appear in the $`2`$–level branching case, but now two dimensions higher than in the $`1`$–level case.
There is an apparent relation between the critical dimension for transience of the motion and the critical dimension for the classical $`t^{1/2}`$–norming of the occupation time fluctuations: for Brownian particle systems without branching, 2 is the critical dimension above which the occupation time fluctuations have the classical norming and it is also the dimension above which the particle motion is transient. In the 1–level branching case, 4 is the critical dimension above which the occupation time fluctuations have the classical norming and it is also the dimension above which the equilibrium clans are transient (in the sense that they eventually leave each bounded region of $`\text{}^d`$ forever (Stöckl and Wakolbinger<sup>(38)</sup>). We shall see that an analogous result holds for the $`2`$–level case.
One of our main objectives is to put these results in a general context for branching systems in locally compact Abelian groups, which clarifies the role played by the Green potential operator $`G`$ of the particle motion and the (operator) powers $`G^k`$ of $`G`$ in relation with the various levels $`k=0,1,2,\mathrm{}`$ of branching. A key role will be played by the level $`k`$ transience and recurrence properties of the motion, $`k=0,1,2,\mathrm{}`$ defined in Section 2. In this paper we will treat only branching levels $`k=0,1,2`$, but the results show a pattern which allows one to guess what the form of the results would be for systems with higher levels of branching. The analysis of 2–level systems is considerably more difficult than that of 1–level systems due to the dependencies among the particles caused by the simultaneous branching of families of particles.
Roughly speaking, the bigger the branching level $`k`$ is, the more long range dependencies are introduced into the system. These dependencies increase the mass fluctuations in the system, whereas a strong spreading out of mass by the particle motion has a smoothing effect on the mass fluctuations.
In the case of finite variance branching (where we will restrict for simplicity to binary branching) it turns out that finiteness of $`G^k`$ corresponds to existence of the $`k`$–level branching equilibrium clans, and then $`G^{k+1}=\mathrm{}`$ corresponds to a long time dependence of the visits of single $`k`$–level clans to bounded sets. In the latter case, a crucial feature of these models is that the growth of the $`(k+1)`$–st power of the operator $`G_t=_0^tT_s𝑑s`$ as $`t\mathrm{}`$, where $`T_t`$ denotes the semigroup of the individual motion, determines the right norming $`a_t`$ for the occupation time fluctuations of the $`k`$–level system, whereas for finite $`G^{k+1}`$ their right norming is the classical $`t^{1/2}`$. In the case of finite $`G^{k+1}`$ the covariance of the limiting Gaussian field of the occupation time fluctuations of the $`k`$–level system contains terms induced by direct ancestry and by level $`j`$–relationship, $`1jk`$.
For infinite variance “$`(1+\beta )`$–branching” $`(0<\beta <1)`$, we have so far results only for the case with $`t^{1/(1+\beta )}`$–norming. Then the norming is determined by the highest level of branching, and the occupation time fluctuations converge to stable random fields.
For the superprocess limits of the 1– and the 2–level branching particle systems the corresponding occupation time fluctuation results are basically analogous to those for the particle systems and even somewhat simpler.
We will focus on two examples of particle motions $`W_t`$: the symmetric $`\alpha `$–stable Lévy process in $`\text{}^d,0<\alpha 2`$ (including Brownian motion, $`\alpha =2`$), and a “hierarchical” random walk in $`\mathrm{\Omega }_N`$, the hierarchical group of order $`N`$, which is a direct sum of a countable number of copies of $`\text{}_{N1}`$, i.e. $`\mathrm{\Omega }_N=\{x=(x_1,x_2,\mathrm{})|x_i\{0,1,\mathrm{},N1\},x_i0`$ except for finitely many i$`\}`$.
For the symmetric $`\alpha `$–stable process in $`\text{}^d`$ and $`k0`$, $`G_t^{k+1}`$ has a power growth in $`t`$ if $`d/(k+1)<\alpha `$, a logarithmic growth if $`\alpha =d/(k+1)`$, and $`G^{k+1}`$ is finite if $`\alpha <d/(k+1)`$.
For the random walk in $`\mathrm{\Omega }_N`$ we consider a probability of jumping a distance $`i`$ proportional to $`(c/N)^{i1}`$, where $`c`$ is a constant such that $`0<c<N`$. (Here, the distance between two elements $`(x_i)`$ and $`(y_i)`$ of $`\mathrm{\Omega }_N`$ is defined to be the highest index $`i`$ for which $`x_i`$ and $`y_i`$ are different). Since the random walk is characterized by this “mobility parameter” $`c`$, we will call it the $`c`$hierarchical random walk in $`\mathrm{\Omega }_N`$. In this case, for $`k0`$, $`G_t^{k+1}`$ has a power growth in $`t`$ if $`c<N^{k/(k+1)}`$, a logarithmic growth if $`c=N^{k/(k+1)}`$, and $`G^{k+1}`$ is finite if $`c>N^{k/(k+1)}`$. In this example the growths of $`G_t`$ and its powers have also oscillating modulations in a logarithmic time scale. However, the oscillations vanish in the cases of logarithmic growth. Oscillatory phenomena have also been observed in another class of random walks on groups which are direct sums of a countable number of copies of a discrete group (Cartwright<sup>(2)</sup>), and in random walks on the Sierpiński graph (Barlow and Perkins<sup>(1)</sup>, Grabner and Woess<sup>(23)</sup> and references therein). A basic reference for random walks on Abelian groups is the paper by Kesten and Spitzer<sup>(29)</sup>.
The values of the parameters $`\alpha `$ (or $`d`$) and $`c`$ which correspond to logarithmic growths of the powers $`G_t^{k+1}`$ will sometimes be referred to as “critical”, since they are boundaries between intervals with different regimes of the longtime behavior of the systems.
We will show that the asymptotics of the occupation time fluctuations of the branching systems are analogous for the two examples. A key role is played by a constant $`\gamma >1`$, which is $`\gamma =d/\alpha 1`$ for the $`\alpha `$-stable system in $`\text{}^d`$, and $`\gamma =\mathrm{log}c/\mathrm{log}(N/c)`$ for the $`c`$–hierarchical system in $`\mathrm{\Omega }_N`$. For $`\gamma (1,0)`$, $`G_t`$ grows like $`t^\gamma `$, for $`\gamma =0`$, $`G_t`$ grows logarithmically, and for $`\gamma >0`$, $`GG_t`$ decays like $`t^\gamma `$. Thus, $`\gamma `$ is an order parameter for the transience/recurrence behavior of the motion, and we will refer to the three above mentioned cases as recurrence of order –$`\gamma `$, critical recurrence, and transience of order $`\gamma `$. It turns out that $`G_t^{k+1}`$ grows like $`t^{k\gamma }`$ if $`k>\gamma `$, grows logarithmically if $`k=\gamma `$, and $`G^{k+1}`$ is finite if $`k<\gamma `$.
For the symmetric $`\alpha `$–stable processes in $`\text{}^d`$, $`\gamma `$ is restricted to $`[(d2)/2,\mathrm{})`$, whereas for the $`c`$–hierarchical random walks in $`\mathrm{\Omega }_N`$, $`\gamma `$ can range over the entire interval $`(1,\mathrm{})`$. In this sense the hierarchical random walks are a richer class of models. Choosing $`c=N^{1\alpha /d}`$, the $`c`$–hierarchical random walk in $`\mathrm{\Omega }_N`$ has the same order of transience/recurrence as the $`\alpha `$–stable process in $`\text{}^d`$, and this allows to think about non–integer dimensions $`d`$.
The idea of using hierarchical systems as models of systems with non–integer dimension was first introduced in the context of statistical physics. A model of ferromagnetic behavior involving the case of
$`N=2`$ is known as Dyson’s hierarchical model and has been used by Sinai<sup>(36)</sup> as a framework in which to carry out a rigorous renormalization group analysis following the ideas of Wilson<sup>(39)</sup>. In the case of ferromagnetism, $`4`$ is the critical dimension and the hierarchical group has been used to study large scale fluctuations near the critical point in $`4\epsilon `$ dimensions. In the case of 1–level branching Brownian motion in $`\text{}^d`$, the dimension $`2`$ is the critical dimension for the persistence/extinction dichotomy. In Dawson and Greven<sup>(7)</sup>, 1–level hierarchical branching random walks (indexed by a sequence $`(c_j)`$ rather than just one parameter $`c`$) have been analyzed in the so called mean–field limit $`N\mathrm{}`$ in the “nearly 2–dimensional analogue”. Since the dimension 4 is the critical dimension for the occupation time fluctuations of 1–level branching Brownian motions and for the long time behavior of 2–level branching Brownian motions in $`\text{}^d`$, it is conceivable that the hierarchical mean–field limit can be used to carry out a similar analysis of these phenomena “near dimension 4”.
The outline of the paper is as follows: Section 2 presents the general results for the particle systems, and the results for the superprocesses are mentioned. Sections 3 and 4 are devoted to the two above mentioned examples of individual motions. Section 4 contains a list of constants and functions that appear in the results of Section 3. The proofs are given in section 5. An Appendix contains definitions and background on 1– and 2–level branching systems, and some tools. 2. GENERAL NOTIONS, RESULTS AND COMMENTS 2.1. The individual motion: powers of the Green potential, strong and weak transience
We consider as a space of sites a locally compact Abelian group $`S`$ with Haar measure $`\rho `$, and as individual particle motion a process $`W_t`$ with stationary independent increments. We assume that for each $`s>0,W_sW_0`$ has a strictly positive symmetric density with respect to $`\rho `$, and that $`W_t`$ has càdlàg paths.
Let us fix some notation. The function spaces $`𝒞_c(S),𝒞_\tau (S),𝒞_c^+(S),𝒞_\tau ^+(S)`$, and the measure spaces $`_\tau (S),𝒩_\tau (S)`$, where $`\tau `$ is a reference function on $`S`$, are defined in the Appendix. For $`\mu _\tau (S)`$ and $`\phi 𝒞_\tau (S)`$, we write $`\mu ,\phi =\phi 𝑑\mu `$, and we denote
$$(\phi ,\psi )_\rho =_S\phi \psi 𝑑\rho ,\phi ,\psi 𝒞_\tau (S).$$
We designate by $`T_t`$ the semigroup of $`W_t`$. Note that $`\rho `$ is $`T_t`$–reversible, i.e., $`(\phi ,T_t\psi )_\rho =(T_t\phi ,\psi )_\rho `$ for all $`\phi ,\psi 𝒞_\tau (S)`$ and $`t>0`$, which implies that $`\rho `$ is $`T_t`$–invariant for each $`t>0`$.
The Green potential operator $`G`$ of the motion is defined by
$$G\phi =_0^{\mathrm{}}T_t\phi 𝑑t,\phi 𝒞_c(S).$$
Together with $`T_t`$, also $`G`$ is self–adjoint with respect to $`\rho `$, so
$$\rho ,(G\phi )(G\psi )=(\phi ,G^2\psi )_\rho .$$
Observe that the (operator) powers of $`G`$ are given by $`G^k\phi ={\displaystyle \frac{1}{(k1)!}}{\displaystyle _0^{\mathrm{}}}t^{k1}T_t\phi 𝑑t,k2.(\mathrm{2.1.1})`$
The quantities $`G^{k+1}(x,B):=G^{k+1}1_B(x),xS,B`$ a measurable subset of $`S`$, $`k=0,1,\mathrm{}`$, have an interpretation in terms of occupation times of a mass flow with continuous birth of mass, which will be helpful later on in the genealogical picture of 1– and 2–level branching systems. Consider the case $`k=1`$ and imagine an initial “parent” unit mass at $`xS`$, which evolves according to the flow $`T_t`$. This parent mass generates its own amount of “daughter” mass at its own site continuously at rate 1, and this daughter mass is again transported by the flow $`T_t`$. Then $`G^2(x,B)`$ is simply the total occupation time of the daughter mass in $`B`$; this is immediate from the semigroup property. To interpret $`G^{k+1}(x,B)`$ in a similar way, imagine an initial “k–level” unit mass at $`x`$ which evolves according to the flow $`T_t`$. For every $`j=k,\mathrm{},1`$, the j–level mass generates its own amount of $`(j1)`$–level mass, which is again transported by the flow $`T_t`$. Then $`G^{k+1}(x,B)`$ results as the total occupation time of 0–level mass in $`B`$.
We use the notation $`||||`$ for the supremum norm. Definition 2.1.1. (a) For $`k0`$, we say that $`W_t`$ is level k transient if
$$G^{k+1}\phi <\mathrm{}\text{ for}\phi 𝒞_c^+(S),$$
and level k recurrent if
$$G^{k+1}\phi \mathrm{}\text{for}\phi 𝒞_c^+(S),\phi \mathrm{\hspace{0.17em}0}.$$
(b) For $`k0`$, we say that $`W_t`$ is level $`k`$ strongly transient if it is level $`k+1`$ transient, and $`W_t`$ is level k weakly transient if it is level $`k`$ transient and level $`k+1`$ recurrent. Note that level $`0`$ transience and recurrence are (because of the assumed irreducibility of $`W_t`$) just the ordinary notions of transience and recurrence, and note also that level $`0`$ strong and weak transience coincide with the notions of strong and weak transience as defined, e.g., in Port and Stone<sup>(33)</sup>. Clearly, level $`k`$ transience implies level $`j`$ transience for $`j<k`$, and level $`k`$ recurrence implies level $`j`$ recurrence for $`j>k`$. In terms of the interpretation given above, level $`k`$ transience (resp. recurrence) means that a $`k`$–level parent unit mass sends up to time infinity a finite (resp. infinite) amount of 0–level daughter mass to any bounded region. Definition 2.1.2. (a) We define the operator
$$G_t\phi =_0^tT_s\phi 𝑑s,\phi 𝒞_\tau (S),t>0,$$
and denote by $`G_t^k`$ the $`k`$–th (operator) power of $`G_t`$, $`k2`$.
(b) For transient motion, we define the the bilinear form
$$R_t(\phi ,\psi )=(\phi ,(GG_t)\psi )_\rho ,\phi ,\psi 𝒞_\tau (S),t>0.$$
Note that for each $`k1`$ and $`\phi 𝒞_c^+(S),\phi 0`$, $`_0^{\mathrm{}}t^{k1}R_t(\phi ,\phi )𝑑t<\mathrm{}`$ ($`\mathrm{resp}.=\mathrm{}`$) if the motion is level $`k`$ transient (resp. level $`k`$ recurrent). Definition 2.1.3. (a) Let $`H`$ and, for each $`t>0,H_t`$ be positive definite bilinear forms on $`𝒞_c(S)`$, and let $`f:[T,\mathrm{})\text{}_+`$ for some $`T>0`$. We write $`H_tf_tH`$ if
$$\frac{1}{f_t}H_t(\phi ,\psi )H(\phi ,\psi )\text{as}t\mathrm{},\phi ,\psi 𝒞_c^+(S).$$
If $`Q_t`$ is a linear operator from $`𝒞_c(S)`$ into $`𝒞_\tau (S)`$, the notation $`Q_tf_tH`$ means that $`H_tf_tH`$ holds for $`H_t(\phi ,\psi ):=(\phi ,Q_t\psi )_\rho `$. (b) We call $`f_t`$ a growth function if it is increasing and $`\underset{t\mathrm{}}{lim}f_t=\mathrm{}`$. Growth functions will be used to characterize the growth of $`G_t`$ and its powers. The growth functions we shall encounter in the examples are of the form $`f_t=t^\zeta h_t`$, where $`\zeta (0,1)`$ and $`h_t`$ is either identically equal to 1 or a slowly oscillating function, or $`f_t=\mathrm{log}t`$. Part (a) of the definition will also be used to characterize the “order of transience” of $`W_t`$ in terms of $`R_t`$ with $`f_t0.`$ 2.2. Main results: Occupation time fluctuations for $`k`$–level critical binary branching particle systems ($`k=0,1,2`$) We consider the following particle systems in $`S`$ with the individual particles moving independently according to the process $`W_t`$:
* 0–level system: The system starts off from a Poisson system with intensity $`\rho `$.
* 1–level branching system: The motion is transient, the system starts off from a Poisson system with intensity $`\rho `$, and the particles undergo critical binary branching at rate $`V`$. We recall that this system has a (infinitely divisible) “Poisson type” equilibrium (in the sense of Liemant et al<sup>(31)</sup>, section 2.3), and we denote by $`R_{\mathrm{}}^1`$ its canonical measure, which we call equilibrium canonical measure. Note that $`R_{\mathrm{}}^1`$ has intensity $`\rho `$ (see the Appendix, (A.1.11)).
* 2–level branching system: The motion is strongly transient, the system starts off from a Poisson system of “2–level particles” with intensity $`R_{\mathrm{}}^1`$, individual particles undergo critical binary branching at rate $`V_1`$ and clans undergo critical binary branching at rate $`V_2`$. Note that $`R_{\mathrm{}}^1`$ is an invariant measure for the $`1`$–level dynamics, just as $`\rho `$ is an invariant measure for the $`0`$–level dynamics (however, $`R_{\mathrm{}}^1`$ is not reversible for the 1–level dynamics).
We refer the reader to Gorostiza<sup>(16)</sup>, and Hochberg and Wakolbinger<sup>(26)</sup> for a detailed description of the dynamics of 2–level branching systems. The necessary background for the present paper is given in the Appendix and should be consulted as the need arises.
For the three particle systems above, $`X_t`$ stands for the empirical measure of the locations of all the particles present at time $`t`$. Thus, $`X_t`$ is a random point measure on $`S`$. In the 2–level case $`X_t`$ corresponds to the aggregated system, i.e., the particles are counted as “1–level particles” independently of what clans (“2–level particles”) they belong to. In each one of these systems the intensity is preserved, i.e., $`EX_t=\rho `$ for all $`t>0`$, as a consequence of the initial conditions and the criticality in the branching cases. We consider the occupation time fluctuation, which is the random signed measure $`Y_t`$ on $`S`$ defined by
$$Y_t=_0^t(X_s\rho )𝑑s,t>0.$$
The following theorems describe the asymptotic distribution of $`Y_t`$ as $`t\mathrm{}`$ for the $`k`$–level systems, $`k=0,1,2,`$ described above. All the convergence assertions are understood to be in distribution as $`t\mathrm{}`$, all the test functions belong to $`𝒞_c^+(S)`$, i.e., all random fields are considered over $`𝒞_c^+(S)`$. The results for the 0–level system are basically known in special cases, but we include them for completeness and because they are the initial step in the multilevel ladder. Theorem 2.2.1. Let $`X_t`$ be the 0–level system.
* If the motion $`W_t`$ is transient, then $`t^{1/2}Y_t`$ converges to a Gaussian field with covariance functional $`2(\phi ,G\psi )_\rho `$.
* If the motion $`W_t`$ is recurrent with $`G_tf_tH`$ for some growth function $`f_t`$, then $`(_0^tf_s𝑑s)^{1/2}Y_t`$ converges to a Gaussian field with covariance functional $`2H(\phi ,\psi )`$.
Theorem 2.2.2. Let $`X_t`$ be the 1–level branching system.
* If $`W_t`$ is strongly transient, then $`t^{1/2}Y_t`$ converges to a Gaussian field with covariance functional
$$(\phi ,(2G+VG^2)\psi )_\rho =2(\phi ,\left(I+\frac{V}{2}G\right)G\psi )_\rho .$$
* If $`W_t`$ is weakly transient with $`G_t^2f_tH`$ for some growth function $`f_t`$, then $`(_0^tf_s𝑑s)^{1/2}Y_t`$ converges to a Gaussian field with covariance functional $`VH(\phi ,\psi )`$.
Theorem 2.2.3. Let $`X_t`$ be the 2–level branching system.
* If $`W_t`$ is level $`1`$ strongly transient and if there exists $`\delta >5/2`$ such that $`T_t\phi =O(t^\delta )\text{as}t\mathrm{},\phi 𝒞_c^+(S),(\mathrm{2.2.1})`$ then $`t^{1/2}Y_t`$ converges to a Gaussian field with covariance functional
$$(\phi ,\left(2G+(V_1+V_2)G^2+\frac{1}{2}V_1V_2G^3\right)\psi )_\rho =2(\phi ,\left(I+\frac{V_2}{2}G\right)\left(I+\frac{V_1}{2}G\right)G\psi )_\rho .$$
* If $`W_t`$ is level $`1`$ weakly transient with $`G_t^2Gf_tH`$ for some growth function $`f_t`$, then $`(_0^tf_s𝑑s)^{1/2}Y_t`$ converges to a Gaussian field with covariance functional $`{\displaystyle \frac{1}{2}}V_1V_2H(\phi ,\psi ).`$
2.3. Comments on the assumptions and the results
1. In Subsection 2.4 we will give conditions on the motion $`W_t`$ which imply the growth assumptions of Theorems 2.2.2(b) and 2.2.3(b). The $`\alpha `$–stable motion fits into this framework.
2. We do not know if condition (2.2.1) holds in general for level 1 strongly transient motion. We will show, however, that it holds for the motions in the examples.
3. We give an explanation of the second moment structure ocurring in Theorem 2.2.2(a). Let $`\underset{¯}{R}_{\mathrm{}}^1`$ be the historical counterpart of the canonical equilibrium measure $`R_{\mathrm{}}^1`$ (which is a time–shift invariant measure on the space of clans ranging from time $`\mathrm{}`$ to time $`+\mathrm{}`$) (see Dawson and Perkins<sup>(12)</sup>, Sections 5.4.3, 5.4.4). Firstly, we observe that the second moment measure of $`R_{\mathrm{}}^1`$ is given by (see the Appendix, (A.1.12)) $`{\displaystyle \mu ,\phi \mu ,\psi R_{\mathrm{}}^1(d\mu )}=(\phi ,\left(I+{\displaystyle \frac{1}{2}}VG\right)\psi )_\rho .(\mathrm{2.3.1})`$ The equality (2.3.1) can be intuitively understood through the backward tree picture (Gorostiza and Wakolbinger<sup>(22)</sup>, Section 4, and references therein, and Appendix, Section 3.A): The intensity measure of the canonical Palm distribution $`(R_{\mathrm{}}^1)_x`$ (with ego at site $`x`$) is
$`{\displaystyle \mu ,\psi (R_{\mathrm{}}^1)_x(d\mu )}`$ $`=`$ $`\psi (x)+E_x\left({\displaystyle _0^{\mathrm{}}}{\displaystyle p_t(W_t,dy)\psi (y)V𝑑t}\right)`$
$`=`$ $`\psi (x)+{\displaystyle \frac{1}{2}}VG\psi (x),(\mathrm{2.3.2})`$
where $`p_t`$ denotes the transition probability of the motion. Then, by the Palm formula (A.2.1) (Appendix), (2.3.1) follows by integrating (2.3.2) with respect to $`\phi (x)\rho (dx)`$. The same reasoning shows that the space–time correlation structure of $`\underset{¯}{R}_{\mathrm{}}^1`$ is given by
$$E_{\underset{¯}{R}_{\mathrm{}}^1}(X_t,\phi X_{t+s},\psi )=(\phi ,\left(I+\frac{1}{2}VG\right)T_s\psi )_\rho .$$
This reveals how the normed second moment measure of the occupation time behaves:
$$\frac{1}{t}E_{\underset{¯}{R}_{\mathrm{}}^1}\left(_0^tX_s,\phi 𝑑s_0^tX_r,\psi 𝑑r\right)2(\phi ,\left(I+\frac{V}{2}G\right)G\psi )_\rho \text{ as }t\mathrm{}.$$
4. Because of the obvious identities
$$\frac{1}{2}(\phi ,G\phi )_\rho =_0^{\mathrm{}}(\phi ,T_{2r}\phi )_\rho 𝑑r=_0^{\mathrm{}}\rho ,(T_r\phi )^2𝑑r,$$
transience of the motion implies persistence of the $`1`$–level branching system (Gorostiza and Wakolbinger<sup>(19)</sup> Corollary 2.2). Therefore the existence of the measure $`R_{\mathrm{}}^1`$, which is the assumed intensity for the Poisson initial condition of the 2–level branching system, is implied by the strong transience assumption on the system.
5. We now explain the second moment structure appearing in Theorem 2.2.3(a). For a level $`1`$ strongly transient motion it can be shown along the same lines as in Gorostiza et al<sup>(17)</sup> that the $`2`$–level branching particle system, started off from the Poisson system of $`1`$–level equilibrium clans, is persistent. Using e.g. the argument of Gorostiza<sup>(16)</sup> (Lemma 4.6), one derives the following expression for the second moment measure of the canonical measure $`Q_{\mathrm{}}`$ of the aggregated 2–level equilibrium with intensity $`\rho `$ (using the fact that $`R_1^{\mathrm{}}`$ is the Poissonization of the canonical measure of the superprocess counterpart, Appendix, (A.1.10)): $`{\displaystyle \mu ,\phi \mu ,\psi Q_{\mathrm{}}(d\mu )}=\phi ,\left(I+{\displaystyle \frac{1}{2}}(V_1+V_2)G+{\displaystyle \frac{1}{4}}V_1V_2G^2\psi \right)_\rho .(\mathrm{2.3.3})`$ Then the second moment structure of the occupation time follows as in the 1–level case. Equality (2.3.3) can also be understood through the backward tree picture: The intensity measure of the Palm distribution $`(Q_{\mathrm{}})_x`$ has the representation (Hochberg and Wakolbinger<sup>(26)</sup>)
$`{\displaystyle \mu ,\psi (Q_{\mathrm{}})_x(d\mu )}`$ $`=`$ $`\psi (x)+E_x[.{\displaystyle _0^{\mathrm{}}}(.{\displaystyle }(V_1+V_2)p_t(W_t,dy)\psi (y)`$
$`+{\displaystyle _0^t}{\displaystyle }V_1p_s(W_t,dz){\displaystyle }V_2p_{ts}(z,dy)\psi (y)ds.)dt.]`$
$`=`$ $`\psi (x)+{\displaystyle \frac{1}{2}}(V_1+V_2)G\psi (x)+{\displaystyle \frac{1}{4}}V_1V_2G^2\psi (x).(\mathrm{2.3.4})`$
The summands on the r.h.s. of (2.3.4) have an interpretation in terms of the genealogy: $`{\displaystyle \frac{1}{2}}V_1G\psi (x)`$ is the contribution of the 1–level relatives, $`{\displaystyle \frac{1}{2}}V_2G\psi (x)`$ is that of the 2–level relatives breaking off directly from the individual trunk, and $`{\displaystyle \frac{1}{4}}V_1V_2G^2\psi (x)`$ is that of the 2–level relatives having been generated by 1–level relatives, after those have broken off from the individual trunk. Now (2.3.3) follows by integrating (2.3.4) with respect to $`\phi (x)\rho (dx)`$.
6. The following table subsumes the covariance kernels appearing in the second moment structures discussed above (Theorems 2.2.2(a) and 2.2.3(a)). Columns 1, 2 and 3 refer to simple motion, 1–level branching and 2–level branching, respectively.
| $`2G`$ aa | aa$`2G+VG^2`$ | aa$`2G+(V_1+V_2)G^2+{\displaystyle \frac{1}{2}}V_1V_2G^3`$ |
| --- | --- | --- |
| | aa$`=2\left(I+{\displaystyle \frac{V}{2}}G\right)G`$aa | aa$`=2\left(I+{\displaystyle \frac{V_2}{2}}G\right)\left(I+{\displaystyle \frac{V_1}{2}}G\right)G`$ |
We observe a relationship between the covariance kernels for the $`1`$– and $`2`$–level cases: For $`V>0`$ and an operator $`Q`$, we define the operator
$$C_V(Q)=\left(I+\frac{V}{2}G\right)Q.$$
Then the $`1`$– and $`2`$–level covariance kernels are given by $`2C_V(G)`$ and $`2C_{V_2}(C_{V_1}(G))`$, respectively. Thus, the $`2`$–level covariance kernel is like the $`1`$–level covariance kernel with $`V`$ replaced by $`V_2`$ and the operator $`G`$ replaced by $`C_{V_1}(G)`$. Recall that $`G`$ represents the expected total occupation time of the motion, and note that $`C_V(G)`$ represents the expected total occupation time of the mass flow of the motion plus a continuous throwing off of mass with intensity $`{\displaystyle \frac{1}{2}}V`$, which also evolves by the same flow. One can then guess that for a 3–level system the covariance kernel would be given by
$$2C_{V_3}(C_{V_2}(C_{V_1}(G)))=2\left(I+\frac{V_3}{2}G\right)\left(I+\frac{V_2}{2}G\right)\left(I+\frac{V_1}{2}G\right)G,$$
and so on for higher levels of branching.
7. By using an argument of Dawson and Perkins<sup>(12)</sup> (Section 5.4.4), one observes that a level $`1`$ transient motion leads to transient equilibrium clans. Indeed, even the expected value $`a`$ of the total future occupation time in a bounded region $`B`$, starting from an equilibrium Palm cluster, is finite: Recall from (2.3.2) that the intensity measure $`\nu (dy)`$ of the Palm distribution $`(R_{\mathrm{}}^1)_0`$ of an equilibrium cluster (with ego at the origin) is
$$\nu (dy)=\delta _0(dy)+\frac{1}{2}VG(0,dy),$$
hence
$$a=\nu (dy)G(y,B)=G(0,B)+\frac{1}{2}VG^2(0,B)<\mathrm{}.$$
8. A level $`2`$ transient motion leads to transient aggregated $`2`$–level equilibrium clans. Indeed, by (2.3.4) the intensity measure $`\sigma (dy)`$ of the Palm distribution $`(Q_{\mathrm{}})_0`$ is
$$\sigma (dy)=\delta _0(dy)+\left(\frac{1}{2}(V_1+V_2)G+\frac{1}{4}V_1V_2G^2\right)(0,dy),$$
hence the expected value of the total future occupation time in a bounded region $`B`$ is
$$\sigma (dy)G(y,B)=G(0,B)+\left(\frac{1}{2}(V_1+V_2)G^2+\frac{1}{4}V_1V_2G^3\right)(0,B)<\mathrm{}.$$
9. For the $`2`$–level system, if the intensity of the initial Poisson distribution is $`\delta _{\delta _x}\rho (dx)`$ (instead of $`R_{\mathrm{}}^1`$), then in the assumption for Theorem 2.2.3(b) the growth of $`G_t^2G`$ is replaced by the growth of $`G_t^3`$, and an analogous result holds.
2.4. Order of transience and recurrence, asymptotics of powers of $`G_t`$, and some special growth functions Definition 2.4.1. (a) Let $``$ denote the class of differentiable functions $`h:\text{}_+\text{}_+`$ that are bounded and bounded away from $`0`$ on $`[T,\mathrm{})`$ for some $`T>0`$. (b) For fixed $`a(0,1)`$, let
$$\stackrel{~}{}_a=\{h|h_t=h_{at}\text{for all}t>0\}$$
(with $`T=0`$). Note that the elements of $`\stackrel{~}{}_a`$ are periodic in a logarithmic scale. Definition 2.4.2. (a) For a given $`\gamma >0`$, we say that $`W_t`$ is transient of order $`\gamma `$ if $`R_tt^\gamma h_tJ`$ for some $`h`$ as in Definition 2.4.1 and some bilinear form $`J`$ as in Definition 2.1.3.
(b) For a given $`\gamma (1,0)`$, we say that $`W_t`$ is recurrent of order –$`\gamma `$ if $`G_tt^\gamma h_tJ`$ for some $`h`$ and $`J`$ as above, and we say that $`W_t`$ is critically recurrent if $`G_t\mathrm{log}tJ`$ for some $`J`$ as above. Clearly, for transience of order $`\gamma `$ we have
$`\text{level}k\text{transience if and only if}`$ $`k<\gamma ,`$
$`\text{level}k\text{recurrence if and only if}`$ $`k\gamma ,`$
and therefore the following lemma holds. Lemma 2.4.1. If $`W_t`$ is transient of order $`\gamma `$, then for each $`k0,W_t`$ is
$`\text{level}k\text{strongly transient if and only if}`$ $`\gamma >k+1,`$
$`\text{level}k\text{weakly transient if and only if}`$ $`k<\gamma k+1.`$
The next lemma shows how the assumptions of Theorems 2.2.2(b) and 2.2.3(b) on the growth of powers of $`G_t`$ are implied by transience of order $`\gamma `$ with $`h1`$. Lemma 2.4.2. Let $`W_t`$ be transient of order $`\gamma `$ with $`R_tt^\gamma J`$.
(1) If $`0<\gamma <1`$, then
$$G_tGt^{1\gamma }H\text{with }H=\frac{1}{1\gamma }J,$$
and
$$G_t^2t^{1\gamma }H\text{with }H=\frac{22^{1\gamma }}{1\gamma }J.$$
(2) If $`\gamma =1`$, then
$$G_tGG_t^2\mathrm{log}tH\text{with }H=J.$$
(3) If $`1<\gamma <2`$, then
$$G_t^2Gt^{2\gamma }H\text{with }H=\frac{22^{2\gamma }}{(2\gamma )(\gamma 1)}J,$$
and
$$G_t^3t^{2\gamma }H\text{with }H=\frac{3^{2\gamma }2^{2\gamma }1}{(2\gamma )(\gamma 1)}J.$$
(4) If $`\gamma =2`$, then
$$G_t^2GG_t^3\mathrm{log}tH\text{with }H=J.$$
More generally, under the assumptions of Lemma 2.4.2 one can show that
$$G_t^\gamma GG_t^{\gamma +1}\mathrm{log}tJ\text{if}\gamma \text{ is an integer,}$$
and
$$G_t^{[\gamma ]+1}Gt^{[\gamma ]+1\gamma }c_\gamma J,G_t^{[\gamma ]+2}t^{[\gamma ]+1\gamma }c_\gamma ^{}J\text{otherwise}$$
(for suitable constants $`c_\gamma ,c_\gamma ^{}`$). But we will use only the cases $`\gamma 2`$ considered in the lemma. For growth functions $`f_t`$ of the form $`f_t=t^\zeta h_t,0<\zeta <1,h\stackrel{~}{}_a`$,which will occur in one of the examples, the normalizations $`(_0^tf_s𝑑s)^{1/2}`$ that appear in the conclusions of Theorems 2.2.1(b), 2.2.2(b) and 2.2.3(b) can be replaced by $`(t^{\zeta +1}\stackrel{~}{h}_t)^{1/2}`$, where $`\stackrel{~}{h}\stackrel{~}{}_a`$, thanks to the following lemma.
Lemma 2.4.3. Let $`h\stackrel{~}{}_a`$. If $`\zeta >1`$, then $`{\displaystyle _1^t}s^\zeta h_s𝑑st^{\zeta +1}\stackrel{~}{h}_t\text{as}t\mathrm{},(\mathrm{2.4.1})`$ where
$$\stackrel{~}{h}_t=\mathrm{log}a_0^{\mathrm{}}a^{r(1+\zeta )}h_{a^rt}𝑑r.$$
Note that $`\stackrel{~}{h}\stackrel{~}{}_a`$. Lemma 2.4.3 applies to the growths of the fluctuations of the 0, 1 and 2–level hierarchical system, but the l.h.s. of (2.4.1) can be computed explicitly in this case. 2.5. Infinite variance branching
In the case of infinite variance branching we consider here only the so called “$`(1+\beta )`$–branching”, ($`0<\beta <1`$), whose offspring generating function is of the form $`s+q(1s)^{1+\beta },s[0,1]`$, for some constant $`q>0`$ (this law belongs to the domain of normal attraction of a stable law with exponent $`1+\beta `$). The picture now is less complete: we have only results for the “classical” $`t^{1/(1+\beta )}`$–norming. In the 1–level case it can be shown (along the lines of Stöckl and Wakolbinger<sup>(38)</sup>) that this regime coincides with that of clan transience. We conjecture that an analogous result holds also in the 2–level case. Note that for $`\beta =1`$ we have the binary branching system considered above. However, we shall see that the results for the finite variance case are not special cases of the ones in this subsection.
We consider the following branching systems:
* 1–level system: The system is as described in Subsection 2.2, except that the particles undergo $`(1+\beta )`$–branching at rate $`V`$. We assume that is is persistent. Hence the system (with Poisson initial condition) converges to an equilibrium with intensity $`\rho `$. (Sufficient conditions for this persistence are given in Gorostiza and Wakolbinger<sup>(22)</sup>, Theorem 2.1). The equilibrium is then a Poisson superposition of “2–level particles”. Again we denote the equilibrium canonical measure by $`R_{\mathrm{}}^1`$.
* 2–level system: The system is as described in Subsection 2.2, except that individual particles undergo $`(1+\beta _1)`$–branching at rate $`V_1`$, clans undergo $`(1+\beta _2)`$–branching at rate $`V_2`$, and it starts off from a Poisson system of “2–level particles” with intensity $`R_{\mathrm{}}^1`$ (the equilibrium canonical measure for the 1–level $`(1+\beta _1)`$–branching system with rate $`V_1`$).
As above, $`Y_t`$ denotes the occupation time fluctuation (of the aggregated system in the 2–level case).
Theorem 2.5.1. Let $`X_t`$ be the 1–level branching system. Assume that $`\rho ,(G\phi )^{1+\beta }<\mathrm{},\phi 𝒞_c^+(S).(\mathrm{2.5.1})`$ Then $`t^{1/(1+\beta )}Y_t`$ converges to a random field $`Z`$ with Laplace functional
$$E\mathrm{exp}\{Z,\phi \}=\mathrm{exp}\left\{\frac{V}{1+\beta }\rho ,(G\phi )^{1+\beta }\right\},\phi 𝒞_c^+(S).$$
Here and in the next theorem the notation $`Z,\phi `$ means the action of the random field $`Z`$ on $`\phi `$. Note that the finite variance case $`\beta =1`$ (Theorem 2.2.2 (a)) is not a special case of Theorem 2.5.1, since this theorem would provide only the second term of the covariance functional. The term $`2(\phi ,G\psi )_\rho `$ in Theorem 2.2.2(a) does not appear in Theorem 2.5.1 because the normalization is now strong enough to kill this contribution to the occupation time fluctuations which comes from individuals related in direct line. Theorem 2.5.2. Let $`X_t`$ be the 2–level branching system. Assume that $`\beta _2<1`$ and $`{\displaystyle \mu ,G\phi ^{1+\beta _2}R_{\mathrm{}}^1(d\mu )}<\mathrm{},\phi 𝒞_c^+(S).(\mathrm{2.5.2})`$ Then $`t^{1/(1+\beta _2)}Y_t`$ converges to a random field $`Z`$ with Laplace functional
$$E\mathrm{exp}\{Z,\phi \}=\mathrm{exp}\left\{\frac{V_2}{1+\beta _2}\mu ,G\phi ^{1+\beta _2}R_{\mathrm{}}^1(d\mu )\right\},\phi 𝒞_c^+(S).$$
Note that also for the 2–level system the finite variance case $`\beta _1=\beta _2=1`$ (Theorem 2.2.3(a)) is not a special case of Theorem 2.5.2.
In the case of $`\alpha `$–stable motions in $`\text{}^d`$, we shall see in the next section that conditions (2.5.1) and (2.5.2) hold for “high” dimensions. 2.6. Occupation time fluctuations of superprocesses
As stated in the Introduction, results analogous to the previous ones hold also for the occupation time fluctuations of the 1– and 2– level superprocesses corresponding to the branching particle systems. The results are simpler because all the mass relationships closer than level $`k`$ do not contribute to the occupation time fluctuation limit of the $`k`$–level branching system. Theorems 2.2.2(a) and 2.2.3(a) hold with only the terms involving $`G^2`$ and $`G^3`$ present in the covariances of the limit random fields, respectively. The proofs are similar to those for the branching particle system and we shall omit them. 3. TWO EXAMPLES OF INDIVIDUAL MOTIONS: SYMMETRIC $`\alpha `$–STABLE PROCESS AND $`c`$–HIERARCHICAL RANDOM WALK 3.1. Transience/recurrence properties of the motions In the first example the particle motion $`W_t`$ is the spherically symmetric $`\alpha `$–stable Lévy process in $`S=\text{}^d`$, $`(0<\alpha 2)`$, and $`\rho `$ is the Lebesgue measure $`\lambda `$. In the second example $`W_t`$ is a continuous–time random walk in the hierarchical group $`S=\mathrm{\Omega }_N`$ and $`\rho `$ is the counting measure $`\nu `$. Due to similarities in the asymptotic behavior of $`G_t`$ and its powers for the two examples, which we will work out in this subsection, the limits of the occupation time fluctuations of the corresponding 1– and 2–level branching systems will also be analogous.
In this subsection we use several constants and functions whose definitions are collected in Section 4 for easy reference.
The analogy between the two motions is exhibited by a constant $`\gamma `$ which we will define in each case.
For the $`\alpha `$–stable process we define $`\gamma ={\displaystyle \frac{d}{\alpha }}1.(\mathrm{3.1.1})`$
Before defining the $`\gamma `$ for the hierarchical random walk we will give some background.
The hierarchical group of order $`N`$ is a countable group defined by
$$\mathrm{\Omega }_N=\{x=(x_1,x_2,\mathrm{})|x_i\{0,1,\mathrm{},N1\},x_i0\text{except for finitely many}i\},$$
with addition componentwise mod$`(N)`$. The hierarchical distance $`||`$ on $`\mathrm{\Omega }_N`$ is defined by
$$|xy|=\mathrm{max}\{i|x_iy_i\}.$$
A discrete–time hierarchical random walk in $`\mathrm{\Omega }_N`$ jumps from $`x`$ to $`y`$ such that $`|xy|=i1`$ with probability $`r_i/N^{i1}(N1)`$, where $`r_1,r_2,\mathrm{}`$ is a probability distribution on $`\{1,2,\mathrm{}\}`$. This type of random walk was introduced by Spitzer<sup>(37)</sup> for $`N=2`$ (the “light bulb random walk”), and by Sawyer and Felsenstein<sup>(35)</sup> for general $`N`$ in the context of genetics models.
The continuous–time analogue of the hierarchical random walk with jump rate $`\sigma >0`$ has transition density
$$p_t(0,x)=\frac{1}{N^i}(\delta _{0i}1)e^{\sigma (1f_i)t}+(N1)\underset{j=i+1}{\overset{\mathrm{}}{}}\frac{1}{N^j}e^{\sigma (1f_j)t}$$
if $`|x|=i0`$, where $`f_j`$ is given by
$$f_j=r_1+\mathrm{}+r_{j1}\frac{r_j}{N1},j1,$$
(note that $`f_0`$ is irrelevant); see e.g. Fleischmann and Greven<sup>(15)</sup> for additional information.
Here we will take $`r_i`$ of the form
$$r_i=\left(\frac{c}{N}\right)^{i1}\left(1\frac{c}{N}\right),i1,$$
where $`c`$ is a constant such that $`0<c<N`$. In this case we have $`f_j=1a^jb,j1,`$ where $`a={\displaystyle \frac{c}{N}},b={\displaystyle \frac{N^2/c1}{N1}},(\mathrm{3.1.2})`$ (note that $`0<a<1,b>1`$), and the transition density becomes $`p_t(0,x)={\displaystyle \frac{1}{N^i}}(\delta _{0i}1)e^{\sigma a^ibt}+(N1){\displaystyle \underset{j=i+1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{N^j}}e^{\sigma a^jbt}(\mathrm{3.1.3})`$ if $`|x|=i0`$.
Since this random walk is characterized by the constant $`c`$ (for fixed $`N`$), we will call it the $`c`$–hierarchical random walk. We shall see that $`c`$ is a mobility constant which plays a similar role to that of $`\alpha `$ for the $`\alpha `$-stable process.
We define $`\gamma ={\displaystyle \frac{\mathrm{log}c}{\mathrm{log}(N/c)}},i.e.,c=N^{\gamma /(\gamma +1)},(\mathrm{3.1.4})`$ The following lemma shows that the constant $`\gamma `$ defined in (3.1.1) for the $`\alpha `$–stable process and in (3.1.4) for the $`c`$–hierarchical random walk corresponds to the order of transience/recurrence parameter (Definition 2.4.2) for the respective processes, and it also shows the analogies for the semigroups $`T_t`$ and for the growth conditions for $`G_t^2`$ and $`G_t^2G`$ that appear as assumptions in Theorem 2.2.2(b) and Theorem 2.2.3(b) for the two motions. Each one of the functions $`h_t`$ appearing in the lemma equals 1 for the $`\alpha `$-stable process, and belongs to $`\stackrel{~}{}_a`$ for the c-hierarchical random walk, with $`a`$ given by (3.1.2), or, equivalently, by (4.2.3). Lemma 3.1.1. Let $`W_t`$ be either the $`\alpha `$–stable process in $`\text{}^d`$ or the $`c`$–hierarchical random walk in $`\mathrm{\Omega }_N`$.
* For all $`\gamma >1`$,
$$(\phi ,T_t\psi )_\rho t^{(\gamma +1)}h_tq_\gamma \rho ,\phi \rho ,\psi .$$
* For $`\gamma >0,W_t`$ is transient of order $`\gamma `$ with
$$R_t(\phi ,\psi )t^\gamma h_tq_\gamma \rho ,\phi \rho ,\psi .$$
* For $`1<\gamma <0,W_t`$ is recurrent of order –$`\gamma `$ with
$$(\phi ,G_t\psi )_\rho t^\gamma h_tq_\gamma \rho ,\phi \rho ,\psi .$$
* For $`\gamma =0,W_t`$ is critically recurrent with
$$(\phi ,G_t\psi )_\rho \mathrm{log}tq_0\rho ,\phi \rho ,\psi .$$
* For $`0<\gamma <1,W_t`$ is weakly transient with
$$(\phi ,G_t^2\psi )_\rho t^{1\gamma }h_tq_\gamma \rho ,\phi \rho ,\psi .$$
* For $`\gamma =1,W_t`$ is weakly transient with
$$(\phi ,G_t^2\psi )_\rho \mathrm{log}tq_1\rho ,\phi \rho ,\psi .$$
* For $`1<\gamma <2,W_t`$ is level 1 weakly transient with
$$(\phi ,G_t^2G\psi )_\rho t^{2\gamma }h_tq_\gamma \rho ,\phi \rho ,\psi .$$
* For $`\gamma =2,W_t`$ is level 1 weakly transient with
$$(\phi ,G_t^2G\psi )_\rho \mathrm{log}tq_2\rho ,\phi \rho ,\psi .$$
The correspondences for the examples are:
For the $`\alpha `$–stable process: $`\gamma ={\displaystyle \frac{d}{\alpha }}1,\rho =\lambda ,h1`$ in all cases,
* $`q_\gamma =\kappa _{d,\gamma }`$,
* $`\alpha <d,q_\gamma ={\displaystyle \frac{\kappa _{d,\gamma }}{\gamma }}`$,
* $`\alpha >d,q_\gamma ={\displaystyle \frac{\kappa _{d,\gamma }}{\gamma }}`$,
* $`\alpha =d,q_0=\kappa _{d,0}`$,
* $`{\displaystyle \frac{d}{2}}<\alpha <d,q_\gamma ={\displaystyle \frac{22^{1\gamma }}{(1\gamma )\gamma }}\kappa _{d,\gamma }`$,
* $`\alpha ={\displaystyle \frac{d}{2}},q_1=\kappa _{d,1}`$,
* $`{\displaystyle \frac{d}{3}}<\alpha <{\displaystyle \frac{d}{2}},q_\gamma ={\displaystyle \frac{22^{2\gamma }}{(2\gamma )(\gamma 1)\gamma }}\kappa _{d,\gamma }`$,
* $`\alpha ={\displaystyle \frac{d}{3}},q_2=\kappa _{d,2}`$.
For the $`c`$–hierarchical random walk: $`\gamma ={\displaystyle \frac{\mathrm{log}c}{\mathrm{log}(N/c)}},\rho =\nu `$,
* $`q_\gamma =\kappa _{N,\gamma },h=h^{(1,\gamma +1)}`$,
* $`c>1,q_\gamma =\kappa _{N,1},h=h^{(1,\gamma )}`$,
* $`c<1,q_\gamma =\kappa _{N,\gamma },h=h^{(2,\gamma )}`$,
* $`c=1,q_0={\displaystyle \frac{\kappa _{N,0}}{\mathrm{log}N}},`$
* $`1<c<N^{1/2},q_\gamma =\kappa _{N,\gamma },h=h^{(3,\gamma 1)}`$,
* $`c=N^{1/2},q_1={\displaystyle \frac{2\kappa _{N,1}}{\mathrm{log}N}}`$,
* $`N^{1/2}<c<N^{2/3},q_\gamma =\kappa _{N,\gamma },h=h^{(3,\gamma 2)}`$,
* $`c=N^{2/3},q_2={\displaystyle \frac{3\kappa _{N,2}}{\mathrm{log}N}}`$.
The constants $`\kappa _{d,\alpha }`$ and $`\kappa _{N,\gamma }`$, and the functions $`h^{(,)}`$ are defined in Section 4: expressions (4.1.1), (4.2.1), (4.2.5), (4.2.6) and (4.2.7). Corollary 3.1.1 (to Lemmas 3.1.1 and 2.4.1).
The $`\alpha `$–stable process is level $`k`$ strongly transient if and only if
$$\alpha <\frac{d}{k+2},$$
and level $`k`$ weakly transient if and only if
$$\frac{d}{k+2}\alpha <\frac{d}{k+1}.$$
The $`c`$–hierarchical random walk is level $`k`$ strongly transient if and only if
$$c>N^{(k+1)/(k+2)},$$
and level $`k`$ weakly transient if and only if
$$N^{k/(k+1)}<cN^{(k+1)/(k+2)}.$$
Corollary 3.1.1 for the $`\alpha `$–stable process with $`k=0`$ is well known (Sato<sup>(34)</sup>).
The powers of the Green potential operator of $`W_t`$ are given in the next lemma.
Lemma 3.1.2. Let $`\gamma >0`$ and $`1j<\gamma +1`$.
* For the $`\alpha `$–stable process, the integral kernel of $`G^j`$ is $`G_{d,\gamma ,j}(x)`$ given by (4.1.2) with $`\alpha j<d`$.
* For the $`c`$–hierarchical random walk, the integral kernel of $`G^j`$ is $`G_{N,\gamma ,j}(x)`$ given by (4.2.2) with $`c>N^{(j1)/j}`$.
The next lemma shows that the condition (2.2.1) in Theorem 2.2.3(a) is fulfilled in both examples with $`\delta =3`$. Lemma 3.1.3. Let $`W_t`$ be either the $`\alpha `$-stable process in $`\text{}^d`$ or the $`c`$–hierarchical random walk in $`\mathrm{\Omega }_N`$. Then level $`k`$ strong transience implies that $`T_t\phi =\mathrm{o}(t^{(k+2)})`$, $`\phi 𝒞_c^+(\text{}^d)`$ (resp. $`C_c^+(\mathrm{\Omega }_N)`$).
We give next upper and lower bounds for the functions $`h`$ and $`\stackrel{~}{h}`$ defined in Section 4, which appear in Lemma 3.1.1 and in Theorem 3.2.1 below for the hierarchical case. Proposition 3.1.1.
| Function | Lower bound | Upper bound |
| --- | --- | --- |
| $`h_t^{(1,\zeta )},\zeta >0`$ | $`{\displaystyle \frac{\mathrm{\Gamma }}{a^\zeta 1}}`$ | $`{\displaystyle \frac{a^\zeta \mathrm{\Gamma }}{a^\zeta 1}}`$ |
| $`h_t^{(2,\zeta )},1<\zeta <0`$ | $`{\displaystyle \frac{a^\zeta \mathrm{\Gamma }}{1a^\zeta }}`$ | $`{\displaystyle \frac{\mathrm{\Gamma }}{1a^\zeta }}`$ |
| $`h_t^{(3,\zeta )},1<\zeta <0`$ | $`{\displaystyle \frac{a^\zeta (22^{\zeta )})\mathrm{\Gamma }}{1a^\zeta }}`$ | $`{\displaystyle \frac{(22^\zeta )\mathrm{\Gamma }}{1a^\zeta }}`$ |
| $`\stackrel{~}{h}_t^{(2,\zeta )},1<\zeta <0`$ | $`{\displaystyle \frac{a^\zeta \mathrm{\Gamma }}{(1a^\zeta )(1\zeta )}}`$ | $`{\displaystyle \frac{\mathrm{\Gamma }}{(1a^\zeta )(1\zeta )}}`$ |
| $`\stackrel{~}{h}_t^{(3,\zeta )},1<\zeta <0`$ | $`{\displaystyle \frac{a^\zeta (22^\zeta )\mathrm{\Gamma }}{(1a^\zeta )(1\zeta )}}`$ | $`{\displaystyle \frac{(22^\zeta )\mathrm{\Gamma }}{(1a^\zeta )(1\zeta )}}`$ |
where $`\mathrm{\Gamma }\mathrm{\Gamma }(\zeta +1)`$. 3.2. Occupation time fluctuations limits We give first the occupation time fluctuation limits for the two examples in the case of finite variance branching. In contrast with the general theorems (Theorems 2.2.1 – 2.2.3) we present the results in a different order, ending up with the classical $`t^{1/2}`$–norming in each case. This is because the emphasis now, in the $`\alpha `$–stable case, is in going from the intermediate to the high dimensions $`d`$.
We denote by $`𝒩`$ a real–valued centered Gaussian random variable whose variance is specified in each case.
Theorem 3.2.1. Let $`W_t`$ be either the $`\alpha `$–stable process in $`\text{}^d`$ or the $`c`$–hierarchical random walk in $`\mathrm{\Omega }_N`$. 0–level:
* For –$`1<\gamma <0,(t^{1\gamma }\stackrel{~}{h}_t^{(2,\gamma )})^{1/2}Y_t`$ converges to $`𝒩\rho `$, where $`𝒩`$ has variance
$$2q_\gamma ^{(0)}.$$
* For $`\gamma =0,(t\mathrm{log}t)^{1/2}Y_t`$ converges to $`𝒩\rho `$, where $`𝒩`$ has variance
$$2q_0^{(0)}.$$
* For $`\gamma >0,t^{1/2}Y_t`$ converges to a Gaussian field with covariance kernel
$$2Q_{\gamma ,1}(x).$$
1–level:
* For $`0<\gamma <1,(t^{2\gamma }\stackrel{~}{h}_t^{(3,\gamma 1)})^{1/2}Y_t`$ converges to $`𝒩\rho `$, where $`𝒩`$ has variance
$$Vq_\gamma ^{(1)}.$$
* For $`\gamma =1,(t\mathrm{log}t)^{1/2}Y_t`$ converges to $`𝒩\rho `$, where $`𝒩`$ has variance
$$Vq_1^{(1)}.$$
* For $`\gamma >1,t^{1/2}Y_t`$ converges to a Gaussian field with covariance kernel
$$2Q_{\gamma ,1}(x)+VQ_{\gamma ,2}(x).$$
2–level:
* For $`1<\gamma <2,(t^{3\gamma }\stackrel{~}{h}_t^{(3,\gamma 2)})^{1/2}Y_t`$ converges to $`𝒩\rho `$, where $`𝒩`$ has variance
$$\frac{V_1V_2}{2}q_\gamma ^{(2)}.$$
* For $`\gamma =2,(t\mathrm{log}t)^{1/2}Y_t`$ converges to $`𝒩\rho `$, where $`𝒩`$ has variance
$$\frac{V_1V_2}{2}q_2^{(2)}.$$
* For $`\gamma >2,t^{1/2}Y_t`$ converges to a Gaussian field with covariance kernel
$$2Q_{\gamma ,1}(x)+(V_1+V_2)Q_{\gamma ,2}(x)+\frac{V_1V_2}{2}Q_{\gamma ,3}(x).$$
The correspondences for the examples are as follows (recall that $`a_t`$ denotes the normalization):
For the $`\alpha `$ stable process: $`\gamma ={\displaystyle \frac{d}{\alpha }}1`$, in this case the functions $`\stackrel{~}{h}^{(,)}`$ are constant, and the constants are included in the $`q_\gamma `$’s. 0–level:
* $`\alpha >d`$, $`a_t=t^{(1d/2\alpha )},q_\gamma ^{(0)}={\displaystyle \frac{\kappa _{d,\gamma }}{(1\gamma )(\gamma )}}={\displaystyle \frac{\kappa _{d,\alpha }^{}}{(2d/\alpha )(1d/\alpha )}}`$.
* $`\alpha =d,a_t=(t\mathrm{log}t)^{1/2},q_0^{(0)}=\kappa _{d,0}=\kappa _{d,d}^{}`$.
* $`\alpha <d,a_t=t^{1/2},Q_{\gamma ,1}(x)=G_{d,\gamma ,1}(x)=G_{d,\alpha ,1}^{}(x)`$.
1–level:
* $`{\displaystyle \frac{d}{2}}<\alpha <d`$, $`a_t=t^{(3/2d/2\alpha )},q_\gamma ^{(1)}={\displaystyle \frac{\kappa _{d,\gamma }(22^{1\gamma })}{(2\gamma )(1\gamma )\gamma }}={\displaystyle \frac{\kappa _{d,\alpha }^{}(22^{2d/\alpha })}{(3d/\alpha )(2d/\alpha )(d/\alpha 1)}}`$.
* $`\alpha ={\displaystyle \frac{d}{2}},a_t=(t\mathrm{log}t)^{1/2},q_1^{(1)}=\kappa _{d,1}=\kappa _{d,d/2}^{}`$.
* $`\alpha <{\displaystyle \frac{d}{2}},a_t=t^{1/2},Q_{\gamma ,j}(x)=G_{d,\gamma ,j}(x)=G_{d,\alpha ,j}^{}(x),j=1,2`$.
2–level:
* $`{\displaystyle \frac{d}{3}}<\alpha <{\displaystyle \frac{d}{2}}`$, $`a_t=t^{(2d/2\alpha )},q_\gamma ^{(2)}={\displaystyle \frac{\kappa _{d,\gamma }(22^{2\gamma })}{(3\gamma )(2\gamma )(\gamma 1)\gamma }}={\displaystyle \frac{\kappa _{d,\alpha }^{}(22^{3d/\alpha })}{(4d/\alpha )(3d/\alpha )(d/\alpha 2)(d/\alpha 1)}}`$.
* $`\alpha ={\displaystyle \frac{d}{3}},a_t=(t\mathrm{log}t)^{1/2},q_2^{(2)}=\kappa _{d,2}=\kappa _{d,d/3}^{}`$.
* $`\alpha <{\displaystyle \frac{d}{3}},a_t=t^{1/2},Q_{\gamma ,j}(x)=G_{d,\gamma ,j}(x)=G_{d,\alpha ,j}^{}(x),j=1,2,3`$.
For the $`c`$–hierarchical random walk: $`\gamma ={\displaystyle \frac{\mathrm{log}c}{\mathrm{log}(N/c)}},c=N^{\gamma /(\gamma +1)}`$. 0–level:
* $`c<1`$, $`a_t=t^{\mathrm{log}(N^{1/2}/c)/\mathrm{log}(N/c)}(\stackrel{~}{h}_t^{(2,\gamma )})^{1/2},q_\gamma ^{(0)}=\kappa _{N,\gamma }=\kappa _{N,c}^{}`$.
* $`c=1,a_t=(t\mathrm{log}t)^{1/2},q_0^{(0)}={\displaystyle \frac{\kappa _{N,0}}{\mathrm{log}a}}={\displaystyle \frac{(N1)^2}{\sigma (N^21)\mathrm{log}N}}`$.
* $`c>1,a_t=t^{1/2},Q_{\gamma ,1}(x)=G_{N,\gamma ,1}(x)=G_{N,c,1}^{}(x)`$.
1–level:
* $`1<c<N^{1/2}`$, $`a_t=t^{\mathrm{log}(N/c^{3/2})/\mathrm{log}(N/c)}(\stackrel{~}{h}_t^{(3,\gamma 1)})^{1/2},q_\gamma ^{(1)}=\kappa _{N,\gamma }=\kappa _{N,c}^{}`$.
* $`c=N^{1/2},a_t=(t\mathrm{log}t)^{1/2},q_1^{(1)}={\displaystyle \frac{\kappa _{N,1}}{\mathrm{log}a^{1/2}}}={\displaystyle \frac{2(N1)^3}{(\sigma (N^{3/2}1))^2\mathrm{log}N}}`$.
* $`c>N^{1/2},Q_{\gamma ,j}(x)=G_{N,\gamma ,j}(x)=G_{N,c,j}^{}(x),j=1,2`$.
2–level:
* $`N^{1/2}<c<N^{2/3}`$, $`a_t=t^{\mathrm{log}(N^{3/2}/c^2)/\mathrm{log}(N/c)}(\stackrel{~}{h}_t^{(3,\gamma 2)})^{1/2},q_\gamma ^{(2)}=\kappa _{N,\gamma }=\kappa _{N,c}^{}`$.
* $`c=N^{2/3},a_t=(t\mathrm{log}t)^{1/2},q_2^{(2)}={\displaystyle \frac{\kappa _{N,2}}{\mathrm{log}a^{1/3}}}={\displaystyle \frac{2(N1)^4}{(\sigma (N^{4/3}1))^3\mathrm{log}N}}`$.
* $`c>N^{2/3},a_t=t^{1/2},Q_{\gamma ,j}(x)=G_{N,\gamma ,j}(x)=G_{N,c,j}^{}(x),j=1,2,3`$.
We turn now to the case of infinite variance branching with $`\alpha `$-stable motion in $`\text{}^d`$. The question is when do the assumptions for Theorems 2.5.1 and 2.5.2 hold. Proposition 3.2.1. For the symmetric $`\alpha `$–stable motion in $`\text{}^d`$, condition (2.5.1) holds if and only if $`d>\alpha \left(1+{\displaystyle \frac{1}{\beta }}\right).`$ Proposition 3.2.2. For the symmetric $`\alpha `$–stable motion in $`\text{}^d`$, condition (2.5.2) holds if and only if $`\beta _2<\beta _1`$ and $`d>\alpha \left(1+{\displaystyle \frac{1}{\beta }}_2\left(1+{\displaystyle \frac{1}{\beta }}_1\right)\right).`$ 3.3. Comments on the results
1. Note that for the 2–level branching Brownian system $`(\alpha =2)`$ the normings $`t^{3/4},(t\mathrm{log}t)^{1/2}`$ and $`t^{1/2}`$ mentioned in the Introduction correspond to dimensions $`d=5,d=6`$ and $`d>6`$, respectively, i.e., two dimensions higher than for the $`1`$–level system. In the Brownian case the critical dimensions for the 1– and 2–level branching systems are $`d=4`$ and $`d=6`$, corresponding to $`\gamma =1`$ and $`\gamma =2`$, respectively. For 1–level Brownian systems, occupation time large deviation results have been obtained by Deuschel and Rosen<sup>(13)</sup> (and references therein).
2. For the $`\alpha `$–stable process in $`\text{}^d`$ the order of transience/recurrence parameter $`\gamma `$ defined in (3.1.1) can take values only in the interval $`[(d2)/2,\mathrm{})`$. On the other hand, for the $`c`$-hierarchical random walk in $`\mathrm{\Omega }_N`$ the possible values of the parameter $`\gamma `$ defined in (3.1.4) range over the whole interval $`(1,\mathrm{})`$, which means that in $`\mathrm{\Omega }_N`$ there is a rich structure of naturally ordered random walks.
3. Note that for both the $`\alpha `$–stable motion and the $`c`$–hierarchical random walk, in all cases (i) and (ii) of Theorem 3.2.1 the occupation time fluctuation limits in different regions of the space $`S`$ are perfectly correlated. An intuitive explanation for this might come from the recurrent visits of each $`k`$–level equilibrium clan, $`k=0,1,2`$, to all bounded regions $`BS`$.
4. Equating the parameters $`\gamma `$ for the $`\alpha `$–stable process (3.1.1) and the $`c`$–hierarchical random walk (3.1.4) we obtain
$$c=N^{1\alpha /d}.$$
For this value of $`c`$, by Lemma 3.1.1 the $`c`$–hierarchical random walk in $`\mathrm{\Omega }_N`$ and the $`\alpha `$–stable process in $`\text{}^d`$ have the same order of transience/recurrence. Consequently, by Theorems 2.2.1, 2.2.2 and 2.2.3 the asymptotics of the occupation time fluctuations are analogous for the corresponding $`k`$–level branching systems, $`k=0,1,2`$. The only differences are in the constants and the kernels of the powers of the Green operators which appear in the fluctuation limits in Theorem 3.2.1. The same observation holds for branching systems of “$`\alpha `$–stable” random walks on the lattice $`\text{}^d`$. In passing we note that $`\alpha `$–stable motions with $`\alpha <2`$ do not have finite moments of order $`\alpha `$, but this plays no role in the asymptotics of the occupation times. The corresponding $`c`$–hierarchical random walks have finite moments of all orders.
5. For the $`c`$–hierarchical random walk with $`c=c_N=\eta N^{k/(k+1)},k0`$, $`\eta >1`$, the powers of the Green potential operator take a simple form in the limit $`N\mathrm{}`$ : all the powers of order $`1jk`$ vanish as $`N\mathrm{}`$, and for the $`(k+1)`$–st power we observe from (4.2.2) and Lemma 3.1.2(b) that
$$\underset{N\mathrm{}}{lim}G_{N,c__N,k+1}(x)=\frac{\eta ^{k+1}}{\sigma ^{k+1}(\eta ^{k+1}1)}(\eta ^{k+1})^{|x|}.$$
In particular, $`lim_N\mathrm{}G_{N,\eta ,1}`$ and $`lim_N\mathrm{}G_{N,\eta N^{1/2},2}`$ have the same spatial asymptotics. This indicates that “near $`\eta =1`$” a similar analysis as was carried out by Dawson and Greven<sup>(7)</sup> for 1- -level branching hierarchical random walks ($`k=0`$) might also be possible for 2–level branching hierarchical random walks ($`k=1`$).
6. The results for the $`\alpha `$–stable motion with the $`t^{1/2}`$–norming can be extended to test functions in $`𝒞_\tau ^+(\text{}^d)`$. For example, for the 2–level system (with $`d>3\alpha `$) we can take $`\tau (x)=(1+|x|^2)^q`$ with $`d/2<q<(d+\alpha )/2`$ (Dawson and Gorostiza<sup>(6)</sup>). Moreover, the results can be extended to convergence of $`𝒮^{}(\text{}^d)`$–valued random fields, where $`𝒮^{}(\text{}^d)`$ is the space of tempered distributions on $`\text{}^d`$, using an argument of Iscoe<sup>(27)</sup>.
7. Similarly to the previous comment, for the $`c`$–hierarchical random walk the results with the $`t^{1/2}`$–norming can be extended to test functions in $`𝒞_\tau ^+(\mathrm{\Omega }_N)`$ with an appropriate function $`\tau `$. For example, for the 2–level system, $`\tau `$ should be a function in $`L^1(\mathrm{\Omega }_N,\nu )`$ such that the function $`x\underset{y}{}\tau (y)(N^2/c^3)^{|xy|}`$ is bounded.
4. DEFINITIONS OF CONSTANTS AND FUNCTIONS FOR THE EXAMPLES 4.1. Notation for $`\alpha `$–stable motion: $`\kappa _{d,\gamma }={\displaystyle \frac{1}{(2\pi )^d}}{\displaystyle _\text{}^d}e^{|x|^{d/(\gamma +1)}}𝑑x=\kappa _{d,\alpha }^{}={\displaystyle \frac{1}{(2\pi )^d}}{\displaystyle _\text{}^d}e^{|x|^\alpha }𝑑x,(\mathrm{4.1.1})`$$`G_{d,\gamma ,j}(x)=C_{d,\gamma ,j}|x|^{d(1j/(\gamma +1))}=G_{d,\alpha ,j}^{}(x)=C_{d,\alpha ,j}^{}|x|^{(dj\alpha )},(\mathrm{4.1.2})`$where
$$C_{d,\gamma ,j}=\frac{\mathrm{\Gamma }\left({\displaystyle \frac{d(\gamma +1j)}{2(\gamma +1)}}\right)}{2^{jd/(\gamma +1)}\pi ^{d/2}\mathrm{\Gamma }\left({\displaystyle \frac{dj}{2(\gamma +1)}}\right)}=C_{d,\alpha ,j}^{}=\frac{\mathrm{\Gamma }\left({\displaystyle \frac{dj\alpha }{2}}\right)}{2^{j\alpha }\pi ^{d/2}\mathrm{\Gamma }\left({\displaystyle \frac{j\alpha }{2}}\right)},$$
and $`j`$ is a positive integer such that $`j<\gamma +1`$, i.e. $`\alpha j<d`$. 4.2. Notation for $`c`$–hierarchical random walk: $`\kappa _{N,\gamma }=(N1)^{\gamma +2}(\sigma (N^{(\gamma +2)/(\gamma +1)}1))^{(\gamma +1)}(\mathrm{4.2.1})`$
$$=\kappa _{N,c}^{}=(N1)^{\mathrm{log}(N^2/c)/\mathrm{log}(N/c)}(\sigma (N^2/c1))^{\mathrm{log}N/\mathrm{log}(N/c)},$$
$`G_{N,\gamma ,j}(x)=C_{N,\gamma ,j}N^{|x|(1j/(\gamma +1))}=G_{N,c,j}^{}(x)=C_{N,c,j}^{}\left({\displaystyle \frac{N^{j1}}{c^j}}\right)^{|x|},(\mathrm{4.2.2})`$ where
$`C_{N,\gamma ,j}=\left({\displaystyle \frac{N1}{\sigma (N^{(\gamma +2)/(\gamma +1)}1)}}\right)^j\left[\delta _{0,|x|}1+{\displaystyle \frac{(N1)N^{j1}}{N^{j\gamma /(\gamma +1)}N^{j1}}}\right]`$
$`=`$ $`C_{N,\gamma ,j}^{}=\left({\displaystyle \frac{N1}{\sigma (N^2/c1)}}\right)^j\left[\delta _{0,|x|}1+{\displaystyle \frac{(N1)N^{j1}}{c^jN^{j1}}}\right],`$
and $`j`$ is a positive integer such that $`j<\gamma +1`$, i.e. $`c>N^{(j1)/j}`$.
Note that $`a`$ and $`b`$ defined in (3.1.2) are also expressed as $`a=N^{1/(\gamma +1)},b={\displaystyle \frac{N^{(\gamma +2)/(\gamma +1)}1}{N1}}.(\mathrm{4.2.3})`$ To simplify notation we write $`\theta =\sigma {\displaystyle \frac{N^{(\gamma +2)/(\gamma +1)}1}{N1}}.(\mathrm{4.2.4})`$ For $`\zeta >0`$, let $`h_t^{(1,\zeta )}={\displaystyle \underset{j=\mathrm{}}{\overset{\mathrm{}}{}}}(\theta a^jt)^\zeta e^{\theta a^jt},t>0,(\mathrm{4.2.5})`$ and for $`1<\zeta <0`$, let $`h_t^{(2,\zeta )}={\displaystyle \underset{j=\mathrm{}}{\overset{\mathrm{}}{}}}(\theta a^jt)^\zeta (1e^{\theta a^jt}),t>0,(\mathrm{4.2.6})`$ $`h_t^{(3,\zeta )}={\displaystyle \underset{j=\mathrm{}}{\overset{\mathrm{}}{}}}(\theta a^jt)^\zeta (1e^{\theta a^jt})^2,t>0,(\mathrm{4.2.7})`$ $`\stackrel{~}{h}_t^{(2,\zeta )}={\displaystyle \underset{j=\mathrm{}}{\overset{\mathrm{}}{}}}(\theta a^jt)^{\zeta 1}(e^{\theta a^jt}1+\theta a^jt),t>0,(\mathrm{4.2.8})`$ $`\stackrel{~}{h}_t^{(3,\zeta )}={\displaystyle \underset{j=\mathrm{}}{\overset{\mathrm{}}{}}}(\theta a^jt)^{\zeta 1}\left(2e^{\theta a^jt}{\displaystyle \frac{1}{2}}e^{2\theta a^jt}+\theta a^jt{\displaystyle \frac{3}{2}}\right),t>0,(\mathrm{4.2.9})`$
with $`a`$ and $`\theta `$ given by (4.2.3) and (4.2.4). Note that all the functions defined in (4.2.5)–(4.2.9) belong to $`\stackrel{~}{}_a`$ for $`a`$ given by (4.2.3). The functions $`\stackrel{~}{h}^{(2,\zeta )}`$ and $`\stackrel{~}{h}^{(3,\zeta )}`$ correspond (asymptotically) to $`h^{(2,\zeta )}`$ and $`h^{(3,\zeta )}`$, respectively, by Lemma 2.4.3, but in this case they are obtained by explicit calculation of the l.h.s. of (2.4.1).
5. PROOFS 5.1. Asymptotics of the powers of $`G_t`$ m Proof of Lemma 2.4.2:
Note that it suffices to do the proofs with $`\phi =\psi `$. Fix $`\phi 𝒞_c^+(S),\phi 0`$, and denote $`J=J(\phi ,\phi )`$ and $`R_t=R_t(\phi ,\phi )`$. By assumption, $`R_tt^\gamma J`$. (1)
$$(\phi ,G_tG\phi )_\rho =(\phi ,_0^t_0^{\mathrm{}}T_{s+u}\phi 𝑑s𝑑u)_\rho =_0^tR_u𝑑uJ\frac{1}{1\gamma }t^{1\gamma },$$
and
$`(\phi ,(G_tGG_t^2)\phi )_\rho =(\phi ,G_t{\displaystyle _t^{\mathrm{}}}T_s\phi 𝑑s)_\rho =(\phi ,{\displaystyle _0^t}{\displaystyle _t^{\mathrm{}}}T_{s+u}\phi 𝑑s𝑑u)_\rho ={\displaystyle _0^t}R_{s+t}𝑑s`$
$`J{\displaystyle _0^t}(s+t)^\gamma 𝑑s=J{\displaystyle \frac{1}{1\gamma }}((2t)^{1\gamma }t^{1\gamma })={\displaystyle \frac{J}{1\gamma }}(2^{1\gamma }1)t^{1\gamma },`$
hence
$$(\phi ,G_t^2\phi )_\rho =(\phi ,(G_t^2G_tG)\phi )_\rho +(\phi ,G_tG\phi )_\rho \frac{J}{1\gamma }(22^{1\gamma })t^{1\gamma }.$$
$`(2)(\phi ,G_tG\phi )_\rho ={\displaystyle _0^t}R_u𝑑uJ{\displaystyle _1^t}u^1𝑑uc\mathrm{log}t,`$
$$(\phi ,(G_tGG_t^2)\phi )_\rho =_0^tR_{s+t}𝑑s$$
$$J_0^t(s+t)^1𝑑sJ(\mathrm{log}(2t)\mathrm{log}t)=o(\mathrm{log}t),$$
hence the assertion for $`G_t^2`$ follows. $`(3)(\phi ,G_t^2G\phi )_\rho ={\displaystyle _0^t}{\displaystyle _0^t}R_{s+u}𝑑r𝑑u`$
$$J_1^t_1^t(s+u)^\gamma 𝑑s𝑑u\frac{J}{(2\gamma )(\gamma 1)}(22^{2\gamma })t^{2\gamma },$$
$$(\phi ,(G_t^2GG_t^3)\phi )_\rho =_0^t_0^tR_{s+u+t}𝑑u𝑑s$$
$$J_0^t_0^t(s+u+t)^\gamma 𝑑u𝑑s\frac{J}{(2\gamma )(\gamma 1)}(3^{2\gamma }+22^{2\gamma }1)t^{2\gamma },$$
hence the assertion for $`G_t^3`$ follows. $`(4)(\phi ,G_t^2G\phi )_\rho ={\displaystyle _0^t}{\displaystyle _0^t}R_{s+u}𝑑r𝑑u`$
$$J_1^t_0^t(s+u)^2𝑑s𝑑uc\mathrm{log}t,$$
$$(\phi ,(G_t^2GG_t^3)\phi )_\rho =_0^t_0^tR_{s+u+t}𝑑u𝑑s$$
$$J_0^t_0^t(s+u+t)^2𝑑s𝑑u=o(\mathrm{log}t),$$
hence the assertion for $`G_t^3`$ follows. $`\mathrm{}`$ Proof of Lemma 2.4.3:
$$_1^ts^\zeta h_s𝑑s=\mathrm{log}a_0^\tau a^{r(1+\zeta )}h_{a^r}𝑑r,$$
where $`\tau ={\displaystyle \frac{\mathrm{log}t}{\mathrm{log}a}}`$. Hence, since $`a<1`$ and $`h`$ is bounded, we have
$`t^{(1+\zeta )}{\displaystyle _1^t}s^\zeta h_s𝑑s`$ $`=`$ $`\mathrm{log}a{\displaystyle _0^\tau }a^{(\tau r)(1+\zeta )}h_{a^r}𝑑r`$
$`=`$ $`\mathrm{log}a{\displaystyle _0^\tau }a^{r(1+\zeta )}h_{a^{(\tau r)}}𝑑r`$
$``$ $`\mathrm{log}a{\displaystyle _0^{\mathrm{}}}a^{r(1+\zeta )}h_{a^rt}𝑑r.`$
$`\mathrm{}`$ 5.2. Main results
We will not include the proofs for the $`0`$–level system (Theorem 2.2.1) because they are simpler versions of the proofs for the branching systems. The proofs for the 1– and 2–level systems follow the idea of the method employed by Iscoe<sup>(27)</sup> for (1–level) superprocesses in $`\text{}^d`$. The particle systems are somewhat harder to deal with than the superprocesses, but the main point is it has been necessary to modify the method in order to deal with the new technical difficulties that arise from the second level branching. We will use the modified approach also for the 1–level system. Proof of Theorems 2.2.2 and 2.5.1 (1–level branching system):
In order to simplify notation we write $`\phi _t=F_t^{1/(1+\beta )}\phi \text{for}\phi 𝒞_c^+(S),\phi 0\text{and}\beta 1,\text{where}F_t={\displaystyle _0^t}f_s𝑑s,(\mathrm{5.2.1})`$
and $`f_s`$ is a growth function. For Theorem 2.2.2(a) and Theorem 2.5.1, $`f_t`$ is interpreted as $`f_t1`$.
We have, by (A.1.1) and (A.1.2) (Appendix),
$$E\mathrm{exp}\left\{F_t^{1/(1+\beta )}_0^tX_s𝑑s,\phi \right\}=\mathrm{exp}\{\rho ,u_{\phi _t}(t)\},$$
where $`u_{\phi _t}(s,x)`$ with values in $`[0,1]`$ is the unique solution of
$`u_{\phi _t}(s)={\displaystyle \frac{V}{1+\beta }}{\displaystyle _0^s}T_{sr}(u_{\phi _t}(r)^{1+\beta })𝑑r+{\displaystyle _0^s}T_{sr}(\phi _t(1u_{\phi _t}(r)))𝑑r(\mathrm{5.2.2})`$
Hence, by $`T_t`$–invariance of $`\rho `$ and $`E_0^tX_s𝑑s,\phi =t\rho ,\phi `$, for the occupation time fluctuation $`Y_t`$ we have
$`E\mathrm{exp}\{F_t^{1/(1+\beta )}Y_t,\phi \}=\mathrm{exp}\{{\displaystyle \frac{V}{1+\beta }}(I_1(t)+I_2(t))+I_3(t))\},(\mathrm{5.2.3})`$ where $`I_1(t)={\displaystyle _0^t}\rho ,w_{\phi _t}(s)^{1+\beta }𝑑s,(\mathrm{5.2.4})`$$`I_2(t)={\displaystyle _0^t}\rho ,u_{\phi _t}(s)^{1+\beta }w_{\phi _t}(s)^{1+\beta }𝑑s,(\mathrm{5.2.5})`$$`I_3(t)={\displaystyle _0^t}\rho ,\phi _tu_{\phi _t}(s)𝑑s,(\mathrm{5.2.6})`$with $`w_\phi (s)(x):={\displaystyle _0^s}T_r\phi (x)𝑑r=G_s\phi (x),xS,\phi 𝒞_c^+(S).(\mathrm{5.2.7})`$
We will prove the following limits as $`t\mathrm{}`$:
For Theorem 2.2.2(a):
$`I_1(t)(\phi ,G^2\phi )_\rho .(\mathrm{5.2.8})`$ For Theorem 2.2.2(b): $`I_1(t)H(\phi ,\phi ).(\mathrm{5.2.9})`$ For Theorem 2.5.1: $`I_1(t)\rho ,(G\phi )^{1+\beta }.(\mathrm{5.2.10})`$ For $`\beta 1`$: $`I_2(t)0.(\mathrm{5.2.11})`$ For $`\beta <1`$: $`I_3(t)0.(\mathrm{5.2.12})`$For Theorem 2.2.2(a): $`I_3(t)(\phi ,G\phi )_\rho .(\mathrm{5.2.13})`$For Theorem 2.2.2(b): $`I_3(t)0.(\mathrm{5.2.14})`$These limits will yield the conclusions of the theorems.
Proof of (5.2.8) and (5.2.9): From (5.2.4) and (5.2.7) we have
$$I_1(t)=F_t^1_0^t\rho ,(G_s\phi )^2𝑑s.$$
By L’Hôpital’s rule, for (5.2.8) we have
$$I_1(t)=t^1_0^t\rho ,(G_s\phi )^2𝑑s(\phi ,G_t^2\phi )_\rho (\phi ,G^2\phi )_\rho ,$$
and for (5.2.9),
$$I_1(t)=F_t^1_0^t\rho ,(G_s\phi )^2𝑑s\frac{1}{f_t}(\phi ,G_t^2\phi )_\rho H(\phi ,\phi ).$$
Proof of (5.2.10): The same as (5.2.8).
Proof of (5.2.11): We rewrite (5.2.5) as
$$I_2(t)=_0^t\rho ,w_{\phi _t}(s)^{1+\beta }[1(u_{\phi _t}(s)/w_{\phi _t}(s))^{1+\beta }]𝑑s.$$
We have from (5.2.2) and (5.2.7) $`u_{\phi _t}(s)w_{\phi _t}(s)={\displaystyle \frac{V}{1+\beta }}{\displaystyle _0^s}T_{sr}(u_{\phi _t}(r)^{1+\beta })𝑑s{\displaystyle _0^s}T_{sr}(\phi _tu_{\phi _t}(r))𝑑r0,(\mathrm{5.2.15})`$hence
$$01(u_{\phi _t}(s)/w_{\phi _t}(s))^{1+\beta }1.$$
Therefore, since the convergence of $`I_1(t)`$ implies uniform integrability, it suffices to show that
$$\underset{t\mathrm{}}{lim}\frac{u_{\phi _t}(s)}{w_{\phi _t}(s)}=1\text{for all}s,\rho a.e.$$
We have from (5.2.7) and (5.2.15)
$$\frac{u_{\phi _t}(s)}{w_{\phi _t}(s)}=1(G_s\phi )^1(J_t(s)+K_t(s)),$$
where
$$J_t(s)=\frac{V}{1+\beta }F_t^{1/(1+\beta )}_0^sT_{sr}(u_{\phi _t}(r)^{1+\beta })𝑑r0,$$
and
$$K_t(s)=_0^sT_{sr}(\phi u_{\phi _t}(r))𝑑r0.$$
By (5.2.15),
$`J_t(s)`$ $``$ $`{\displaystyle \frac{V}{1+\beta }}F_t^{1/(1+\beta )}{\displaystyle _0^s}T_{sr}(w_{\phi _t}(r)^{1+\beta })𝑑r`$
$`=`$ $`{\displaystyle \frac{V}{1+\beta }}F_t^{\beta /(1+\beta )}{\displaystyle _0^s}T_{rs}((G_r\phi )^{1+\beta })𝑑r0.`$
Similarly, $`K_t(s)0`$, so (5.2.11) is proved. Proof of (5.2.12): We have from (5.2.6) and (5.2.15)
$$I_3(t)t^{2/(1+\beta )}_0^t(\phi ,G_s\phi )_\rho 𝑑s=t^{(\beta 1)/(1+\beta )}t^1_0^t(\phi ,G_s\phi )_\rho 𝑑s,$$
and the result follows since $`(\phi ,G_s\phi )_\rho (\phi ,G\phi )_\rho <\mathrm{}`$ as $`s\mathrm{}`$. Proof of (5.2.13): We rewrite (5.2.6) as
$$I_3(t)=t^1_0^t\rho ,\phi \stackrel{~}{u}_{\phi _t}(s)𝑑s,$$
where $`\stackrel{~}{u}_{\phi _t}(s):=t^{1/2}u_{\phi _t}(s)`$. Since
$$t^1_0^t\rho ,\phi G_s\phi 𝑑s(\phi ,G\phi )_\rho ,$$
it suffices to prove that
$$A_t:=t^1_0^t\rho ,\phi \stackrel{~}{u}_{\phi _t}(s)𝑑st^1_0^t\rho ,\phi G_s\phi 𝑑s0$$
We have, from (5.2.2)
$$\stackrel{~}{u}_{\phi _t}(s)=\frac{V}{2}t^{1/2}_0^sT_{sr}(\stackrel{~}{u}_{\phi _t}(r)^2)𝑑r+G_s\phi t^{1/2}_0^sT_{sr}(\phi \stackrel{~}{u}_{\phi _t}(r))𝑑r,$$
hence
$$|A_t|t^{3/2}(\frac{V}{2}_0^t_0^s\rho ,\phi T_{sr}(\stackrel{~}{u}_{\phi _t}(r)^2)drds.+_0^t_0^s\rho ,\phi T_{sr}(\phi \stackrel{~}{u}_{\phi _t}(r))drds.),$$
but, from (5.2.2), $`\stackrel{~}{u}_{\phi _t}(r)G_r\phi `$ (since $`f_t1`$), so
$`|A_t|`$ $``$ $`t^{3/2}\left({\displaystyle \frac{V}{2}}{\displaystyle _0^t}{\displaystyle _0^s}\rho ,\phi T_{sr}(G_r\phi )^2𝑑r𝑑s+{\displaystyle _0^t}{\displaystyle _0^s}\rho ,\phi T_{sr}(\phi G_r\phi )𝑑r𝑑s\right)`$
$`=t^{3/2}({\displaystyle \frac{V}{2}}{\displaystyle _0^t}\rho ,\phi G_{tr}(G_r\phi )^2dr+{\displaystyle _0^t}\rho ,\phi G_{tr}(\phi G_r\phi )dr.),(\mathrm{5.2.16})`$therefore
$$|A_t|t^{1/2}\rho ,\frac{V}{2}\phi G(G\phi )^2+\phi G(\phi G\phi ).$$
Now, by strong transience,
$$\rho ,\phi G(G\phi )^2G\phi G^2\phi \rho ,\phi <\mathrm{},$$
and by transience,
$$\rho ,\phi G(\phi G\phi )G\phi ^2\rho ,\phi <\mathrm{},$$
hence the result follows. Proof of (5.2.14): We can follow the same argument used for (5.2.13) replacing $`\phi `$ by $`\phi F_t^{1/2}`$. Both terms on the r.h.s. of (5.2.16) can be shown to converge to $`0`$ by L’Hôpital’s rule and the assumptions.
It follows from (5.2.3)–(5.2.6) and the limits (5.2.8)–(5.2.14) that
$`E\mathrm{exp}\{(F_t)^{1/(1+\beta )}Y_t,\phi \}`$
$`\{\begin{array}{cc}\mathrm{exp}\left\{{\displaystyle \frac{V}{2}}(\phi ,G^2\phi )_\rho +(\phi ,G\phi )_\rho \right\}\hfill & \text{for Theorem 2.2.2(a),}\hfill \\ \mathrm{exp}\left\{{\displaystyle \frac{V}{2}}H(\phi ,\phi )\right\}\hfill & \text{for Theorem 2.2.2(b),}\hfill \\ \mathrm{exp}\left\{{\displaystyle \frac{V}{1+\beta }}\rho ,(G\phi )^{1+\beta }\right\}\hfill & \text{for Theorem 2.5.1}\hfill \end{array}`$
as $`t\mathrm{}`$.
Finally, the convergence of the bilateral Laplace functional implies the weak convergence of $`Y_t`$ as $`t\mathrm{}`$ (Iscoe<sup>(27)</sup>, pp. 106–107 and 112). $`\mathrm{}`$
Proof of Theorems 2.2.3 and 2.5.3 (2–level branching system):
We will follow the same steps for the proof of the 1–level case, but now some of them are harder. The problem is that, while the test functions $`\phi 𝒞_c^+(S)`$ and the measure $`\rho `$ for the 1–level system are not so difficult to work with, for the 2–level system the test functions $`\mu \mu ,\phi `$ and the measure $`R_{\mathrm{}}^1`$ (which now plays the role of $`\rho `$) raise new technical questions that are not easy to deal with, in particular involving the third moments of $`R_{\mathrm{}}^1`$. The background on the 2–level system in the Appendix should be consulted at this point.
We continue to use the notation $`\phi _t`$ introduced in (5.2.1), but now with $`\beta =\beta _2`$, and we put $`f_t1`$ for Theorem 2.2.3(a) and Theorem 2.5.2.
For $`\phi 𝒞_c^+(S)`$ we have, by (A.1.13) and (A.1.14) (Appendix),
$$E\mathrm{exp}\left\{F_t^{1/(1+\beta _2)}_0^tX_s𝑑s,\phi \right\}=\mathrm{exp}\{R_{\mathrm{}}^1,𝐮_{\phi _t}(t)\},$$
where $`𝐮_{\phi _t}(s)={\displaystyle \frac{V_2}{1+\beta _2}}{\displaystyle _0^s}U_{sr}(𝐮_{\phi _t}(r)^{1+\beta _2})𝑑r+{\displaystyle _0^s}U_{sr}(,\phi _t(1𝐮_{\phi _t}(r)()))𝑑r.(\mathrm{5.2.17})`$Hence, by $`U_t`$–invariance of $`R_{\mathrm{}}^1`$ and $`E_0^tX_s𝑑s,\phi =t\rho ,\phi `$, $`E\mathrm{exp}\{F_t^{1/(1+\beta _2)}Y_t,\phi \}=\mathrm{exp}\{{\displaystyle \frac{V_2}{1+\beta _2}}(I_1(t)+I_2(t))+I_3(t))\},(\mathrm{5.2.18})`$ where $`I_1(t)={\displaystyle _0^t}R_{\mathrm{}}^1,𝐰_{\phi _t}(s)^{1+\beta _2}𝑑s,(\mathrm{5.2.19})`$$`I_2(t)={\displaystyle _0^t}R_{\mathrm{}}^1,𝐮_{\phi _t}(s)^{1+\beta _2}𝐰_{\phi _t}(s)^{1+\beta _2}𝑑s,(\mathrm{5.2.20})`$$`I_3(t)={\displaystyle _0^t}R_{\mathrm{}}^1,,\phi _t𝐮_{\phi _t}(s)()𝑑s,(\mathrm{5.2.21})`$with, by (A.1.4), $`𝐰_\phi (s)(\mu ):={\displaystyle _0^s}U_r(,\phi )(\mu )𝑑r={\displaystyle _0^s}\mu ,T_r\phi 𝑑r=\mu ,G_s\phi ,\mu _\tau (S).(\mathrm{5.2.22})`$
We will prove the following limits as $`t\mathrm{}`$:
For Theorem 2.2.3(a): $`I_1(t)(\phi ,G^2\phi )_\rho +\frac{V_1}{2}(\phi ,G^3\phi )_\rho .(\mathrm{5.2.23})`$For Theorem 2.2.3(b): $`I_1(t){\displaystyle \frac{V_1}{2}}H(\phi ,\phi ).(\mathrm{5.2.24})`$For Theorem 2.5.2: $`I_1(t)R_{\mathrm{}}^1,,G\phi ^{1+\beta _2}.(\mathrm{5.2.25})`$For $`\beta _21`$: $`I_2(t)0.(\mathrm{5.2.26})`$For $`\beta _2<1`$: $`I_3(t)0.(\mathrm{5.2.27})`$For Theorem 2.2.3(a): $`I_3(t)(\phi ,G\phi )_\rho +{\displaystyle \frac{V_1}{2}}(\phi ,G^2\phi )_\rho .(\mathrm{5.2.28})`$For Theorem 2.2.3(b): $`I_3(t)0.(\mathrm{5.2.29})`$Proof of (5.2.23) and (5.2.24): We have from (5.2.19) and (5.2.22)
$$I_1(t)=F_t^1_0^tR_{\mathrm{}}^1,,G_s\phi ^2𝑑s.$$
By L’Hôpital’s rule and using (A.1.12) we have, for (5.2.23),
$`I_1(t)`$ $`=`$ $`t^1{\displaystyle _0^t}R_{\mathrm{}}^1,,G_s\phi ^2𝑑sR_{\mathrm{}}^1,,G_t\phi ^2`$
$`=`$ $`(\phi ,G_t^2\phi )_\rho +{\displaystyle \frac{V_1}{2}}(\phi ,G_t^2G\phi )_\rho (\phi ,G^2\phi )_\rho +{\displaystyle \frac{V_1}{2}}(\phi ,G^3\phi )_\rho ,`$
and for (5.2.24),
$`I_1(t)`$ $`=`$ $`F_t^1{\displaystyle _0^t}R_{\mathrm{}}^1,,G_s\phi ^2𝑑s`$
$``$ $`{\displaystyle \frac{1}{f_t}}R_{\mathrm{}}^1,,G_t\phi ^2`$
$`=`$ $`{\displaystyle \frac{1}{f_t}}\left((\phi ,G_t^2\phi )_\rho +{\displaystyle \frac{V_1}{2}}(\phi ,G_t^2G\phi )_\rho \right)`$
$`{\displaystyle \frac{V_1}{2}}H(\phi ,\phi ).`$
Proof of (5.2.25): The same as (5.2.23).
Proof of (5.2.26): We rewrite (5.2.20) as
$`I_2(t)`$ $`=`$ $`{\displaystyle _0^t}R_{\mathrm{}}^1,𝐰_{\phi _t}(s)^{1+\beta _2}𝐮_{\phi _t}(s)^{1+\beta _2}𝑑s`$
$`=`$ $`{\displaystyle _0^t}R_{\mathrm{}}^1,𝐰_{\phi _t}(s)^{1+\beta }[1(𝐮_{\phi _t}(s)/𝐰_{\phi _t}(s))^{1+\beta _2}]𝑑s.`$
Since $`01(𝐮_{\phi _t}(s)/𝐰_{\phi _t}(s))^{1+\beta _2}1`$ by (5.2.2) and (A.1.16), and since $`I_1(t)`$ converges, it suffices to prove that
$$\underset{t\mathrm{}}{lim}\frac{𝐮_{\phi _t}(s)}{𝐰_{\phi _t}(s)}=1\text{for all}s,R_{\mathrm{}}^1a.e.$$
We have from (5.2.22) and (A.1.16)
$$\frac{𝐮_{\phi _t}(s)(\mu )}{𝐰_{\phi _t}(s)(\mu )}=1\mu ,G_s\phi ^1(J_t(s)+K_t(s))(\mu )0,$$
where
$$J_t(s)(\mu )=\frac{V_2}{1+\beta _2}F_t^{1/(1+\beta _2)}_0^sU_{sr}(𝐮_{\phi _t}(r)^{1+\beta _2})(\mu )𝑑r0$$
and
$$K_t(s)(\mu )=_0^sU_{sr}(,\phi 𝐮_{\phi _t}(r)())(\mu )𝑑r0.$$
By (A.1.17) and (5.2.22)
$`J_t(s)(\mu )`$ $``$ $`{\displaystyle \frac{V_2}{1+\beta _2}}F_t^{1/(1+\beta _2)}{\displaystyle _0^s}U_{sr}(𝐰_{\phi _t}(r)^{1+\beta _2})(\mu )𝑑r`$
$`=`$ $`{\displaystyle \frac{V_2}{1+\beta _2}}F_t^{\beta _2/(1+\beta _2)}{\displaystyle _0^s}U_{sr}(,G_r\phi ^{1+\beta _2})(\mu )𝑑r0.`$
Similarly, $`K_t(s)(\mu )0`$, so the result follows. Proof of (5.2.27): We have from (5.2.21), (5.2.22) and (A.1.12)
$`I_3(t)t^{2/(1+\beta _2)}{\displaystyle _0^t}R_{\mathrm{}}^1,,\phi ,G_s\phi 𝑑s`$
$`=t^{(\beta _21)/(1+\beta _2)}t^1{\displaystyle _0^t}((\phi ,G_s\phi )_\rho +{\displaystyle \frac{V_1}{2}}(\phi ,G_sG\phi )_\rho ))ds.`$
Since
$$(\phi ,G_s\phi )_\rho +\frac{V_1}{2}(\phi ,G_sG\phi )(\phi ,G\phi )_\rho +\frac{V_1}{2}(\phi ,G^2\phi )_\rho \text{as}s\mathrm{},$$
the result follows. Proof of (5.2.28): We rewrite (5.2.21) as
$$I_3(t)=t^1_0^tR_{\mathrm{}}^1,,\phi \stackrel{~}{𝐮}_{\phi _t}(s)()𝑑s,$$
where $`\stackrel{~}{𝐮}_{\phi _t}(s)=t^{1/2}𝐮_{\phi _t}(s)`$. Since
$$t^1_0^tR_{\mathrm{}}^1,,\phi ,G_s\phi (\phi ,G\phi )_\rho +\frac{V_1}{2}(\phi ,G^2\phi )_\rho ,$$
by (A.1.12), it suffices to prove that
$$A_t:=t^1_0^tR_{\mathrm{}}^1,,\phi \stackrel{~}{𝐮}_{\phi _t}(s)()𝑑st^1_0^tR_{\mathrm{}}^1,,\phi ,G_s\phi 𝑑s0.$$
We have from (A.1.16)
$$\stackrel{~}{𝐮}_{\phi _t}(s)(\mu )=\frac{V_2}{2}t^{1/2}_0^sU_{sr}(\stackrel{~}{𝐮}_{\phi _t}(r)^2)(\mu )𝑑r+\mu ,G_s\phi t^{1/2}_0^sU_{sr}(,\phi \stackrel{~}{𝐮}_{\phi _t}(r)())(\mu )𝑑r,$$
hence
$`|A_t|`$ $``$ $`t^{3/2}({\displaystyle \frac{V_2}{2}}{\displaystyle _0^t}{\displaystyle _0^s}{\displaystyle }R_{\mathrm{}}^1(d\mu )\mu ,\phi U_{sr}(\stackrel{~}{𝐮}_{\phi _t}(r)^2)(\mu )drds`$
$`+{\displaystyle _0^t}{\displaystyle _0^s}{\displaystyle }R_{\mathrm{}}^1(d\mu )\mu ,\phi U_{sr}(,\phi \stackrel{~}{𝐮}_{\phi _t}(r))(\mu )drds).`$
Now $`\stackrel{~}{𝐮}_{\phi _t}(r)(\mu )\mu ,G_r\phi `$ (since $`f_t1`$). Note that this estimate is not too rough because $`\stackrel{~}{𝐮}_{\phi _t}(r)(\mu )\mu ,G_r\phi `$ as $`t\mathrm{}`$. Hence
$$|A_t|const.(H_1(t)+H_2(t)),$$
where
$`H_1(t)`$ $`=`$ $`t^{3/2}{\displaystyle _0^t}{\displaystyle _0^s}{\displaystyle R_{\mathrm{}}^1(d\mu )\mu ,\phi U_{sr}(,G_r\phi ^2)(\mu )𝑑r𝑑s},`$
$`H_2(t)`$ $`=`$ $`t^{3/2}{\displaystyle _0^t}{\displaystyle _0^s}{\displaystyle R_{\mathrm{}}^1(d\mu )\mu ,\phi U_{sr}(,\phi ,G_r\phi )(\mu )𝑑r𝑑s}.`$
We will show that $`H_1(t)0`$. The proof that $`H_2(t)0`$ is similar.
Using (A.1.8) we have
$$H_1(t)const.\underset{j=1}{\overset{3}{}}J_j(t),$$
where
$`J_1(t)`$ $`=`$ $`t^{3/2}{\displaystyle _0^t}{\displaystyle _0^s}{\displaystyle R_{\mathrm{}}^1(d\mu )\mu ,\phi \mu ,T_{sr}G_r\phi ^2𝑑r𝑑s},`$
$`J_2(t)`$ $`=`$ $`t^{3/2}{\displaystyle _0^t}{\displaystyle _0^s}{\displaystyle R_{\mathrm{}}^1(d\mu )\mu ,\phi \mu ,T_{sr}(G_r\phi )^2𝑑r𝑑s},`$
$`J_3(t)`$ $`=`$ $`t^{3/2}{\displaystyle _0^t}{\displaystyle _0^s}{\displaystyle R_{\mathrm{}}^1(\mu )\mu ,\phi _0^{sr}\mu ,T_u(T_{sru}G_r\phi )^2𝑑u𝑑r𝑑s}.`$
By (A.1.12) and (A.1.13) we obtain
$$J_1(t)const.\underset{j=1}{\overset{5}{}}K_{1,j}(t),J_2(t)const.\underset{j=1}{\overset{2}{}}K_{2,j}(t),J_3(t)const.\underset{j=1}{\overset{2}{}}K_{3,j}(t),$$
where
$`K_{1,1}(t)`$ $`=`$ $`t^{3/2}{\displaystyle _0^t}{\displaystyle _0^s}\rho ,\phi (T_{sr}G_r\phi )^2𝑑r𝑑s,`$
$`K_{1,2}(t)`$ $`=`$ $`t^{3/2}{\displaystyle _0^t}{\displaystyle _0^s}\rho ,\phi T_{sr}G_r\phi T_{sr}GG_r\phi 𝑑r𝑑s,`$
$`K_{1,3}(t)`$ $`=`$ $`t^{3/2}{\displaystyle _0^t}{\displaystyle _0^s}\rho ,\phi G(T_{sr}G_r\phi )^2𝑑r𝑑s,`$
$`K_{1,4}(t)`$ $`=`$ $`t^{3/2}{\displaystyle _0^t}{\displaystyle _0^s}{\displaystyle _0^{\mathrm{}}}\rho ,\phi GT_u(T_uT_{sr}G_r\phi )^2𝑑u𝑑r𝑑s,`$
$`K_{1,5}(t)`$ $`=`$ $`t^{3/2}{\displaystyle _0^t}{\displaystyle _0^s}{\displaystyle _0^{\mathrm{}}}\rho ,T_{sr}G_r\phi GT_u(T_u\phi T_uT_{sr}G_r\phi )𝑑u𝑑r𝑑s,`$
$`K_{2,1}(t)`$ $`=`$ $`t^{3/2}{\displaystyle _0^t}{\displaystyle _0^s}\rho ,\phi T_{sr}(G_r\phi )^2𝑑r𝑑s,`$
$`K_{2,2}(t)`$ $`=`$ $`t^{3/2}{\displaystyle _0^t}{\displaystyle _0^s}\rho ,\phi GT_{sr}(G_r\phi )^2𝑑r𝑑s,`$
$`K_{3,1}(t)`$ $`=`$ $`t^{3/2}{\displaystyle _0^t}{\displaystyle _0^s}{\displaystyle _0^{sr}}\rho ,\phi T_u(T_{sru}G_r\phi )^2𝑑u𝑑r𝑑s,`$
$`K_{3,2}(t)`$ $`=`$ $`t^{3/2}{\displaystyle _0^t}{\displaystyle _0^s}{\displaystyle _0^{sr}}\rho ,\phi GT_u(T_{sru}G_r\phi )^2𝑑u𝑑r𝑑s.`$
We will show that each of these terms converges to $`0`$ as $`t\mathrm{}`$. Recall that $`G^j\phi <\mathrm{},j=1,2,3`$, for $`\phi 𝒞_c^+(S)`$.
$$K_{1,1}(t)G\phi t^{3/2}_0^t\rho ,\phi G_sG\phi 𝑑sG\phi G^2\phi \rho ,\phi t^{1/2}0.$$
$$K_{1,2}(t)G^2\phi t^{3/2}_0^t\rho ,\phi G_sG\phi 𝑑sG^2\phi ^2\rho ,\phi t^{1/2}0.$$
$$K_{1,3}(t)G\phi t^{3/2}_0^t\rho ,\phi G_sG^2\phi 𝑑sG\phi G^3\phi \rho ,\phi t^{1/2}0.$$
$`K_{1,4}(t)`$ $`=`$ $`t^{3/2}{\displaystyle _0^t}{\displaystyle _0^s}{\displaystyle _0^{\mathrm{}}}\rho ,\phi GT_u(T_uT_rG_{sr}\phi )^2𝑑u𝑑r𝑑s`$
$``$ $`const.t^{1/2}{\displaystyle _0^t}{\displaystyle _0^{\mathrm{}}}\rho ,\phi GT_u(T_{u+r}G_{tr}\phi )^2𝑑u𝑑r\text{(by l’Hôpital)}`$
$``$ $`const.t^{1/2}{\displaystyle _0^t}{\displaystyle _0^{\mathrm{}}}\rho ,\phi GT_u(T_{u+r}G_{tr}\phi T_{u+r}T_{tr}\phi )𝑑u𝑑r\text{(by l’Hôpital)}`$
$``$ $`const.t^{1/2\delta }{\displaystyle _0^t}{\displaystyle _0^{\mathrm{}}}\rho ,\phi GT_u(T_{u+r}G_{tr}\phi (t+u)^\delta T_{t+u}\phi )𝑑u𝑑r`$
$``$ $`const.t^{3/2\delta }G^3\phi \rho ,\phi \text{(by (2.2.1))}`$
$`0.`$
$`K_{1,5}(t)`$ $`=`$ $`t^{3/2}{\displaystyle _0^t}{\displaystyle _0^s}{\displaystyle _0^{\mathrm{}}}\rho ,T_{sr}G_r\phi T_u(T_u\phi T_uT_rG_{sr}\phi )𝑑u𝑑r𝑑s`$
$``$ $`t^{3/2}{\displaystyle _0^t}{\displaystyle _0^s}{\displaystyle _0^{\mathrm{}}}\rho ,\phi G^2T_u(T_u\phi T_{u+r}G_{sr}\phi )𝑑u𝑑r𝑑s`$
$``$ $`const.t^{1/2}{\displaystyle _0^t}{\displaystyle _0^{\mathrm{}}}\rho ,\phi G^2T_u(T_u\phi T_{u+r}G_{tr}\phi )𝑑u𝑑r\text{(by l’Hôpital)}`$
$``$ $`const.t^{1/2}{\displaystyle _0^t}{\displaystyle _0^{\mathrm{}}}\rho ,\phi G^2T_u(T_u\phi T_{u+r}T_{tr}\phi )𝑑u𝑑r\text{(by l’Hôpital)}`$
$``$ $`const.t^{1/2\delta }{\displaystyle _0^t}{\displaystyle _0^{\mathrm{}}}\rho ,\phi G^2T_u(T_u\phi (t+u)^\delta T_{t+u}\phi )𝑑u𝑑r`$
$``$ $`const.t^{3/2\delta }G^3\phi \rho ,\phi \text{(by (2.2.1))}`$
$`0.`$
$`K_{2,1}(t)0`$, similarly to $`K_{1,1}(t)0.`$
$`K_{2,2}(t)0`$, similarly to $`K_{1,3}(t)0`$.
$`K_{3,1}(t)`$ $``$ $`G\phi t^{3/2}{\displaystyle _0^t}{\displaystyle _0^s}\rho ,\phi rT_rG\phi 𝑑r𝑑s`$
$``$ $`const.G\phi G^3\phi \rho ,\phi t^{1/2}0.`$
$`K_{3,2}(t)`$ $`=`$ $`t^{2/3}{\displaystyle _0^t}{\displaystyle _0^s}{\displaystyle _0^r}\rho ,\phi GT_u(T_{ru}G_{sr}\phi )^2𝑑u𝑑r𝑑s`$
$``$ $`const.t^{1/2}{\displaystyle _0^t}{\displaystyle _0^r}\rho ,\phi GT_u(T_{ru}G_{tr}\phi )^2𝑑u𝑑r\text{(by l’Hôpital)}`$
$``$ $`const.t^{1/2}{\displaystyle _0^t}{\displaystyle _0^r}\rho ,\phi GT_u(T_{ru}G_{tr}\phi T_{ru}\phi )𝑑u𝑑r\text{(by l’Hôpital)}`$
$`=`$ $`const.t^{1/2}{\displaystyle _0^t}{\displaystyle _0^r}\rho ,\phi GT_{ru}(T_uG_{tr}\phi T_uT_{tr}\phi )𝑑u𝑑r`$
$`=`$ $`const.t^{1/2}{\displaystyle _0^r}{\displaystyle _0^{tr}}\rho ,\phi GT_{tru}(T_uG_r\phi T_{u+r}\phi )𝑑u𝑑r`$
$``$ $`const.(M_1H)+M_2(t)),`$
where (by l’Hôpital),
$$M_1(t)=t^{3/2}_0^t\rho ,\phi G(T_{tr}G_r\phi T_t\phi )𝑑r$$
and, since $`{\displaystyle \frac{d}{dt}}GT_t\phi =T_t`$,
$$M_2(t)=t^{3/2}_0^t_0^{tr}\rho ,\phi T_{tru}(T_uG_r\phi T_{u+r}\phi )𝑑u𝑑r.$$
Now,
$`M_1(t)`$ $`=`$ $`t^{3/2\delta }{\displaystyle _0^t}\rho ,\phi G(T_{tr}G_r\phi t^\delta T_t\phi dr`$
$``$ $`const.t^{5/2\delta }G^2\phi \rho ,\phi \text{(by (2.2.1))}`$
$`0,`$
$`|M_2(t)|`$ $``$ $`t^{3/2}{\displaystyle _0^t}{\displaystyle _0^r}\rho ,\phi T_{ru}(T_uG_{tr}\phi T_{t(ru)}\phi dudr`$
$`=`$ $`t^{3/2}{\displaystyle _0^t}{\displaystyle _0^r}\rho ,\phi T_u(T_{ru}G_{tr}\phi T_{tu}\phi )𝑑u𝑑r`$
$`=`$ $`t^{3/2}{\displaystyle _0^t}\rho ,\phi T_u\left({\displaystyle _u^t}T_{ru}G_{tr}\phi 𝑑rT_{tu}\phi \right)𝑑u`$
$``$ $`G^2\phi t^{5/2}\rho ,\phi T_t\phi `$
$`=`$ $`G^2\phi t^{5/2\delta }\rho ,\phi t^\delta T_t\phi `$
$``$ $`const.G^2\phi \rho ,\phi t^{5/2\delta }\text{(by (2.2.1))}`$
$`0.`$
Finally, the weak convergence of $`Y_t`$ as $`t\mathrm{}`$ follows as in the $`1`$–level case. $`\mathrm{}`$ 5.3. Examples Proof of Lemma 3.1.1: $`\alpha `$–stable process:
All the proofs for the $`\alpha `$–stable case can be done by using the self–similarly of the transition probability $`p_t`$. We will prove only (c) and (d) to exemplify. $`\text{(c)}G_t\phi (x)={\displaystyle _0^t}{\displaystyle _\text{}^d}p_s(xy)\phi (y)𝑑y𝑑s={\displaystyle _0^t}s^{d/\alpha }{\displaystyle _\text{}^d}p_1(s^{1/\alpha }(xy))\phi (y)𝑑y𝑑s`$
$$=t^{d/\alpha +1}_0^1r^{d/\alpha }_\text{}^dp_1(t^{1/\alpha }r^{1/\alpha }(xy))\phi (y)𝑑y𝑑r,$$
hence
$$\underset{t\mathrm{}}{lim}t^{1d/\alpha }G_t\phi (x)=\frac{1}{1d/\alpha }p_1(0)_\text{}^d\phi (y)𝑑y,$$
and (by Fourier transform)
$$p_1(0)=\frac{1}{(2\pi )^d}_\text{}^de^{|z|^\alpha }𝑑z.$$
$`\text{(d)}G_t^{}\phi (x)={\displaystyle _\text{}^d}p_t(xy)\phi (y)𝑑y=t^1{\displaystyle _\text{}^d}p_1(t^{1/\alpha }(xy))\phi (y)𝑑y,`$ hence
$$\underset{t\mathrm{}}{lim}tG_t^{}\phi (x)=p_1(0)_\text{}^d\phi (y)𝑑y,$$
with
$$p_1(0)=\frac{1}{(2\pi )^d}_\text{}^de^{|z|^d}𝑑z,$$
and the result follows by l’Hôpital’s rule. Hierarchical random walk:
Recall that $`a^{\gamma +1}={\displaystyle \frac{1}{N}}`$, by (4.2.3), and the definitions of $`\theta `$ and the functions $`h`$ in Subsection 4.2 (see (4.2.4)–(4.2.9)).
(a) By (3.1.3), up to a summand which converges to 0 exponentially fast as $`t\mathrm{},p_t(0,x)`$ equals
$`(N1){\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{N^j}}e^{\theta a^jt}`$ $`=`$ $`(N1)(\theta t)^{(\gamma +1)}{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}e^{\theta a^jt}(\theta a^jt)^{\gamma +1}`$
$`=`$ $`q_\gamma t^{(\gamma +1)}\left(h_t^{(1,\gamma +1)}{\displaystyle \underset{j0}{}}e^{\theta a^jt}(\theta a^jt)^{\gamma +1}\right).`$
To conclude the proof of (a) if suffices to observe that $`_{j0}e^{\theta a^jt}(\theta a^jt)^{\gamma +1}0`$ as $`t\mathrm{}`$. Indeed for all $`t`$ and all negative integers $`j`$ we have
$$(\theta a^jt)^{\gamma +1}e^{\theta a^jt}const.\frac{1}{\theta a^jt},$$
hence $`\underset{j0}{}e^{\theta a^jt}(\theta a^jt)^{\gamma +1}`$ is majorized by $`{\displaystyle \frac{1}{t}}`$ times a convergent geometric series. (b) Because of (a) it suffices to compute
$`{\displaystyle _t^{\mathrm{}}}s^{(\gamma +1)}h_s^{(1,\gamma +1)}𝑑s`$ $`=`$ $`{\displaystyle _t^{\mathrm{}}}{\displaystyle \underset{j=\mathrm{}}{\overset{\mathrm{}}{}}}(\theta a^j)^{\gamma +1}e^{\theta a^js}ds`$
$`=`$ $`{\displaystyle \underset{j=\mathrm{}}{\overset{\mathrm{}}{}}}(\theta a^j)^\gamma e^{\theta a^jt}=t^\gamma h_t^{(1,\gamma )}.`$
(c) Since $`G_t(0,x)=_0^tp_s(0,x)𝑑s`$, we infer from (a) that
$`G_t(0,x)`$ $``$ $`{\displaystyle _0^t}s^{(\gamma +1)}h_s^{(1,\gamma +1)}𝑑s`$
$`=`$ $`q_\gamma {\displaystyle _0^t}{\displaystyle \underset{j=\mathrm{}}{\overset{\mathrm{}}{}}}(\theta a^j)^{\gamma +1}e^{\theta a^js}ds`$
$`=`$ $`q_\gamma {\displaystyle \underset{j\mathrm{}}{\overset{\mathrm{}}{}}}(\theta a^j)^\gamma (1e^{\theta a^jt})`$
$`=`$ $`q_\gamma t^\gamma h_t^{(2,\gamma )}.`$
(d) Since $`a={\displaystyle \frac{1}{N}}`$ for $`\gamma =0`$, up to a summand which is uniformly bounded in $`t`$, $`G_t(0,x)`$ equals
$$\frac{N1}{\theta }\underset{j=1}{\overset{\mathrm{}}{}}(1e^{\theta a^jt}).$$
Since $`a^{j+1}a^ya^i`$ for $`y[j,j+1]`$, then $`{\displaystyle _1^{\mathrm{}}}(1e^{a^yt})𝑑y{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}(1e^{a^jt}){\displaystyle _1^{\mathrm{}}}(1e^{a^ya^1t})𝑑y.(\mathrm{5.3.1})`$Now,
$$_1^{\mathrm{}}(1e^{a^yt})𝑑y=\frac{1}{\mathrm{log}a}_0^a\frac{1e^{zt}}{z}𝑑z=\frac{1}{\mathrm{log}a}_0^{at}\frac{1e^r}{r}𝑑r,$$
and by l’Hôpital’s rule,
$$_0^{at}\frac{1e^r}{r}𝑑r\mathrm{log}t.$$
Hence it follows from (5.3.1) that
$$\underset{j=1}{\overset{\mathrm{}}{}}(1e^{a^jt})\frac{\mathrm{log}t}{\mathrm{log}a},$$
which implies the result. (e) Since $`G_t^2(0,x)=_0^t_0^tp_{u+v}(0,x)𝑑u𝑑v`$, we infer from (a) that
$`G_t^2(0,x)`$ $``$ $`q_\gamma {\displaystyle \underset{j=\mathrm{}}{\overset{\mathrm{}}{}}}(\theta a^j)^{\gamma +1}{\displaystyle _0^t}{\displaystyle _0^t}e^{\theta a^j(u+v)}𝑑u𝑑v`$
$`=`$ $`q_\gamma {\displaystyle \underset{j=\mathrm{}}{\overset{\mathrm{}}{}}}(\theta a^j)^{\gamma 1}(1e^{\theta a^jt})^2`$
$`=`$ $`q_\gamma t^{1\gamma }h_t^{(3,\gamma 1)}.`$
(f) Since $`a={\displaystyle \frac{1}{N^{1/2}}}`$ for $`\gamma =1`$, up to a summand which is uniformly bounded in $`t`$, $`G_t^2(0,x)`$ equals
$$\frac{N1}{\theta ^2}\underset{j=1}{\overset{\mathrm{}}{}}(1e^{\theta a^jt})^2.$$
The same argument used for (d) shows that
$$\underset{j=1}{\overset{\mathrm{}}{}}(1e^{a^jt})^2\frac{\mathrm{log}t}{\mathrm{log}a},$$
and the result follows.
(g) Since $`G_t^2G(0,x)=_0^{\mathrm{}}_0^t_0^tp_{s+u+v}(0,x)𝑑u𝑑v𝑑s`$, we infer from (a) that
$`G_t^2G(0,x)`$ $``$ $`q_\gamma {\displaystyle \underset{j=\mathrm{}}{\overset{\mathrm{}}{}}}(\theta a^j)^{\gamma +1}{\displaystyle _0^{\mathrm{}}}{\displaystyle _0^t}{\displaystyle _0^t}e^{\theta a^j(s+u+v)}𝑑u𝑑v𝑑s`$
$`=`$ $`q_\gamma {\displaystyle \underset{j=\mathrm{}}{\overset{\mathrm{}}{}}}(\theta a^j)^{\gamma 2}(1e^{\theta a^jt})^2`$
$`=`$ $`q_\gamma t^{2\gamma }h_t^{(3,\gamma 2)}.`$
(h) Since $`a={\displaystyle \frac{1}{N^{1/3}}}`$ for $`\gamma =2`$, up to a summand which is uniformly bounded in $`t`$, $`G_t^2G(0,x)`$ equals
$$\frac{N1}{\theta ^3}\underset{j=1}{\overset{\mathrm{}}{}}(1e^{\theta a^jt})^2,$$
and the proof is the same as for (f). $`\mathrm{}`$ Proof of Lemma 3.1.2: By the observation before Lemma 2.4.1, $`G^j\phi <\mathrm{}`$ for $`𝒞_c^+(S)`$ if and only if $`j<\gamma +1`$, i.e. $`\alpha j<d`$ for the $`\alpha `$–stable case, and $`c>N^{(j1)/j}`$ for the $`c`$–hierarchical case.
The expression for $`G^j`$ can be obtained by formula (2.1.1). For the $`c`$–hierarchical case the form of the semigroup $`T_t`$ is explicit from the transition probability $`p_t(0,x)`$ given in (3.1.3). For the $`\alpha `$–stable case the transition probability $`p_t(0,x)`$ is not known explicitly for general $`\alpha `$, but its Fourier transform,
$$_\text{}^de^{ixz}p_t(0,x)𝑑x=e^{t|z|^\alpha },$$
can be used for the proof. $`\mathrm{}`$ Proof of Lemma 3.1.3: We sketch the main idea of the proof. By Lemma 3.1.1(a), $`T_tt^{(\gamma +1)}`$. By Lemma 2.4.1, $`\gamma >k+1`$. Hence $`T_t=o(t^{(k+2)}).\mathrm{}`$ Proof of Proposition 3.1.1: Let $`h_t`$ denote any of the functions defined in (4.2.5) – (4.2.7). We write $`h_t`$ as $`h_t=h_t^{()}+h_t^{(+)}`$, where $`h_t^{()}`$ and $`h_t^{(+)}`$ stand for the sums $`\underset{j}{}`$ with $`j<0`$ and $`j0`$, respectively. Since $`lim_t\mathrm{}h_t^{()}=0`$ (as in the proof of Lemma 3.1.1 for the hierarchical case), and $`h_t`$ is periodic in a logarithmic scale, in order to prove that
$$L_1\underset{t}{inf}h_t\underset{t}{sup}h_tL_2,$$
for some positive constants $`L_1`$ and $`L_2`$, it suffices to show that $`L_1\underset{t\mathrm{}}{lim\; inf}h_t^{(+)}\underset{t\mathrm{}}{lim\; sup}h_t^{(+)}L_2.(\mathrm{5.3.2})`$
We will prove (5.3.2) for $`h_t^{(1,\zeta )}`$. Using the formula
$$q^j=\frac{q\mathrm{log}q}{q1}_j^{j+1}q^y𝑑y,q>0,q1,$$
and $`a^{j+1}a^ya^j`$ for $`jyj+1`$ (since $`0<a<1`$), we obtain
$$\frac{a^\zeta \mathrm{log}a^\zeta }{a^\zeta 1}_0^{\mathrm{}}a^{y\zeta }e^{a^ya^1t}𝑑y\underset{j=0}{\overset{\mathrm{}}{}}a^{j\zeta }e^{a^jt}\frac{a^\zeta \mathrm{log}a^\zeta }{a^\zeta 1}_0^{\mathrm{}}a^{y\zeta }e^{a^yt}𝑑y.$$
We have
$$_0^{\mathrm{}}a^{y\zeta }e^{a^yt}𝑑y=\frac{1}{\mathrm{log}a}_0^1z^{\zeta 1}e^{zt}𝑑z=\frac{t^\zeta }{\mathrm{log}a}_0^tr^{\zeta 1}e^r𝑑r,$$
and since $`_0^tr^{\zeta 1}e^r𝑑r\mathrm{\Gamma }(\zeta )`$ as $`t\mathrm{}`$, putting these results together we obtain
$$\frac{\mathrm{\Gamma }(\zeta +1)}{a^\zeta 1}\underset{t\mathrm{}}{lim\; inf}t^\zeta \underset{j=0}{\overset{\mathrm{}}{}}a^{j\zeta }e^{a^jt}\underset{t\mathrm{}}{lim\; sup}t^\zeta \underset{j=0}{\overset{\mathrm{}}{}}a^{j\zeta }e^{a^jt}\frac{a^\zeta \mathrm{\Gamma }(\zeta +1)}{a^\zeta 1},$$
which finishes the proof.
This method can be used for the other functions $`h_t`$ as well, with slightly more elaborate calculations. For $`\stackrel{~}{h}_t^{(2,\zeta )}`$ and $`\stackrel{~}{h}_t^{(3,\zeta )}`$ we need to use the fact that the functions $`xe^x1+x`$ and $`x2e^x\frac{1}{2}e^{2x}+x\frac{3}{2}`$, respectively, are increasing. However, we can obtain bounds for $`\stackrel{~}{h}_t^{(2,\zeta )}`$ and $`\stackrel{~}{h}_t^{(3,\zeta )}`$ from the bounds for $`h_t^{(2,\zeta )}`$ and $`h_t^{(3,\zeta )}`$, simply by dividing them by $`1\zeta `$. This is clear from the form of $`\stackrel{~}{h}_t`$ in Lemma 2.4.3. $`\mathrm{}`$ Proof of Theorem 3.2.1: The proof is a direct application of Theorems 2.2.1, 2.2.1, 2.2.3, Lemmas 3.1.1, 3.1.2, 3.1.3, and Corollary 3.1.1. $`\mathrm{}`$ 5.4. Conditions for the results on infinite variance branching results Proof of Porposition 3.2.1: We have to show that $`d>\alpha \left(1+{\displaystyle \frac{1}{\beta }}\right)(\mathrm{5.4.1})`$ is necessary and sufficient for condition (2.5.1) of Theorem 2.5.1.
That (5.4.1) implies (2.5.1) is proved by Iscoe<sup>(27)</sup>. For the converse, note that
$$G1_B(x)k|x|^{(d\alpha )}\text{for }|x|2,$$
where $`B`$ denotes the unit ball centered at the origin and $`k`$ is some positive constant. Hence, if $`d\alpha (1+1/\beta ),(G1_B)^{1+\beta }`$ is not $`\lambda `$–integrable. $`\mathrm{}`$
For the proof of Proposition 3.2.2 we need some preliminary results.
Let $`R_{\mathrm{}}^1`$ and $`R_{\mathrm{}}`$ be the canonical measure of the equilibrium of the particle system (started off in the Poisson system $`\mathrm{\Pi }_\lambda `$ with intensity $`\lambda `$) and that of the superprocess (started off in $`\lambda `$), respectively (Appendix, Subsection A.1).
Lemma 5.4.1. Assume that $`\beta _2<\beta _1`$ and let $`\phi :^d[0,\mathrm{}]`$. Then
(a)
$$_{_\tau (S)}\nu ,\phi ^{1+\beta _2}R_{\mathrm{}}(d\nu )_{_\tau (S)}\nu ,\phi ^{1+\beta _2}R_{\mathrm{}}^1(d\nu ).$$
b) If $`\phi `$ is $`\lambda `$-integrable, then
$$_{_\tau (S)}\nu ,\phi ^{1+\beta _2}R_{\mathrm{}}^1(d\nu )C<\mathrm{},$$
where the constant $`C`$ depends only on $`d,\alpha ,\beta _1,\beta _2`$ and $`\lambda ,\phi `$.
Proof: (a) We have, by Jensen’s inequality and (A.1.10),
$`{\displaystyle \nu ,\phi ^{1+\beta _2}R_{\mathrm{}}\left(d\nu \right)}`$ $`=`$ $`{\displaystyle \left(\mu ,\phi \mathrm{\Pi }_\nu \left(d\mu \right)\right)^{1+\beta _2}R_{\mathrm{}}\left(d\nu \right)}`$
$``$ $`{\displaystyle \mu ,\phi ^{1+\beta _2}\mathrm{\Pi }_\nu \left(d\mu \right)R_{\mathrm{}}\left(d\nu \right)}`$
$`=`$ $`{\displaystyle \mu ,\phi ^{1+\beta _2}R_{\mathrm{}}^1\left(d\nu \right)}.`$
(b) By the Palm formula (A.2.1) it suffices to show for some $`\delta >\beta _2`$,
$$\nu ,\phi ^\delta (R_{\mathrm{}}^1)_x(d\nu )<C<\mathrm{},$$
where the constant $`C`$ does not depend on $`x`$. We can choose $`\delta (\beta _2,\beta _1)`$ such that $`d\beta _2/\alpha >1`$.
We will use the tree representation of $`(R_{\mathrm{}}^1)_x^{red}`$ given in (A.3.2), and we denote $`Z_{t,i}=X_t^{W_t^x,i},\phi `$. Since $`\left(_{j=1}^na_j\right)^\delta _{j=1}^na_j^\delta `$, for all nonnegative sequences $`(a_j)`$ and $`0<\delta <1`$, we have
$`{\displaystyle \nu ,\phi ^\delta \left(R_{\mathrm{}}^1\right)_x^{red}\left(d\nu \right)}`$ $`=`$ $`E\left(\left({\displaystyle _0^{\mathrm{}}}\left({\displaystyle \underset{i=1}{\overset{N_t}{}}}Z_{t,i}\right)\pi (dt)\right)^\delta \right)`$
$`=`$ $`E(E\left[\left({\displaystyle _0^{\mathrm{}}}({\displaystyle \underset{i=1}{\overset{N_t}{}}}Z_{t,i}biggr)\pi (dt)\right)^\delta |\pi ,W^x]\right)`$
$``$ $`E\left(E\left[\left({\displaystyle _0^{\mathrm{}}}({\displaystyle \underset{i=1}{\overset{N_t}{}}}Z_{t,i}biggr)^\delta \pi (dt)\right)|\pi ,W^x\right]\right)`$
$`=`$ $`E({\displaystyle _0^{\mathrm{}}}E\left[\left({\displaystyle \underset{i=1}{\overset{N_t}{}}}Z_{t,i}\right)^\delta \right|W^x]\pi (dt)).`$
Now, by the self–similarity of the $`\alpha `$–stable process,
$`E\left[Z_{t,i}|W^x\right]`$ $`=`$ $`{\displaystyle \phi \left(y\right)p_t(W_t^x,y)\lambda \left(dy\right)}`$
$`=`$ $`{\displaystyle \phi \left(yW_t^x\right)p_t(0,y)\lambda \left(dy\right)}`$
$`=`$ $`{\displaystyle \frac{1}{t^{d/\alpha }}}{\displaystyle \phi \left(yW_t^x\right)p_1(0,yt^{1/\alpha })𝑑y}`$
$``$ $`K{\displaystyle \frac{1}{t^{d/\alpha }}}\lambda ,\phi ,`$
for some real constant $`K`$ not depending on $`W^x`$ and $`\phi `$. Hence, by Hölder’s inequality,
$`E\left[\left({\displaystyle \underset{i=1}{\overset{N_t}{}}}Z_{t,i}\right)^\delta |W^x\right]`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}E\left[\left({\displaystyle \underset{i=1}{\overset{n}{}}}Z_{t,i}\right)^\delta |W^x\right]P[N_t=n]`$
$`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}n^\delta E[\left({\displaystyle \frac{1}{n}}{\displaystyle \underset{i=1}{\overset{n}{}}}Z_{t,i}\right)^\delta |W^xbiggr]P[N_t=n]`$
$``$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}n^\delta \left(E\left[{\displaystyle \frac{1}{n}}{\displaystyle \underset{i=1}{\overset{n}{}}}Z_{t,i}|W^x\right]\right)^\delta P\left[N_t=n\right]`$
$``$ $`\left(K\lambda ,\phi \right)^\delta t^{d\delta /\alpha }{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}n^\delta P\left[N_t=n\right].`$
For $`(1+\beta _1)`$–branching and $`\delta <\beta _1`$,
$$\underset{n=0}{\overset{\mathrm{}}{}}n^\delta P\left[N_t=n\right]<\mathrm{}.$$
On the other hand, since $`d\delta /\alpha >1`$, we have
$$E\left(_1^{\mathrm{}}t^{d\delta /\alpha }\pi \left(dt\right)\right)<\mathrm{}.$$
Putting these results together finishes the proof. $`\mathrm{}`$
Corollary 5.4.1. In the setting of Lemma 5.4.1(b),
$$\nu ,\phi ^{\beta _2}(R_{\mathrm{}})_x(d\nu )<\mathrm{},x^d.$$
Proof: Let $`\psi `$ be strictly positive and $`\lambda `$-integrable. Then, by the Palm formula (A.2.1),
$`{\displaystyle \nu ,\phi ^{\beta _2}(R_{\mathrm{}})_x(d\nu )\psi (x)\lambda (dx)}`$ $`=`$ $`{\displaystyle \nu ,\phi ^{\beta _2}\nu ,\psi R_{\mathrm{}}(d\nu )}`$
$``$ $`{\displaystyle \nu ,\phi +\psi ^{1+\beta _2}R_{\mathrm{}}(d\nu )}<\mathrm{}.`$
This shows that the assertion of the corollary holds for $`\lambda `$–almost all $`x`$. The shift–invariance of the system implies that it is in fact true for all $`x^d`$. $`\mathrm{}`$
Lemma 5.4.2.
$$(R_{\mathrm{}}^1)_x=\delta _{\delta _x}_{_\tau (S)}\mathrm{\Pi }_\nu ()(R_{\mathrm{}})_x(d\nu )\text{for}\lambda \text{almost all }x.$$
Proof: Since $`R_{\mathrm{}}^1`$ has intensity measure $`\lambda `$, by (A.1.11), the assertion is obtained from the following chain of equalities, where we use the Palm formula (A.2.1) for $`\mathrm{\Pi }_\nu `$ and for $`R_{\mathrm{}}`$, the fact that $`(\mathrm{\Pi }_\nu )_x=\delta _{\delta _x}\mathrm{\Pi }_\nu `$, and (A.1.10),
$`{\displaystyle }`$ $`f(x)`$ $`F(\mu )(\delta _{\delta _x}\mathrm{\Pi }_\nu )(d\mu )(R_{\mathrm{}})_x(d\nu )\lambda (dx)`$
$`=`$ $`{\displaystyle f(x)F(\mu )(\delta _{\delta _x}\mathrm{\Pi }_\nu )(d\mu )\nu (dx)R_{\mathrm{}}(d\nu )}`$
$`=`$ $`{\displaystyle F(\mu )f(x)\mu (dx)\mathrm{\Pi }_\nu (d\mu )R_{\mathrm{}}(d\nu )}`$
$`=`$ $`{\displaystyle F(\mu )f(x)\mu (dx)R_{\mathrm{}}^1(d\mu )}.`$
$`\mathrm{}`$
Proof of Proposition 3.2.2: We have to prove that $`\beta _2<\beta _1<1(\mathrm{5.4.2})`$and $`d>\alpha \left(1+{\displaystyle \frac{1}{\beta _2}}\left(1+{\displaystyle \frac{1}{\beta _1}}\right)\right).(\mathrm{5.4.3})`$ are necessary and sufficient for condition (2.5.2) of Theorem 2.5.2.
1. Assume that (5.4.2) and (5.4.3) hold. Because of Lemma 3.1.2 and (4.1.2), we have to show that
$$\left(\frac{\phi (y)}{|xy|^{d\alpha }}𝑑y\mu (dx)\right)^{1+\beta _2}R_{\mathrm{}}^1(d\mu )<\mathrm{},\phi 𝒞_c^+(\text{}^d).$$
First we observe that (Iscoe<sup>(27)</sup>, Lemma 5.3) $`{\displaystyle \frac{\phi (y)}{|xy|^{d\alpha }}𝑑y}const.(1|x|)^{d+\alpha }.(\mathrm{5.4.4})`$Hence from Lemma 5.4.1(b) (denoting by $`B_r`$ the ball with radius $`r`$ centered at $`0`$) we obtain
$$\left(_{B_1}_^d\frac{\phi (y)}{|xy|^{d\alpha }}𝑑y\mu (dx)\right)^{1+\beta _2}R_{\mathrm{}}^1(d\mu )<\mathrm{},$$
and it suffices to show that
$$A:=\left(_^d1_{B_1^c}(x)|x|^{d+\alpha }\mu (dx)\right)^{1+\beta _2}R_{\mathrm{}}^1(d\mu )<\mathrm{}.$$
Using the Palm formula (A.2.1) and Lemma 5.4.2, we have
$`A`$ $`=`$ $`{\displaystyle 1_{B_1^c}(x)|x|^{d+\alpha }\left(1_{B_1^c}(z)|z|^{d+\alpha }\mu (dz)\right)^{\beta _2}(R_{\mathrm{}}^1)_x(d\mu )\lambda (dx)}`$
$`=`$ $`{\displaystyle 1_{B_1^c}(x)|x|^{d+\alpha }\left(1_{B_1^c}(x)|x|^{d+\alpha }+1_{B_1^c}(z)|z|^{d+\alpha }\mu (dz)\right)^{\beta _2}}`$
$`\mathrm{\Pi }_\nu (d\mu )(R_{\mathrm{}})_x(d\nu )\lambda (dx)`$
$``$ $`{\displaystyle 1_{B_1^c}(x)|x|^{d+\alpha }(|x|^{d+\alpha })^{\beta _2}\lambda (dx)}`$
$`+{\displaystyle 1_{B_1^c}(x)|x|^{d+\alpha }\left(1_{B_1^c}(z)|z|^{d+\alpha }\mu (dz)\right)^{\beta _2}\mathrm{\Pi }_\nu (d\mu )(R_{\mathrm{}})_x(d\nu )\lambda (dx)}.`$
Since $`\alpha +\beta _2(d+\alpha )<0`$ by (5.4.3), the first term on the r.h.s. is finite.
Using Hölder’s inequality, the second term can be bounded by
$$B:=_{B_1^c}|x|^{d+\alpha }\left(|z|^{d+\alpha }\nu (dz)\right)^{\beta _2}(R_{\mathrm{}})_x(d\nu )\lambda (dx),$$
and by shift–invariance,
$$B=_{B_1^c}|x|^{d+\alpha }\left(|zx|^{d+\alpha }\nu (dz)\right)^{\beta _2}(R_{\mathrm{}})_0(d\nu )\lambda (dx).$$
The scaling property of $`(R_{\mathrm{}})_0`$ (Dawson and Perkins<sup>(11)</sup>, Theorem 6.7) yields
$`B`$ $`=`$ $`{\displaystyle _{B_1^c}}|x|^{d+\alpha }|x|^{\alpha /\beta _1}{\displaystyle \left(|z|x|x|^{d+\alpha }\nu (dz)\right)^{\beta _2}(R_{\mathrm{}})_0(d\nu )\lambda (dx)}`$
$`=`$ $`{\displaystyle _{B_1^c}}|x|^{d+\alpha +\alpha /\beta _1+\beta _2(d+\alpha )}{\displaystyle \left(\left|z\frac{x}{|x|}\right|^{d+\alpha }\nu (dz)\right)^{\beta _2}(R_{\mathrm{}})_0(d\nu )\lambda (dx)}.`$
By isotropy of $`(R_{\mathrm{}})_0`$, the integral w.r. to $`(R_{\mathrm{}})_0`$ does not depend on $`x/|x|`$. For an arbitrary fixed $`x_0^d`$ we put $`e=x_0/|x_0|`$ and we obtain
$$B=_{B_1^c}|x|^{d+\alpha +\alpha /\beta _1+\beta _2(d+\alpha )}\lambda (dx)\left(|ze|^{d+\alpha }\nu (dz)\right)^{\beta _2}(R_{\mathrm{}})_0(d\nu ).$$
Since $`\alpha +\alpha /\beta _1+\beta _2(d+\alpha )<0`$ by (5.4.3), the integral w.r. to $`\lambda `$ is finite. Hence we will be done if we can show that the integral w.r. to $`(R_{\mathrm{}})_0`$ is finite as well.
We will show that
$$C_1:=\left(_{B_2}|ze|^{d+\alpha }\nu (dz)\right)^{\beta _2}(R_{\mathrm{}})_0(d\mu )$$
and
$$C_2:=\left(_{B_2^c}|ze|^{d+\alpha }\nu (dz)\right)^{\beta _2}(R_{\mathrm{}})_0(d\mu )$$
are finite.
Let
$$g(z)=1_{B_2}(z)|ze|^{d+\alpha }.$$
By Corollary 5.4.1 we have
$$C_1=\left(g(z)\nu (dz)\right)^{\beta _2}(R_{\mathrm{}})_0(d\mu )<\mathrm{},$$
since $`\lambda ,g<\mathrm{}`$. On the other hand, for $`zB_2^c`$, $`|ze|\frac{1}{2}|z|.`$ Therefore
$$C_2const.\left(_{B_2^c}|z|^{d+\alpha }\nu (dz)\right)^{\beta _2}(R_{\mathrm{}})_0(d\nu ).$$
To show the finiteness of the latter integral we will decompose the random variable $`_{B_2^c}|z|^{d+\alpha }\nu (dz)`$ into a sum of terms whose $`L^{\beta _2}`$–norms add up to something finite. Put $`D_k:=B_{2^{k+1}}\backslash B_{2^k}`$, $`k=0,1,\mathrm{}`$ Then, again by the scaling property of $`(R_{\mathrm{}})_0`$ we have
$`{\displaystyle \left(_{D_k}|z|^{d+\alpha }\nu (dz)\right)^{\beta _2}(R_{\mathrm{}})_0(d\nu )}`$
$`=(2^k)^{\alpha /\beta _1}{\displaystyle \left(_{D_0}(2^k|z|)^{d+\alpha }\nu (dz)\right)^{\beta _2}(R_{\mathrm{}})_0(d\nu )}`$
$`2^{k\left(\alpha /\beta _1+\beta _2(d+\alpha )\right)}{\displaystyle (\nu (D_0))^{\beta _2}(R_{\mathrm{}})_0(d\nu )}.`$
The $`L^{\beta _2}`$–norm of the $`k`$–th summand is thus bounded by $`const.2^{k\left(\alpha /\beta _1\beta _2d+\alpha \right)}`$, which is summable since $`\alpha /\beta _1\beta _2d+\alpha <0`$ due to (5.4.3).
2. To prove the converse we assume that (2.5.2) holds.
(a) Assume that $`\beta _2\beta _1`$, $`\beta _1<1`$. We will show that, for $`\phi C_c(S)`$, $`\phi 0`$, $`\phi 0`$,
$$\mu ,G\phi ^{1+\beta _2}R_{\mathrm{}}^1(d\mu )=\mathrm{}.$$
First note that by the Palm formula (A.2.1), $`{\displaystyle \mu ,G\phi ^{1+\beta _2}R_{\mathrm{}}^1(d\mu )}={\displaystyle G\phi (x)\mu ,G\phi ^{\beta _2}(R_{\mathrm{}}^1)_x(d\mu )\lambda (dx)}.(\mathrm{5.4.5})`$Choose $`\psi :^d_+`$ such that $`G\phi \psi (x)`$ provided that $`|x|1`$. Then the r.h.s. of (5.4.5) is bounded above by
$$1_{\{|x|1\}}G\phi (x)\mu ,\psi ^{\beta _2}(R_{\mathrm{}}^1)_0(d\mu )\lambda (dx).$$
By the tree representation (A.3.2) of $`(R_{\mathrm{}}^1)_0`$ we have $`{\displaystyle \mu ,\psi ^{\beta _2}(R_{\mathrm{}}^1)_0(d\mu )}E\left(\left({\displaystyle \underset{i=1}{\overset{N}{}}}X_\sigma ^{W_\sigma ^0,i},\psi \right)^{\beta _2}\right),(\mathrm{5.4.6})`$where $`\sigma `$ is exponentially distributed with parameter $`V_1`$, and $`N`$ is distributed like any of the $`N_t`$. Since $`EN^{\beta _2}=\mathrm{}`$, and since, conditioned on $`\sigma `$ and $`W^0`$, the random variables $`X_\sigma ^{W_\sigma ^0,i},\psi `$ are i.i.d. with positive expectation, it follows from the law of large numbers that the r.h.s. of (5.4.6) is infinite.
(b) Now assume that $`1\beta _1>\beta _2`$, and suppose that $`\phi (x)1`$ for $`|x|1`$. Then, for some $`k>0`$,
$$G\phi (x)k|x|^{(d\alpha )}\text{ if }|x|2.$$
Assume that $`\mu ,G\phi ^{1+\beta _2}R_{\mathrm{}}^1(d\mu )<\mathrm{}`$. Then, by Lemma 5.4.1(a),
$$\nu ,G\phi ^{1+\beta _2}R_{\mathrm{}}^1(d\nu )<\mathrm{}.$$
Therefore, by the Palm formula (A.2.1) we have
$`\mathrm{}`$ $`>`$ $`{\displaystyle \nu ,G\phi ^{1+\beta _2}R_{\mathrm{}}^1(d\nu )}`$
$``$ $`k^{^{1+\beta _2}}{\displaystyle 1_{\{|x|1\}}|x|^{d+\alpha }\left(1_{\{|z|1\}}|z|^{d+\alpha }\nu (dz)\right)^{\beta _2}(R_{\mathrm{}})_x(d\nu )\lambda (dx)}.`$
By shift–invariance this equals
$$k^{^{1+\beta _2}}1_{\{|x|1\}}|x|^{d+\alpha }\left(1_{\{|zx|1\}}|zx|^{d+\alpha }\nu (dz)\right)^{\beta _2}(R_{\mathrm{}})_0(d\nu )\lambda (dx).$$
The scaling property of $`(R_{\mathrm{}})_0`$ permits to rewrite the inner integral (with $`e(x)=x/|x|`$) as
$$|x|^{\alpha /\beta _1}\left(1_{\{|z|x|x|\mathrm{\hspace{0.17em}1}\}}|z|x|x|^{d+\alpha }\nu (dz)\right)^{\beta _2}(R_{\mathrm{}})_0(d\nu )$$
$$=|x|^{\alpha /\beta _1+\beta _2(d+\alpha )}\left(1_{\{|x||ze(x)|1\}}|ze(x)|^{d+\alpha }\nu (dz)\right)^{\beta _2}(R_{\mathrm{}})_0(d\nu ).$$
For $`|x|1`$, the latter integral is bounded below by
$$\left(1_{\{|z|1\}}|2z|^{d+\alpha }\nu (dz)\right)^{\beta _2}(R_{\mathrm{}})_0(d\nu )>0.$$
Now,
$$|x|^{d+\alpha +\alpha /\beta _1+\beta _2(d+\alpha )}1_{\{|x|1\}}\lambda (dx)<\mathrm{}$$
implies that $`\alpha +\alpha /\beta _1+\beta _2(d+\alpha )<0`$, or equivalently, (5.4.3) holds. $`\mathrm{}`$
APPENDIX A.1. Background on 1- and 2-level branching systems We consider particle systems in a locally compact Abelian group $`S`$ with Haar measure $`\rho `$. Recall that $`T_t`$ denotes the semigroup of the particle motion and $`G`$ the corresponding Green operator.
Let $`𝒞(S)`$ denote the space of bounded continuous functions on $`S,𝒞_0(S)`$ the subspace of functions vanishing at infinity, and $`𝒞_c(S)`$ that of functions with compact support. For a strictly positive function $`\tau 𝒞_0(S)`$, let
$$𝒞_\tau (S)=\{\phi 𝒞(S):\phi \tau ^1𝒞_0(S)\}$$
with the norm $`\phi _\tau =\phi \tau ^1`$. We assume that $`\tau `$ is such that $`tT_t\phi `$ is a continuous curve in $`(𝒞_\tau (S),||||_\tau )`$ for each $`\phi 𝒞_\tau (S)`$. For example, in the case of the $`\alpha `$-stable motion in $`\text{}^d`$ we may take $`\tau (x)=(1+|x|^2)^q`$ with $`d/2<q<(d+\alpha )/2`$ (Dawson and Gorostiza<sup>(6)</sup>). The subspaces of non–negative functions are indicated with the superscript ‘+’, e.g. $`𝒞_\tau ^+(S)`$. Let $`_\tau (S)`$ denote the space of non–negative Radon measures $`\mu `$ on $`S`$ such that $`\mu ,\tau <\mathrm{}`$, endowed with the smallest topology which makes the maps $`\mu \mu ,\phi `$ continuous for all $`\phi 𝒞_c^+(S)\{\tau \}`$. We assume that $`\rho _\tau (S)`$. The subspace of $`_\tau (S)`$ of integer–valued measures is designated by $`𝒩_\tau (S)`$.
The Laplace functional of the occupation time of the 1–level branching particle system $`X_t`$ with $`(1+\beta )`$–branching at rate $`V`$ and started off from a Poisson system with intensity $`\rho `$ is given by $`E\text{exp}\left\{{\displaystyle _0^t}X_s𝑑s,\phi \right\}=\text{exp}\{\rho ,u_\phi (t)\},\phi 𝒞_\tau ^+(S),(A\mathrm{.1.1})`$ where $`u_\phi (x,t)`$ with values in $`[0,1]`$ is the unique solution of the non–linear evolution equation $`u_\phi (t)={\displaystyle \frac{V}{1+\beta }}{\displaystyle _0^t}T_{ts}(u_\phi (s)^{1+\beta })𝑑s+{\displaystyle _0^t}T_{ts}(\phi (1u_\phi (s)))𝑑s.(A\mathrm{.1.2})`$ This is shown by the same argument of Theorem 5 in Gorostiza and López–Mimbela<sup>(18)</sup> (formulas (4.8) and (4.9)). It follows that $`u_\phi (t){\displaystyle _0^t}T_{ts}(\phi (1u_\phi (s)))𝑑sG_t\phi .(A\mathrm{.1.3})`$ Let $`U_t`$ denote the semigroup of the 1–level branching particle system. We have $`U_t(,\phi )(\mu )=\mu ,T_t\phi ,\phi 𝒞_\tau ^+(S),\mu 𝒩_\tau (S),(A\mathrm{.1.4})`$ and, if $`\beta =1`$,
$`U_t(,\phi ,\psi )(\mu )`$ $`=`$ $`\mu ,T_t\phi \mu ,T_t\psi +\mu ,T_t(\phi \psi )T_t\phi T_t\psi `$
$`+`$ $`V{\displaystyle _0^t}\mu ,T_s(T_{ts}\phi T_{ts}\psi )𝑑s,\phi ,\psi 𝒞_\tau ^+(S),\mu 𝒩_\tau (A\mathrm{.1.5})`$
The formulas (A.1.4) and (A.1.5) can be derived by martingale methods from the Markov property of the system (see e.g. Gorostiza and Rodrigues<sup>(20)</sup> for explicit calculations of this type). In particular, $`U_t(,\phi )(\delta _x)=T_t\phi (x),(A\mathrm{.1.6})`$
$$U_t(,\phi ,\psi )(\delta _x)=T_t(\phi \psi )(x)+V_0^tT_s(T_{ts}\phi T_{ts}\psi )(x)ds.(A\mathrm{.1.7})$$
We have from (A.1.4)
$$U_t(,\phi ,\psi )(\mu )\mu ,T_t\phi \mu ,T_t\psi +\mu ,T_t(\phi \psi )+V_0^t\mu ,T_s(T_{ts}\phi T_{ts}\psi )𝑑s,$$
$`\phi ,\psi 𝒞_\tau ^+(S).(A\mathrm{.1.8})`$ If the 1–level branching system is persistent, it has a “Poisson type” equilibrium (in the sense of Liemant et al<sup>(31)</sup>, section 2.3), which is an infinitely divisible random element of $`_\tau (S)`$. Its canonical measure, which is a measure on $`_\tau (S)`$, is denoted by $`R_{\mathrm{}}^1`$. A sufficient condition for persistence is
$$_0^{\mathrm{}}\rho ,(T_t\phi )^{1+\beta }𝑑t<\mathrm{},\phi 𝒞_c^+(S)$$
(Gorostiza and Wakolbinger<sup>(22)</sup>, Theorem 2.1).
For each $`t>0`$, the random measure $`X_t`$ is infinitely divisible and its canonical measure $`R_t`$ has the form $`R_t^1={\displaystyle _S}P[X_t^x()]\rho (dx),(A\mathrm{.1.9})`$ where $`X_t^x`$ corresponds to the branching system starting with a single ancestor in $`x`$ at time $`0`$ (Gorostiza and Wakolbinger<sup>(21)</sup>, formula (3.1), Liemant et al<sup>(31)</sup>).
The measure $`R_{\mathrm{}}^1`$ is the “Poissonization” of the equilibrium canonical measure $`R_{\mathrm{}}`$ of the superprocess counterpart of the particle system, i.e., $`R_{\mathrm{}}^1={\displaystyle _{_\tau (S)}}\mathrm{\Pi }_\nu ()R_{\mathrm{}}(d\nu ),(A\mathrm{.1.10})`$ where $`\mathrm{\Pi }_\nu `$ is the distribution of a Poisson random measure on $`S`$ with intensity measure $`\nu `$. Indeed, in Gorostiza et al<sup>(19)</sup> it is shown that the distibution $`L_t^1`$ of the branching particle system $`X_t`$ is a Cox process, i.e.,
$$L_t^1=_{_\tau (S)}\mathrm{\Pi }_\nu ()L_t(d\nu ),$$
where $`L_t`$ is the distribution of the superprocess counterpart of $`X_t`$. By continuity and the assumed persistence, this relation carries over to $`t=\mathrm{}`$:
$$L_{\mathrm{}}^1=_{_\tau (S)}\mathrm{\Pi }_\nu ()L_{\mathrm{}}(d\nu ).$$
Together with the Lévy-Khinchin formula (Kallenberg<sup>(28)</sup>), this implies the following chain of equalities for each $`\phi C_c(S)`$:
$`\text{exp}\left\{{\displaystyle R_{\mathrm{}}^1(d\mu )(1e^{\mu ,\phi })}\right\}`$ $`=`$ $`{\displaystyle e^{\mu ,\phi }L_{\mathrm{}}^1\left(d\mu \right)}`$
$`=`$ $`{\displaystyle e^{\mu ,\phi }\mathrm{\Pi }_\nu \left(d\mu \right)L_{\mathrm{}}\left(d\nu \right)}`$
$`=`$ $`{\displaystyle e^{\nu ,1e^\phi }L_{\mathrm{}}\left(d\nu \right)}`$
$`=`$ $`\text{exp}\left\{{\displaystyle R_{\mathrm{}}\left(d\nu \right)(1e^{\nu ,1e^\phi })}\right\}`$
$`=`$ $`\text{exp}\left\{{\displaystyle R_{\mathrm{}}\left(d\nu \right)\left(1e^{\mu ,\phi }\mathrm{\Pi }_\nu (d\mu )\right)}\right\}`$
$`=`$ $`\text{exp}\left\{{\displaystyle R_{\mathrm{}}\left(d\nu \right)\mathrm{\Pi }_\nu \left(d\mu \right)(1e^{\mu ,\phi })}\right\},`$
which yields (A.1.10).
Let $`_\tau ^2(S)`$ denote the space of Radon measures $`\underset{¯}{\mu }`$ on $`_\tau (S)`$ such that $`\underset{¯}{\mu },,\tau <\mathrm{}`$, where
$$\underset{¯}{\mu },F=_{_\tau (S)}F(\nu )\underset{¯}{\mu }(d\nu )$$
(sometimes we use the notation on the r.h.s in order to avoid confusion). We have that $`R_{\mathrm{}}^1_\tau ^2(S)`$, $`R_{\mathrm{}}^1`$ is invariant (but not reversible) for the 1–level dynamics (Liemant et al<sup>(31)</sup>, Chapter 2, Dawson and Perkins<sup>(11)</sup>), and it has intensity $`\rho `$ in the sense that $`R_{\mathrm{}}^1,,\phi =\rho ,\phi ,\phi 𝒞_\tau ^+(S).(A\mathrm{.1.11})`$ If $`\beta =1`$, then $`R_{\mathrm{}}^1`$ has finite moments of all orders and the second and third moments are given by
$`R_{\mathrm{}}^1,,\phi ,\psi `$ $`=`$ $`\rho ,\phi \psi +{\displaystyle \frac{V}{2}}\rho ,\phi G\psi ,\phi ,\psi 𝒞_c^+(S),`$ (A.1.13)
$`R_{\mathrm{}}^1,,\phi ,\psi ,\zeta =\rho ,\phi \psi \zeta +{\displaystyle \frac{V}{2}}\rho ,\phi \psi G\zeta +\phi \zeta G\psi +\psi \zeta G\phi `$
$`+{\displaystyle \frac{V^2}{2}}\rho ,{\displaystyle _0^{\mathrm{}}}[\phi GT_t(T_t\psi T_t\zeta )+\psi GT_t(T_t\phi T_t\zeta )+\zeta GT_t(T_t\phi T_t\psi )]𝑑t,\phi ,\psi ,\zeta 𝒞_c^+(S).`$
See Subsection A.4 for a proof.
Note also, from (A.1.9) and $`T_t`$–invariance of $`\rho `$, that for each $`t>0`$, $`R_t^1,,\phi =\rho ,\phi ,quad\phi 𝒞_\tau ^+(S).(A\mathrm{.1.14})`$
We pass now to the 2–level branching system. A “2–level particle” is an element $`\mu `$ of $`𝒩_\tau (S)`$ of the form $`\mu =_{i=1}^n\delta _{x_i}`$. A “clan” is the progeny under the 1–level dynamics (i.e., individual particles undergoing $`(1+\beta )`$–branching at rate $`V_1`$) of a family of particles which constitute an initial 2–level particle. Clans undergo $`(1+\beta _2)`$–branching at rate $`V_2`$. Assuming persistence of the 1–level system, the 2–level system starts off from a Poisson system of “2–level particles” with intensity measure $`R_{\mathrm{}}^1`$. The empirical measures of the 2–level system take values in $`_\tau ^2(S)`$. Restricting to test functions on $`_\tau (S)`$ of the form $`\mu \mu ,\phi `$, $`\phi 𝒞_c^+(S)`$, amounts to considering the aggregated system, i.e., we consider the empirical measure of all the point masses disregarding which 2–level particles they belong to. Note that the moments of $`R_{\mathrm{}}^1`$ in (A.1.11), (A.1.12), (A.1.13) correspond to the aggregated Poisson system. The empirical measures $`X_t`$ of the aggregation of the 2–level system take values in $`𝒩_\tau (S)`$.
The same argument used to obtain the Laplace functional of the 1–level system can be used for the 2–level system. Hence, analogously to (A.1.1), the Laplace functional of the occupation time of the 2–level system is given by $`E\text{exp}\left\{{\displaystyle _0^t}X_s𝑑s,\phi \right\}=\text{exp}\{R_{\mathrm{}}^1,𝐮_\phi (t)\},\phi 𝒞_\tau ^+(S),(A\mathrm{.1.15})`$ where $`𝐮_\phi `$ with values in $`[0,1]`$ is the unique solution of the non-linear evolution equation $`𝐮_\phi (t)={\displaystyle \frac{V_2}{1+\beta _2}}{\displaystyle _0^t}U_{ts}(𝐮_\phi (s)^{1+\beta _2})𝑑s+{\displaystyle _0^t}U_{ts}(,\phi (1𝐮_\phi (s)))𝑑s.(A\mathrm{.1.16})`$ If follows from (A.1.4) and (A.1.16) that $`𝐮_\phi (t)(\mu ){\displaystyle _0^t}U_s(,\phi )(\mu )𝑑s={\displaystyle _0^t}\mu ,T_s\phi 𝑑s=\mu ,G_t\phi ,\mu _\tau (S).(A\mathrm{.1.17})`$
A.2. The Palm formula Let $`M`$ be a measure on $`_\tau (S)`$ whose intensity measure
$$\mathrm{\Lambda }_M:=_{_\tau (S)}\nu ()M(d\nu )$$
is locally finite, i.e. $`\mathrm{\Lambda }_H,\phi <\mathrm{}`$ for all $`\phi 𝒞_c(S)`$. The Palm measures of $`M`$ are a family $`(M_x)_{xS}`$ of probability measures on $`_\tau (S)`$ which satisfy $`{\displaystyle _{_\tau (S)}}\nu ,\phi F(\nu )M(d\nu )={\displaystyle _S}\phi (x){\displaystyle _{_\tau (S)}}F(\nu )M_x(d\nu )\mathrm{\Lambda }_M(dx)(A\mathrm{.2.1})`$for all measurable $`F:_\tau (S)_+`$ and $`\phi :S_+`$. If $`M`$ is supported by $`\{\mu _\tau (S)|\mu \text{ is }\{0,1,2,\mathrm{},\mathrm{}\}`$–valued$`\}`$, then
$$M_x(\{\mu |\mu (x)1\})=1$$
for $`\mathrm{\Lambda }_M`$–almost all $`x`$, and in this case the reduced Palm measures $`M_x^{red}`$, $`xS`$, are defined by $`M_x^{red}=M_x(\{\mu \delta _x()\}).(A\mathrm{.2.2})`$
A.3. Tree representations of the Palm measures of $`R_t^1`$ and $`R_{\mathrm{}}^1`$ The Palm measures of $`R_t^1`$ and of the equilibrium canonical measure $`R_{\mathrm{}}^1`$ described in Subsection A.1 have a representation in terms of a backward tree which we recall here.
Lemma A.3.1. (Gorostiza and Wakolbinger<sup>(21)</sup>, Theorem 2.3). Let $`W^x`$ be a random path of the motion process starting in $`xS`$, let $`\pi `$ be a random Poisson configuration on $`_+`$ with intensity $`V`$, and for each $`yS`$, $`r>0`$ and $`i=1,2,\mathrm{}`$, let $`X_r^{y,i}`$ be a branching particle system arising from one ancestor at site $`y`$ and developing over time $`r`$. Let $`N_r`$ be an integer–valued random variable with generating function $`1(1+\beta )q(1s)^\beta `$ (see Subsection 2.5). Assume all these objects are independent. Then the particle systems $`\mathrm{\Phi }_x^t:={\displaystyle _0^t}{\displaystyle \underset{i=1}{\overset{Nr}{}}}X_r^{W_r^x,i}()\pi (dr)(A\mathrm{.3.1})`$and $`\mathrm{\Phi }_x^{\mathrm{}}:={\displaystyle _0^{\mathrm{}}}{\displaystyle \underset{i=1}{\overset{Nr}{}}}X_r^{W_r^x,i}()\pi (dr)(A\mathrm{.3.2})`$have distributions $`(R_t^1)_x^{\text{red}}`$ and $`(R_{\mathrm{}}^1)_x^{\text{red}}`$, respectively, for $`\rho `$–almost all $`xS`$.
It follows immediately from (A.3.1), (A.3.2), (A.1.11), (A.1.14) and the Palm formula (A.2.1) that all the moments of $`R_t^1`$ increase to those of $`R_{\mathrm{}}^1`$ as $`t\mathrm{}`$.
A.4. Second and third moments of $`R_{\mathrm{}}^1`$ Proof of (A.1.12) and (A.1.13): The proof of can be carried out directly by using the explicit form of the Laplace transform of $`R_{\mathrm{}}^1`$ (as given, e.g., in Gorostiza and Wakolbinger<sup>(21)</sup>, Theorem 3.3). Here we give an argument which uses the structure of $`R_t^1`$ in (A.1.9) and the monotone convergence of the moments of $`R_t^1`$ mentioned above.
Let us introduce the notation $`A_{t,x}(\phi ,\psi )={\displaystyle _0^t}T_s(T_{ts}\phi T_{ts}\psi )(x)𝑑s,(A\mathrm{.4.1})`$$`B_{t,x}(\phi ,\psi ,\zeta )={\displaystyle _0^t}T_s\left(T_{ts}\phi \left({\displaystyle _0^{ts}}T_{tsr}\left(T_r\psi T_r\zeta \right)𝑑r\right)\right)(x)𝑑s.(A\mathrm{.4.2})`$
The following formulae, which can be obtained either by differentiating the Laplace functional of $`X_t^x`$ or by means of a tree representation of the Palm measures of the distribution of $`X_t^x`$ (similar to (A.3.1)), are well known (e.g. Klenke<sup>(30)</sup>, Lemma 3.1),
$`E[X_t^x,\phi X_t^x,\psi ]=T_t(\phi \psi )(x)+VA_{t,x}(\phi ,\psi ),`$ (A.4.3)
$`E[X_t^x,\phi X_t^x,\psi X_t^x,\zeta ]`$
$`=T_t(\phi \psi \zeta )(x)+V(A_{t,x}(\phi ,\psi \zeta )+A_{t,x}(\psi ,\phi \zeta )+A_{t,x}(\zeta ,\phi \psi ))`$
$`+V^2(B_{t,x}(\phi ,\psi ,\zeta )+B_{t,x}(\psi ,\phi ,\zeta )+B_{t,x}(\zeta ,\phi ,\psi )).`$ (A.4.4)
Note that (A.4.3) is just a rewriting of (A.1.7). We need the following lemma.
Lemma A.4.1.
(a) $`{\displaystyle _S}A_{t,x}(\phi ,\psi )\rho (dx){\displaystyle \frac{1}{2}}\rho ,\phi G\psi `$ as $`t\mathrm{}`$. (b) $`{\displaystyle _S}B_{t,x}(\phi ,\psi ,\zeta )\rho (dx){\displaystyle \frac{1}{2}}\rho ,\phi {\displaystyle _0^{\mathrm{}}}GT_r(T_r\psi T_r\zeta )𝑑r`$ as $`t\mathrm{}`$.
Proof: (a) $`{\displaystyle _S}A_{t,x}(\phi ,\psi )\rho (dx)=\rho ,{\displaystyle _0^t}(T_s\phi T_s\psi )𝑑s=\rho ,\phi {\displaystyle _0^t}T_{2s}(\psi )𝑑s{\displaystyle \frac{1}{2}}\rho ,\phi G\psi `$ as $`t\mathrm{}.`$
(b)
$`{\displaystyle _S}B_{t,x}(\phi ,\psi ,\zeta )\rho (dx)`$ $`=`$ $`\rho ,{\displaystyle _0^t}T_s\phi \left({\displaystyle _0^s}T_{sr}(T_r\psi T_r\zeta )𝑑r\right)𝑑s`$
$`=`$ $`\rho ,\phi {\displaystyle _0^t}{\displaystyle _0^s}T_{2sr}(T_r\psi T_r\zeta )𝑑r𝑑s`$
$`=`$ $`\rho ,\phi {\displaystyle _0^t}{\displaystyle _r^t}T_{2sr}(T_r\psi T_r\zeta )𝑑s𝑑r`$
$`=`$ $`{\displaystyle \frac{1}{2}}\rho ,\phi {\displaystyle _0^t}{\displaystyle _r^{2tr}}T_u(T_r\psi T_r\zeta )𝑑u𝑑r`$
$`=`$ $`{\displaystyle \frac{1}{2}}\rho ,\phi {\displaystyle _0^t}{\displaystyle _0^{2(tr)}}T_vT_r(T_r\psi T_r\zeta )𝑑v𝑑r`$
$``$ $`{\displaystyle \frac{1}{2}}\rho ,\phi {\displaystyle _0^{\mathrm{}}}GT_r(T_r\psi T_r\zeta )𝑑r\text{as}t\mathrm{}.`$
$`\mathrm{}`$ Combining Lemma A.4.1 with (A.1.9), (A.4.3), (A.4.4), and using the above mentioned monotonicity of the moments of $`R_t^1`$, the proof of (A.1.12) and (A.1.13) is complete. $`\mathrm{}`$
ACKNOWLEDGEMENTS
This research was carried out mainly during visits to The Fields Institute (Toronto, Canada), the Johann Wolfgang Goethe University (Frankfurt, Germany), and the Center for Mathematical Research (CIMAT, Guanajuato, Mexico). We express our gratitude for the hospitality of these institutions. We thank Dr. Mike Porter, who helped us by his computer analysis to unravel the oscillations of the hierarchical random walk. =====
REFERENCES
1. Barlow, M. T. and Perkins, E. A. (1988). Brownian motion on the Sierpiński gasket, Probab. Th. Rel. Fields 79 543–623.
2. Cartwright, D. I. (1988). Random walks on direct sums of discrete groups, J. Theor. Probab. 1 341–356.
3. Collet, P. and Eckmann, J-P. (1978). A renormalization group analysis of the hierarchical model in statistical mechanics, Lecture Notes in Physics 74, Springer-Verlag, Berlin-New York.
4. Cox, J. T. and Griffeath, D. (1984). Large deviations for Poisson systems of independent random walks, Probab. Th. Rel. Fields 66 543-558.
5. Cox, J. T. and Griffeath, D. (1985). Occupation times for critical branching Brownian motions, Ann. Probab. 13 1108-1132.
6. Dawson D. A. and Gorostiza, L. G. (1990). Generalized solutions of a class of nuclear–space–valued stochastic evolution equations, Appl. Math. Optim. 22 241–263.
7. Dawson, D. A. and Greven, A. Multiple space–time scale analysis for interacting branching models, Electron. J. Probab. 1, no. 14, 84 pp.
8. Dawson, D. A. and Ivanoff, G. (1978). Branching diffusions and random measures, in Branching Processes, Joffe, A. Ney P. (editors), M. Dekker, New York, 61–103.
9. Dawson,D. A. and Hochberg, K. J. (1991). A multilevel branching model, Adv. Appl. Prob. 23 701–715.
10. Dawson,D. A., Hochberg, K. J. and Vinogradov, V. (1996). High–density limits of hierarchically structured branching–diffusing populations, Stoch. Proc. Appl. 62 191–222.
11. Dawson, D. A. and Perkins, E. (1991). Historical Processes. Memoirs of the AMS 454, Providence, R. I.
12. Dawson, D. A. and Perkins, E. (1999). Measure–valued processes and renormalization of branching particle systems, in Stochastic Partial Differential Equations: Six Perspectives, Carmona, R. A. and Rozovskii, B. (editors), Math. Surveys and Monographs 64 45–106, AMS.
13. Deuschel, J. D. and Rosen, J. (1998). Occupation time large deviations for critical branching Brownian motion, super–Brownian motion and related processes, Ann. Probab. 26 602–643.
14. Deuschel, J. D. and Wang, K. (1994). Large deviations for the occupation time functional of a Poisson system of independent Brownian particles, Stoch. Proc. Appl. 52 183-209.
15. Fleischmann, K. and Greven, A. (1994). Diffusive clustering in an infinite system of hierarchically interacting diffusions, Probab. Th. Rel. Fields 98 517–566.
16. Gorostiza, L.G. (1996). Asymptotic fluctuations and critical dimension for a two-level branching system, Bernoulli 2 109-132.
17. Gorostiza, L. G. Hochberg, K. and Wakolbinger, A. (1995). Persistence of a critical super-2 process. J. Appl. Prob. 32 534-540.
18. Gorostiza, L. G. and López–Mimbela, J. A. (1994). An occupation time approach for convergence of measure–valued processes, and the death process of a branching system. Stat. Prob. Lett. 21 59–67.
19. Gorostiza, L. G., Roelly–Coppoletta, S. and Wakolbinger, A. (1990). Sur la persistence du processus de Dawson–Watanabe stable, in Séminaire de Probabilités XXIV, Azéma, J, Meyer, P. A. and Yor, M. (editors). Lect. Notes Math. 1426, Springer–Verlag, Berlin, 275–281.
20. Gorostiza, L. G. and Rodrigues, E.R. (1999). A stochastic model for transport of particulate matter in air: an asymptotic analysis, Acta Applicandae Mathematicae (in press).
21. Gorostiza, L. G. and Wakolbinger, A. (1991). Persistence criteria for a class of critical branching particle systems in continuous time, Ann. Probab. 19 266–288.
22. Gorostiza, L. G. and Wakolbinger, A. (1994). Long time behavior of critical particle systems and applications, in Measure–Valued Processes, Stochastic Partial Differential Equations, and Interacting Systems, Dawson, D. A. (editor), CRM Proc. and Lecture Notes 5, AMS, 119–137.
23. Grabner, P. J. and Woess, W. (1997). Functional iterations and periodic oscillations for simple random walk on the Sierpiński gasket, Stoch. Proc. Appl. 69 127–138.
24. Greven, A. and Hochberg, K. J. (1999). New behavioral patterns for two–level branching systems, in Stochastic Models, Gorostiza, L. G. and Ivanoff, B. G. (editors), Conference Procedings Series, CMS (in press).
25. Hochberg, K. J. (1995). Hierarchically structured branching populations with spatial motion. Rocky Mountain J. Math. 25 269–283.
26. Hochberg, K. J. and Wakolbinger, A. (1995). Non-persistence of two-level branching particle systems in low dimensions, in Stochastic Partial Differential Equations, Etheridge, A. (editor), London Mathem. Soc. Lecture Note Series 216, Cambridge Univ. Press, 126–140.
27. Iscoe, I. (1986). A weighted occupation time for a class of measure-valued branching processes, Probab. Th. Rel. Fields 71 85 – 116.
28. Kallenberg, O. (1983). Random Measures, 3rd. Edition, Akademie–Verlag, Berlin, Academic Press, New York.
29. Kesten, H. and Spitzer, F. (1965). Random walks on countably infinite Abelian groups, Acta Math. 114 237–265.
30. Klenke, A. (1997). Multiple scale analysis of clusters in spacial branching models. Ann. Probab. 25, 1670–1711.
31. Liemant, A., Matthes, K. and Wakolbinger, A. (1998). Equilibrium Distributions of Branching Processes, Akademie-Verlag, Berlin and Kluwer, Dordrecht.
32. Méléard, S. and Roelly, S. (1992). An ergodic result for critical spatial branching systems, Stochastic Analysis and Related Topics, Progress in Probability 31, Birkhauser, Boston, 333–341.
33. Port, S. C. and Stone, C. J. (1971). Infinitely divisible processes and their potential theory (First Part), Ann. Inst. Fourier 21 (2) 157–275.
34. Sato, K. (1996). Criteria of weak and strong transience for Lévy processes, in Probability Theory and Mathematical Statistics, Proceedings of the Seventh Japan-Russia Symposium, World Scientific, Singapore, 438–449.
35. Sawyer, S. and Felsenstein, J. (1983). Isolation by distance in a hierarchically clustered population, J. Appl. Prob. 20 1–10.
36. Sinai, Ya. G. (1982). Theory of Phase Transitions: Rigorous results. Pergamon Press.
37. Spitzer, F. (1964). Principles of Random Walk. Van Nostrand, Princeton.
38. Stöckl, A. and Wakolbinger, A. (1994). On clan-recurrence and -transience in time stationary branching Brownian particle systems, in Measure–Valued Processes, Stochastic Partial Differential Equations, and Interacting Systems, Dawson, D. A. (editor), CRM Proc. and Lecture Notes 5, AMS, 213–219.
39. Wilson, K. (1976). The renormalization group and block spins, in Proceedings of the International Conference on Statistical Physics, Pál, L. and Szépfalusy, P. (editors), North Holland, Amsterdam.
40. Wu, Y. (1994). Asymptotic behavior of two level measure branching processes, Ann. Probab. 22 854–874.
|
warning/0005/hep-th0005272.html
|
ar5iv
|
text
|
# A Consistent Noncommutative Field Theory: the Wess-Zumino Model
## I INTRODUCTION
Noncommutative geometry has been receiving a great deal of attention in the context of string/M-theory. Initially it appeared as a possible compactification manifold of space-time and led to the appearance of noncommutative quantum field theories on noncommutative tori . More recently it was shown that the dynamics of a D-brane in the presence of a $`B`$-field can, in certain limits, be described by a deformed gauge field theory in terms of Moyal products on space-time. Since this field theory arose from a coherent truncation of a string theory it is expected that deformed field theories are consistent by themselves. This motivated an intensive investigation of noncommutative quantum field theories on four dimensional Euclidean and Minkowski spaces. Scalar fields , gauge fields and supersymmetric theories have been studied. Some two dimensional models have also been analyzed .
A distinct characteristic of a class of noncommutative quantum field theories is the mixing of ultraviolet (UV) and infrared (IR) divergences reminiscent of the UV/IR connection of string theory. For the $`\varphi _4^4`$ massive scalar field there is an infrared quadratic singularity in the propagator at the one loop level, which jeopardizes the perturbative formulation of the theory. On the other hand, the theory has been proved to remain ultraviolet renormalizable up two loops although this does not seem to hold at all orders . Also, models involving complex scalar fields are not always renormalizable not even at one-loop approximation . Therefore, it is relevant to understand the renormalizability properties of noncommutative field theories to find out whether they are consistent.
It has been suggested that, due to the absence of quadratic divergences in their commutative version, noncommutative supersymmetric theories may remain ultraviolet renormalizable . The superspace formulation has already been accomplished at the classical level . However, at the quantum level only one loop results have been reported for supersymmetric gauge theories. As in the commutative case only logarithmic divergences show up .
This paper is dedicated to show that the noncommutative Wess-Zumino model in four dimensions is a consistent quantum field theory in the sense of being ultraviolet renormalizable and free of the IR/UV mixing at any arbitrary order of perturbation. This happens even though the scalar potential of the noncommutative Wess-Zumino model belongs to the class of non-renormalizable theories discussed in . It is a potential typical of a F-term $`\varphi ^{}\varphi ^{}\varphi \varphi `$ (while a D-term induces $`\varphi ^{}\varphi \varphi ^{}\varphi `$) but nevertheless supersymmetry still eliminates all quadratic divergences. This is at the root of the renormalizability of the model.
Noncommutative field theories containing just scalar and fermion fields, as is the case in the Wess-Zumino model, are constructed from the usual Lagrangian by replacing the ordinary product by the Moyal product of fields, i. e. $`ABAB`$. The Moyal product is noncommutative and obeys the rule
$`{\displaystyle 𝑑x\varphi _1(x)\varphi _2(x)\mathrm{}\varphi _n(x)}=`$ (1)
$`{\displaystyle \frac{d^4k_i}{(2\pi )^4}(2\pi )^4\delta (k_1+k_2+\mathrm{}+k_n)\stackrel{~}{\varphi }_1(k_1)\stackrel{~}{\varphi }_2(k_2)\mathrm{}\stackrel{~}{\varphi }_n(k_n)\mathrm{exp}(i\underset{i<j}{}k_ik_j)},`$ (2)
where $`\stackrel{~}{\varphi }_i`$ is the Fourier transform of the field $`\varphi _i`$, the index $`i`$ being used to distinguish different fields. In (2) we have introduced the notation $`ab=1/2a^\mu b^\nu \mathrm{\Theta }_{\mu \nu }`$, where $`\mathrm{\Theta }_{\mu \nu }`$ is the anti-symmetric constant matrix characterizing the noncommutativity of the underlying space. We shall assume from now on that $`\mathrm{\Theta }_{0i}=0`$ in order to evade causality and unitarity problems .
To represent Feynman amplitudes one could either use a double line notation, as the one introduced by ’t Hooft for matrix models, or single lines, which demands the symmetrization of the kernel (2) over the arguments of fields of the same kind. In this work we adopt the second systematics.
The paper is organized as follows. In Section II we present and discuss general aspects of the noncommutative Wess-Zumino model. The one loop analysis is performed in Section III, while in Section IV we demonstrate the renormalizability of the model to all orders of perturbation theory. Section V contains some final comments and the conclusions.
## II THE NONCOMMUTATIVE WESS-ZUMINO MODEL
In four dimensional Minkowski space-time the Wess-Zumino model is defined by the Lagrangian density
$``$ $`=`$ $`{\displaystyle \frac{1}{2}}A(^2)A+{\displaystyle \frac{1}{2}}B(^2)B+{\displaystyle \frac{1}{2}}\overline{\psi }(i\overline{)}m)\psi +{\displaystyle \frac{1}{2}}F^2+{\displaystyle \frac{1}{2}}G^2+mFA+mGB+`$ (3)
$`g(FA^2FB^2+2GAB\overline{\psi }\psi Ai\overline{\psi }\gamma _5\psi B),`$ (4)
where $`A`$ is a scalar field, $`B`$ is a pseudo scalar field, $`\psi `$ is a Majorana spinor field and $`F`$ and $`G`$ are, respectively, scalar and pseudoscalar auxiliary fields. By extending the above model to a noncommutative space one is led to the Lagrangian density
$``$ $`=`$ $`{\displaystyle \frac{1}{2}}A(^2)A+{\displaystyle \frac{1}{2}}B(^2)B+{\displaystyle \frac{1}{2}}\overline{\psi }(i\overline{)}m)\psi +{\displaystyle \frac{1}{2}}F^2+{\displaystyle \frac{1}{2}}G^2+mFA+mGB+`$ (5)
$`g(FAAFBB+GAB+GBA\overline{\psi }\psi A\overline{\psi }i\gamma _5\psi B).`$ (6)
It should be noticed that there is only one possible extension of the cubic term $`2GAB`$, to the noncommutative case, which preserves supersymmetry. It should also be emphasized that the noncommutative supersymmetry transformations are identical to the commutative ones since they are linear in the fields and no Moyal products are, therefore, involved. Hence, the extension of the theory to the noncommutative case does not alter the form of the Ward identities, which in turn implies that all fields have vanishing vacuum expectation values.
The elimination of the auxiliary fields through their corresponding equations of motion turns the bilinear terms in the Lagrangian Eq.(6) into the standard mass terms. On the other hand, the cubic terms produce quartic interactions which, in terms of a complex field $`\varphi =A+iB`$, can be cast as $`\varphi ^{}\varphi ^{}\varphi \varphi `$. This potential belongs to a class of non-renormalizable potentials, as discussed in . As it will be shown below, supersymmetry saves the day turning the theory into a renormalizable one.
The Lagrangian (6) was also written using the superspace formalism in . However, we will work with components fields in order to trace the effects of noncommutativity in the divergent Feynman integrals.
The propagators for the $`A`$ and $`F`$ fields are (see Fig. 1)
$`\mathrm{\Delta }_{AA}(p)`$ $`=`$ $`\mathrm{\Delta }(p){\displaystyle \frac{i}{p^2m^2+iϵ}},`$ (8)
$`\mathrm{\Delta }_{FF}(p)`$ $`=`$ $`p^2\mathrm{\Delta }(p),`$ (9)
$`\mathrm{\Delta }_{AF}(p)`$ $`=`$ $`\mathrm{\Delta }_{FA}(p)=m\mathrm{\Delta }(p),`$ (10)
whereas the propagators involving the $`B`$ and $`G`$ fields have identical expression (i.e., they are obtained by replacing $`A`$ by $`B`$ and $`F`$ by $`G`$). For the $`\psi `$ field we have
$$S(p)=\frac{i}{\overline{)}pm}.$$
(11)
The analytical expressions associated to the vertices are:
$`FA^2\text{vextex:}`$ $`ig\mathrm{cos}(p_1p_2),`$ (13)
$`FB^2\text{vextex:}`$ $`ig\mathrm{cos}(p_1p_2),`$ (14)
$`GAB\text{vertex:}`$ $`2ig\mathrm{cos}(p_1p_2),`$ (15)
$`\overline{\psi }\psi A\text{vertex:}`$ $`ig\mathrm{cos}(p_1p_2),`$ (16)
$`\overline{\psi }\psi B\text{vertex:}`$ $`ig\gamma _5\mathrm{cos}(p_1p_2).`$ (17)
Due to the oscillating factors provided by the cosines some of the integrals constructed with the above rules will be finite but in general divergences will survive, the degree of superficial divergence for a generic 1PI graph $`\gamma `$ being
$$d(\gamma )=4I_{AF}I_{BF}N_AN_B2N_F2N_G\frac{3}{2}N_\psi ,$$
(18)
where $`N_𝒪`$ denotes the number of external lines associated to the field $`𝒪`$ and $`I_{AF}`$ and $`I_{BF}`$ are the numbers of internal lines associated to the indicated mixed propagators. In all cases we will regularize the divergent Feynman integrals by using the supersymmetric regularization method proposed in .
## III THE ONE LOOP APPROXIMATION
It is straightforward to verify that, at the one loop level, all the tadpoles contributions add up to zero. This confirms the statement made in the previous section concerning the validity of the Ward identities.
Let us now examine the contributions to the self-energy of the $`A`$ field. The corresponding graphs are those shown in Fig 2$`a`$2$`e`$. In that figure diagrams $`a,b`$ and $`c`$ are quadratically divergent whereas graphs $`d`$ and $`e`$ are logarithmically divergent. We shall first prove that the quadratic divergences are canceled. In fact, we have that
$`\mathrm{\Gamma }_{\text{2}ac}(AA)`$ $`=`$ $`g^2{\displaystyle \frac{d^4k}{(2\pi )^4}\mathrm{cos}^2(kp)\{4k^2+4k^22\mathrm{T}\mathrm{r}[(\overline{)}k+\overline{)}p+m)(\overline{)}k+m)]\}}`$ (19)
$`\times \mathrm{\Delta }(k+p)\mathrm{\Delta }(k),`$ (20)
where the terms in curly brackets correspond to the graphs $`a`$, $`b`$ and $`c`$, respectively. After calculating the trace we obtain
$$\mathrm{\Gamma }_{\text{2}ac}(AA)=8g^2\frac{d^4k}{(2\pi )^4}(pk+m^2)\mathrm{cos}^2(kp)\mathrm{\Delta }(k)\mathrm{\Delta }(k+p).$$
(21)
This last integral is, at most, linearly divergent. However, the would be linearly divergent term vanishes by symmetric integration thus leaving us with an integral which is, at most, logarithmically divergent. Adding to Eq.(21) the contribution of the graphs 2$`d`$ and 2$`e`$ one arrives at
$$\mathrm{\Gamma }_{\text{2}ae}(AA)=8g^2\frac{d^4k}{(2\pi )^4}\mathrm{cos}^2(pk)(pk)\mathrm{\Delta }(k)\mathrm{\Delta }(k+p).$$
(22)
To isolate the divergent contribution to $`\mathrm{\Gamma }_{\text{2}ae}(AA)`$ we Taylor expand the coefficient of $`\mathrm{cos}^2(pk)`$ with respect to the variable $`p`$ around $`p=0`$, namely,
$`8g^2{\displaystyle }{\displaystyle \frac{d^4k}{(2\pi )^4}}\mathrm{cos}^2(pk)t^{(1)}(p)\left[(pk)\mathrm{\Delta }(k)\mathrm{\Delta }(k+p)\right].|_{p=0}`$ (23)
$`=`$ $`16g^2{\displaystyle \frac{d^4k}{(2\pi )^4}\mathrm{cos}^2(pk)\frac{(pk)^2}{(k^2m^2)^3}},`$ (24)
where $`t^{(r)}(p)`$ denotes the Taylor operator of order $`r`$. Since $`\mathrm{cos}^2(kp)=(1+\mathrm{cos}(2kp))/2`$ the divergent part of (24) is found to read
$$\mathrm{\Gamma }_{Div}(AA)=2g^2p^2\frac{d^4k}{(2\pi )^4}\frac{1}{(k^2m^2)^2}iI_\xi g^2p^2,$$
(25)
where the subscript $`\xi `$ remind us that all integrals are regularized through the procedure indicated in . In the commutative Wess-Zumino model this divergence occurs with a weight twice of the above. As usual, it is eliminated by the wave function renormalization $`A=Z^{1/2}A_r`$, where $`A_r`$ denotes the renormalized $`A`$ field. Indeed, it is easily checked that with the choice $`Z=1I_\xi g^2`$ the contribution (25) is canceled.
We turn next into analyzing the term containing $`\mathrm{cos}(2kp)`$ in (24). For small values of $`p`$ it behaves as $`p^2\mathrm{ln}(p^2/m^2)`$. Thus, in contradistinction to the nonsupersymmetric $`\varphi _4^4`$ case , there is no infrared pole and the function actually vanishes at $`p=0`$.
One may check that at one-loop the $`B`$ field self-energy is the same as the self-energy for the $`A`$ field, i. e., $`\mathrm{\Gamma }(BB)=\mathrm{\Gamma }(AA)`$. Therefore the divergent part of $`\mathrm{\Gamma }(BB)`$ will be eliminated if we perform the same wave function renormalization as we did for the $`A`$ field, $`B=Z^{1/2}B_r`$. We also found that the mixed two point Green functions do not have one-loop radiative corrections, $`\mathrm{\Gamma }(AF)=\mathrm{\Gamma }(BG)=0`$.
The one-loop corrections to the two point of the auxiliary field $`F`$ are depicted in Fig 3. The two graphs give identical contributions leading to the result
$$\mathrm{\Gamma }(FF)=4g^2\frac{d^4k}{(2\pi )^4}\mathrm{cos}^2(kp)\mathrm{\Delta }(k)\mathrm{\Delta }(k+p),$$
(26)
whose divergent part is
$$\mathrm{\Gamma }_{Div}(FF)=2g^2\frac{d^4k}{(2\pi )^4}\frac{1}{(k^2m^2)^2}=iI_\xi g^2,$$
(27)
involving the same divergent integral of the two point functions of the basic fields. It can be controlled by the field renormalization $`F=Z^{1/2}F_r`$, as in the case of $`A`$ and $`B`$. Analogous reasoning applied to the auxiliary field $`G`$ leads to the conclusion that $`G=Z^{1/2}G_r`$. However, things are different as far as the term containing $`\mathrm{cos}(2kp)`$ is concerned. It diverges as $`\mathrm{ln}(p^2/m^2)`$ as $`p`$ goes to zero. Nevertheless, this is a harmless singularity in the sense that its multiple insertions in higher order diagrams do not produce the difficulties pointed out in .
Let us now consider the corrections to the self-energy of the spinor field $`\psi `$ which are shown in Fig. 4. The two contributing graphs give
$`\mathrm{\Gamma }(\psi \overline{\psi })`$ $`=`$ $`4g^2{\displaystyle \frac{d^4k}{(2\pi )^4}\mathrm{cos}^2(kp)\mathrm{\Delta }(k)\mathrm{\Delta }(k+p)[(\overline{)}k+m)\gamma _5(\overline{)}k+m)\gamma _5]}`$ (28)
$`=`$ $`8g^2{\displaystyle \frac{d^4k}{(2\pi )^4}\mathrm{cos}^2(kp)\overline{)}k\mathrm{\Delta }(k)\mathrm{\Delta }(k+p)},`$ (29)
so that for the divergent part we get $`\mathrm{\Gamma }_{Div}(\psi \overline{\psi })=ig^2\overline{)}pI_\xi `$ leading to the conclusion that the spinor field presents the same wave function renormalization of the bosonic fields, i. e., $`\psi =Z^{1/2}\psi _r`$. As for the term containing $`\mathrm{cos}(2kp)`$ it behaves as $`\overline{)}p\mathrm{ln}(p^2/m^2)`$ and therefore vanishes as $`p`$ goes to zero.
The one-loop superficially (logarithmically) divergent graphs contributing to the three point function of the $`A`$ field are shown in Fig 5. The sum of the amplitudes corresponding to the graphs 5$`a`$ and 5$`b`$ is
$`\mathrm{\Gamma }_{5a+5b}(AAA)`$ $`=`$ $`96ig^3m{\displaystyle \frac{d^4k}{(2\pi )^4}(kp_2)^2\mathrm{\Delta }(k)\mathrm{\Delta }(k+p_3)\mathrm{\Delta }(kp_2)\mathrm{cos}(kp_1+p_3p_1)}`$ (30)
$`\times \mathrm{cos}(p_2k)\mathrm{cos}(p_3k),`$ (31)
while its divergent part is found to read
$$\mathrm{\Gamma }_{5a+5bDiv}(AAA)=24ig^3m\mathrm{cos}(p_3p_1)\frac{d^4k}{(2\pi )^4}(k)^2(\mathrm{\Delta }(k))^3.$$
(32)
The divergent part of the graph 5$`c`$, nonetheless, gives a similar contribution but with a minus sign so that the two divergent parts add up to zero. Thus, up to one-loop the three point function $`\mathrm{\Gamma }(AAA)`$ turns out to be finite. Notice that a nonvanishing result would spoil the renormalizability of the model. The analysis of $`\mathrm{\Gamma }(ABB)`$ follows along similar lines and with identical conclusions. Furthermore, it is not difficult to convince oneself that $`\mathrm{\Gamma }(FAA)`$, $`\mathrm{\Gamma }(FBB)`$ and $`\mathrm{\Gamma }(GAB)`$ are indeed finite.
As for $`\mathrm{\Gamma }(A\psi \overline{\psi })`$ we notice that superficially divergent contributions arise from the diagrams depicted in Figs 6$`a`$ and 6$`b`$. In particular, diagram 6$`a`$ yields
$`\mathrm{\Gamma }_{6a}(A\psi \overline{\psi })`$ $`=`$ $`8ig^3{\displaystyle \frac{d^4k}{(2\pi )^4}\mathrm{\Delta }(k)\mathrm{\Delta }(p_2+k)\mathrm{\Delta }(kp_1)(\overline{)}p_2+\overline{)}k+m)(\overline{)}k\overline{)}p_1+m)}`$ (33)
$`\times \mathrm{cos}(kp_3p_1p_3)\mathrm{cos}(kp_1)\mathrm{cos}(kp_2),`$ (34)
while 6$`b`$ gives
$`\mathrm{\Gamma }_{6b}(A\psi \overline{\psi })`$ $`=`$ $`8ig^3{\displaystyle \frac{d^4k}{(2\pi )^4}\mathrm{\Delta }(k)\mathrm{\Delta }(p_2+k)\mathrm{\Delta }(kp_1)\gamma _5(\overline{)}p_2+\overline{)}k+m)(\overline{)}k\overline{)}p_1+m)\gamma _5}`$ (35)
$`\times \mathrm{cos}(kp_3p_1p_3)\mathrm{cos}(kp_1)\mathrm{cos}(kp_2),`$ (36)
so that the sum of the two contributions is also finite. The same applies for $`\mathrm{\Gamma }(B\psi \overline{\psi })`$.
We therefore arrive at another important result, namely, that there is no vertex renormalization at the one loop level. This parallels the result of the commutative Wess-Zumino model.
To complete the one-loop analysis we must examine the four point functions. Some of the divergent diagrams contributing to $`\mathrm{\Gamma }(AAAA)`$ are depicted in Fig. 7$`ac`$. The analytical expression associated with the graph 7$`a`$ is
$`\mathrm{\Gamma }_{7a}(AAAA)`$ $`=`$ $`16g^4{\displaystyle \frac{d^4k}{(2\pi )^4}k^2\mathrm{\Delta }(k)\mathrm{\Delta }(k+p_1)(k+p_1+p_3)^2\mathrm{\Delta }(k+p_1+p_3)\mathrm{\Delta }(p_2k)}`$ (37)
$`\times \mathrm{cos}(kp_1)\mathrm{cos}(kp_2)\mathrm{cos}[(k+p_1)p_3]\mathrm{cos}[(kp_2)p_4].`$ (38)
There are five more diagrams of this type, which are obtained by permuting the external momenta $`p_2,p_3`$ and $`p_4`$ while keeping $`p_1`$ fixed. Since we are interested in the (logarithmic) divergence associated with this diagram, we set all the external momenta to zero in the propagators but not in the arguments of the cosines. This yields
$`\mathrm{\Gamma }_{7aDiv}(AAAA)`$ $`=`$ $`16g^4{\displaystyle \frac{d^4k}{(2\pi )^4}(k^2)^2(\mathrm{\Delta }(k))^4}`$ (39)
$`\times \mathrm{cos}(kp_1)\mathrm{cos}(kp_2)\mathrm{cos}[(k+p_1)p_3]\mathrm{cos}[(kp_2)p_4].`$ (40)
Adopting the same procedure for the other five graphs we notice that the corresponding contributions are pairwise equal. The final result is therefore
$`\mathrm{\Gamma }_{Div}(AAAA)=32g^4{\displaystyle \frac{d^4k}{(2\pi )^4}(k^2)^2(\mathrm{\Delta }(k))^4\mathrm{cos}(kp_1)}`$ (41)
$`\times [\mathrm{cos}(kp_2)\mathrm{cos}[(k+p_1)p_3]\mathrm{cos}[(kp_2)p_4]+p_3p_4+p_2p_4].`$ (42)
There is another group of six diagrams, Fig 7$`b`$, which are obtained from the preceding ones by replacing the propagators of $`A`$ and $`F`$ fields by the propagator of the $`B`$ and $`G`$ fields, respectively. The net effect of adding these contributions is, therefore, just to double the numerical factor in the right hand side of the above formula.
Besides the two groups of graphs just mentioned, there are another six graphs with internal fermionic lines. A representative of this group has been drawn in Fig 7$`c`$. It is straightforward to verify that because of the additional minus sign due to the fermionic loop, there is a complete cancellation with the other contributions described previously. The other four point functions may be analyzed similarly with the same result that no quartic counterterms are needed.
## IV Absence of Mass and Coupling Constant Renormalization to all Orders of Perturbation Theory
In the previous section we proved that up to one loop the noncommutative Wess-Zumino model is renormalizable and only requires a common wave function renormalization. Here, we shall prove that no mass and coupling constant counterterms are needed at any finite order of perturbation theory. As in the commutative case, our proof relies heavily on the Ward identities.
We start by noticing that from Eq.(2) it follows that
$$d^4y\frac{\delta }{\delta 𝒪(y)}d^4x\underset{n\mathrm{factors}}{\underset{}{𝒪(x)𝒪(x)\mathrm{}𝒪(x)}}=nd^4x\underset{n1\mathrm{factors}}{\underset{}{𝒪(x)𝒪(x)\mathrm{}𝒪(x)}}.$$
(43)
In turns, this enables one to find
$$\frac{}{m}Z(J)=\frac{m}{2g}\frac{\delta Z(J)}{\delta J_F(y)}d^4y\frac{iZ(J)}{2g}J_A(y)d^4y,$$
(44)
which looks formally identical to the corresponding relation in the commutative case . Here, $`Z(J)`$ is the Green function generating functional and $`J_𝒪`$ is the external source associated to the field $`𝒪`$. By collectively denoting the fields by $`\varphi `$, $`Z(J)`$ can be cast as
$$Z(J)=D\varphi \mathrm{exp}i\left(S+d^4xJ\varphi \right),$$
(45)
where $`S=d^4x`$ and $``$ is the regularized Lagrangian.
In terms of the 1PI generating functional $`\mathrm{\Gamma }(R)`$ the identity (44) becomes
$$\frac{}{m}\mathrm{\Gamma }[R]=\frac{m}{2g}R_F(y)d^4y+\frac{1}{2g}\frac{\delta \mathrm{\Gamma }[R]}{\delta R_A(y)}d^4y.$$
(46)
By taking the functional derivative with respect to $`R_F`$ and then putting all $`R`$’s equal to zero we obtain
$$m=\mathrm{\Gamma }(FA)|_{p^2=0}=Z^1\mathrm{\Gamma }_r(FA)|_{p^2=0},$$
(47)
where $`\mathrm{\Gamma }_r(AF)`$ is the renormalized 1PI Green function of the indicated fields. We take as normalization conditions those specified in . Specifically, $`\mathrm{\Gamma }_r(FA)|_{p^2=0}=m_r`$, where $`m_r`$ is taken to be the renormalized mass. Hence, $`m_r=Zm`$ implying that there is no additive mass renormalization. Through similar steps one also finds that $`g_r=Z^{3/2}g`$, where $`g_r`$ is the renormalized coupling constant. This implies the absence of coupling constant counterterms.
We stress the fact that, by exploiting the Ward identities, we have succeeded in generalizing to all orders of perturbation theory the one loop result concerned with the absence of counterterms different from those already present in the original Lagrangian.
## V Conclusions
After extending the Wess-Zumino model to the noncommutative Minkowski space, we succeeded in demonstrating, to all orders of perturbation, that the theory is free of nonintegrable infrared singularities and renormalizable. Thus, this model provides an example of a fully consistent noncommutative quantum field theory.
It shares some properties with the Wess-Zumino model. The quadratic and linear divergences are absent. Furthermore, only a wave function renormalization is needed to make the theory finite. Also, all fields exhibit the same mass as is the case in any ordinary supersymmetric theory.
On the other hand, one should notice that the commutative Wess-Zumino model can not be recovered from the noncommutative one at the limit of vanishing deformation. In fact, the limit of vanishing deformation does not exist because of logarithmic singularities.
A very important feature of the noncommutative Wess-Zumino model is that all vertices were deformed in the same way. This was essential to split the amplitudes into planar and non planar contributions in a uniform way so that the renormalizability properties of the Wess-Zumino model is always present in the planar sector. The reason for the deformed vertices to be the same is the presence of the auxiliary fields. With them all interactions are cubic. The elimination of the auxiliary fields produces cubic and quartic interactions and the vertices will be deformed in different ways. Of course, the renormalizability properties will be the same without the auxiliary fields but, surely, more difficult to prove. Supersymmetric gauge theories have cubic and quartic vertices even in the presence of auxiliary fields. We expect that the renormalizability proof will be much more difficult unless further simplification arise. Studies in this direction are in progress.
ACKNOWLEDGMENTS
This work was partially supported by Fundação de Amparo à Pesquisa do Estado de São Paulo (FAPESP) and Conselho Nacional de Desenvolvimento Científico e Tecnológico (CNPq).
|
warning/0005/nlin0005069.html
|
ar5iv
|
text
|
# Absence of supersensitivity to small input signals in generalized on–off systems
## Abstract
It has recently been shown that nonlinear skew product dynamical systems with invariant subspaces which are capable of displaying on–off intermittency can show supersensitivity to small input signals.
Here we show that this supersensitivity is absent for more general dynamical systems with non–skew product structure, capable of displaying a generalized form of on–off intermittency, and is therefore in this sense fragile. This absence of supersensitivity is of importance in view of the fact that dynamical systems are generically expected to be of non–skew product nature.
Many dynamical systems of interest possess symmetries or constraints which force the presence of invariant subspaces. A great deal of effort has recently gone into the study of such systems (see e.g. ). A sub–class of these models, namely those with skew product structure (and normal parameters<sup>*</sup><sup>*</sup>*Parameters that leave the dynamics on the invariant manifold unchanged are called normal, otherwise they are referred to as non–normal.), have been shown to be capable of producing a number of novel modes of behavior, including on–off intermittency and bubbling .
Recently, Zhou and Lai have shown that systems of this type can display supersensitivity, in the sense that small constant or time varying inputs to the system can induce extremely large responses. The authors further claim that with an additional odd symmetry condition, this supersensitivity is robust to addition of noise. Such supersensitivity could be of importance in many fields, including the study of synchronization of coupled chaotic systems and the design of sensor devices .
The results on on–off intermittent systems reported by these authors can all be described within the framework of skew product systems. Generically, however, one would expect typical dynamical systems to have non–skew product structure (with non–normal parameters). Systems of this type have recently been studied and have been found to be capable of displaying a number of additional novel dynamical modes of behavior, absent in skew product systems, including a generalization of on–off intermittency, referred to as in–out intermittency .
The easiest way to characterise in–out intermittency is by contrasting it with on–off intermittency. Briefly, let $`M_I`$ be the invariant subspace and $`A`$ the attractor which exhibits either on–off or in–out intermittency. If the intersection $`A_0=AM_I`$ is a minimal attractor, then we have on–off intermittency, whereas if $`A_0`$ is not a minimal attractor, then we have in–out intermittency. In the latter case there can be different invariant sets in $`A_0`$ associated with attraction and repulsion transverse to $`A_0`$, hence the name in–out. Another crucial difference between the two is that, as opposed to on–off intermittency, in the case of in–out intermittency the minimal attractors in the invariant subspaces do not necessarily need to be chaotic and hence the trajectories can (and often do) shadow a periodic orbit in the ‘out’ phases .
Our aim here is to find out whether this type of supersensitivity, observed in on–off intermittent systems, persists in more general non–skew product systems which are capable of displaying in–out intermittent behaviour.
A simple class of maps that can model both on–off and in–out types of intermittency is given by
$$x_{n+1}=F(x_n,y_n,𝐚),y_{n+1}=G(x_n,y_n,𝐚),$$
(1)
where $`G(x_n,0,𝐚)=0`$, the variables $`x_n`$ and $`y_n`$ represent the dynamics within the invariant submanifold ($`y=0`$) and the transverse distance to it respectively and $`𝐚^𝐦`$ are the control parameters of the system. A special subset of these systems, with skew product form over the dynamics in $`x`$, can be written as
$$x_{n+1}=F(x_n,𝐚),y_{n+1}=G(x_n,y_n,𝐚).$$
(2)
By considering a skew product system of type (2), Zhou and Lai modelled the motion near the invariant submanifold $`y=0`$, using a Fokker–Planck equation. In this way they were able to predict that the sensitivity $`S`$ of the map in the neighbourhood of a blowout bifurcation, leading to on–off intermittency, is given by
$$S=\frac{y}{p}=\frac{\tau }{p\mathrm{ln}(\tau /p)},$$
(3)
where $`y`$ is the average of the transverse variable $`y`$, $`p`$ is the input signal and $`\tau `$ is the threshold below which $`y`$ goes through a laminar phase. They were able to confirm this prediction numerically.
To study whether non–skew product (in–out intermittent) systems can also display supersensitivity, we considered a particular example of the map (1) in the form
$`x_{n+1}`$ $`=`$ $`rx_n(1x_n)+sx_ny_n^2,`$ (4)
$`y_{n+1}`$ $`=`$ $`\nu e^{bx_n}y_n+ay_n^3,`$ (5)
where $`r[0,4]`$ and $`(s,\nu ,a,b)`$ are the control parameters. Note that for $`s=0`$, the map (4) has the skew product form (2) and for fixed $`r`$, the parameters $`s`$, $`a`$, $`b`$ or $`\nu `$ are normal. Thus depending upon the choice of its parameters, this map is capable of displaying both on–off and in–out types of intermittency, Note also that this map possesses the odd symmetry condition $`G(y)=G(y)`$ that was found in to be required for the robustness of supersensitivity with respect to noise. Also the transverse Lyapunov exponent $`\lambda _T`$ for this map can be readily calculated to be
$$\lambda _T=\mathrm{ln}\nu +bx_r,$$
(6)
where $`x_r`$ is the average of the variable $`x_n`$ for an initial condition on the invariant submanifold $`y=0`$.
To study the effects of an input signal on in–out systems, we considered a variant of this map given by
$`x_{n+1}`$ $`=`$ $`rx_n(1x_n)+sx_ny_n^2,`$ (7)
$`y_{n+1}`$ $`=`$ $`\nu e^{bx_n}y_n+ay_n^3+p,`$ (8)
where the real parameter $`p`$ models the effects of a small input signal.
To begin with, we made a comparative numerical study of the sensitivity of in–out and on–off systems to input signals $`p`$, using (7). Fig. 1 shows a comparative study of the bursting behaviour in the two cases close to, but below, their blowout points. As can be seen from the comparison of the lower panels, there is very little bursting in the in–out case.
To further demonstrate this relative insensitivity in the in–out case, we made a study of the sensitivity $`S`$ of the systems close to their blowout points, as a function of the input signal $`p`$. This is shown in Fig. 2, which again demonstrates a distinct absence of supersensitivity for the in–out case, specially for the lower input signal levels. Furthermore, it shows a saturation in sensitivity in the in–out case for input signals $`p\begin{array}{c}<\\ \end{array}10^7`$.
These results indicate a clear lack of supersensitivity in the in–out case to small constant input signals.
To understand this qualitative difference between the on–off and in–out cases, we briefly recall a number of differences between the two cases, relevant to our discussion here. In the case of on–off, the attraction and the ejection of the orbits near the invariant manifold are brought about by a single chaotic attractor in $`M_I`$. Thus for the values of the control parameter close to but below the blowout point, the chaotic attractor in $`M_I`$ becomes transversally weakly attracting, but there can be repelling orbits within this attractor that are transversally unstable, leading to bubbling and allowing the orbit to access the lower and upper boundaries frequently. The system thus becomes sensitive to small inputs, producing large bursts and hence supersensitivity.
For the in–out case, on the other hand, the ‘in’ and ‘out’ phases are governed by two separate invariant sets in $`M_I`$: a chaotic saddle and a periodic attractor respectively. Thus for the values of control parameter ($`b`$ in our case) above the blowout value, the chaotic saddle in the invariant submanifold is transversally attracting whereas the periodic attractor in $`M_I`$ is transversally unstable with a positive transverse Lyapunov exponent $`\lambda _T`$. As a result, an orbit drawn towards the invariant submanifold by the chaotic saddle is thus ejected by the transversally unstable periodic attractor, leading to in–out intermittency. On the other hand, for the values of control parameter $`b`$ just below the blowout value, the unique periodic attractor in the invariant submanifold becomes transversally stable (with $`\lambda _T<0`$), while the chaotic saddle still remains transversally attracting. As a result orbits drawn towards the invariant submanifold by the chaotic saddle get attracted to the periodic orbit there (see Fig. 3 for a schematic depiction).
We shall now show that the presence of a small input signal $`p`$ leaves the above dynamical picture essentially unaltered, apart from displacing the location of the periodic attractor off the previous invariant submanifold. There are two ways to see this. Firstly, for small values of the input signal $`p𝒪(1)`$, the periodic orbit is expected to persist by continuity. We numerically confirmed that the period 12 orbit involved in the in–out intermittency studied here (see for details) does indeed survive for small values of $`p`$, albeit shifted slightly off the invariant submanifold $`M_I`$ (see the left panels of Fig. 1).
Alternatively, we can estimate the transverse location of the displaced periodic orbit. To do this we recall that we are interested in small displacements from $`M_I`$, which implies that as a first approximation we may ignore higher order dependence on $`y`$. We therefore approximate the map (7) by
$`x_{n+1}`$ $`=`$ $`rx_n(1x_n)+sx_ny_n^2,`$ (9)
$`y_{n+1}`$ $``$ $`\nu e^{bx_n}y_n+p,`$ (10)
where the second order term in $`y`$ has been kept in the $`x`$ map in order to ensure the essential overall non–skew product structure of the system.
The period 12 attractor involved in this case has $`x_n`$ values satisfying $`x_{n+12}=x_n`$ and $`y_n`$ values given by
$`y_0`$ $`=`$ $`p,`$
$`y_1`$ $`=`$ $`\left(\nu e^{bx_0}+1\right)p,`$
$`y_2`$ $`=`$ $`\left(\nu ^2e^{b(x_0+x_1)}+\nu e^{bx_1}+1\right)p,`$
$`\mathrm{}`$
$`y_{12}`$ $`=`$ $`\left(\nu ^{12}e^{b(x_0+x_1+\mathrm{}+x_{11})}+\mathrm{}+\nu e^{bx_{11}}+1\right)p,`$
$`y_{13}`$ $`=`$ $`(\nu ^{13}e^{b(2x_0+x_1+\mathrm{}+x_{11})}`$
$`+\nu ^{12}e^{b(x_0+x_1+\mathrm{}+x_{11})}+\mathrm{}+\nu e^{bx_0}+1)p,`$
$`\mathrm{}`$
$`y_n`$ $`=`$ $`{\displaystyle \underset{j=0}{\overset{n}{}}}\left(\nu ^je^{b\left[\underset{i=1}{\overset{12}{}}\frac{j+11(imod12)}{12}x_{(i+n1)mod12}\right]}\right)p,`$
where $``$ denotes the integer part.
The above expressions for $`y_n`$ change periodically (with period 12), depending on the initial $`x`$. The asymptotic average value of $`y`$ can then be approximated by
$`\underset{n\mathrm{}}{lim}y`$ $``$ $`{\displaystyle \underset{j=0}{\overset{n}{}}}\left(\nu ^je^{b\left[\underset{i=0}{\overset{11}{}}\frac{j}{12}x_i\right]}\right)p,`$ (11)
$`=`$ $`\left(1\nu e^{\frac{b}{12}\underset{i=0}{\overset{11}{}}x_i}\right)^1p`$ (12)
$`=`$ $`{\displaystyle \frac{p}{1\nu e^{\lambda _T\mathrm{ln}\nu }}}={\displaystyle \frac{p}{1e^{\lambda _T}}},`$ (13)
where we have used (6). Interestingly this enables us to find the sensitivity $`S`$ as a function of $`\lambda _T`$
$$S=\frac{1}{1e^{\lambda _T}},$$
(14)
which is independent of $`p`$, thus explaining the saturation in sensitivity $`S`$ observed in the in–out case in Fig. 2, in clear contrast to expression for the on–off sensitivity given by (3) .
It now remains to show that apart from the above shift off the submanifold, the periodic attractor remains essentially intact. To see this, recall that the effect of a non–zero $`p`$ on $`x`$ is given by
$$x_{n+1}F^n(x_1,𝐚)+sx_ny^2,$$
(15)
where $`F^n(x_1,𝐚)`$ represents the $`x`$ component of the n<sup>th</sup> iterate of the map (2). Using (11) this gives
$$|s|x_ny^2|s|𝒪(1)\left[\frac{p}{1e^{\lambda _T}}\right]^2.$$
(16)
Now for input signals $`p𝒪(1)`$ and for the parameter regimes chosen here, $`1e^{\lambda _T}\sqrt{p}`$, which implies
$$|s|x_ny^2p,$$
(17)
showing that to this approximation the $`p`$-induced variations in $`x`$ are extremely small (relative to $`p`$), hence providing a good indication that the periodic attractor remains essentially intact.
The above arguments and results demonstrate the qualitative differences between the responses of the on–off and the in–out dynamics to small input signals. In particular, the survival of the periodic orbit in the latter case acts to trap the incoming orbits and therefore blocks the possibility of supersensitivity in this case. We expect this picture to be common and thus supersensitivity to be absent in the generic non–skew product (in–out) settings.
To further substantiate this finding, we calculated the average blowout variable $`y`$ in system (7) for fixed input signals, as a function of $`\lambda _T`$. The results are summarised in Fig. 4, which show that for $`\lambda _T>0`$, both cases are relatively insensitive to input signals, whereas for $`\lambda _T<0`$, the on–off case is much more sensitive to input signals than the in–out case, with the latter dependence in very good agreement with our prediction (11).
Finally we calculated the dependence of the sensitivity $`S`$ for the in–out case as a function of $`\lambda _T`$, with different input signals, as a function of $`\lambda _T`$. The results are shown in Fig. 5, together with our predicted expression (14), which show excellent agreement.
To summarise, we have argued that the supersensitivity found in for the case of skew product systems with on–off intermittency is absent in the more general setting of non–skew product systems, capable of displaying in–out intermittency. We have substantiated this claim both analytically and through extensive numerical simulations. We have also checked that the absence of supersensitivity in the in–out case remains robust to changes in both the input signal (of the form $`p\mathrm{sin}(2\pi x)`$) as well as to unbiased noise in the transverse direction (of the order of the input signal).
The absence of supersensitivity for systems displaying in–out intermittency is important, particularly given that dynamical systems are generically expected to be of non–skew product type.
EC is supported by a PPARC fellowship and RT benefited from PPARC UK Grant No. L39094.
|
warning/0005/math0005168.html
|
ar5iv
|
text
|
# On some automorphisms of the set of effects on Hilbert space
## Abstract.
The set of all effects on a Hilbert space has an affine structure (it is a convex set) as well as a multiplicative structure (it can be equipped with the so-called Jordan triple product). In this paper we describe the corresponding automorphisms of that set.
###### Key words and phrases:
Effect algebra, Hilbert space, operator algebra, factor, automorphism This research was supported from the following sources:
1) Hungarian National Foundation for Scientific Research (OTKA), Grant No. T–030082 F–019322,
2) A grant from the Ministry of Education, Hungary, Reg. No. FKFP 0304/1997
Let $`H`$ be a complex Hilbert space. Denote by $`B(H)`$ the $`C^{}`$-algebra of all bounded linear operators on $`H`$. The operator interval $`[0,I]`$ of all positive operators in $`B(H)`$ which are bounded by the identity $`I`$ is called the Hilbert space effect algebra. This has important applications in quantum mechanics. The effect algebra $`[0,I]`$ can be equipped with several algebraic operations. For example, one can define a partial addition on it. Namely, if $`A,B[0,I]`$ and $`A+B[0,I]`$, then one can set $`AB=A+B`$. This structure has been investigated in several papers (see and the references therein). Moreover, on $`[0,I]`$ there is a natural partial ordering $``$ which comes from the usual ordering between the self-adjoint operators on $`H`$ and one can also define the operation of the so-called orthocomplementation by $`{}_{}{}^{}:AIA`$. The set of all effects on $`H`$ equipped with this ordering and orthocomplementation has been studied for example in . Next, $`[0,I]`$ is clearly a convex subset of the linear space $`B(H)`$. So, one can consider the operation of convex combinations on it. The set of all effects with this structure has been investigated in , for instance. Finally, as for a mutliplicative operation on $`[0,I]`$, note that in general $`A,B[0,I]`$ does not imply that $`AB[0,I]`$. However, we all the time have $`ABA[0,I]`$. This multiplication which is a nonassociative operation and sometimes called Jordan triple product also appears in infinite dimensional holomorphy as well as in connection with the geometrical properties of $`C^{}`$-algebras.
Because of the importance of effect algebras, it is a natural problem to study the isomorphisms of the mentioned structures. The aim of this paper is to contribute to these investigations.
The automorphisms of $`[0,I]`$ with the operation of partial addition were described in . The automorphisms of the Hilbert space effect algebra equipped with the partial ordering $``$ and the orthocomplementation were characterized in (see also ). The isomorphisms of $`[0,I]`$ as a convex subset of $`B(H)`$ were investigated in . However, in that paper the authors considered such affine functions (maps preserving convex combinantions) which are homogenous for the scalars in $`[0,1]`$. This means that they supposed that their affine bijections have the additional property that they send 0 to 0. As a corollary of our first theorem we describe the affine isomorphisms of $`[0,I]`$ without this extra condition. In the second theorem we determine the automorphisms of $`[0,I]`$ equipped with the Jordan triple product. It is worth mentioning that, as it turns out from our results, the linear and multiplicative structures of $`[0,I]`$ are very closely related to each other.
Let us fix the notation and definitions that we shall use throughout the paper. So, $`B(H)`$ and $`B_s(H)`$ denote the $`C^{}`$-algebra of all bounded linear operators on $`H`$ and the $`JB^{}`$-algebra of all bounded self-adjoint operators on $`H`$, respectively. A self-adjoint idempotent $`P`$ in $`B(H)`$ is called a projection. A von Neumann algebra $`𝒜`$ on $`H`$ is said to be a factor if its center is trivial, that is, it equals $`I`$ ($`I`$ is the identity on $`H`$). Define the set $`(𝒜)`$ of all effects in $`𝒜`$ by $`(𝒜)=[0,I]𝒜`$. If $`_1,_2`$ are \*-algebras over the complex field, then a linear map $`\psi :_1_2`$ satisfying $`\psi (A)^{}=\psi (A^{})`$ $`(A_1)`$ is called
* a Jordan \*-homomorphism if $`\psi (A)^2=\psi (A^2)`$ $`(A_1)`$;
* a \*-homomorphism if $`\psi (A)\varphi (B)=\psi (AB)`$ $`(A,B_1)`$;
* a \*-antihomomorphism if $`\psi (A)\psi (B)=\psi (BA)`$ $`(A,B_1)`$.
If $`X,Y`$ are linear spaces over $``$ and $`CX`$ is a convex set, then the function $`\psi :CY`$ is called affine if it satisfies
$$\psi (\lambda x+(1\lambda )y)=\lambda \psi (x)+(1\lambda )\psi (y)$$
for every $`x,yC`$ and $`\lambda [0,1]`$.
Our first result determines the affine automorphisms of $`(𝒜)`$ for any factor $`𝒜`$.
###### Theorem 1.
Let $`𝒜`$ be a factor. If $`\varphi :(𝒜)(𝒜)`$ is a bijective affine function, then there is an either \*-automorphism or \*-antiautomorphism $`\mathrm{\Phi }`$ of $`𝒜`$ such that
$$\varphi (A)=\mathrm{\Phi }(A)(A(𝒜))$$
or
$$\varphi (A)=\mathrm{\Phi }(IA)(A(𝒜)).$$
###### Proof.
First we recall the following fact whose proof requires only trivial calculations. Let $`\varphi :(𝒜)X`$ be an affine function with $`\varphi (0)=0`$ where $`X`$ is a linear space. Define
$$\mathrm{\Phi }_1(A)=\{\begin{array}{cc}0\hfill & \text{ if }A=0\hfill \\ A\varphi \left(\frac{A}{A}\right)\hfill & \text{ if }A0,0A𝒜.\hfill \end{array}$$
Next let
$$\mathrm{\Phi }_2(A)=\mathrm{\Phi }_1(A^+)\mathrm{\Phi }_1(A^{})(A^{}=A𝒜),$$
where $`A^+`$ and $`A^{}`$ are the positive part and the negative part of $`A`$, respectively. That is,
$$A^+=(1/2)(|A|+A)\text{ and }A^{}=(1/2)(|A|A).$$
Finally, set
$$\mathrm{\Phi }(A)=\mathrm{\Phi }_2(\mathrm{Re}A)+i\mathrm{\Phi }_2(\mathrm{Im}A)(A𝒜),$$
where $`\mathrm{Re}A`$ and $`\mathrm{Im}A`$ denote the real part and the imaginary part of $`A`$, respectively. Then $`\mathrm{\Phi }:𝒜X`$ is the unique linear extension of $`\varphi `$ from $`(𝒜)`$ to $`𝒜`$.
Let $`\varphi :(𝒜)(𝒜)`$ be an affine function. We assert that $`\varphi `$ is continuous in the norm topology. To see this, consider the affine function
$$\psi :A\varphi (A)\varphi (0)$$
on $`(𝒜)`$ which sends 0 to 0. Since its unique linear extension $`\mathrm{\Psi }:𝒜𝒜`$ has the property that $`\mathrm{\Psi }(A)+\varphi (0)(𝒜)`$ for every $`A(𝒜)`$, we deduce that
$$\mathrm{\Psi }(A)\mathrm{\Psi }(A)+\varphi (0)+\varphi (0)2(A(𝒜)).$$
Clearly, every element $`A`$ of the unit ball of $`𝒜`$ can be written as $`A=A_1A_2+i(A_3A_4)`$ for some $`A_1,A_2,A_3,A_4(𝒜)`$. It follows that $`\mathrm{\Psi }`$ is bounded on the unit ball of $`𝒜`$. This implies that $`\mathrm{\Psi }`$ and hence $`\varphi `$ are norm continuous.
Let now $`\varphi :(𝒜)(𝒜)`$ be an affine bijection. Then $`\varphi `$ and it inverse are norm-continuous. Moreover, $`\varphi `$ obviously preserves the extreme points of $`(𝒜)`$ which are exactly the projections in $`𝒜`$.
We claim that $`\varphi (0)`$ is either $`0`$ or $`I`$. Let $`P0,I`$ be a projection in $`𝒜`$. If every projection in $`𝒜`$ commutes with $`P`$, then we obtain that every element of $`𝒜`$ commutes with $`P`$ which, $`𝒜`$ being a factor, would imply that $`P`$ is a scalar multiple of the identity but this is a contradiction. So, we can choose a projection $`Q`$ in $`𝒜`$ which does not commute with $`P`$. Considering the operator $`U=I2Q𝒜`$ we get a unitary element in $`𝒜`$ which does not commute with $`P`$. So, we have $`PUPU^{}`$. In any von Neumann algebra the unitary group is arcwise connected. Therefore, there is an arc within the set of all projections in $`𝒜`$ connecting $`P=IPI^{}`$ to $`UPU^{}`$. To sum up, every nontrivial projection in $`𝒜`$ can be connected by an arc within the set of all projections to another projection different from the first one. It is trivial that $`0`$ and $`I`$ can be connected only to themselves. Since $`\varphi `$ is a homeomorphism of the set of all projections in $`𝒜`$, we deduce that $`\varphi `$ sends nontrivial projections to nontrivial projections and hence we have either $`\varphi (0)=0`$ or $`\varphi (0)=I`$. Clearly, we can assume without loss of generality that $`\varphi (0)=0`$ (otherwise, we consider the transformation $`AI\varphi (A)`$). Let $`\mathrm{\Phi }`$ be the unique linear extension of $`\varphi `$ onto $`𝒜`$. We already know that $`\mathrm{\Phi }`$ is a bounded linear transformation which sends projections to projections. It is a standard algebraic argument to verify that $`\mathrm{\Phi }`$ is a Jordan \*-homomorphism (see \[1, Remark 2.2\] and use the spectral theorem of self-adjoint operators together with the continuity of $`\mathrm{\Phi }`$).
We assert that $`\mathrm{\Phi }`$ is bijective. If $`\mathrm{\Phi }(A)=0`$, then we see that $`\mathrm{\Phi }_2(\mathrm{Re}A)=\mathrm{\Phi }_2(\mathrm{Im}A)=0`$. Let $`B,C`$ denote the positive and negative parts of $`\mathrm{Re}A`$, respectively, From $`\mathrm{\Phi }_2(\mathrm{Re}A)=0`$ we infer that $`\mathrm{\Phi }_1(B)=\mathrm{\Phi }_1(C)`$. Supposing that $`B,C0`$, this means that
$$B\varphi \left(\frac{B}{B}\right)=C\varphi \left(\frac{C}{C}\right).$$
Using the homogenity of $`\varphi `$ for the scalars in $`[0,1]`$, we conclude that
$$\varphi \left(\frac{B}{B+C}\right)=\varphi \left(\frac{C}{B+C}\right).$$
Since $`\varphi `$ is injective, it follows that $`B=C`$ which gives us that $`\mathrm{Re}A=0`$. Similarly, one can check that $`\mathrm{Im}A=0`$ is also true, so we have $`A=0`$. Therefore, $`\mathrm{\Phi }`$ is injective. Since the range of $`\mathrm{\Phi }`$ is a linear subspace of $`𝒜`$ which contains $`(𝒜)`$ (recall that $`\mathrm{\Phi }`$ is an extension of $`\varphi `$), it follows that $`\mathrm{\Phi }`$ is surjective. So, $`\mathrm{\Phi }`$ is a Jordan \*-automorphism of $`𝒜`$.
It is well-known that every factor is a prime algebra. This means that for any $`A,B𝒜`$, the equality $`A𝒜B=\{0\}`$ implies that $`A=0`$ or $`B=0`$. Now, a classical theorem of Herstein on Jordan homomorphisms applies to obtain that $`\mathrm{\Phi }`$ is either a \*-automorphism or a \*-antiautomorphism of $`𝒜`$. This completes the proof of the theorem. ∎
Taking into account the form of \*-automorphisms and \*-antiautomorphisms of $`B(H)`$, we immediately have the following corollary.
###### Corollary 2.
Let $`\varphi :[0,I][0,I]`$ be a bijective affine function. Then there is an either unitary or antiunitary operator $`U`$ on $`H`$ such that
$$\varphi (A)=UAU^{}(A[0,I])$$
or
$$\varphi (A)=U(IA)U^{}(A[0,I]).$$
Our next result describes the automorphisms of $`[0,I]`$ equipped with the Jordan triple product.
###### Theorem 3.
Suppose that $`dimH3`$. Let $`\varphi :[0,I][0,I]`$ be a bijective function satisfying
$$\varphi (ABA)=\varphi (A)\varphi (B)\varphi (A)(A,B[0,I]).$$
Then $`\varphi `$ is of the form
$$\varphi (A)=UAU^{}(A[0,I]),$$
where $`U`$ is either a unitary or an antiunitary operator on $`H`$.
###### Proof.
First observe that $`\varphi `$ sends projections to projections. Indeed, if $`PB(H)`$ is a projection, then we have $`\varphi (P)=\varphi (P)^3`$. Since $`\varphi (P)`$ is a positive operator, by the spectral mapping theorem we obtain that $`\sigma (\varphi (P))\{0,1\}`$ and this proves that $`\varphi (P)`$ is a projection.
We next show that $`\varphi `$ preserves the partial ordering $``$ among the projections. Let $`P,QB(H)`$ be projections and suppose that $`PQ`$. Then we have $`PQP=P`$ which yields $`\varphi (P)=\varphi (P)\varphi (Q)\varphi (P)`$. This implies that $`\varphi (P)\varphi (Q)`$ is an idempotent. On the other hand, since $`\varphi (P)`$ and $`\varphi (Q)`$ are projections, the norm of their product is not greater than 1. So, $`\varphi (P)\varphi (Q)`$ is a contractive idempotent. It is well-known that this implies that $`\varphi (P)\varphi (Q)`$ is a projection and hence, due to the self-adjointness, it follows that $`\varphi (P)`$ and $`\varphi (Q)`$ are commuting. Hence, we can compute
$$\varphi (P)=\varphi (P)\varphi (Q)\varphi (P)=\varphi (Q)\varphi (P)\varphi (P)=\varphi (Q)\varphi (P)$$
which yields that $`\varphi (P)\varphi (Q)`$. Since $`\varphi ^1`$ has the same properties as $`\varphi `$, it follows that $`\varphi `$ preserves the ordering $``$ in both directions. In particular, we obtain that $`\varphi (0)=0`$, $`\varphi (I)=I`$ and that $`\varphi `$ preserves the nonzero minimal projections, that is, the rank-one projections on $`H`$.
We claim that $`\varphi `$ preserves also the orthocomplementation on the set of projections. To see this, we first show that $`\varphi `$ preserves the mutual orthogonality. Let $`P,QB(H)`$ be projections such that $`PQ=0`$. Then we have $`0=\varphi (PQP)=\varphi (P)\varphi (Q)\varphi (P)`$ which implies that
$$0=\varphi (P)\varphi (Q)\varphi (Q)\varphi (P)=\varphi (P)\varphi (Q)(\varphi (P)\varphi (Q))^{}.$$
This gives us that $`\varphi (P)\varphi (Q)=0`$. It follows that $`\varphi (P)+\varphi (IP)`$ is a projection, say $`\varphi (Q)`$. Since $`\varphi (P),\varphi (IP)\varphi (Q)`$ and $`\varphi `$ preserves the ordering in both directions, we infer that $`P,IPQ`$. This gives us that $`Q=I`$ and, hence, $`\varphi (P)+\varphi (IP)=I`$. Therefore, $`\varphi `$ preserves the orthocomplementation on the set of all projections. The form of such transformations, that is, the form of all bijections of the set of all projections on a Hilbert space with dimension not less than 3 which preserve the order in both directions and the orthocomplementation, is well-known (see, for example, ). Namely, there is an either unitary or antiunitary operator $`U`$ on $`H`$ such that
$$\varphi (P)=UPU^{}$$
for all projections $`P`$ on $`H`$.
We next prove that $`\varphi (\lambda P)=\lambda \varphi (P)`$ for every $`\lambda [0,1]`$ and every rank-one projection $`P`$. In fact, in that case we can compute
$$\varphi (\lambda P)=\varphi (P(\lambda P)P)=\varphi (P)\varphi (\lambda P)\varphi (P)=f_P(\lambda )\varphi (P)$$
for some scalar $`f_P(\lambda )[0,1]`$ which follows from the fact that $`\varphi (P)`$ is of rank one. We assert that $`f_P`$ is a multiplicative function. If $`\mu [0,1]`$, then we have
$$f_P(\lambda ^2\mu )\varphi (P)=\varphi (\lambda ^2\mu P)=\varphi ((\lambda P)(\mu P)(\lambda P))=$$
$$\varphi (\lambda P)\varphi (\mu P)\varphi (\lambda P)=f_P(\lambda )^2f_P(\mu )\varphi (P)$$
which implies that $`f_P(\lambda ^2\mu )=f_P(\lambda )^2f_P(\mu )`$. Choosing $`\mu =1`$, it follows that $`f_P(\lambda ^2)=f_P(\lambda )^2`$. We next obtain that $`f_P(\lambda ^2\mu )=f_P(\lambda ^2)f_P(\mu )`$. Since this holds for every $`\lambda ,\mu [0,1]`$, we conclude that $`f_P`$ is multiplicative. We now claim that $`f_P`$ does not depend on the rank-one projection $`P`$. If $`P,Q`$ are rank-one projections which are not mutually orthogonal, then $`PQP0`$ and we have
$$f_Q(\lambda ^2)\varphi (PQP)=f_Q(\lambda ^2)\varphi (P)\varphi (Q)\varphi (P)=\varphi (P)\varphi (\lambda ^2Q)\varphi (P)=$$
$$\varphi (P(\lambda ^2Q)P)=\varphi ((\lambda P)Q(\lambda P))=$$
$$\varphi (\lambda P)\varphi (Q)\varphi (\lambda P)=f_P(\lambda ^2)\varphi (PQP).$$
This gives us that $`f_P=f_Q`$. If $`P,Q`$ are mutually orthogonal, then there is a rank-one projection $`R`$ such that $`PRP0`$ and $`RQR0`$. Thus we have $`f_P=f_R=f_Q`$. So, there is a multiplicative function $`f:[0,1][0,1]`$ such that
$$\varphi (\lambda P)=f(\lambda )\varphi (P)$$
for every $`\lambda [0,1]`$ and rank-one projection $`P`$ on $`H`$. We show that $`f`$ is also additive on $`[0,1]`$. To see this, for any unit vector $`xH`$ denote by $`P_x`$ the rank-one projection onto the linear subspace of $`H`$ spanned by $`x`$. Let $`x,yH`$ be mutually orthogonal unit vectors and $`\lambda ,\mu [0,1]`$ such that $`\lambda ^2+\mu ^2=1`$. Then $`z=\lambda x+\mu y`$ is a unit vector. We compute $`\varphi (P_z(P_x+P_y)P_z)`$ in two different ways. On the one hand, since $`P_z(P_x+P_y)P_z=P_z`$, we have $`\varphi (P_z(P_x+P_y)P_z)=\varphi (P_z)`$. On the other hand, we compute
$$\varphi (P_z(P_x+P_y)P_z)=\varphi (P_z)\varphi (P_x+P_y)\varphi (P_z)=\varphi (P_z)(\varphi (P_x)+\varphi (P_y))\varphi (P_z)=$$
$$\varphi (P_z)\varphi (P_x)\varphi (P_z)+\varphi (P_z)\varphi (P_y)\varphi (P_z)=\varphi (P_zP_xP_z)+\varphi (P_zP_yP_z)=$$
$$\varphi (\lambda ^2P_z)+\varphi (\mu ^2P_z)=(f(\lambda ^2)+f(\mu ^2))\varphi (P_z)$$
where we have used the fact that $`\varphi `$ is orthoadditive on the set of all projections (this follows from the form of $`\varphi `$ on that set). Therefore, we have $`f(\lambda ^2)+f(\mu ^2)=1=f(\lambda ^2+\mu ^2)`$. By multiplicativity, we obtain the additivity of $`f`$. We claim that $`f`$ is in fact the identity on $`[0,1]`$. Since $`f`$ maps into $`[0,1]`$, one can easily check that $`f`$ is monotone increasing. Moreover, as $`f(1)=1`$, the additivity of $`f`$ implies that $`f(r)=r`$ for every rational number $`r`$ in $`[0,1]`$. If $`\lambda ]0,1[`$ is arbitrary, then approximating $`\lambda `$ by rationals $`r,s`$ from below and above, respectively, by the monotonity we can infer that $`f(\lambda )=\lambda `$.
We already know the form of $`\varphi `$ on the set of all projections. It is easy to see that without loss of generality we can assume that $`\varphi (P)=P`$ holds for every projection $`P`$ and we then have to prove that $`\varphi `$ is the identity on the whole interval $`[0,I]`$. But this is now easy. Indeed, let $`A[0,I]`$. Pick an arbitary rank-one projection $`P=P_x`$, where $`xH`$ is a unit vector. Then we compute
$$P\varphi (A)P=\varphi (PAP)=\varphi (Ax,xP)=Ax,x\varphi (P)=Ax,xP=PAP.$$
Since $`P`$ was arbitrary, we obtain $`\varphi (A)=A`$ for every $`A[0,I]`$. This completes the proof of the theorem. ∎
Since the Jordan algebra $`B_s(H)`$ of all self-adjoint operators also plays very important role in the mathematical foundations of quantum mechanics, we were tempted to determine the automorphisms of the set $`B_s(H)`$ equipped with the Jordan triple product. Observe that the following theorem has the interesting corollary that every such automorphism is automatically linear, so one can say that the linear structure of $`B_s(H)`$ is completely determined by its multiplicative Jordan triple structure. We remark that the question when a multiplicative function is necessarily additive was investigated for associative rings (recall that our structure is highly nonassociative) in the purely algebraic setting (see ).
###### Theorem 4.
Suppose that $`dimH3`$. Let $`\varphi :B_s(H)B_s(H)`$ be a bijective function (linearity is not assumed) satisfying
$$\varphi (ABA)=\varphi (A)\varphi (B)\varphi (A)(A,BB_s(H)).$$
Then there is an either unitary or antiunitary operator $`U`$ on $`H`$ such that either
$$\varphi (A)=UAU^{}(AB_s(H))$$
or
$$\varphi (A)=UAU^{}(AB_s(H)).$$
###### Proof.
We first prove that $`\varphi (I)`$ is either $`I`$ or $`I`$. Since
$$\varphi (I)\varphi (A)\varphi (I)=\varphi (A)$$
for every $`AB_s(H)`$ and $`\varphi `$ is surjective, it follows that $`\varphi (I)^2=I`$. Therefore, we have
$$\varphi (I)\varphi (A)=\varphi (I)\varphi (A)\varphi (I)\varphi (I)=\varphi (A)\varphi (I).$$
Since this holds for every $`AB_s(H)`$, by the surjectivity of $`\varphi `$, it follows that $`\varphi (I)`$ is in the center of $`B(H)`$ and, consequently, $`\varphi (I)`$ is a scalar. This yields that either $`\varphi (I)=I`$ or $`\varphi (I)=I`$. Clearly, the function $`\varphi `$ is a bijective mapping of $`B_s(H)`$ satisfying the equation appearing in the statement. So, without loss of generality we can assume that $`\varphi (I)=I`$.
We prove that $`\varphi `$ sends projections to projections. If $`P`$ is a projection, then $`\varphi (P)`$ is self-adjoint and we have $`\varphi (P)^2=\varphi (P)\varphi (I)\varphi (P)=\varphi (PIP)=\varphi (P)`$ which shows that $`\varphi (P)`$ is an idempotent.
Now, we can follow the argument in the proof of Theorem 3. One can verify that $`\varphi `$ preserves the partial ordering $``$ in both directions and the orthocomplementation on the set of all projections. So, we have an either unitary or antiunitary operator $`U`$ on $`H`$ such that
$$\varphi (P)=UPU^{}$$
for every projection $`P`$ on $`H`$. One can check that for every rank-one projection $`P`$ there exists a function $`f_P:`$ such that $`\varphi (\lambda P)=f_P(\lambda )\varphi (P)`$ $`(\lambda )`$. We next obtain that $`f_P(\lambda ^2\mu )=f_P(\lambda )^2f_P(\mu )`$ and choosing $`\mu =1`$, this gives us that $`f_P(\lambda ^2)=f_P(\lambda )^2`$. In particular, by the injectivity of $`f_P`$, from
$$f_P(\lambda )^2=f_P((\lambda )^2)=f_P(\lambda ^2)=f_P(\lambda )^2$$
we deduce that $`f_P(\lambda )=f_P(\lambda )`$. Since $`f_P(\lambda ^2\mu )=f_P(\lambda ^2)f_P(\mu )`$, we get that $`f_P`$ is multiplicative. One can next show that $`f_P`$ does not depend on the rank-one projection $`P`$. So, there is a multiplicative function $`f:`$ such that
$$\varphi (\lambda P)=f(\lambda )\varphi (P)$$
for every real number $`\lambda `$ and rank-one projection $`P`$ on $`H`$. As for the additivity of $`f`$, just as in the proof of our previous theorem we get that
$$f(t)=f(t\lambda ^2)+f(t\mu ^2)$$
for every real $`t`$, where $`\lambda ^2+\mu ^2=1`$. To show that $`f`$ is additive, it is enough to verify that $`f(1)=f(t)+f(1t)`$ for every real $`t`$. If $`t[0,1]`$, then we already know this. If $`t[0,1]`$, say $`t<0`$, then we can refer to the equality
$$f(1t)=f\left((1t)\frac{1}{1t}\right)+f\left((1t)\frac{t}{1t}\right)$$
what is known to be valid since the numbers $`\frac{1}{1t},\frac{t}{1t}`$ belong to $`[0,1]`$ and their sum is 1. So, we have
$$f(1t)=f(1)+f(t)=f(1)f(t).$$
Consequently, we obtain that $`f:`$ is an injective multiplicative and additive function. This means that $`f`$ is a nontrivial ring endomorphism of $``$. It is well-known that $`f`$ is necessarily the identity (anyway, this can be proved quite similarly to the corresponding part of the proof of Theorem 3). Finally, one can complete the proof of the statement just as in the case of our previous theorem. ∎
|
warning/0005/cond-mat0005185.html
|
ar5iv
|
text
|
# Disorder from disorder in a strongly frustrated transverse field Ising chain
## I Introduction
The study of quantum versions of classically frustrated magnets has been a subject of interest at least since the work of Anderson and Fazekas on the possibility of a quantum disordered state for the triangular lattice Heisenberg antiferromagnet. Following their work, and especially after the early suggestion of Anderson that the cuprate superconductors derive their special properties from the proximity of a spin liquid state, there has been a considerable amount of work on quantum Heisenberg models on various frustrated lattices such as the kagomé.
One can think of quantum Heisenberg models as (classical) Ising models perturbed by a transverse (XY) exchange—indeed, this was the strategy followed by Anderson and Fazekas in their analysis of the triangular lattice problem. Phrased in this fashion, the problem becomes one of the effects of introducing a quantum dynamics into a highly degenerate ground state manifold of the Ising system—a procedure of evident interest on account of the singular effects of any perturbation. The simplest instance of this more general problem is the introduction of a transverse magnetic field, which has been used in several contexts previously to generate the simplest quantum statistical mechanics; for a review see Ref. which also reviews some work on some one-dimensional systems with competing interactions.
A number of geometrically frustrated quantum Ising systems have been studied recently by Moessner, Chandra and one of the present authors who have reported instances of “order by disorder” in which a non-trivial ordering pattern is selected by the quantum fluctuations as well as one of “disorder by disorder” (on the kagomé lattice) in which the Anderson-Fazekas scenario of a disordered quantum state constructed out of a disordered classical manifold is realized. In this paper, we report studies of a one-dimensional, geometrically frustrated chain which also exhibits disorder by disorder. This is interesting in itself; additionally, through a technical advance, we were able to provide strong evidence for the quantum disordered state by means of a strong coupling perturbation expansion, which we anticipate will be a useful technique in studying higher-dimensional frustrated quantum Ising systems.
## II The model and its possible ordering
We consider the Hamiltonian
$$H=J\underset{ij}{}\sigma _i^z\sigma _j^z+\mathrm{\Gamma }\underset{i}{}\sigma _i^x$$
(1)
where $`J>0`$ is the antiferromagnetic Ising exchange, the sum in the first term runs over the bonds of the triangular chain in Figure 1, $`\mathrm{\Gamma }`$ is the strength of the transverse field and the $`\sigma ^a`$ are the Pauli operators for $`S=1/2`$. (As the true spin operators are $`S^a=\sigma ^a/2`$, when $`\mathrm{}=1`$, our $`J`$ is really one quarter the exchange and our $`\mathrm{\Gamma }`$ is half the physical transverse field.)
The ground states of the system without the transverse field can be obtained by minimizing the energy of each triangle separately and this in turn requires that we choose the one bond in each that will remain unsatisfied (the “bad” bond). The choice of the bad bonds is totally independent from triangle to triangle (the system is a bad bond paramagnet); a typical state is exhibited in Fig 2. With free boundary conditions, the number of such states is thus $`2\times 3^N`$ where $`N`$ is the number of triangles and the extra factor of $`2`$ accounts for the two Ising reversed states that give rise to the same bond configuration. This extensive entropy is accompanied by short-ranged correlations and we find that the spin-spin correlation function averaged over the ground state manifold takes the form
$`S_{\alpha t}^zS_{\beta t}^z=S_{\alpha b}^zS_{\beta t}^z`$ $`=`$ $`(1/3)^{\beta \alpha +1}`$ (2)
$`S_{\alpha b}^zS_{\beta b}^z=S_{\alpha t}^zS_{\beta b}^z`$ $`=`$ $`(1/3)^{\beta \alpha }`$ (3)
where $`\alpha t`$ and $`\alpha b`$ denote the top and bottom spins in unit cell $`\alpha `$. Due to the lack of inversion symmetry in our choice of unit cell, these forms hold for $`\beta >\alpha `$; they can be extended to cover $`\beta \alpha `$ by inspecting the lattice and noting its symmetries. Evidently, the system is disordered at all temperatures.
We turn now to the quantum problem posed by the inclusion of the transverse field.
### A Analysis at $`\mathrm{\Gamma }J`$
At $`T=0`$, an infinitesimally small $`\mathrm{\Gamma }`$ is a singular perturbation which will lift the macroscopic degeneracy of the classical problem. Following Ref. , we will attempt to identify the physics of this regime variationally. To this end we identify the configuration (unique, up to global symmetry operations) that maximizes the number of “flippable” spins, which is shown in Figure 3. A flippable spin in a given configuration is one which can be reversed without violating the ground state constraint. As a flippable spin can be polarized along $`\widehat{x}`$ at no cost in exchange, we expect a state that maximizes their number to be especially favored by the transverse field. As not all flippable spins are independently flippable, we have two options. We may construct the staggered state in which we polarize the maximal set of independently flippable spins along the top row,
$$|s=(_{\mathrm{top}}|x)|\mathrm{}_{\mathrm{bottom}};$$
(4)
here $`|x`$ is the state with a spin pointing down along the x axis. Alternatively we may construct the uniform state in which we polarize all the flippable spins but correct for the conflicts by projecting onto the ground state manifold by the action of the projector $`P`$ which eliminates those configurations that do not have exactly one bad bond per triangle:
$`|u`$ $`=`$ $`P\{(_{\mathrm{top}}|x)`$ (6)
$`(||x||x|\mathrm{})_{\mathrm{bottom}}\}.`$
The uniform state inherits all the symmetries of the maximally flippable configuration but the staggered state does not. For the wave functions as written, we can readily evaluate the energies and we find that the staggered state has the same energy $`(\mathrm{\Gamma }+J)/2`$ per site as the uniform state. We can also write down the magnetization (up to translations and Ising reversal) in the staggered state:
$$\sigma ^z(\alpha )=m(\alpha )=\{\begin{array}{cc}0\hfill & \mathrm{top}\hfill \\ (\frac{1}{2})^\alpha \hfill & \mathrm{bottom}\hfill \end{array}$$
(7)
and the uniform state
$$m(\alpha )=\{\begin{array}{cc}\frac{1}{10}\hfill & \mathrm{top}\hfill \\ \frac{1}{10}+(\frac{2}{5})^\alpha \hfill & \mathrm{bottom}.\hfill \end{array}$$
(8)
Needless to say, this analysis being purely variational is only indicative of what kind of “order from disorder” the system might exhibit at $`\mathrm{\Gamma }J`$. At issue is whether the connectivity of the set of configurations is too high (e.g. as reflected in the disorder in the classical problem) to permit a localization of the ground state wavefunction near the maximally flippable configuration. As advertised we will, instead, find that the system remains disordered.
### B Analysis at $`\mathrm{\Gamma }J`$
One of the attractive features of transverse field Ising models, emphasized in Ref. , is that they exhibit a gapped, paramagnetic phase at $`\mathrm{\Gamma }J`$. This enables a natural perturbative expansion for the problem, in which we perturb in the Ising exchange about the purely transverse field paramagnet. A second attractive feature, exhibited by the well known Trotter-Suzuki procedure is that the quantum partition function (with the transverse field present) has the form of the classical partition function of copies of the chain (now without the transverse field) coupled ferromagnetically in the imaginary time direction which can then be analyzed from its “high temperature” phase via the standard Landau-Ginzburg-Wilson (LGW) procedure. It is not difficult to show that these are attempts to search for the same instability. We will pursue the perturbative expansion a bit later; here, we will attempt to identify the leading instability from the paramagnetic phase by the LGW step of diagonalizing the bond (adjacency) matrix. (The equivalence to the first step in the perturbation expansion will be evident when we come to it.)
As the ferromagnetic interaction in the imaginary time direction is trivially incorporated by including a quadratic dispersion about $`k_\tau =0`$, we focus on the problem in the chains. To diagonalize the quadratic form,
$$=\underset{ij}{}S_iS_j$$
(9)
we introduce Fourier variables with respect to the unit cell location,
$`S_{\alpha t}`$ $`=`$ $`{\displaystyle \underset{k}{}}e^{ik\alpha }S_{kt}`$ (10)
$`S_{\alpha b}`$ $`=`$ $`{\displaystyle \underset{k}{}}e^{ik\alpha }S_{kb}.`$ (11)
and find that the Hamiltonian takes the form
$``$ $`=`$ $`{\displaystyle \underset{k>0}{}}\left[\begin{array}{cc}S_{kt}^{}\hfill & \hfill S_{kb}^{}\end{array}\right]\left[\begin{array}{cc}0\hfill & \hfill (1+e^{ik})\\ (1+e^{ik})\hfill & \hfill 2\mathrm{cos}k\end{array}\right]\left[\begin{array}{c}S_{kt}\\ S_{kb}\end{array}\right].`$ (17)
Further diagonalization yields the eigenvalues
$$\lambda _\pm =(\epsilon 1)\pm \sqrt{\epsilon ^2+1}$$
(18)
where we have defined $`\epsilon =1+\mathrm{cos}k`$. Evidently, $`\lambda _{}`$ is minimized when one has $`\epsilon =0k=\pi `$. The corresponding eigenvector is
$$\stackrel{}{v}_{\mathrm{min}}=\left[\begin{array}{c}0\\ 1\end{array}\right]$$
(19)
Hence, the spatial dependence of the lowest eigenmode, which is the candidate for ordering, is
$$\varphi _\pi =\mathrm{cos}\pi \alpha \left[\begin{array}{c}0\\ 1\end{array}\right]=(1)^\alpha \left[\begin{array}{c}0\\ 1\end{array}\right].$$
(20)
Thus, as $`\mathrm{\Gamma }/J`$ is lowered—which corresponds to lowering the temperature in the classical representation—the LGW analysis suggests a transition in the $`d=2`$ Ising universality class to a state with the top row disordered and bottom row ordered antiferromagnetically, which is the staggered state that we had constructed earlier.
### C Mean Field Theory
As we are interested in the suppression of ordering by fluctuations due to the frustration, it is useful to have a mean-field estimate for the location of the transition. To this end we consider the possibility of ordering into the staggered state. In this case only the spins on the bottom row see an effective field and the problem (in mean field theory) is identical to that of the purely one-dimensional transverse field Ising chain. This yields a critical coupling
$$x_c=(J/\mathrm{\Gamma })_c=1/2$$
(21)
and a staggered magnetization
$$m=\sqrt{1(2x)^2}$$
(22)
in the ordered phase. The magnetization of the top row is always zero. Better estimates can be obtained by treating the top row in mean field theory and using the known exact results for the bottom row. In this fashion we obtain $`x_c=1`$ and $`m(xx_c)^{1/8}`$.
## III Exact diagonalization
### A Modified Lanczos method
We now turn to our numerical studies of the model, beginning with exact diagonalization of the Hamiltonian for finite lattices with periodic boundary conditions. Exact diagonalization unfortunately imposes a demand on computer memory which rises exponentially with the size of the system considered. Consequently, the systems discussed in this report are limited to at most 8 unit cells (16 spins).
In the modified Lanczos method that we apply, one begins with a state $`|0`$ expected to have nonzero overlap with the ground state of $``$. For the present problem we used the trivial (fully $`x`$-polarized) $`J=0`$ ground state, which is a $`q=0`$ state. One generates a state orthogonal to $`|0`$ by acting upon $`|0`$ with $``$ and using Gram-Schmidt orthogonalization to generate the state
$$|1=|0\frac{0||0}{0|0}$$
(23)
One then diagonalizes $``$ in the resulting two-dimensional subspace. The lowest eigenvalue is the improved estimate for the ground state energy, and the corresponding eigenvector is the improved estimate for the ground state wave function. This process is iterated to convergence.
A slight modification of the procedure used in obtaining the ground state energy and wave function can be employed to find the energy gap and first excited state wave function. One begins with an initial guess with nonzero overlap with the lowest excited state wave function and orthogonal to the ground state, and then proceeds as described above for the ground state. The initial guess for the excited state was chosen to be one of the excited states in the $`J=0`$ system, with a single flipped spin at a particular location. Although the state is not one of definite momentum, many iterations of the modified Lanczos procedure converge to a state of definite momentum in which $`k=\pi `$, as expected from the analysis in Section II. Although successive excited state estimates should in principle be orthogonal to the ground state, rounding errors cause an admixture with the ground state to occur at each iteration. To prevent a convergence toward the ground state, Gram-Schmidt orthogonalization is used at each step to remove any component of the ground state which may have entered. Although one does not have the exact ground state at a particular iteration, convergence can still be achieved by performing Gram-Schmidt orthogonalization with the most recent estimate for the ground state wave function.
### B Excitation gap
Figure 4 displays the gap to the first excited state computed for a variety of system sizes ($`N`$ denoting the number of unit cells) and values of $`x=J/\mathrm{\Gamma }`$. For convenience, the domain $`[0,\mathrm{}]`$ of $`x`$ is compressed into the interval $`[0,1]`$ via the transformation $`y=x/(x+1)`$. For most values of $`y`$, convergence of the gap with increasing system size is quite rapid, in itself an indication of the presence of a short correlation length. In fact, for $`y<2/3`$, the gaps for systems with $`6`$ and $`8`$ unit cells are virtually indistinguishable. In that regime, it is clear that the system remains disordered (in the same phase as the pure transverse-field model $`y=0`$). For larger $`y`$ the situation is less clear-cut, so we next consider a more sensitive test for the existence of a continuous phase transition. (We note that our largest system size yields a ground state energy of $`\frac{1}{2}\left[J+(1.11)\mathrm{\Gamma }\right]`$ per spin at $`\mathrm{\Gamma }J`$, 11% lower than the energy for the variational states considered earlier.)
### C Cumulants
Following Binder, we examine the cumulants
$$r=\frac{|\mathrm{\Psi }|^4}{|\mathrm{\Psi }|^2^2}$$
(24)
In equation 24, $`\mathrm{\Psi }`$ is the order parameter of interest. Since any ordering is expected to be accompanied by antiferromagnetic ordering of the bottom row, a sensible choice of $`\mathrm{\Psi }`$ seems to be the staggered magnetic moment operator,
$$M_s=\underset{\alpha }{}(1)^\alpha \sigma _{\alpha b}^z$$
(25)
on the bottom row.
It is expected that for systems of size $`L`$, in the vicinity of a critical point, $`r`$ obeys the scaling form
$$r=f((xx_c)L^{1/\nu })$$
(26)
If a critical point exists, the cumulants computed for various system sizes should intersect at the same point for $`x=x_c`$.
Cumulants for systems containing $`4`$, $`6`$, and $`8`$ unit cells are displayed in Figure 5. There is apparently no intersection of the cumulants for different system sizes for $`x20`$ (and, by extrapolation, for even larger values of $`x`$). The behavior of the cumulants is consistent with the absence of a critical point and therefore suggests that the system is disordered for all values of $`J/\mathrm{\Gamma }`$. We note that in a disordered phase with a Gaussian distributed order parameter, the cumulant will approach the value 3. At $`x=0`$ the cumulant can be shown analytically to have the value $`34/N`$ for $`N`$ unit cells; evidently the finite size effects are greater at large $`x`$, where the ground state constraint produces correlations, but even there the steady growth with system size is consistent with a thermodynamic limit value of 3.
## IV Perturbation expansion generation and analysis
### A Cluster expansions for the elementary excitation spectra
A perturbation expansion for the energy gap (and the full elementary excitation spectrum, incidentally) about the $`J=0`$ limit was constructed to 14th order in $`J/\mathrm{\Gamma }`$ using a cluster expansion technique. The ideas underlying this calculation are set out briefly in Ref. , and described in more detail in Ref. . However, this particular calculation is unusual (but not unique, see Müller-Hartmann et al.) in one interesting respect: it is carried out for a system with more than one branch of elementary excitations which are degenerate in the unperturbed limit, but are not degenerate by symmetry and hence are not degenerate for $`J>0`$. That such calculations can be done by cluster expansion methods was mentioned briefly in Ref. . Here we supply some further details which are important in making the connection between the the immediate product of the cluster expansion, which is a matrix containing power expansions that describe the motion of elementary excitations in real space, and the desired end product, the power series expansions for the excitation spectra in wave vector space.
For the purposes of this section we assign the sites at the base of the triangles the coordinates $`(n,0)`$, and those at the tips of the triangles $`(n+\frac{1}{2},1)`$, with $`n`$ running over the integers. The unperturbed Hamiltonian is the transverse field $`_n[\sigma _{(n,0)}^x+\sigma _{(n+\frac{1}{2},1)}^x]`$ and the perturbation is the Ising exchange $`_n[\sigma _{(n,0)}^z\sigma _{(n+1,0)}^z+\sigma _{(n,0)}^z\sigma _{(n+\frac{1}{2},1)}^z+\sigma _{(n+1,0)}^z\sigma _{(n+\frac{1}{2},1)}^z]`$. The perturbation expansion is a power series in $`x=J/\mathrm{\Gamma }`$.
For the unperturbed system, the ground state consists of all spins “down” (pointing along $`x`$) and the elementary excitations are single spins which have been flipped “up”. Note that there are two degenerate branches of such excitations, flat bands at $`\omega =2`$. (One could consider an unperturbed Hamiltonian of the more general form $`_n\left[\sigma _{(n,0)}^z+\zeta \sigma _{(n,1)}^z\right]`$, which has the same symmetry as the triangular chain but lacks the degeneracy. For $`\zeta 1`$, then, perturbations expansions for each of the branches of excitations could be constructed separately; however, there would be energy denominators $`(1\zeta )`$ in abundance that would lead to poorly converging series for $`\zeta `$ near 1.)
In our calculations, with $`\zeta =1`$, we treat all of the elementary excitations on equal footing. We keep track of three distinct types of matrix elements in the effective Hamiltonian, those that connect two base sites (type $`A`$), those that connect two tip sites (type $`B`$), and those that connect base to tip sites (type $`C`$); let us call $`t_d^X`$ the matrix element (which is a series expansion in $`\lambda `$) in the effective Hamiltonian of type $`X`$ that couples sites at horizontal distance $`d`$. The effective Hamiltonian has the following structure near the diagonal
$$\left(\begin{array}{cccccccc}\mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \\ \mathrm{}& t_0^A& t_{1/2}^C& t_1^A& t_{3/2}^C& t_2^A& t_{5/2}^C& \mathrm{}\\ \mathrm{}& t_{1/2}^C& t_0^B& t_{1/2}^C& t_1^B& t_{3/2}^C& t_2^B& \mathrm{}\\ \mathrm{}& t_1^A& t_{1/2}^C& t_0^A& t_{1/2}^C& t_1^A& t_{3/2}^C& \mathrm{}\\ \mathrm{}& t_{3/2}^C& t_1^B& t_{1/2}^C& t_0^B& t_{1/2}^C& t_1^B& \mathrm{}\\ \mathrm{}& t_2^A& t_{3/2}^C& t_1^A& t_{1/2}^C& t_0^A& t_{1/2}^C& \mathrm{}\\ \mathrm{}& t_{5/2}^C& t_2^B& t_{3/2}^C& t_1^B& t_{1/2}^C& t_0^B& \mathrm{}\\ & \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right).$$
(27)
This matrix acts on a column-vector describing the spin-flip amplitude at each site, listed in order of its $`x`$ coordinate. This matrix is readily diagonalized by plane waves with a two-site basis. One obtains the two branches of the excitation spectrum in the form
$$ϵ_\pm =\frac{1}{2}\left(F_A+F_B\pm \sqrt{(F_AF_B)^2+F_C^2}\right)$$
(28)
with
$$F_A(q)=\underset{n0}{}\stackrel{~}{t}_n^A\mathrm{cos}nq,$$
(29)
likewise for $`F_B(q)`$, and
$$F_C(q)=\underset{n0}{}\stackrel{~}{t}_{(2n+1)/2}^C\mathrm{cos}(2n+1)q/2.$$
(30)
The relationship between the $`t`$s and $`\stackrel{~}{t}`$s is as follows. Take any two adjacent rows of the effective Hamiltonian, and count how many $`t_d^X`$ there are for a given $`X`$ and $`d`$: that number is the ratio $`\stackrel{~}{t}_d^X/t_d^X`$. Thus $`\stackrel{~}{t}_0^A=t_0^A`$, $`\stackrel{~}{t}_{1/2}^C=4t_{1/2}^C`$, and so forth. It is the $`\stackrel{~}{t}`$s that come directly out of our computer programs and so it is convenient to express the functions that appear in Eq. (28) in those terms.
We should note that our 14th order calculation of the excitation spectrum involved weight calculations for 1355 graphs of up to 15 spins. We did not attempt to classify the graphs topologically. The implementation of the weight calculations traded off considerable efficiency for safety (guarding against a variety of possible coding errors), and the calculations were carried out on a modest machine (333 MHz Pentium II with 256 megabytes of RAM), so it should be feasible to evaluate several more terms of the perturbation expansion if desired.
### B Analysis of the energy gap series
One is faced immediately with an interesting choice in analyzing the series for the gap, $`ϵ_{}`$. One can either extrapolate the series for $`F_A`$, $`F_B`$ and $`F_C`$ individually and insert the results into the formula (28), or evaluate that formula for the series and extrapolate $`ϵ_{}`$ directly.
However, the lowest excited state is always found at $`q=\pi `$, and in this case the need to choose vanishes because $`F_C(\pi )0`$ and thus $`ϵ_{}(\pi )=\mathrm{min}(F_A(\pi ),F_B(\pi ))=F_A(\pi )`$. In Table I we display $`F_A(\pi )`$ and $`F_B(\pi )`$. The complete set of $`\stackrel{~}{t}_D^X`$ are available from the authors on request.
To analyze the gap series, we applied the transformation $`y=x/(x+1)`$, and constructed Padé approximants to the transformed series. We use the notation $`P_M^{}^M`$ to denote an approximant with a polynomial of order $`M`$ in the numerator and order $`M^{}`$ in the denominator. Figure 6 displays several approximants which use all the terms of the series (and hence $`M+M^{}=14`$).
A salient feature of Figure 6 is the consistency of the Padé approximants over a large fraction of the domain. The approximants displayed in the plot are consistent to within $`1\%`$ for $`y`$ as large as $`0.75`$. Note that this corresponds to a value of $`x=3`$, well in excess of the mean field estimate of $`x_c=1/2`$ and the $`d=1`$ transverse field Ising model estimate of $`x_c=1`$. Even without further analysis these indicate the disordering effects of the strong frustration in the problem.
Figure 6 also displays the excitation gap as computed via the modified Lanczos algorithm. Evident in the graph is the excellent agreement of the approximants and the exact diagonalization results over the region of consistency among the Padé approximants.
#### 1 Two Point Approximants and The Global Phase Portrait
Although the region of validity of the analyzed gap series extends well into the domain of strong exchange coupling, the Padé approximants plotted in Figure 6 tell an ambiguous story. In particular, it is not clear from the displayed approximants whether or not the excitation gap vanishes for large $`J/\mathrm{\Gamma }`$ and hence it is unclear whether the triangular chain is ordered in the $`J\mathrm{\Gamma }`$ limit. The reader should note though that even those Padé approximants that indicate a vanishing of the gap, do so for values of $`x`$ in excess of 11.
Biasing of the Padé approximants is one means of extending the results of the gap series analysis to larger $`y`$ values. A simple way to bias the approximants is by means of “two-point” approximants. In particular, one would hope that by biasing the value of the gap at $`y=1`$ the resulting approximants would be accurate over the entire range $`0y1`$. However, since it is precisely the value of the gap at $`y=1`$ (let us denote that $`\mathrm{\Delta }(1)`$) that we know least well, such a procedure appears to beg the question. What we have done, then, is to construct two-point approximants (with various orders of numerators and denominator) for a range of $`\mathrm{\Delta }(1)`$ values, and observed behavior. Two of the biased approximants, $`P_9^5`$ and $`P_8^6`$, show a high degree of consistency with each other for a relatively small range of $`\mathrm{\Delta }(1)`$ values. In fact, the range over which the maximum discrepancy between the two is less than $`1\%`$ is confined between $`\mathrm{\Delta }(1)=.15`$ and $`\mathrm{\Delta }(1)=.2`$. For values of $`\mathrm{\Delta }(1)`$ between $`.15`$ and $`.20`$, $`P_9^5`$ and $`P_8^6`$ are in good enough agreement that the maximum disagreement between the two over the entire $`y`$ domain falls below $`1\%`$. Outside the $`\mathrm{\Delta }(1)`$ range indicated above, the maximum disagreement between $`P_9^5`$ and $`P_8^6`$ exceeds $`1\%`$. Figure 7 displays the biased approximants corresponding to the limits of $`\mathrm{\Delta }(1)`$ discussed above. No matter what the choice of $`\mathrm{\Delta }(1)`$, none of the other two point approximants are in agreement with each other or with $`P_9^5`$ and $`P_8^6`$ to an extent which approaches the level of agreement of $`P_9^5`$ and $`P_8^6`$ for $`\mathrm{\Delta }(1)`$ values between $`.15`$ and .$`.20`$.
The biased $`P_9^5`$ approximants displayed in the figure are in reasonable, though not perfect agreement with the Lanzcos gap values. The approximants shown are biased at $`\mathrm{\Delta }(1)=.15`$ and $`\mathrm{\Delta }(1)=.2`$, the bounds of the region in which $`P_9^5`$ and $`P_8^6`$ satisfy the $`1\%`$ consistency criterion discussed above. The agreement is best for small values of $`y`$ but deteriorates as $`y`$ approaches 1. However, even at $`y=1`$, averaging the two approximants yields a prediction for $`\mathrm{\Delta }(1)`$ of $`.175`$, a value which differs from the Lanczos prediction of $`.132`$ (albeit without any attempt at scaling with system size) by only $`30\%`$. The approximants echo the exact diagonalization results in indicating the absence of order for all values of $`J/\mathrm{\Gamma }`$.
## V Summary
Results for the energy gap and cumulants from exact diagonalization of finite systems indicate that the transverse field, triangular chain Ising model is disordered, even though several analytic approaches (from both the small and large transverse-field limits) suggest that a particular ordered phase could exist for small transverse field. We note that the lowest lying excited state is consistent with the ordering analysis, indicating that the fluctuations generated by the ground state constraint at small $`\mathrm{\Gamma }/J`$ are too strong to allow order to set in. Hence we get “disorder by disorder” instead of “order by disorder”. Estimates of the gap obtained by direct Padé approximants to 14th order perturbation expansions about the strong transverse field limit produce excellent results over a modest range of the Ising exchange coupling to transverse field ratio that exceeds the natural estimates for a critical value by a factor of 3 and do not indicate ordering for values of the ratio as big as 11. Globally reasonable results have been obtained by biasing the approximants and applying a consistency criterion and these indicate a lack of ordering at any value of $`J/\mathrm{\Gamma }`$. This suggests that an application of this method to systems of greater experimental interest, such as the higher dimensional kagomé and pyrochlore lattices, may be a viable approach to deducing the global phase behavior of such systems. We should note that in $`d>1`$ we expect the series method to be increasingly competitive with exact diagonalization as the latter technique has to contend with much stronger finite size effects at computationally feasible system sizes.
###### Acknowledgements.
We would like to thank Roderich Moessner for collaboration in the initial stages of the project and for much useful advice in its course. We would also like to thank him for a careful reading of the manuscript. This work was supported by the National Science Foundation through grants DMR 94-57928 (MPG) and DMR 99-78074 (DJP and SLS) as well as by the Alfred P. Sloan Foundation and the David and Lucille Packard Foundation (SLS).
|
warning/0005/gr-qc0005116.html
|
ar5iv
|
text
|
# Dynamics of spatially homogeneous locally rotationally symmetric solutions of the Einstein-Vlasov equations
## 1 Introduction
The most popular matter content by far in the study of spatially homogeneous cosmological models is a perfect fluid with linear equation of state (see e.g., the book ). It is important to know if the results obtained for this class are structurally stable if we change the matter content. Thus it is of interest to investigate other types of sources. Here we will consider certain diagonalizable locally rotationally symmetric (LRS) spatially homogeneous models with collisionless matter. This class of models was previously studied in the case of massless particles in . Here we will focus on the case with massive particles. We will recast Einstein’s field equations into a form so that one part of the boundary of the state space for the massive case can be identified with the state space for the massless case while another part can be identified with the state space for the corresponding dust equations. (In addition other parts of the boundary have the interpretation of state spaces associated with certain models with distributional matter.) It will be shown that these boundary submanifolds are intimately connected with the early and late time behaviour of the LRS massive collisionless gas models respectively.
The results of our analysis can be summarized as follows. Consideration is restricted to models of Bianchi types I, II and III. This is enough to display a large variety of phenomena. At early times, i.e. close to the initial singularity, the dynamics of solutions with massive particles mimics closely the dynamics for the corresponding symmetry type with massless particles. In particular there are solutions whose behaviour near the singularity is quite different from that of any fluid model of any of these Bianchi types. At late times, i.e. in a phase of unlimited expansion, the general picture is that the dynamics resembles that of a dust model. This is proved for Bianchi types I and II. For type III the results are consistent with dust-like asymptotics but we were not able to prove that this is what happens. If kinetic theory with massive particles always behaved like dust at late times this would provide a justification of the use of a fluid model in that regime.
The outline of the paper is as follows. In section 2 we derive the dynamical system. Sections 3, 4 and 5 analyse the models of types I, II and III respectively, with the main results being stated in Theorems 2.1, 3.1 and 4.1. In section 6 we conclude with some remarks and speculations. An appendix contains some information about dynamical systems which is applied frequently in the paper.
## 2 A dynamical systems formulation
We will consider LRS models for which the metric can be written in the form
$$ds^2=dt^2+g_{11}(t)(\theta ^1)^2+g_{22}(t)((\theta ^2)^2+(\theta ^3)^2),$$
(1)
where $`\theta ^i`$ are suitable one-forms describing the various symmetry types. The energy-momentum tensor $`T_{ij}`$ for the Einstein-Vlasov system with massive particles is assumed to be diagonal and is described by
$`\rho `$ $`=`$ $`{\displaystyle f_0(v_i)(m^2+g^{11}(v_1)^2+g^{22}((v_2)^2+(v_3)^2))^{1/2}(detg)^{1/2}𝑑v_1𝑑v_2𝑑v_3},`$
$`p_i`$ $`=`$ $`{\displaystyle f_0(v_i)g^{ii}(v_i)^2(m^2+g^{11}(v_1)^2+g^{22}((v_2)^2+(v_3)^2))^{1/2}(detg)^{1/2}𝑑v_1𝑑v_2𝑑v_3},`$ (2)
where $`\rho `$ is the energy density and $`p_i=T^i_i`$ the pressure components of the energy-momentum tensor. The function $`f_0`$ is determined at some fixed time $`t_0`$ by $`f_0(v_i)=f(t_0,v_i)`$ where $`f`$ is the phase space density of particles. The covariant components $`v_i`$ are independent of time. The function $`f_0`$ satisfies the condition $`f_0(v_1,v_2,v_3)=F(v_1,(v_2)^2+(v_3)^2)`$.
Some further technical conditions will be imposed on $`f_0`$. It is assumed to be non-negative and have compact support. It is also assumed that the support does not intersect the coordinate planes $`v_i=0`$. A function $`f_0`$ with this property will be said to have split support. The reason for the assumption of split support will be seen later. In the following it will always be assumed without further comment that the data considered have split support. It follows from the assumptions already made that $`f_0(x_i)=f_0(x_i)`$ for $`i=2,3`$. It will be assumed that this also holds for $`i=1`$ and functions $`f_0`$ with this property will be called reflection-symmetric. This ensures that the form of the phase space density of particles is compatible with a diagonal metric and, in particular, that the energy-momentum tensor is diagonal. For the symmetry types to be considered in the following it then follows that the entire system consisting of geometry and matter is invariant under three commuting reflections. For this reason, solutions where the metric is diagonal and $`f_0`$ has the symmetry properties just mentioned will be called reflection-symmetric. A solution is said to be isotropic if $`f_0(v_1,v_2,v_3)=F((v_1)^2+(v_2)^2+(v_3)^2)`$ and if $`g_{11}g_{22}g_{33}`$ for all time.
The momentum constraints are automatically satisfied for these models. Only the Hamiltonian constraint and the evolution equations are left. Instead of considering a set of second order equations in terms of e.g., $`a`$ and $`b`$, where
$$a^2=g_{11},b^2=g_{22},$$
(3)
we will reformulate these equations as a first order system of ODEs by introducing a new set of variables. The mean curvature $`\mathrm{tr}k`$ (where $`k_{ij}`$ is the second fundamental form) is given by
$$\mathrm{tr}k=(a^1da/dt+2b^1db/dt).$$
(4)
A new dimensionless time coordinate $`\tau `$ is defined by $`\frac{1}{3}_{t_0}^t\mathrm{tr}k(t)𝑑t`$ for some arbitrary fixed time $`t_0`$. (We will follow the conventions in . The time variable thus differs by a factor 3 from the one in ). In the following a dot over a quantity denotes its derivative with respect to $`\tau `$. The Hubble variable $`H`$ is given by $`H=\mathrm{tr}k/3`$. Now define the following dimensionless variables:
$`z`$ $`=`$ $`m^2/(a^2+2b^2+m^2),`$
$`s`$ $`=`$ $`b^2/(b^2+2a^2),`$
$`M_2`$ $`=`$ $`\sigma _2(a^2/b^4)(\mathrm{tr}k)^2,`$
$`M_3`$ $`=`$ $`3\sigma _3b^2(\mathrm{tr}k)^2,`$
$`\mathrm{\Sigma }_+`$ $`=`$ $`3(b^1db/dt)(\mathrm{tr}k)^11,`$ (5)
where $`\sigma _2`$ is 1 for Bianchi types II, VIII, IX and 0 for Bianchi types I, III and the Kantowski-Sachs (KS) models. The coefficient $`\sigma _3`$ is $`1`$ for types III and VIII. It is $`1`$ for KS and type IX and 0 for types I, II <sup>1</sup><sup>1</sup>1The motivation for the variable $`s`$ comes from more general diagonal models where it is convenient to introduce variables of the type $`s_i=g^{ii}/(g^{11}+g^{22}+g^{33})`$. $`s`$ is simply $`s_1`$ in the case when $`g^{22}=g^{33}`$.. These variables lead to a decoupling of the equation for the only remaining dimensional variable $`H`$ (or equivalently $`\mathrm{tr}k`$)
$$\dot{H}=(1+q)H,$$
(6)
where the deceleration parameter $`q`$ is given by
$$q=2\mathrm{\Sigma }_+^2+\frac{1}{2}\mathrm{\Omega }(1+R).$$
(7)
The quantity $`R`$ is defined by
$$R=(p_1+2p_2)/\rho ,$$
(8)
where
$`p_1/\rho `$ $`=`$ $`(1z)sg_1/h,`$
$`p_2/\rho `$ $`=`$ $`\frac{1}{2}(1z)(1s)g_2/h,`$
$`g_{1,2}`$ $`=`$ $`{\displaystyle f_0(v_i)(v_{1,2})^2[z+(1z)(s(v_1)^2+\frac{1}{2}(1s)((v_2)^2+(v_3)^2))]^{1/2}𝑑v_1𝑑v_2𝑑v_3},`$
$`h`$ $`=`$ $`{\displaystyle f_0(v_i)[z+(1z)(s(v_1)^2+\frac{1}{2}(1s)((v_2)^2+(v_3)^2))]^{1/2}𝑑v_1𝑑v_2𝑑v_3}.`$ (9)
The assumption of split support ensures that the function $`R(s,z)`$ is a smooth ($`C^{\mathrm{}}`$) function of its arguments. The related quantity $`R_+`$ defined by
$$R_+=(p_2p_1)/\rho .$$
(10)
is a smooth function of $`s`$ and $`z`$ for the same reason.
The normalized energy density $`\mathrm{\Omega }=\rho /(3H^2)`$ is determined by the Hamiltonian constraint and, in units where $`G=1/8\pi `$, is given by
$$\mathrm{\Omega }=1\mathrm{\Sigma }_+^2M_2M_3.$$
(11)
The assumption of a distribution of massive particles with non-negative mass leads to inequalities for $`R`$, $`R_+`$ and $`\mathrm{\Omega }`$. Firstly, $`0R1`$ with $`R=0`$ only when $`z=1`$ and $`R=1`$ only when $`z=0`$. Secondly, $`RR_+\frac{1}{2}R`$ with $`R_+=\frac{1}{2}R`$ for $`s=0`$ and $`R_+=R`$ for $`s=1`$. Thirdly $`\mathrm{\Omega }0`$. Using these inequalities in equation (7) in turn results in $`0q2`$ for Bianchi types I, II, III and VIII (i.e., the same inequality as for causal perfect fluids, see ).
The remaining dimensionless coupled system is:
$`\dot{\mathrm{\Sigma }}_+`$ $`=`$ $`(2q)\mathrm{\Sigma }_+𝒮_++\mathrm{\Omega }R_+,`$
$`\dot{s}`$ $`=`$ $`6s(1s)\mathrm{\Sigma }_+,`$
$`\dot{z}`$ $`=`$ $`2z(1z)(1+\mathrm{\Sigma }_+3\mathrm{\Sigma }_+s),`$
$`\dot{M}_2`$ $`=`$ $`2(q4\mathrm{\Sigma }_+)M_2,`$
$`\dot{M}_3`$ $`=`$ $`2(q\mathrm{\Sigma }_+)M_3,`$ (12)
where $`𝒮_+`$ is given by
$$𝒮_+=4M_2M_3.$$
(13)
There are a variety of submanifolds corresponding to different symmetry types:
$`S_\mathrm{I}`$ $`:`$ $`M_2=M_3=0,`$
$`S_{\mathrm{II}}`$ $`:`$ $`M_2>0,M_3=0,`$
$`S_{\mathrm{III}}`$ $`:`$ $`M_2=0,M_3>0,`$
$`S_{\mathrm{KS}}`$ $`:`$ $`M_2=0,M_3<0,`$
$`S_{\mathrm{VIII}}`$ $`:`$ $`M_2>0,M_3>0,(1s)M_3=6M_2s`$
$`S_{\mathrm{IX}}`$ $`:`$ $`M_2>0,M_3<0,(1s)M_3=6M_2s.`$ (14)
The relationship between the various models can be visualized in a symmetry reduction diagram given in Fig. 1 (a collective treatment of the corresponding vacuum models from a Hamiltonian perspective and with the aim of quantizing the models was given in ). Note that while this diagram accurately reflects the relationship of the geometry in the different cases, the relationship of the matter content is more subtle when types VIII or IX are involved. This complication does not occur for the Bianchi types studied in detail in this paper and will therefore not be discussed further here.
Note that a non-negative energy density implies that $`\mathrm{\Omega }0`$, which in turn implies that our variables are bounded for types I,II,III and VIII, since $`M_2`$ and $`M_3`$ are non-negative and since by definition $`z`$ and $`s`$ are bounded. These models expand indefinitely. The KS and type IX models are recollapsing models and since $`H`$ becomes zero at the point of maximal expansion, the Hubble-normalized variables blow up at this point. However, one can find other variables that are bounded along the lines found in . Neither are the above variables ‘optimal’ for the other LRS models. One can adapt to the particular mathematical features these models exhibit. However, we choose to use the above formulation since the present variables are easier to interpret physically and are naturally generalizable to more general non-LRS models. For simplicity we will from now on study Bianchi types I,II, and III.
It is of interest to note that the metric functions $`a,b`$ are expressible in terms of $`s,z`$ in the massive case. The relations are
$$a^2=z(m^2s(1z))^1,b^2=2z(m^2(1s)(1z))^1.$$
(15)
In addition to the symmetry submanifolds there are a number of other boundary submanifolds:
$`z`$ $`=`$ $`0,1,`$
$`s`$ $`=`$ $`0,1,`$
$`\mathrm{\Omega }`$ $`=`$ $`0.`$ (16)
The submanifold $`z=0`$ corresponds to the massless case. The submanifold $`z=1`$ leads to a decoupling of the $`s`$-equation, leaving a system identical to the corresponding dust equations. The submanifolds $`s=0,s=1`$ correspond to problems with $`f_0`$ being a distribution while $`\mathrm{\Omega }=0`$ constitutes the vacuum submanifold with test matter. Apart from these solutions there exists an isotropic solution in Bianchi type I characterized by $`\mathrm{\Sigma }_+=R_+=0`$ and a constant value for $`s`$ that depends on the function $`f_0`$.
Including these boundaries yields compact state spaces for types I,II and III. In order to apply the standard theory of dynamical systems the coefficients must be $`C^1`$ on the entire compact state space $`G`$ of a given model. This is necessary even for uniqueness. In the present case it suffices to show that $`R`$ and $`R_+`$ are $`C^1`$ on $`G`$, i.e, that they are $`C^1`$ for $`s,z`$ when $`0s1,0z1`$. As has already been pointed out, this follows from the assumption of split support, which even implies the analogous statement with $`C^1`$ replaced by $`C^{\mathrm{}}`$. It would be possible to get $`C^1`$ regularity under the weaker assumption that $`f_0`$ vanishes as fast as a sufficiently high power of the distance to the coordinate planes. We have not, however, examined in detail how large the power would have to be since this is of little relevance to our main concerns in this paper.
Of key importance is the existence of a monotone function in the ‘massive’ interior part of the state space:
$`M`$ $`=`$ $`(s(1s)^2)^{1/3}z(1z)^1,`$
$`\dot{M}`$ $`=`$ $`2M.`$ (17)
Note that the volume $`ab^2`$ is proportional to $`M^{3/2}`$. This monotone function rules out any interior $`\omega `$\- and $`\alpha `$-limit sets and forces these sets to lie on the $`s=0`$, $`s=1`$, $`z=0`$ or $`z=1`$ parts of the boundary.
## 3 Type I models
It is natural to start investigating the type I system since it is a submanifold of the state space of all other symmetry types. The physical state space, $`G`$, of these models is given by the region in $`𝐑^3`$ defined by the inequalities $`1\mathrm{\Sigma }_+1`$, $`0s1`$ and $`0z1`$.
To understand the dynamics of the type I models, it is necessary to determine the stationary points and their stability. The coordinates, in terms of $`(\mathrm{\Sigma }_+,s,z)`$, of the various stationary points are the following: $`(0,s_0,0)`$,$`(\frac{1}{2},0,0)`$, $`(1,0,0)`$,$`(1,1,0)`$,$`(1,1,0)`$,$`(1,0,1)`$,$`(1,1,1)`$,$`(1,1,1)`$, where $`s_0`$ is a particular constant value of $`s`$ depending on the function $`f_0`$ (see ). These points are called $`P_1,\mathrm{}P_8`$. (Note that they are numbered differently in the massless case compared to those in .) In addition there exist two lines of equilibrium points, $`(1,0,K),(0,F,1)`$, denoted by $`L_1,L_2`$, where $`K`$ and $`F`$ are constant values. The points $`P_1,P_2,P_4,P_6,P_7,P_8`$ are hyperbolic saddles while $`P_5`$ is degenerate, with one zero eigenvalue. The point $`P_3`$ is a hyperbolic source. The line $`L_1`$ is a transversally hyperbolic saddle while the line $`L_2`$ is a transversally hyperbolic sink. (For an explanation of this terminology we refer to the appendix.)
The state space together with equilibrium points and separatrix orbits is depicted in Fig. 2.
The main result in this section is the following theorem:
Theorem 3.1 If a smooth non-vacuum reflection-symmetric LRS solution of Bianchi type I of the Einstein-Vlasov equations for massive particles is represented as a solution of (2) with $`M_2=M_3=0`$ then for $`\tau \mathrm{}`$ it converges to a point of the line $`L_2`$. For $`\tau \mathrm{}`$ there exists
(i) a single (isotropic) solution that converges to $`P_1`$ and
(ii) a one-parameter set of solutions lying on the unstable manifold of $`P_2`$ and
(iii) all remaining solutions belong to a two-parameter set (the generic case) of solutions converging to $`P_3`$.
This will be proved in a series of lemmas. We refer to for terminology from the theory of dynamical systems.
Lemma 3.1 There exist open neighbourhoods $`U_1`$ and $`U_2`$ of the point $`P_3`$ and the line $`L_2`$ respectively such that:
(i) if a solution belongs to $`U_1`$ at any time it belongs to $`U_1`$ at all earlier times and its $`\alpha `$-limit set consists of the point $`P_3`$ alone.
(ii) if a solution belongs to $`U_2`$ at any time it belongs to $`U_2`$ at all later times and its $`\omega `$-limit point consists of a single point of the line $`L_2`$.
Proof Part (i) follows from the fact that $`P_3`$ is a hyperbolic source and the Hartman-Grobman theorem. Part (ii) follows from the fact that $`L_2`$ is a transversally hyperbolic sink and the reduction theorem (, Theorem A1).
As a step towards analysing the dynamics of the full system we determine the $`\omega `$-limit points of solutions of the dynamical system on the parts of the boundary of $`G`$ defined by $`s=0`$ and $`s=1`$. This information will later be combined with the monotone function $`M`$ when determining the $`\omega `$-limit sets of solutions of the full system. In the case of the $`\alpha `$-limit sets the monotone function alone accomplishes the same thing.
Lemma 3.2 A solution of the restriction of the system to the part of the boundary of $`G`$ defined by $`s=1`$ for which neither $`z`$ nor $`\mathrm{\Sigma }_+`$ take on one of their limiting values has the endpoint of $`L_2`$ with $`s=1`$ as its $`\omega `$-limit set.
Proof If $`\mathrm{\Sigma }_+0`$ at any time, then $`\mathrm{\Sigma }_+`$ is decreasing. The rate of decrease remains uniform as long as $`\mathrm{\Sigma }_+`$ does not tend to zero. It follows that after a finite time $`\mathrm{\Sigma }_+`$ must be strictly less than $`1/2`$. On the other hand, $`z`$ is monotone increasing in the region $`\mathrm{\Sigma }_+<1/2`$ and the rate of increase remains uniform as long as $`z`$ does not tend to one. It follows that $`z1`$ as $`\tau \mathrm{}`$. If $`\mathrm{\Sigma }_+`$ tends to zero in this limit then the conclusion of the lemma holds. Otherwise $`\mathrm{\Sigma }_+`$ must become negative at some time. Thus it can be seen that any $`\omega `$-limit points satisfy $`z=1`$ and $`1\mathrm{\Sigma }_+0`$. From part (ii) of Lemma 3.1 it follows that any solution which enters $`U_2`$ has the desired $`\omega `$-limit set. Since the $`\omega `$-limit set is a union of orbits, it is possible as a consequence to exclude the points with $`z=1`$ and $`1<\mathrm{\Sigma }_+<0`$ from the $`\omega `$-limit set. To complete the proof of the lemma it remains only to exclude the point $`P_8`$ from the $`\omega `$-limit set. This point is a hyperbolic saddle of the restriction of the system to $`s=1`$ and so it follows from the discussion in the appendix and what has been proved already that it cannot belong to the $`\omega `$-limit set. For if $`P_8`$ belonged to the $`\omega `$-limit set points of its stable and unstable manifolds would also have to do so, and this has already been ruled out.
Lemma 3.3 A solution of the restriction of the system to the part of the boundary of $`G`$ defined by $`s=0`$ for which neither $`z`$ nor $`\mathrm{\Sigma }_+`$ take on one of their limiting values has the endpoint of $`L_2`$ with $`s=0`$ as its $`\omega `$-limit set.
Proof Along any solution of this system $`z`$ is monotone increasing on the part of the state space of the restricted dynamical system with $`\mathrm{\Sigma }_+1`$ and $`z(1z)0`$. Hence, by the monotonicity principle (see ), any $`\omega `$-limit point must satisfy $`z=1`$ or $`\mathrm{\Sigma }_+=1`$. However $`\mathrm{\Sigma }_+`$ is increasing for $`\mathrm{\Sigma }_+`$ close to but not equal to $`1`$. Hence there can be no $`\omega `$-limit points with $`\mathrm{\Sigma }_+=1`$. It follows that $`z`$ tends to one as $`\tau \mathrm{}`$ for any solution and any $`\omega `$-limit point satisfies $`z=1`$. From Lemma 3.1, any solution which enters $`U_2`$ has the desired $`\omega `$-limit set. Arguing as in the proof of Lemma 3.3 allows points with $`1<\mathrm{\Sigma }_+<0`$ and $`0<\mathrm{\Sigma }_+<1`$ to be excluded. The point $`P_6`$, which is a hyperbolic saddle of the restricted system, can be eliminated in the same way as was done in the case of $`P_8`$ in the proof of Lemma 3.2 using the results of the discussion in the appendix. Finally, the non-existence of $`\omega `$-limit points with $`\mathrm{\Sigma }_+=1`$, already mentioned above, shows that the endpoint of the line $`L_1`$ cannot lie in the $`\omega `$-limit set.
Lemma 3.4 If a solution lies in the interior of $`G`$, then unless it lies on the unstable manifold of $`P_1`$ or $`P_2`$ its $`\alpha `$-limit set consists of the point $`P_3`$.
Proof Consider a solution in the interior of $`G`$ which does not lie on the unstable manifold of $`P_1`$ or $`P_2`$. If it intersects $`U_1`$ then by Lemma 3.1 its $`\alpha `$-limit set consists of the point $`P_3`$. There can be no other $`\alpha `$-limit points in $`U_1`$. Because the function $`M`$ tends to zero along the solution as $`\tau \mathrm{}`$ the $`\alpha `$-limit set must be contained in the surface $`z=0`$. Recall that the surface $`z=0`$ corresponds to the case of massless particles which was analysed completely in . (Note that the stationary points were numbered differently in that paper.) Consider the boundary of the surface $`z=0`$. Arguing as in the proof of Lemma 3.2, the lines joining $`P_3`$ to $`P_2`$ and $`P_4`$ can be excluded from the $`\alpha `$-limit set. The discussion of the appendix and the fact that $`P_4`$ is a hyperbolic saddle with stable manifold $`\mathrm{\Sigma }_+=1`$ and unstable manifold the line connecting $`P_4`$ to $`P_5`$ can be used to exclude that line and the point $`P_4`$ itself. The line connecting $`P_5`$ to the endpoint of the line $`L_1`$ can be excluded in an analogous way, noting that the non-hyperbolic point $`P_5`$ is also covered by the discussion of the appendix. The point $`P_5`$ is also excluded by this argument. Applying the reduction theorem allows the line joining the endpoint of the line $`L_1`$ to $`P_2`$ to be excluded together with the endpoint of $`L_1`$. At this stage we can also exclude the point $`P_2`$ itself, using the results of the appendix again and the fact that by assumption the solution does not lie on the unstable manifold of $`P_2`$. Thus the only point of the boundary of the set $`z=0`$ which can belong to the $`\alpha `$-limit set is $`P_3`$. Now suppose that a point of the interior of the surface belongs to the $`\alpha `$-limit set. If it is a point of the unstable manifold of $`P_2`$ then $`P_2`$ also belongs to the $`\alpha `$-limit set, in contradiction to what has just been proved. If it is some other point other than $`P_1`$ then, using the fact that the $`\alpha `$-limit set is a union of orbits and Theorem 3.1 of , it follows that $`P_3`$ belongs to the $`\alpha `$-limit set and we obtain a contradiction again. Finally, if it were $`P_1`$ then the results of the appendix would imply that other points of the interior would belong to the $`\alpha `$-limit set, and this has just been ruled out.
Lemma 3.5 The $`\omega `$-limit point of each solution in the interior of $`G`$ is a point of the line $`L_2`$.
Proof Note first that the function $`M`$ goes to infinity along any such solution as $`\tau \mathrm{}`$. It follows that any $`\omega `$-limit point must satisfy $`z=1`$, $`s=0`$ or $`s=1`$. If the solution enters the set $`U_2`$ then by part (ii) of Lemma 3.1 the $`\omega `$-limit set is as claimed. There are no other $`\omega `$-limit points of any solution in $`U_2`$. Consider now the evolution of $`\mathrm{\Sigma }_+`$ on the surface $`z=1`$. It either increases from $`1`$ to $`0`$ or decreases from $`1`$ to $`0`$. Since the $`\omega `$-limit set is a union of orbits, we conclude that no point of the interior of the surface $`z=1`$ or its boundary lines $`s=0`$ and $`s=1`$ other than the points of the line $`L_2`$ can belong to the $`\omega `$-limit set. Using once more the fact that the $`\omega `$-limit set is a union of orbits, it is possible to exclude the interior of the surface $`s=1`$ from the $`\omega `$-limit set by Lemma 3.2 and the interior of $`s=0`$ by Lemma 3.3. Now all remaining possibilities other than points on $`L_2`$ will be excluded successively. The nature of the line $`L_1`$ as a transversely hyperbolic saddle suffices to eliminate it, as well as the lines joining it to $`P_8`$ and $`P_2`$. The point $`P_3`$, being a hyperbolic source, is clearly ruled out, and with it the lines joining it to $`P_2`$ and $`P_6`$. Further applications of the results of the appendix rule out the remaining lines, namely those joining $`P_8`$ to $`P_5`$, $`P_5`$ to $`P_4`$, $`P_4`$ to $`P_7`$ and $`P_7`$ to $`P_6`$. It follows that the $`\omega `$-limit set is contained in the line $`L_2`$. Applying the reduction theorem then shows that the $`\omega `$-limit set is a single point of $`L_2`$.
The results of Lemma 3.4 and Lemma 3.5 together imply Theorem 3.1.
Theorem 3.1 has been formulated entirely in terms of the dynamical systems picture. It should, however, be pointed out that this allows asymptotic expansions for all quantities of geometrical or physical interest near the singularity or in an expanding phase to be obtained if desired. For example, in an expanding phase in type I the following expansions can be derived:
$`\mathrm{\Sigma }_+`$ $`=`$ $`\alpha t^1+o(t^1)`$ (18)
$`s`$ $`=`$ $`s_0\frac{4}{9}s_0(1s_0)t^1+o(t^1)`$ (19)
$`z`$ $`=`$ $`1\beta t^{4/3}+o(t^{4/3})`$ (20)
$`H`$ $`=`$ $`\frac{2}{3}t^1+O(t^{7/3})`$ (21)
$`\mathrm{\Omega }`$ $`=`$ $`1\alpha t^2+o(t^2)`$ (22)
$`\rho `$ $`=`$ $`\frac{4}{9}t^2\frac{4}{3}\alpha ^2t^4+o(t^4)`$ (23)
$`p_1`$ $`=`$ $`O(t^{10/3})`$ (24)
Here $`\alpha `$ and $`\beta `$ are constants depending on the solution. It should be emphasized that these are not just formal expansions, but rigorous results which emerge from the dynamical systems analysis.
A particular consequence of Theorem 3.1 is that all LRS type I models isotropize at late times. This was already proved by other means in , where it was also shown that non-LRS models of Bianchi type I isotropize and have dust-like behaviour for $`\tau \mathrm{}`$.
## 4 Type II models
The physical state space, $`G`$, of the LRS type II models is given by the region in $`𝐑^4`$ defined by the inequalities $`M_20`$, $`0s1`$, $`0z1`$, and $`1\mathrm{\Sigma }_+^2M_20`$.
The coordinates, in terms of $`(\mathrm{\Sigma }_+,s,z,M_2)`$, of the various stationary points are the following: $`(0,s_0,0,0)`$,$`(\frac{1}{2},0,0,0)`$, $`(1,0,0,0)`$,$`(1,1,0,0)`$,$`(1,1,0,0)`$,$`(1,0,1,0)`$,$`(1,1,1,0)`$,$`(1,1,1,0)`$,
$`(\frac{1}{5},1,0,\frac{6}{25})`$,$`(\frac{1}{8},0,1,\frac{3}{64})`$,$`(\frac{1}{8},1,1,\frac{3}{64})`$, where $`s_0`$ is the same particular constant value of $`s`$ that appeared in the previous type I section. These points are called $`P_1,\mathrm{}P_{11}`$ (note that they are numbered differently than in , in the massless case). In addition there exist two lines of equilibrium points, $`(1,0,K,0),(0,F,1,0)`$, denoted by $`L_1,L_2`$, where $`K`$ and $`F`$ are constants. The first eight stationary points and the two lines correspond to points and lines of the same name in the Bianchi I system and their coordinates are obtained by appending a zero to those of the Bianchi I points. The points $`P_1,P_2,P_3,P_4,P_6,P_7,P_8,P_9,P_{10}`$ are hyperbolic saddles while $`P_5`$ is degenerate, with one zero eigenvalue. The point $`P_{11}`$ is a hyperbolic sink with two real and two complex eigenvalues. The lines $`L_1`$ and $`L_2`$ are transversally hyperbolic saddles.
To prove results about the global properties of solutions it is helpful to use certain monotone functions. The first is defined for $`s<1`$ by
$$Z_1=(2s/(1s))^{4/3}M_2$$
(25)
This is obtained by rewriting the function whose time derivative was calculated in equation (23) of in terms of the variables of this paper and observing that it remains monotone in the massive case. It satisfies $`\dot{Z}_1=2qZ_1`$. The second is obtained by combining $`Z_1`$ with the monotone function $`M`$ available for all the Bianchi types considered in this paper. Let $`Z_2=Z_1M^2=2^{4/3}s^2M_2z^2(1z)^2`$ for $`z>0`$. It satisfies $`\dot{Z}_2=2(q2)Z_2`$. The function $`Z_1`$ is defined on the part of the Bianchi II state space where $`s1`$ and monotonically increasing except where it vanishes. This is clear if $`q0`$. If $`q=0`$ it follows that $`\mathrm{\Sigma }_+=0`$ and $`M_2=1`$ and at points satisfying these conditions $`\dot{\mathrm{\Sigma }}_+0`$. The function $`Z_2`$ is defined on the part of the Bianchi II state space where $`z>0`$ and is monotonically decreasing except on the set where it vanishes, since $`q=2`$ implies $`Z_2=0`$.
Theorem 4.1 If a smooth non-vacuum reflection-symmetric LRS solution of Bianchi type II of the Einstein-Vlasov equations for massive particles is represented as a solution of (2) with $`M_3=0`$, then for $`\tau \mathrm{}`$ it converges to $`P_{11}`$. For $`\tau \mathrm{}`$ there exists
(i) a one-parameter set of solutions converging to the unstable manifold of $`P_1`$ and
(ii) a three-parameter set of all remaining solutions converging to the heteroclinic cycle on the $`z=0`$ submanifold, consisting of the orbits connecting the $`z=0`$ endpoint of the line $`L_1`$ to $`P_5`$, $`P_5`$ to $`P_4`$, $`P_4`$ to $`P_3`$ on the type I boundary and $`P_3`$ to the $`z=0`$ endpoint of the line $`L_1`$ via the vacuum boundary.
Lemma 4.1 If a solution belongs to the interior of the type II state space then any $`\alpha `$-limit point satisfies $`z=0`$ and $`sM_2=0`$. Any $`\omega `$-limit point satisfies $`s=1`$ and $`(z1)M_2=0`$.
Proof From the evolution equation for $`M`$ it follows that $`z=0`$ for any $`\alpha `$-limit point and that for any $`\omega `$-limit point $`z=1`$, $`s=0`$ or $`s=1`$. Next the monotonicity principle will be applied to the functions $`Z_1`$ and $`Z_2`$. Applying it to $`Z_1`$ on the region where $`Z_10`$ shows that for any $`\alpha `$-limit point $`s=0`$ or $`M_2=0`$. It also shows that there are no $`\omega `$-limit points with $`s=0`$. Combining this with the information obtained already shows that any $`\omega `$-limit point satisfies $`z=1`$ or $`s=1`$. If $`z1`$ then it follows from the monotonicity principle applied to $`Z_2`$ that $`M_2=0`$ for any $`\omega `$-limit point. The monotonicity of $`Z_1`$ then implies that $`s1`$ as $`\tau \mathrm{}`$.
Lemma 4.2 Consider the dynamical system obtained by restricting the type II system to the plane defined by the conditions $`s=1`$ and $`z=1`$. If a solution belongs to the interior of the state space for this restricted system then it converges to $`P_{11}`$ as $`\tau \mathrm{}`$.
Proof The restricted dynamical system is identical with that for type II dust solutions. In it was proved by using a monotone function derived by Hamiltonian methods that for $`\tau \mathrm{}`$ the dust solutions satisfy $`\mathrm{\Sigma }_+\frac{1}{8}`$ and $`M_2\frac{3}{64}`$. Hence it can be concluded that the solution approaches $`P_{11}`$ as $`\tau \mathrm{}`$.
Lemma 4.3 If a solution lies in the interior of the type II state space then unless it lies on the unstable manifold of $`P_1`$ (and this does occur) the $`\alpha `$-limit set consists of the heteroclinic cycle described in the statement of Theorem 4.1.
Proof By Lemma 4.1 we know that any $`\alpha `$-limit point satisfies $`z=0`$. Moreover it satisfies $`M_2=0`$ or $`s=0`$. The situation is very similar to that in the massless case treated in and the proof may be taken over rather directly. It is only necessary to pay attention to the fact that it is the nature of the stationary points in the full massive Bianchi II state space which must be taken into account and that the notation is different.
Suppose that the $`\alpha `$-limit set contains a point with $`s=0`$ and $`M_20`$. Then by Lemma 4.3 of it contains the endpoint of $`L_1`$ and either $`P_2`$ or $`P_3`$. On the other hand, if it contains a point with $`M_2=0`$ then this belongs to the massless Bianchi I state space. Then it must contain one of the points $`P_1`$, $`P_2`$, $`P_3`$, $`P_4`$, $`P_5`$ or the endpoint of $`L_1`$. To prove the lemma we may assume that the solution does not lie on the unstable manifold of $`P_1`$. If $`P_1`$ nevertheless belonged to the $`\alpha `$-limit set then this set would have to include points belonging to the unstable manifold of $`P_1`$ other than $`P_1`$ itself. But these satisfy neither $`s=0`$ nor $`M_2=0`$ and so this gives a contradiction. Thus under the given assumptions the $`\alpha `$-limit set does not contain $`P_1`$. If the $`\alpha `$-limit set contained a point with $`M_2=0`$, $`|\mathrm{\Sigma }_+|<1`$ and $`0<s<1`$ it would contain $`P_1`$ (in its role as $`\omega `$-limit set for Bianchi type I solutions), leading once more to a contradiction. If the $`\alpha `$-limit set contained $`P_2`$ then by the results of the appendix it would contain $`P_1`$, which is also not possible. Applying Lemma 4.3 of again allows points with $`M_20`$ which are not on the vacuum boundary to be excluded from the $`\alpha `$-limit set. The straight lines joining $`P_2`$ to $`P_3`$ and the endpoint of $`L_1`$ are excluded as well. The conclusion is that the $`\alpha `$-limit set is contained in the heteroclinic cycle mentioned in the statement of Theorem 4.1. It remains to show that it is the whole heteroclinic cycle. This is straightforward to do using the results of the appendix.
Lemma 4.4 If a solution lies in the interior of the type II state space then it converges to the point $`P_{11}`$ as $`\tau \mathrm{}`$.
Proof Consider any $`\omega `$-limit point with $`z1`$. Then by Lemma 4.1 this point satisfies $`s=1`$ and $`M_2=0`$. Any nearby $`\omega `$-limit points must also satisfy these conditions. If any of these limit points satisfied $`z=0`$ then $`P_4`$ and $`P_5`$ would be $`\omega `$-limit points of the given solution. Using the saddle point properties of these points then shows that $`P_7`$ and $`P_8`$ belong to the $`\omega `$-limit set. Repeating the same argument shows that the endpoint of $`L_1`$ with $`s=1`$ is an $`\omega `$-limit point. The fact that this point is a transversely hyperbolic saddle implies that its unstable manifold in the hyperplane $`s=1`$ is contained in the $`\omega `$-limit set. By Lemma 4.2 the $`\omega `$-limit set also contains $`P_{11}`$. Since $`P_{11}`$ is a hyperbolic sink this contradicts the assumption $`z1`$. Thus we conclude that the entire $`\omega `$ limit set is contained in the plane defined by the equations $`s=1`$ and $`z=1`$. The argument just given rules out the possibility of $`\omega `$-limit points with $`M_2=0`$. Applying Lemma 4.2 once more shows that the only possible $`\omega `$-limit point which does not lie on the vacuum boundary is $`P_{11}`$. Finally the fact that $`P_7`$ and $`P_8`$ are hyperbolic saddles can be used to rule out points of the vacuum boundary, thus completing the proof.
## 5 Type III models
The physical state space, $`G`$, of the LRS type III models is given by the region in $`𝐑^4`$ defined by the inequalities $`M_30`$, $`0s1`$, $`0z1`$, and $`1\mathrm{\Sigma }_+^2M_30`$.
The coordinates, in terms of $`(\mathrm{\Sigma }_+,s,z,M_3)`$, of the various stationary points are the following: $`(0,s_0,0,0)`$,$`(\frac{1}{2},0,0,0)`$, $`(1,0,0,0)`$,$`(1,1,0,0)`$,$`(1,1,0,0)`$,$`(1,0,1,0)`$,$`(1,1,1,0)`$,$`(1,1,1,0)`$,
$`(\frac{1}{2},0,0,\frac{3}{4})`$,$`(\frac{1}{2},0,1,\frac{3}{4})`$, where $`s_0`$ the same particular constant value of $`s`$ that appeared in the previous type I section. These points are called $`P_1,\mathrm{}P_{10}`$ (note that they are numbered differently than in , in the massless case). In addition there exist three lines of equilibrium points, $`(1,0,K,0,0),(0,F,1,0,0),(\frac{1}{2},1,z_0,\frac{3}{4})`$, denoted by $`L_1,L_2,L_3`$, where $`K,F`$ and $`z_0`$ are constants. The first eight stationary points and the first two lines correspond to points and lines of the same name in the Bianchi I system and their coordinates are obtained by appending a zero to those of the Bianchi I points. The points $`P_1,P_2,P_4,P_6,P_7,P_8,P_9`$ are hyperbolic saddles while $`P_5`$ and $`P_{10}`$ are degenerate, with one zero eigenvalue each. The point $`P_3`$ is a hyperbolic source. The lines $`L_1`$ and $`L_2`$ are transversally hyperbolic saddles while the line $`L_3`$ is degenerate with two zero and two negative eigenvalues.
To prove global results about the global properties of solutions it is useful to note the existence of the following bounded monotone function
$`\stackrel{~}{M}_3`$ $`=`$ $`M_3(2\mathrm{\Sigma }_+)^2,`$
$`\dot{\stackrel{~}{M}_3}`$ $`=`$ $`2\stackrel{~}{M}_3[(12\mathrm{\Sigma }_+)^2+\mathrm{\Omega }(R+R_+)](2\mathrm{\Sigma }_+)^1.`$ (26)
Theorem 5.1 If a smooth non-vacuum reflection-symmetric LRS solution of Bianchi type III of the Einstein-Vlasov equations for massive particles is represented as a solution of (2) with $`M_2=0`$, then for $`\tau \mathrm{}`$ it converges to a point of the line $`L_3`$ with $`z>0`$. For $`\tau \mathrm{}`$ there exists
(i) a one-parameter set of solutions lying on the unstable manifold of $`P_1`$ and
(ii) a two-parameter set of solutions lying on the unstable manifold of $`P_2`$ and
(iii) all remaining solutions converge to $`P_3`$.
In all these solutions the scale factor $`a`$ is monotone increasing at late times.
Lemma 5.1 If a solution belongs to the interior of the type III state space any $`\alpha `$-limit point satisfies $`M_3=0`$. Any $`\omega `$-limit point satisfies $`\mathrm{\Sigma }_+=\frac{1}{2}`$, $`M_3=\frac{3}{4}`$ and $`s=1`$.
Proof The continuous function $`\stackrel{~}{M}_3`$ on the state space must have a maximum and since its gradient never vanishes this maximum can only be attained at points with $`M_3=1\mathrm{\Sigma }_+^2`$. Computing the derivative of $`\stackrel{~}{M}_3`$ along the curve in the $`(M_3,\mathrm{\Sigma }_+)`$ plane defined by this relation shows that the maximum value is $`\frac{1}{3}`$ and that it is attained when $`\mathrm{\Sigma }_+=\frac{1}{2}`$ and $`M_3=\frac{3}{4}`$. Now we apply the monotonicity principle. Let $`S`$ be the part of the Bianchi III state space obtained by removing the points with $`M_3=0`$ and those with $`\mathrm{\Sigma }_+\frac{1}{2}=M_3\frac{3}{4}=0`$. This is an invariant set for the dynamical system. It will now be shown that $`\stackrel{~}{M}_3`$ is strictly increasing along solutions on this set. If $`\mathrm{\Sigma }_+\frac{1}{2}`$ or if $`\mathrm{\Omega }(R+R_+)0`$ then this follows immediately from (5). If $`\mathrm{\Sigma }_+=\frac{1}{2}`$ and $`\mathrm{\Omega }(R+R_+)=0`$ then
$$\dot{\mathrm{\Sigma }}_+=\frac{3}{4}(M_3\frac{3}{4})$$
(27)
This completes the proof that $`\stackrel{~}{M}_3`$ is strictly increasing on $`S`$. The monotonicity principle then shows that any point in the $`\alpha `$-limit set must be in the complement of $`S`$ and such that such that $`\stackrel{~}{M}_3`$ does not take on its maximum value on $`\overline{S}`$ there. Hence $`M_3=0`$ there. It also shows that any point in the $`\omega `$-limit set must be in the complement of $`S`$ and that $`\stackrel{~}{M}_3`$ does not take on its minimum value there. Hence in the latter case $`\mathrm{\Sigma }_+=\frac{1}{2}`$ and $`M_3=\frac{3}{4}`$. It follows from this that $`\mathrm{\Sigma }_+\frac{1}{2}`$ as $`\tau \mathrm{}`$ and the equation for $`s`$ then implies that $`s1`$.
Lemma 5.2 A solution which belongs to the interior of the type III state space converges to a point of the line $`L_3`$ with $`z>0`$ as $`\tau \mathrm{}`$.
Proof Because of the result of Lemma 5.1 it only remains to prove that $`z`$ tends to a positive limit as $`\tau \mathrm{}`$. Note first that the evolution equation for $`s`$ implies an equation of the form $`(d/d\tau )(1s)=(1s)F`$ where $`F=6s\mathrm{\Sigma }_+`$. As $`\tau `$ tends to infinity $`F3`$ and a simple comparison argument proves that $`1s(\tau )=O(e^{(3+ϵ)\tau })`$ as $`\tau \mathrm{}`$. In particular, $`1s`$ decays exponentially to zero at late times. The evolution equations imply that $`\mathrm{\Omega }`$ satisfies the equation:
$$\dot{\mathrm{\Omega }}/\mathrm{\Omega }=(\mathrm{\Sigma }_+\frac{1}{2})[(3R)\mathrm{\Sigma }_++\frac{3}{2}(1+R)]2\mathrm{\Sigma }_+(R_++R)(M_3\frac{3}{4})(1+R)$$
(28)
Note that $`R_++R0`$ so that the second term on the right hand side is negative. However it is exponentially small at late times since it contains a factor $`(1s)`$ when expressed in terms of the matter quantities. In particular $`\dot{\mathrm{\Omega }}/\mathrm{\Omega }0`$ as $`\tau \mathrm{}`$ and $`\mathrm{\Omega }^1=O(e^{ϵ\tau })`$ for any $`ϵ>0`$. This means that $`\mathrm{\Omega }`$ converges to zero slower than any exponential. In other words, $`\mathrm{\Omega }e^{ϵ\tau }`$ tends to infinity for any $`ϵ>0`$. Suppose that $`\mathrm{\Sigma }_+\frac{1}{2}`$ for some solution at some time. Then $`M_3\frac{3}{4}`$ and the first and third terms in the expression for $`\dot{\mathrm{\Omega }}/\mathrm{\Omega }`$ are positive at late times while the second term is negative. It will now be shown that the third term decays slower than any exponential and thus must eventually dominate the second term. For
$$\mathrm{\Omega }=(\frac{1}{4}\mathrm{\Sigma }_+^2)+(\frac{3}{4}M_3)(\frac{3}{4}M_3)$$
(29)
It follows that $`\dot{\mathrm{\Omega }}/\mathrm{\Omega }>0`$ at late times as long as $`\mathrm{\Sigma }_+>\frac{1}{2}`$. Since it follows from Lemma 5.3 that $`\mathrm{\Omega }0`$ as $`\tau \mathrm{}`$ it follows that for any time $`\tau _0`$ for which $`\mathrm{\Sigma }_+(\tau _0)>\frac{1}{2}`$ there exists a time $`\tau >\tau _0`$ with $`\mathrm{\Sigma }_+=\frac{1}{2}`$. When $`\mathrm{\Sigma }_+=\frac{1}{2}`$ then
$$\dot{\mathrm{\Sigma }}_+=(M_3\frac{3}{4})[\frac{3}{4}(1R)\frac{1}{4}(R+R_+)]$$
(30)
Now it follows from the evolution equation for $`z`$ that $`1z`$ cannot approach zero faster than, for instance, $`e^\tau `$ and the same is then true of $`1R`$. It can be concluded that the first term in the square bracket on the right hand side of (30) dominates the second at late times. Hence $`\frac{1}{2}\mathrm{\Sigma }_+`$ must be negative at late times, which in turns implies that $`z`$ is increasing and that it must tend to a positive limit.
Lemma 5.3 If a solution lies in the interior of the type III state space then unless it lies on the unstable manifold of $`P_1`$ or $`P_2`$ (and both of these cases occur) the $`\alpha `$-limit set consists of the point $`P_3`$.
Proof Note first that it can be concluded as in the proof of Lemma 3.4 that any $`\alpha `$-limit point satisfies $`z=0`$. Thus, applying Lemma 5.1, it can be identified with a point of the state space for massless type I solutions. Now it is possible to proceed further following the method of proof of Lemma 3.4. Consider the boundary of the state space for massless type I solutions. The point $`P_3`$, being a hyperbolic source in the type III state space, can be excluded as an $`\alpha `$-limit point of a solution of type III. It is then possible to successively exclude points of the boundary as in the proof of Lemma 3.4. The facts which need to be used are that all $`\alpha `$-limit points satisfy $`M_3=0`$ and $`z=1`$ and that the points $`P_4`$, $`P_5`$ and the endpoint of $`L_1`$ are a hyperbolic saddle, a non-hyperbolic saddle topologically equivalent to a hyperbolic one and a transversely hyperbolic saddle, respectively. At this stage it can be concluded that all $`\alpha `$-limit points of solutions of type III are either $`P_1`$, $`P_2`$ or points of the unstable manifold of $`P_1`$. For all other points of the interior of the massless type I state space lie on solutions which converge to the hyperbolic source $`P_3`$ in the past time direction, and so are excluded. It remains to examine what happens in a neighbourhood of the points $`P_1`$ and $`P_2`$, which are both hyperbolic saddles. The unstable manifold of $`P_2`$ in the type III state space is three-dimensional and so there are solutions which converge to $`P_2`$ as $`\tau \mathrm{}`$. Any other type III solutions which had $`P_2`$ as an $`\alpha `$-limit point would have to have $`\alpha `$-limit points on the stable manifold of $`P_2`$, which has already been excluded. Hence solutions of type III which do not converge in the past to $`P_2`$ cannot have $`P_2`$ or a point of its unstable manifold as $`\alpha `$-limit points. Thus the only remaining possibility is that solutions lie on the unstable manifold of $`P_1`$ and converge to that point in the past. Since the unstable manifold is two-dimensional, solutions of this kind exist.
The results of Lemma 5.2 and Lemma 5.3 together imply all the results of Theorem 5.1 except the last directly. The statement about the scale factor $`a`$ follows from the fact, derived in the course of the proof of Lemma 5.2, that $`\mathrm{\Sigma }_+<\frac{1}{2}`$ at late times.
## 6 Concluding remarks
In this paper we studied the dynamics of solutions of the Einstein-Vlasov equations which are locally rotationally symmetric, reflection-symmetric and of Bianchi types I, II and III. The initial singularities are of four types. There are isotropic singularities which, in the dynamical systems description used in this paper, are those which converge to the point $`P_1`$ as $`\tau \mathrm{}`$. The general theory of isotropic singularities developed by Anguige and Tod implies as a very special case the occurrence of isotropic singularities in Bianchi models with collisionless matter and information about how many there are. They only developed the theory for massless particles and so in order to apply to the situations considered here it would have to be generalized to the massive case. There are barrel singularities which occur in types I and III but not in type II. In the dynamical systems picture these are the solutions which converge to $`P_2`$ as $`\tau \mathrm{}`$. Fluid models with corresponding symmetries never have barrel singularities and so this is a peculiarity of collisionless matter, both in the case of massive particles studied here and that of massless particles studied in . There is the generic case in types I and III, which concerns solutions which develop from an open dense set of initial data for each of these Bianchi types. These solutions have a cigar singularity and converge to $`P_3`$ as $`\tau \mathrm{}`$. Finally, there are the generic solutions of type II, which have an oscillatory initial singularity.
As far as the late time behaviour is concerned, it is tempting to speculate that behaving like a dust model at late times in an expanding phase may be a general feature of solutions of the Einstein-Vlasov equations with massive particles. We know of no counterexample to this. For the solutions of types I and II treated in this paper it has been proved to be true. For type III the situation appears to be delicate and the occurrence of degenerate stationary points of the dynamical system may require an application of centre manifold theory in order to determine details of the asymptotics. A possible criterion for detecting cases where there may be trouble is as follows. Consider a dust solution which is a candidate for the asymptotic state of solutions of the Einstein-Vlasov equations. If each eigenvalue of the second fundamental form of the homogeneous hypersurfaces, when divided by the mean curvature, is bounded below by a positive constant in the dust solution then it is a strong candidate. Otherwise difficulties are to be expected. This criterion gives a positive recommendation for types I and II and a warning for type III. Thus at least for the models investigated in this paper it is a good guide. Using the information on dust models in chapter 6 of it also gives a positive recommendation for types I, II and VI<sub>0</sub> without the need to restrict to the LRS case.
In this paper a dynamical system has been set up for all LRS Bianchi models of class A as well as for Kantowski-Sachs models and the type III models, which are of class B. We expect that techniques similar to those used here can be applied to analyse Kantowski-Sachs models and LRS models of type VIII and IX. An important feature of all these LRS models is that the Vlasov equation can be solved exactly. This is also true of general Bianchi type I models. Some limited results on the dynamics of Bianchi type I solutions of the Einstein-Vlasov equations which are reflection-symmetric but not necessarily LRS were proved in . A heuristic analysis of reflection-symmetric type I models was given using Hamiltonian techniques in , where there are also interesting remarks on the general Bianchi I case. It would be very desirable to have a mathematically rigorous implementation of the ideas of .
What can be done in cases where the Vlasov equation cannot be solved exactly? If, as already speculated above, the late time evolution resembles that of a dust solution and if the dust solutions are asymptotically LRS then it may be possible to give a good approximation to the solution of the Vlasov equation in that regime. There is one drawback of this idea as a general tool for Bianchi class A models. Unfortunately there are no LRS spacetimes of Bianchi type VI<sub>0</sub>. In the case of fluids there exists a special class of Bianchi VI<sub>0</sub> spacetimes which is often characterized by the rather opaque statement that $`n_\alpha ^\alpha =0`$. These spacetimes do have a simple geometric characterization which will now be explained. Every Bianchi class A spacetime has a discrete group of isometries whose generators simultaneously reverse two of the invariant one-forms on the group. The special class of Bianchi VI<sub>0</sub> solutions can be characterized by the existence of an additional isometry which reverses just one of the one-forms. It is possible to consider solutions of the Einstein-Vlasov equations with the corresponding type of symmetry. We are, however, not aware that the Vlasov equation can be solved exactly in these special spacetimes. If it could then this might fill the apparent gap in the strategy just suggested.
The oscillatory behaviour observed near the singularity in type II models appears at first sight to indicate that collisionless matter does not fit into the analysis of general spacetime singularities by Belinskii, Khalatnikov and Lifshitz . On the other hand, the fact that in the analysis of Misner using a time-dependent potential we see the phenomenon of walls moving too fast to be caught suggests that the oscillations might go away in general models. This issue requires further work. It could turn out that collisonless matter generically becomes negligible near the singularity, as originally stated for fluids in .
To conclude, we mention some further interesting open problems. What happens in the case of a model with two species of particles, one massive and one massless? Of course this could be thought of as a simple cosmological model incorporating both baryonic matter and the microwave background photons. It is related to the two-fluid models which have been analysed in . Mixtures of fluids and kinetic theory could also be considered. We have seen that the Einstein equations with collisionless matter as source may behave very differently from the Einstein equations with a fluid source at early times (and also at late times in the massless case). Under what circumstances are there intermediate stages of the evolution with collisionless matter which can be well described by a fluid? Since the point $`P_1`$ is a saddle there are obviously solutions which approach this point and then go away again but is there more that can be said about this issue? What can be said about inhomogeneous models? In Rein analysed the behaviour at early times of solutions of the Einstein-Vlasov equations with spherical, plane and hyperbolic symmetry and massive particles. He identified open subsets of initial data for these symmetry types with a singularity resembling the generic LRS solutions of types I and III. There is an overlap between the results of and those of the present paper. It could be illuminating to attempt a common generalization of these. In any case, it is clear that one of the central challenges of the future in the study of cosmological solutions of the Einstein-Vlasov equations, or indeed the Einstein equations coupled to any type of matter fields, is to develop techniques which apply to inhomogeneous problems. A thorough understanding of the homogeneous case is likely to be an invaluable guide in addressing it.
## Appendix A Appendix
In this appendix some general procedures which are useful in determining limit sets of solutions of dynamical systems will be outlined. Let $`\gamma `$ be an orbit of a dynamical system and $`p`$ a stationary point. We will discuss only $`\omega `$-limit sets, but corresponding statements about $`\alpha `$-limit sets follow immediately by reversing the direction of time. We consider the following three statements which may or may not be true for given choices of $`\gamma `$ and $`p`$.
1. $`p`$ is an $`\omega `$-limit point of $`\gamma `$
2. $`\gamma `$ lies on the stable manifold of $`p`$
3. there are $`\omega `$-limit points of $`\gamma `$ different from $`p`$ which are arbitrarily close to $`p`$ and lie on the unstable manifold of $`p`$
4. there are $`\omega `$-limit points of $`\gamma `$ different from $`p`$ which are arbitrarily close to $`p`$ and lie on the stable manifold of $`p`$
In the body of the paper we frequently use certain relations among the statements above which hold under various assumptions on the nature of the stationary point $`p`$. Whatever the stationary point, it is always true that the statement 1. is implied by any of the statements 2., 3. or 4. This is a consequence of the elementary fact that the $`\omega `$-limit set is closed. Now suppose that $`p`$ is a hyperbolic stationary point. In this case, if 1. is true and 2. is false then both 3. and 4. are true. This follows from Lemma A1 of . Combining these statements we see that for a hyperbolic stationary point there are two mutually exclusive cases under which 1. can hold. Either $`\gamma `$ lies on the unstable manifold of $`p`$ or the $`\omega `$-limit set contains points of both the stable and unstable manifolds of $`p`$ arbitrarily close to $`p`$. In particular, if $`p`$ is a hyperbolic source then it cannot be in the $`\omega `$-limit set of $`\gamma `$ and if $`p`$ is a hyperbolic sink and $`p`$ is in the $`\omega `$-limit set of $`\gamma `$ it is the whole $`\omega `$-limit set. If we already have some a priori information about where $`\omega `$-limit points can lie (due, for instance, to the existence of a monotone function) then this gives more information about where the points on the stable and unstable manifolds whose existence is guaranteed by the general statements above can lie.
Next we consider the case of transversally hyperbolic stationary points. Suppose that $`p`$ belongs to a manifold of stationary points of dimension $`d`$. (Only the case $`d=1`$ occurs in this paper.) These points have a zero eigenvalue of multiplicity $`d`$. If all other eigenvalues have non-vanishing real parts then the stationary point is called transversally hyperbolic. (Depending on the signs of the eigenvalues the manifold of stationary points is called a transversally hyperbolic source, sink or saddle.) By the reduction theorem (Theorem A1 of ) each of these points lies on an invariant manifold and the restriction of the flow to each invariant manifold is topologically equivalent to that near a hyperbolic stationary point. The arguments for a hyperbolic stationary point adapt easily to give analogous statements for transversally hyperbolic stationary points. In applying these results we can essentially ignore the directions along the manifold of stationary points.
Finally we consider certain other non-hyperbolic stationary points. A result of the type we need was proved in Lemma A2 of but we would like to formulate the statement in a more transparent way here. Consider an isolated stationary point $`p`$ with a trivial stable manifold and a one-dimensional centre manifold. Using the reduction theorem we see that the unstable manifold divides a neighbourhood of $`p`$ into two parts on each of which the restriction of the dynamical system is topologically equivalent to the restriction of a dynamical system with a hyperbolic stationary point. Whether the latter system has a saddle or a source depends on a certain sign condition. This condition may be different for the two halves. In the dynamical systems considered in this paper the only example of this is provided by the point $`P_5`$. Only one of the halves belongs to the physical part of the state space and in that half the sign is such that a saddle is obtained. The result of these considerations is that for the arguments in this paper $`P_5`$ may be treated just as if it had been a hyperbolic saddle, with the centre manifold taking over the role of the trivial stable manifold.
|
warning/0005/hep-th0005197.html
|
ar5iv
|
text
|
# Acknowledgements
## Acknowledgements
The work of S.O. has been supported in part by Japan Society for the Promotion of Science and that of S.D.O. by CONACyT (CP, ref.990356 and grant 28454E). S.O. thanks T. Sakai for useful discussions in Yukawa Institute of Theoretical Physics.
|
warning/0005/hep-th0005057.html
|
ar5iv
|
text
|
# Contents
## Chapter 1 Introduction
Physical phenomena occurring in high energy physics are analysed in terms of ‘particles’, arising as asymptotic configurations of elementary entities in scattering experiments. These particles are characterized by certain specific intrinsic properties, which are expressed by quantum numbers whose integration in the framework of a consistent and complete theoretical description is an aim of quantum field theory. The usual theoretical description of particles goes back to the famous analysis by Wigner of the irreducible representations of the Poincaré group . He gives a complete classification of all these representations, which are labelled by two parameters $`m`$ and $`s`$. It is assumed that a particle pertains to a specific representation of this group, in which case the parameters $`m`$ and $`s`$ are interpreted as its intrinsic mass and spin, respectively. However, this approach to a theoretical description of mass and spin is not universally applicable. There are quantum field theories in which particles coupled to particles of zero rest mass cannot be described in terms of eigenstates of the mass operator. An example is quantum electrodynamics where charged particles are inevitably accompanied by soft photons. It is an open question, known as the infraparticle problem , how mass and spin of a particle are to be described in the framework of quantum field theory. Moreover, standard collision theory does not work in these cases.
A closer analysis of quantum electrodynamics shows that the infraparticle problem is connected with Gauss’ law . An outline of the underlying mechanism, following arguments of Buchholz in , may be appropriate at this point. Due to Gauss’ law, the charge of a physical state can be determined by measuring the electromagnetic field at asymptotic spacelike distances. These measurements do not interfere with those performed within bounded regions; therefore, being a $`c`$-number, the asymptotic field configuration is a superselection rule of the theory. Its dependence on the state of motion of the charged particle implies that there exists a multitude of superselection sectors and that the Lorentz symmetry is broken. Consequently, charged particles cannot be described according to Wigner’s theory.
The present thesis proposes a novel approach to the concept of particles, elaborating some of the ideas of Buchholz’ which he introduced in . In a model-independent framework, especially without excluding massless states and without assuming asymptotic completeness of the theory, an approach of Araki and Haag to scattering theory is reconsidered. Chapter 2 introduces the concept of detectors to be used in this work and investigates the suitable topologies that the corresponding algebraic structures are furnished with. A basic ingredient here is the interplay between locality and the spectrum condition. In Chapter 3 we pass to the dual point of view and analyse the resulting continuous functionals. Then, on physical grounds, a certain subclass is distinguished, arising as asymptotic limits of certain functionals constructed from physical states of bounded energy. These limits exhibit properties of singly localized systems (particles). The limiting procedure to be presented here is able to directly reproduce charged systems, in contrast to the LSZ-theory where charge-carrying unobservable operators are necessary.
The representations induced by these asymptotic functionals (the particle weights) are highly reducible, so the obvious task is to work out a disintegration theory in terms of irreducible representations (pure particle weights). This will be done in Chapters 4 and 6. The approach of Chapter 4 makes use of the standard decomposition theory for representations of $`C^{}`$-algebras. To be able to apply this theory, the mathematical structures under consideration have to be adapted to its needs. Great care is taking to ensure that the resulting irreducible representations have all the properties allowing for their interpretation as representatives of elementary particles. As demonstrated by Buchholz , it is then possible to classify the pure particle weights according to their spin and mass even in the case of charged systems. This shows that the notion of particle weights provides a promising approach to the aforementioned infraparticle problem. In Chapter 5 a compactness criterion due to Fredenhagen and Hertel is used to impose certain restrictions on the phase space of quantum field theory. The additional information is used to demonstrate that the particle weight representations of Chapter 4 are locally normal. This implies that one does not lose essential information about the physical systems in the course of the constructions needed to adapt the problem at hand to the needs of spatial disintegration. Chapter 6, again drawing on the mentioned compactness criterion, presents the first steps in an alternative approach to disintegration: Choquet theory. Chapter 7 gives a brief summary.
##### Assumptions of Local Quantum Physics
We collect here the main structural postulates upon which Local Quantum Physics is built in the abstract setting of the algebraic approach , principally in order to fix notation.
* The basis of the present investigations is a net
$$𝒪𝔄(𝒪)$$
(1.1a)
of $`C^{}`$-algebras, which are indexed by the bounded regions $`𝒪`$ in space-time $`^{s+1}`$ and which are *concrete* in the sense that they all belong to the algebra of bounded operators $`𝔅()`$ on a certain Hilbert space $``$. The so-called quasi-local algebra $`𝔄`$ is the $`C^{}`$-inductive limit of the net (1.1a) (cf. \[11, Definition 2.63\]):
$$𝔄\underset{𝒪}{\overset{C^{}}{}}𝔄(𝒪)\text{.}$$
(1.1b)
* On the $`C^{}`$-algebra $`𝔄`$ the symmetry transformations in the inhomogeneous Lorentz group, the Poincaré group $`𝖯_+^{}=𝖫_+^{}^{s+1}`$, are implemented via a strongly continuous group of automorphisms:
$$𝖯_+^{}(\mathrm{\Lambda },x)\alpha _{(\mathrm{\Lambda },x)}\mathrm{Aut}𝔄\text{.}$$
(1.2)
* The net (1.1a) is subject to the following conditions:
+ Isotony: For any two bounded regions $`𝒪_1`$ and $`𝒪_2`$ in $`^{s+1}`$
$$𝒪_1𝒪_2𝔄(𝒪_1)𝔄(𝒪_2)\text{.}$$
(1.3a)
+ Locality: If the bounded regions $`𝒪_1`$ and $`𝒪_2`$ are spacelike separated, i. e., $`𝒪_1`$ belongs to the spacelike complement of $`𝒪_2`$, formally $`𝒪_1𝒪_2^{}`$, then
$$𝔄(𝒪_1)𝔄(𝒪_2)^{}\text{,}$$
(1.3b)
where the prime in (1.3b) denotes the commutant in $`𝔅()`$.
+ Relativistic Covariance: For arbitrary bounded regions $`𝒪`$ and arbitrary transformations $`(\mathrm{\Lambda },x)𝖯_+^{}`$ there hold the relations
$$𝔄(\mathrm{\Lambda }𝒪+x)=\alpha _{(\mathrm{\Lambda },x)}\left(𝔄(𝒪)\right)\text{.}$$
(1.3c)
* The subgroup $`^{s+1}`$ of translations in $`𝖯_+^{}`$ is implemented on $`𝔄`$ by a strongly continuous unitary group, i. e., one which is continuous with respect to the strong-operator topology. These unitaries can be expressed through the (unbounded) generators $`P^\mu `$, $`\mu =1\text{,}\mathrm{}\text{,}s+1`$, of space-time translations according to
$$U(x)=\mathrm{exp}(iP^\mu x_\mu )\text{,}$$
(1.4a)
and, by virtue of (1.2), one has for any $`x^{s+1}`$
$$\alpha _x(A)=U(x)AU(x)^{}\text{,}A𝔄\text{.}$$
(1.4b)
The joint spectrum of the generators $`P^\mu `$, expressed by the pertinent spectral resolution $`E(.)`$ in terms of projections in $`𝔄^{\prime \prime }`$, is supposed to lie in the closed forward light cone
$$\overline{V}_+\{p^{s+1}:pp=p^\mu p_\mu 0\}\text{.}$$
This assumption is known under the term ‘positive-energy representation.’
* Physical states are represented by normalized positive linear functionals on the quasi-local algebra $`𝔄`$, which are normal, i. e., continuous with respect to the $`\sigma `$-weak topology that $`𝔄`$ inherits from $`𝔅()`$. The set of all physical states $`\omega `$ is denoted by $`𝒮`$; it is in one-to-one correspondence to the entirety of all density matrices, the positive trace-class operators in $`𝔅()`$ with unit trace, via
$$\omega (A)=\mathrm{Tr}(\rho _\omega A)\text{,}A𝔄\text{,}$$
(1.5a)
where $`\rho _\omega `$ denotes the unique operator of the above kind. The fact that a physical state $`\omega `$ possesses energy-momentum in the Borel set $`\mathrm{\Delta }^{s+1}`$ is expressed by the condition
$$\omega \left(E(\mathrm{\Delta })\right)=\mathrm{Tr}\left(E(\mathrm{\Delta })\rho _\omega E(\mathrm{\Delta })\right)=1\text{.}$$
(1.5b)
The corresponding subset of $`𝒮`$ is written $`𝒮(\mathrm{\Delta })`$.
At this point, for the sake of clarity, a few remarks concerning topological notions seem advisable. The norm topology on $`𝔄`$ is sometimes called the uniform topology and leaves no room for a possible misunderstanding. The situation is more complicated in case of the term ‘strong continuity:’
* An automorphism group $`\{\alpha _g:g𝒢\}\mathrm{Aut}𝔄`$ on the $`C^{}`$-algebra $`𝔄`$, $`𝒢`$ a topological group, is called strongly continuous if the mapping
$$𝒢g\alpha _g(A)𝔄$$
is continuous for arbitrary $`A𝔄`$ with respect to the initial topology of the group $`𝒢`$ and with respect to the uniform topology of $`𝔄`$.
* A unitary group $`\{U(g):g𝒢\}𝔅()`$, $`𝒢`$ again a topological group, is called strongly continuous if the mapping
$$𝒢gU(g)𝔅()$$
is continuous with respect to the topology of $`𝒢`$ and with respect to the strong-operator topology on $`𝔅()`$.
The term ‘$`\sigma `$-weak topology’ is used to denote the locally convex topology on the algebra $`𝔅()`$ that is defined through the family of seminorms
$$𝒬_{\{\varphi _n,\psi _n\}}:𝔅()_0^+A𝒬_{\{\varphi _n,\psi _n\}}(A)\left|\underset{n=1}{\overset{\mathrm{}}{}}(\varphi _n,A\psi _n)\right|\text{,}$$
where the sequences $`\{\varphi _n\}_n`$ and $`\{\psi _n\}_n`$ of vectors in the Hilbert space $``$ are subject to the conditions $`_{n=1}^{\mathrm{}}\varphi _n^2<\mathrm{}`$ and $`_{n=1}^{\mathrm{}}\psi _n^2<\mathrm{}`$. This designation is synonymous with ‘ultra-weak topology.’ Mappings which are continuous with respect to this topology are called normal.
## Chapter 2 Localizing Operators and Spectral Seminorms
The results presented in Chapters 2 and 3 have been worked out in close collaboration with Detlev Buchholz, whose ideas, as set out in , constituted the foundation. Their somewhat complicated presentation is the author’s responsibility. The particle concept to be set forth in the sequel is motivated by the experimental situation encountered in high energy physics where certain physical systems show up as ‘particles,’ being traced by specific measuring devices called ‘detectors.’ The common characteristic of these physical systems is that they are localized in the course of the measuring process. Haag and Kastler stated in their fundamental article on algebraic quantum field theory that ‘…ultimately all physical processes are analyzed in terms of geometric relations of unresolved phenomena,’ emphasizing localization as the very nature of all measurements. To represent the experimental set-up in the framework of the algebraic approach to local quantum physics elements of the quasi-local algebra $`𝔄`$ have to be singled out first that exhibit properties of particle detectors.
### 2.1 The Algebra of Detectors
As argued by Araki and Haag a particle detector $`C𝔄`$ should be insensitive to the vacuum $`\mathrm{\Omega }`$: $`C\mathrm{\Omega }=0`$. In view of the actual experimental situation one can be more specific, noting that a minimal energy, depending on the detector used, has to be deposited to produce a signal. In the present thesis we shall therefore work with a smaller class of operators: the algebraic representatives corresponding to a particle counter are to annihilate all physical states with bounded energy below a specific threshold, to be precise. Now, on account of the Reeh–Schlieder-Theorem, this feature is incompatible with locality since an algebra pertaining to a region $`𝒪`$ with non-void causal complement $`𝒪^{}`$ does not contain any operator annihilating states of bounded energy (cf. ). As a consequence, the operators which comply with the above annihilation property cannot be strictly local; instead their localization has to be weakened. This is done in a way that resembles the introduction of rapidly decreasing functions on $`^n`$: the operators in question are not contained in a local algebra, but they are almost local in the sense of the following definition (‘quasilocal of infinite order’ is the designation used in ).
###### Definition 2.1 (Almost Locality).
Let $`𝒪_r\{(x^0,𝒙)^{𝒔+\mathit{1}}:|𝒙^\mathit{0}|+|𝒙|<𝒓\}`$, $`r>0`$, denote the double cone (standard diamond) with basis $`𝑶_𝒓\{𝒙^𝒔:|𝒙|<𝒓\}`$. An operator $`A𝔄`$ is called almost local if there exists a net $`\{A_r𝔄(𝒪_r):r>0\}`$ of local operators such that
$$\underset{r\mathrm{}}{lim}r^kAA_r=0$$
(2.1)
for any $`k_0`$. The set of almost local operators is a -subalgebra of $`𝔄`$ denoted by $`𝔄_𝒮`$.
###### Remark.
* Let $`A`$ and $`B`$ be almost local operators with approximating nets of local operators $`\{A_r𝔄(𝒪_r):r>0\}`$ and $`\{B_r𝔄(𝒪_r):r>0\}`$, respectively. Then, since $`𝒪_r`$ and $`𝒪_r+2𝒙`$ are spacelike separated for $`r|𝒙|`$ so that the associated algebras $`𝔄(𝒪_r)`$ and $`𝔄(𝒪_r+2𝒙)`$ commute, the following estimate holds for any $`𝒙^𝒔\{\mathit{0}\}`$
$$[\alpha _{2𝒙}(A),B]2\left(AA_{|𝒙|}B+AA_{|𝒙|}BB_{|𝒙|}+ABB_{|𝒙|}\right)$$
(2.2a)
The right-hand side of this inequality is bounded and falls off more rapidly than any power of $`|𝒙|^\mathit{1}`$, therefore the continuous mapping $`^s𝒙[\alpha _𝒙(𝑨),𝑩]`$ turns out to be integrable:
$$_^sd^sx[\alpha x(A),B]<\mathrm{}\text{.}$$
(2.2b)
* The approximating net of local operators $`\{A_r𝔄(𝒪_r):r>0\}`$ for $`A𝔄_𝒮`$ can be used to construct a second approximating net $`\{A_r^{}𝔄(𝒪_r):r>0\}`$ with the additional property $`A_r^{}A`$ for any $`r>0`$, which at the same time is subject to the inequality $`AA_r^{}2AA_r`$ and thus satisfies condition (2.1) for almost locality. Estimates of this kind will later on turn out to be important in solving the problem of existence of uniform bounds for integrals of the form (2.2b), evaluated for sequences or even nets of almost local operators. With approximating nets of local operators of this special kind the estimate (2.2a) can be improved for arbitrary $`A\text{,}B𝔄_𝒮`$ to yield
$$[\alpha _{2𝒙}(A),B]2\left(AA_{|𝒙|}B+ABB_{|𝒙|}\right)\text{,}𝒙^𝒔\{\mathit{0}\}\text{.}$$
(2.2c)
The feature of annihilating states of bounded energy below a certain threshold is called vacuum annihilation property in the sequel and finds its rigorous mathematical expression in the following definition.
###### Definition 2.2 (Vacuum Annihilation Property).
An operator $`A𝔄`$ is said to have the vacuum annihilation property if, in the sense of operator-valued distributions, the mapping
$$^{s+1}x\alpha _x(A)U(x)AU(x)^{}𝔄$$
(2.3)
has a Fourier transform with compact support $`\mathrm{\Gamma }`$ contained in the complement of the forward light cone $`\overline{V}_+`$. The collection of all vacuum annihilation operators is a subspace of $`𝔄`$ denoted $`𝔄_{\text{ann}}`$.
###### Remark.
The support of the Fourier transform of (2.3) is precisely the energy-momentum transfer of $`A`$, and the energy-threshold for the annihilation of states depends on the distance $`d(\mathrm{\Gamma },\overline{V}_+)`$ between $`\mathrm{\Gamma }`$ and $`\overline{V}_+`$. Let $`\mathrm{\Gamma }_0`$ be a closed subset of $`^{s+1}`$ and let $`\stackrel{~}{𝔄}(\mathrm{\Gamma }_0)`$ denote the set of all operators $`A𝔄`$ having energy-momentum transfer $`\mathrm{\Gamma }_A\mathrm{\Gamma }_0`$. Then $`\stackrel{~}{𝔄}(\mathrm{\Gamma }_0)`$ is easily seen to be a uniformly closed subspace of $`𝔄`$, invariant under space-time translations.
The construction of a subalgebra $``$ in $`𝔄`$ containing self-adjoint operators to be interpreted as representatives of particle detectors is accomplished in three steps (Definitions 2.32.5), starting with a subspace $`𝔏_0𝔄`$ consisting of operators which, in addition to the properties mentioned above, are infinitely often differentiable with respect to the automorphism group $`\{\alpha _{(\mathrm{\Lambda },x)}:(\mathrm{\Lambda },x)𝖯_+^{}\}`$ (cf. Definition A.12 in Appendix A).
###### Definition 2.3.
The almost local vacuum annihilation operators $`L_0𝔄`$ which are infinitely often differentiable with respect to the group $`\{\alpha _{(\mathrm{\Lambda },x)}:(\mathrm{\Lambda },x)𝖯_+^{}\}`$ constitute a subspace $`𝔄_𝒮𝔄_{\text{ann}}𝒟^{(\mathrm{})}(𝔄)`$ of $`𝔄`$. The intersection of this set with all the pre-images of $`𝔄_𝒮`$ under arbitrary products of partial derivations $`\delta ^{k_1}\mathrm{}\delta ^{k_N}`$ for any $`N`$ and any $`1k_id_𝖯`$, $`d_𝖯`$ the dimension of $`𝖯_+^{}`$, is again a linear space denoted $`𝔏_0`$. Explicitly, $`𝔏_0`$ consists of all almost local vacuum annihilation operators which are infinitely often differentiable, having *almost local* partial derivatives of any order.
###### Remark.
* The space $`𝔏_0`$ is stable under the action of the Poincaré group. This means that $`\alpha _{(\mathrm{\Lambda },x)}(𝔏_0)=𝔏_0`$ for any $`(\mathrm{\Lambda },x)𝖯_+^{}`$. Due to the properties of Fourier transformation, $`\alpha _{(\mathrm{\Lambda },x)}(L_0)`$ has energy-momentum transfer in $`\mathrm{\Lambda }\mathrm{\Gamma }`$ if $`L_0𝔏_0(\mathrm{\Gamma })𝔏_0\stackrel{~}{𝔄}(\mathrm{\Gamma })`$; the adjoint $`L_{0}^{}{}_{}{}^{}`$ of this $`L_0`$ belongs to $`\stackrel{~}{𝔄}(\mathrm{\Gamma })`$.
* Furthermore $`𝔏_0`$ is invariant under differentiation: The partial derivatives are almost local and infinitely often differentiable operators by definition, and, as uniform limits of vacuum annihilation operators, they inherit the energy-momentum transfer of these so that they belong to $`𝔄_{\text{ann}}`$, too.
A huge number of elements of $`𝔏_0`$ can be constructed by regularizing almost local operators with respect to rapidly decreasing functions on the Poincaré group. The semi-direct product Lie group $`𝖯_+^{}=𝖫_+^{}^{s+1}`$ is unimodular by \[45, Proposition II.29 and Corollary\] since $`𝖫_+^{}`$ is a simple thus semisimple Lie group \[36, Proposition I.1.6\]. So let $`\mu `$ be the Haar measure on $`𝖯_+^{}`$ and $`A𝔄_𝒮`$, then the operator
$$A(F)=𝑑\mu (\mathrm{\Lambda },x)F(\mathrm{\Lambda },x)\alpha _{(\mathrm{\Lambda },x)}(A)$$
(2.4)
belongs to $`𝔏_0(\mathrm{\Gamma })`$ if the infinitely differentiable function $`F`$ is rapidly decreasing on the subgroup $`^{s+1}`$ and compactly supported on $`𝖫_+^{}`$, i. e. $`F𝒮_0\left(𝖯_+^{}\right)=𝒮_0\left(𝖫_+^{}^{s+1}\right)`$ in the notation introduced in , and has the additional property that the Fourier transforms of the partial functions $`F_\mathrm{\Lambda }(.)F(\mathrm{\Lambda },.)`$ have common support in the compact set $`\mathrm{\Gamma }\mathrm{}\overline{V}_+`$ for any $`\mathrm{\Lambda }𝖫_+^{}`$.
The following definition specifies a left ideal $`𝔏`$ of the algebra $`𝔄`$.
###### Definition 2.4.
Let $`𝔏`$ denote the linear span of all operators $`L𝔄`$ of the form $`L=AL_0`$ where $`A𝔄`$ and $`L_0𝔏_0`$; i. e.
$$𝔏𝔄𝔏_0=\mathrm{span}\left\{AL_0:A𝔄,L_0𝔏_0\right\}\text{.}$$
Then $`𝔏`$ is a left ideal of $`𝔄`$, called the ‘left ideal of localizing operators.’
By their very construction, the elements of $`𝔏`$ annihilate the vacuum and all states of bounded energy below a certain threshold that depends on the minimum of $`d(\mathrm{\Gamma }_i,\overline{V}_+)`$, $`i=1\text{,}\mathrm{}\text{,}N`$, with respect to all representations $`L=_{i=1}^NA_iL_i𝔏`$, where $`\mathrm{\Gamma }_i`$ is the energy-momentum transfer of $`L_i`$. The algebra of operators whose self-adjoint elements are to be interpreted as representatives of particle detectors is laid down in the next definition.
###### Definition 2.5.
Let $``$ denote the linear span of all operators $`C𝔄`$ which can be represented in the form $`C=L_{1}^{}{}_{}{}^{}L_2`$ with $`L_1\text{,}L_2𝔏`$; i. e.
$$𝔏^{}𝔏=\mathrm{span}\left\{L_{1}^{}{}_{}{}^{}L_2:L_1,L_2𝔏\right\}\text{.}$$
Then $``$ is a -subalgebra of $`𝔄`$, called the ‘algebra of detectors.’
###### Remark.
The algebra $``$ is smaller than that used by Araki and Haag in . It is not closed in the uniform topology of $`𝔄`$ and does not contain a unit.
### 2.2 Spectral Seminorms on the Algebra of Detectors
The analysis of physical states is performed by use of the algebra of detectors $``$. In a state $`\omega `$ of bounded energy $`E`$ we expect to encounter a finite number of localization centres, since the triggering of a detector $`C`$ requires a minimal energy $`ϵ`$ to be deposited, the number $`N`$ of localization centres being equal to or less than $`E/ϵ`$. Now, according to this heuristic picture, placing the counter $`C`$ for given time $`t`$ at every point $`𝒙^𝒔`$ and adding up the corresponding expectation values $`\omega \left(\alpha _{(t,𝒙)}(C)\right)`$ should result in the finite integral
$$_^sd^sx\left|\omega \left(\alpha _{(t,𝒙)}(C)\right)\right|<\mathrm{}\text{.}$$
(2.5)
As a matter of fact, the operators $`C`$ turn out to have the property (2.5) as was shown by Buchholz in . For the sake of completeness and to demonstrate how phase-space properties of the theory (localization in space combined with energy-bounds) enter the present investigation, we give an elaborate proof.
###### Proposition 2.6.
Let $`E(.)`$ be the spectral resolution of the space-time translations $`U(x)`$, $`x^{s+1}`$, and let $`L_0𝔏_0`$ have energy-momentum transfer $`\mathrm{\Gamma }`$ in a convex subset of $`\mathrm{}\overline{V}_+`$. Then for any bounded Borel set $`\mathrm{\Delta }^{s+1}`$ the net of operator-valued Bochner integrals indexed by compact $`𝐊^𝐬`$,
$$\begin{array}{cc}\hfill Q_{\mathrm{\Delta },𝑲}^{(L_{0}^{}{}_{}{}^{}L_0)}& E(\mathrm{\Delta })Kd^sx\alpha x(L_{0}^{}{}_{}{}^{}L_0)E(\mathrm{\Delta })\hfill \\ & =Kd^sxE(\mathrm{\Delta })\alpha x(L_{0}^{}{}_{}{}^{}L_0)E(\mathrm{\Delta })\text{,}\hfill \end{array}$$
is $`\sigma `$-strongly convergent as $`𝐊^𝐬`$ and the limit $`Q_\mathrm{\Delta }^{(L_{0}^{}{}_{}{}^{}L_0)}𝔅()^+`$ satisfies the estimate
$$Q_\mathrm{\Delta }^{(L_{0}^{}{}_{}{}^{}L_0)}N(\mathrm{\Delta },\mathrm{\Gamma })_^sd^sx[\alpha x(L_0),L_{0}^{}{}_{}{}^{}]$$
(2.6)
for suitable $`N(\mathrm{\Delta },\mathrm{\Gamma })`$, depending on $`\mathrm{\Delta }`$ and $`\mathrm{\Gamma }`$. Moreover the mapping
$$𝒙𝑬(𝜟)\alpha 𝒙(𝑳_{\mathit{0}}^{}{}_{}{}^{}𝑳_\mathit{0})𝑬(𝜟)$$
is integrable with respect to the $`\sigma `$-weak topology on $`𝔅()`$ and its integral coincides with the operator $`Q_\mathrm{\Delta }^{(L_{0}^{}{}_{}{}^{}L_0)}`$:
$$Q_\mathrm{\Delta }^{(L_{0}^{}{}_{}{}^{}L_0)}=\sigma \mathrm{weak}_^sd^sxE(\mathrm{\Delta })\alpha x(L_{0}^{}{}_{}{}^{}L_0)E(\mathrm{\Delta })\text{.}$$
###### Proof.
$`\mathrm{\Delta }`$ being a bounded Borel set, the same is true of its closure $`\overline{\mathrm{\Delta }}`$, so that, due to compactness and convexity of $`\mathrm{\Gamma }`$, there exists a number $`n`$ for which the relation $`(\overline{\mathrm{\Delta }}+\mathrm{\Gamma }_n)\overline{V}_+=\mathrm{}`$ is satisfied, where $`\mathrm{\Gamma }_n`$ denotes the sum $`\mathrm{\Gamma }_n\mathrm{\Gamma }+\mathrm{}+\mathrm{\Gamma }`$ with $`n`$ terms. The spectrum condition then entails:
$$E(\overline{\mathrm{\Delta }}+\mathrm{\Gamma }_n)=0\text{.}$$
(2.7)
Note, that in the derivation of this result compactness of $`\mathrm{\Gamma }`$ is needed to ensure that the distance between $`\mathrm{\Gamma }`$ and $`\overline{V}_+`$ is positive; other shapes of $`\mathrm{\Gamma }`$ are possible as long as convexity and the condition $`d(\mathrm{\Gamma },\overline{V}_+)>0`$ are preserved, e. g. wedges in $`\mathrm{}\overline{V}_+`$. For arbitrary $`𝒙_\mathit{1}\text{,}\mathrm{}\text{,}𝒙_𝒏^𝒔`$ all the operators $`\alpha _{𝒙_𝒊}(L_0)`$, $`i=1\text{,}\mathrm{}\text{,}n`$, belong to $`\stackrel{~}{𝔄}(\mathrm{\Gamma })`$ whilst their product $`_{i=1}^n\alpha _{𝒙_𝒊}(L_0)`$ is an element of $`\stackrel{~}{𝔄}(\mathrm{\Gamma }_n)`$, hence by (2.7)
$$\underset{i=1}{\overset{n}{}}\alpha _{𝒙_𝒊}(L_0)E(\mathrm{\Delta })=E(\overline{\mathrm{\Delta }}+\mathrm{\Gamma }_n)\underset{i=1}{\overset{n}{}}\alpha _{𝒙_𝒊}(L_0)E(\mathrm{\Delta })=0\text{.}$$
(2.8)
Now, \[15, Lemma 2.2\] states that for any $`B𝔅()`$ and any $`k`$
$$P_kKd^sx\alpha x(B^{}B)P_k(k1)\underset{\mathrm{\Psi }}{sup}\left(_{𝑲𝑲}d^sx[\alpha x(B),B^{}]\mathrm{\Psi }\right)\text{,}$$
(2.9)
where $`P_k`$ is the orthogonal projection onto the intersection of the kernels of $`k`$-fold products $`_{i=1}^k\alpha _{𝒚_𝒊}(B)`$ for arbitrary $`𝒚_\mathit{1}\text{,}\mathrm{}\text{,}𝒚_𝒌^𝒔`$, $`𝑲^𝒔`$ is compact and the supremum extends over all unit vectors $`\mathrm{\Psi }P_{k1}`$. According to (2.8) $`E(\mathrm{\Delta })P_n`$ if we take $`BL_0`$, so that the following estimate, uniform in $`𝑲`$, is a consequence of (2.9) combined with almost locality of $`L_0`$ (cf. (2.2b)):
$$Q_{\mathrm{\Delta },𝑲}^{(L_{0}^{}{}_{}{}^{}L_0)}=E(\mathrm{\Delta })Kd^sx\alpha x(L_{0}^{}{}_{}{}^{}L_0)E(\mathrm{\Delta })(n1)_^sd^sx[\alpha x(L_0),L_{0}^{}{}_{}{}^{}]\text{.}$$
(2.10)
The positive operators $`\{Q_{\mathrm{\Delta },𝑲}^{(L_{0}^{}{}_{}{}^{}L_0)}:𝑲^𝒔\text{compact}\}`$ thus constitute an increasing net which is bounded in $`𝔅()^+`$. According to \[11, Lemma 2.4.19\] this net has a least upper bound in $`𝔅()^+`$, which is its $`\sigma `$-strong limit $`Q_\mathrm{\Delta }^{(L_{0}^{}{}_{}{}^{}L_0)}`$ and satisfies
$$Q_\mathrm{\Delta }^{(L_{0}^{}{}_{}{}^{}L_0)}(n1)_^sd^sx[\alpha x(L_0),L_{0}^{}{}_{}{}^{}]\text{.}$$
(2.11)
For $`N(\mathrm{\Delta },\mathrm{\Gamma })n1`$ this is the desired estimate (2.6).
The $`\sigma `$-weak topology of $`𝔅()`$ is induced by the positive normal functionals of the space $`𝔅()_{}^+`$, so that integrability of $`𝒙𝑬(𝜟)\alpha 𝒙(𝑳_{\mathit{0}}^{}{}_{}{}^{}𝑳_\mathit{0})𝑬(𝜟)`$ in the $`\sigma `$-weak topology is implied by integrability of the functions
$$𝒙\left|\psi \left(𝑬(𝜟)\alpha 𝒙(𝑳_{\mathit{0}}^{}{}_{}{}^{}𝑳_\mathit{0})𝑬(𝜟)\right)\right|=\psi \left(𝑬(𝜟)\alpha 𝒙(𝑳_{\mathit{0}}^{}{}_{}{}^{}𝑳_\mathit{0})𝑬(𝜟)\right)$$
for any $`\psi 𝔅()_{}^+`$. Now, given any compact subset $`𝑲`$ of $`^s`$, there holds the estimate
$$\begin{array}{c}Kd^sx\left|\psi \left(E(\mathrm{\Delta })\alpha x(L_{0}^{}{}_{}{}^{}L_0)E(\mathrm{\Delta })\right)\right|=Kd^sx\psi \left(E(\mathrm{\Delta })\alpha x(L_{0}^{}{}_{}{}^{}L_0)E(\mathrm{\Delta })\right)\hfill \\ \hfill =\psi \left(Kd^sxE(\mathrm{\Delta })\alpha x(L_{0}^{}{}_{}{}^{}L_0)E(\mathrm{\Delta })\right)\psi Q_{\mathrm{\Delta },𝑲}^{(L_{0}^{}{}_{}{}^{}L_0)}\psi Q_\mathrm{\Delta }^{(L_{0}^{}{}_{}{}^{}L_0)}\text{,}\end{array}$$
and, as a consequence of the Monotone Convergence Theorem \[26, II.2.7\], the functions $`𝒙\left|\psi \left(𝑬(𝜟)\alpha 𝒙(𝑳_{\mathit{0}}^{}{}_{}{}^{}𝑳_\mathit{0})𝑬(𝜟)\right)\right|`$ indeed turn out to be integrable for any $`\psi 𝔅()_{}^+`$. Thus the integral of the mapping $`𝒙𝑬(𝜟)\alpha 𝒙(𝑳_{\mathit{0}}^{}{}_{}{}^{}𝑳_\mathit{0})𝑬(𝜟)`$ with respect to the $`\sigma `$-weak topology exists (cf. \[26, II.6.2\]) and, through an application of Lebesgue’s Dominated Convergence Theorem \[26, II.5.6\], is seen to be the $`\sigma `$-weak limit of the net of operators $`Q_{\mathrm{\Delta },𝑲}^{(L_{0}^{}{}_{}{}^{}L_0)}`$ which coincides with the $`\sigma `$-strong limit $`Q_\mathrm{\Delta }^{(L_{0}^{}{}_{}{}^{}L_0)}`$ established above. Formally
$$Q_\mathrm{\Delta }^{(L_{0}^{}{}_{}{}^{}L_0)}=\sigma \mathrm{weak}_^sd^sxE(\mathrm{\Delta })\alpha x(L_{0}^{}{}_{}{}^{}L_0)E(\mathrm{\Delta })\text{,}$$
which is the last of the above assertions. ∎
###### Proposition 2.7.
Suppose that $`\mathrm{\Delta }^{s+1}`$ is a bounded Borel set.
* Let $`L𝔏`$ be arbitrary, then the net of operators for compact $`𝑲^𝒔`$
$$\begin{array}{cc}\hfill Q_{\mathrm{\Delta },𝑲}^{(L^{}L)}& E(\mathrm{\Delta })Kd^sx\alpha x(L^{}L)E(\mathrm{\Delta })\hfill \\ & =Kd^sxE(\mathrm{\Delta })\alpha x(L^{}L)E(\mathrm{\Delta })\text{,}\hfill \end{array}$$
converges $`\sigma `$-strongly to $`Q_\mathrm{\Delta }^{(L^{}L)}𝔅()^+`$ in the limit $`𝑲^𝒔`$. Moreover the mapping $`𝒙𝑬(𝜟)\alpha 𝒙(𝑳^{}𝑳)𝑬(𝜟)`$ is integrable with respect to the $`\sigma `$-weak topology on $`𝔅()`$ and satisfies
$$Q_\mathrm{\Delta }^{(L^{}L)}=\sigma \mathrm{weak}_^sd^sxE(\mathrm{\Delta })\alpha x(L^{}L)E(\mathrm{\Delta })\text{.}$$
* Let $`C`$ be arbitrary, then the net of operators indexed by compact $`𝑲^𝒔`$
$$\begin{array}{cc}\hfill Q_{\mathrm{\Delta },𝑲}^{(C)}& E(\mathrm{\Delta })Kd^sx\alpha x(C)E(\mathrm{\Delta })\hfill \\ & =Kd^sxE(\mathrm{\Delta })\alpha x(C)E(\mathrm{\Delta })\hfill \end{array}$$
is $`\sigma `$-strongly convergent to $`Q_\mathrm{\Delta }^{(C)}𝔅()`$ for $`𝑲^𝒔`$. In addition to this the mapping $`𝒙𝑬(𝜟)\alpha 𝒙(𝑪)𝑬(𝜟)`$ is integrable with respect to the $`\sigma `$-weak topology on $`𝔅()`$ and the integral is given by
$$Q_\mathrm{\Delta }^{(C)}=\sigma \mathrm{weak}_^sd^sxE(\mathrm{\Delta })\alpha x(C)E(\mathrm{\Delta })\text{.}$$
Furthermore
$$sup\{_^sd^sx\left|\varphi \left(E(\mathrm{\Delta })\alpha x(C)E(\mathrm{\Delta })\right)\right|:\varphi 𝔅()_{,1}\}<\mathrm{}\text{.}$$
(2.12)
###### Remark.
Note, that relation (2.12) is a sharpened version of (2.5) which, based on heuristic considerations, was the starting point of the present investigation.
###### Proof.
* By partition of unity (cf. \[40, Satz 8.1\]), applied to elements of $`𝔏_0`$ which have arbitrary energy-momentum transfer in $`\mathrm{}\overline{V}_+`$, any $`L𝔏`$ can be written as a finite sum $`L=_{j=1}^mA_jL_j`$ where the $`A_j`$ belong to $`𝔄`$ and the operators $`L_j𝔏_0`$ have energy-momentum transfer in compact and convex subsets $`\mathrm{\Gamma }_j`$ of $`\mathrm{}\overline{V}_+`$. Since
$$L^{}L2^{m1}\left(\underset{1jm}{sup}A_j^2\right)\underset{j=1}{\overset{m}{}}L_{j}^{}{}_{}{}^{}L_j\text{,}$$
we infer
$$Q_{\mathrm{\Delta },𝑲}^{(L^{}L)}2^{m1}\left(\underset{1jm}{sup}A_j^2\right)\underset{j=1}{\overset{m}{}}Q_{\mathrm{\Delta },𝑲}^{(L_{j}^{}{}_{}{}^{}L_j)}\text{,}$$
so that by (2.10) the increasing net $`\{Q_{\mathrm{\Delta },𝑲}^{(L^{}L)}:𝑲^𝒔\mathrm{compact}\}`$ turns out to be bounded, having a least upper bound in $`𝔅()^+`$ that is its $`\sigma `$-strong limit $`Q_\mathrm{\Delta }^{(L^{}L)}`$. Making again use of the above order relation for $`L^{}L`$ one arrives at
$$\psi \left(E(\mathrm{\Delta })\alpha x(L^{}L)E(\mathrm{\Delta })\right)2^{m1}\left(\underset{1jm}{sup}A_j^2\right)\underset{j=1}{\overset{m}{}}\psi \left(E(\mathrm{\Delta })\alpha x(L_{j}^{}{}_{}{}^{}L_j)E(\mathrm{\Delta })\right)$$
for any $`\psi 𝔅()_{}^+`$ and any $`𝒙^𝒔`$, where the right-hand side of this relation is integrable as was shown in the proof of Proposition 2.6. Then the reasoning applied there establishes the $`\sigma `$-weak integrability of $`𝒙𝑬(𝜟)\alpha 𝒙(𝑳^{}𝑳)𝑬(𝜟)`$ together with the relation
$$Q_\mathrm{\Delta }^{(L^{}L)}=\sigma \mathrm{weak}_^sd^sxE(\mathrm{\Delta })\alpha x(L^{}L)E(\mathrm{\Delta })\text{.}$$
* Consider $`C_0=L_{1}^{}{}_{}{}^{}L_2`$ with $`L_1\text{,}L_2𝔏`$. By polarization
$$C_0=\frac{1}{4}\underset{k=0}{\overset{3}{}}i^k(L_1+i^kL_2)^{}(L_1+i^kL_2)=\frac{1}{4}\underset{k=0}{\overset{3}{}}i^kL_{}^{(k)}{}_{}{}^{}L^{(k)}\text{,}$$
where $`L^{(k)}L_1+i^kL_2𝔏`$ for $`k=0\text{,}\mathrm{}\text{,}3`$, and according to (i)
$$\begin{array}{c}Q_{\mathrm{\Delta },𝑲}^{(C_0)}=E(\mathrm{\Delta })Kd^sx\alpha x(C_0)E(\mathrm{\Delta })\hfill \\ \hfill =\frac{1}{4}\underset{k=0}{\overset{3}{}}i^k\left(E(\mathrm{\Delta })Kd^sx\alpha x\left(L_{}^{(k)}{}_{}{}^{}L^{(k)}\right)E(\mathrm{\Delta })\right)=\frac{1}{4}\underset{k=0}{\overset{3}{}}i^kQ_{\mathrm{\Delta },𝑲}^{\left(L_{}^{(k)}{}_{}{}^{}L^{(k)}\right)}\end{array}$$
converges $`\sigma `$-strongly to
$$Q_\mathrm{\Delta }^{(C_0)}\frac{1}{4}\underset{k=0}{\overset{3}{}}i^kQ_\mathrm{\Delta }^{\left(L_{}^{(k)}{}_{}{}^{}L^{(k)}\right)}\text{.}$$
(2.13)
Now, let $`\varphi `$ be a normal functional on $`𝔅()`$. By polar decomposition (cf. \[54, Theorem III.4.2(i), Proposition III.4.6\]) there exist a partial isometry $`V𝔅()`$ and a positive normal functional $`|\varphi |`$ subject to the relation $`|\varphi |=\varphi `$, such that $`\varphi (.)=|\varphi |(.V)`$, allowing for the following estimate ($`𝒙^𝒔`$ arbitrary):
$$\begin{array}{c}2\left|\varphi \left(E(\mathrm{\Delta })\alpha x(C_0)E(\mathrm{\Delta })\right)\right|=2\left||\varphi |\left(E(\mathrm{\Delta })\alpha x(L_{1}^{}{}_{}{}^{}L_2)E(\mathrm{\Delta })V\right)\right|\hfill \\ \hfill 2\sqrt{|\varphi |\left(E(\mathrm{\Delta })\alpha x(L_{1}^{}{}_{}{}^{}L_1)E(\mathrm{\Delta })\right)}\sqrt{|\varphi |\left(V^{}E(\mathrm{\Delta })\alpha x(L_{2}^{}{}_{}{}^{}L_2)E(\mathrm{\Delta })V\right)}\\ \hfill =\underset{\lambda >0}{inf}\left(\lambda ^1|\varphi |\left(E(\mathrm{\Delta })\alpha x(L_{1}^{}{}_{}{}^{}L_1)E(\mathrm{\Delta })\right)+\lambda |\varphi |\left(V^{}E(\mathrm{\Delta })\alpha x(L_{2}^{}{}_{}{}^{}L_2)E(\mathrm{\Delta })V\right)\right)\text{,}\end{array}$$
where we made use of the fact that $`2\sqrt{ab}=inf_{\lambda >0}(\lambda ^1a+\lambda b)`$ for any $`a\text{,}b0`$. Now, from the first part of this Proposition we infer that it is possible to integrate the above expression over all of $`^s`$ to get for any $`\lambda >0`$ the estimate
$$2_^sd^sx\left|\varphi \left(E(\mathrm{\Delta })\alpha x(C_0)E(\mathrm{\Delta })\right)\right|\lambda ^1\varphi Q_\mathrm{\Delta }^{(L_{1}^{}{}_{}{}^{}L_1)}+\lambda \varphi Q_\mathrm{\Delta }^{(L_{2}^{}{}_{}{}^{}L_2)}\text{.}$$
Note, that the normal functionals $`\varphi `$ and the $`\sigma `$-weak integrals commute due to \[26, Proposition II.5.7 adapted to integrals in locally convex spaces\]. Taking the infimum with respect to $`\lambda `$ one finally arrives at
$$_^sd^sx\left|\varphi \left(E(\mathrm{\Delta })\alpha x(C_0)E(\mathrm{\Delta })\right)\right|\varphi Q_\mathrm{\Delta }^{(L_{1}^{}{}_{}{}^{}L_1)}^{1/2}Q_\mathrm{\Delta }^{(L_{2}^{}{}_{}{}^{}L_2)}^{1/2}\text{.}$$
(2.14)
This relation is valid for any normal functional in $`𝔅()_{}`$, so that the $`\sigma `$-weak integrability of $`𝒙𝑬(𝜟)\alpha 𝒙(𝑪_\mathit{0})𝑬(𝜟)`$ is established, the relation
$$Q_\mathrm{\Delta }^{(C_0)}=\sigma \mathrm{weak}_^sd^sxE(\mathrm{\Delta })\alpha x(C_0)E(\mathrm{\Delta })$$
(2.15)
being an immediate consequence (cf. the proof of Proposition 2.6). Another fact implied by inequality (2.14) is the estimate
$$sup\{_^sd^sx\left|\varphi \left(E(\mathrm{\Delta })\alpha x(C_0)E(\mathrm{\Delta })\right)\right|:\varphi 𝔅()_{,1}\}Q_\mathrm{\Delta }^{(L_{1}^{}{}_{}{}^{}L_1)}^{1/2}Q_\mathrm{\Delta }^{(L_{2}^{}{}_{}{}^{}L_2)}^{1/2}\text{.}$$
(2.16)
Since any $`C`$ is a linear combination of operators of the form $`C_0`$, the above relations (2.13) through (2.16) are easily generalized to establish the second part of the Proposition. ∎
The preceding result suggests the introduction of topologies on the left ideal $`𝔏`$ and on the -algebra $``$, respectively, using specific seminorms indexed by bounded Borel subsets $`\mathrm{\Delta }`$ of $`^{s+1}`$.
###### Definition 2.8.
* The left ideal $`𝔏`$ is equipped with a family of seminorms $`q_\mathrm{\Delta }`$ via
$$q_\mathrm{\Delta }(L)Q_\mathrm{\Delta }^{(L^{}L)}^{1/2}\text{,}L𝔏\text{.}$$
(2.17a)
* The -algebra $``$ is furnished with seminorms $`p_\mathrm{\Delta }`$ by assigning
$$p_\mathrm{\Delta }(C)sup\{_^sd^sx\left|\varphi \left(E(\mathrm{\Delta })\alpha x(C)E(\mathrm{\Delta })\right)\right|:\varphi 𝔅()_{,1}\}\text{,}C\text{.}$$
(2.17b)
* The completions of the locally convex (Hausdorff) spaces $`(𝔏,𝔗_q)`$ and $`(,𝔗_p)`$ arising from topologization by these seminorms are denoted $`(\overline{𝔏},\overline{𝔗}_q)`$ and $`(\overline{},\overline{𝔗}_p)`$, respectively. Accordingly, the complete locally convex subspace of $`\overline{𝔏}`$ generated by $`𝔏_0`$ is designated as $`(\overline{𝔏}_0,\overline{𝔗}_q)`$.
* The completions of the locally convex spaces $`(𝔏,𝔗_q^u)`$ and $`(,𝔗_p^u)`$ arising from topologization by all the seminorms $`q_\mathrm{\Delta }`$ and $`p_\mathrm{\Delta }`$, respectively, together with the initial uniform (norm) topology inherited from the quasi-local algebra $`𝔄`$ are denoted $`(𝔄_𝔏,\overline{𝔗}_q^u)`$ and $`(𝔄_{},\overline{𝔗}_p^u)`$.
###### Remark.
* Let $`𝔅()_{}^+`$ denote the positive cone in $`𝔅()_{}`$, then for any $`L𝔏`$
$$q_\mathrm{\Delta }(L)^2=sup\{_^sd^sx\omega \left(E(\mathrm{\Delta })\alpha x(L^{}L)E(\mathrm{\Delta })\right):\omega 𝔅()_{,1}^+\}\text{,}$$
(2.17c)
a formulation that will frequently be used.
* The seminorm properties of $`q_\mathrm{\Delta }`$ and $`p_\mathrm{\Delta }`$ are easily checked. To establish the subadditivity of $`q_\mathrm{\Delta }`$ one has to observe that
$$\begin{array}{c}q_\mathrm{\Delta }(L_1+L_2)^2q_\mathrm{\Delta }(L_1)^2+\underset{\lambda >0}{inf}\left[\lambda ^1q_\mathrm{\Delta }(L_1)^2+\lambda q_\mathrm{\Delta }(L_2)^2\right]+q_\mathrm{\Delta }(L_2)^2\hfill \\ \hfill =q_\mathrm{\Delta }(L_1)^2+2q_\mathrm{\Delta }(L_1)q_\mathrm{\Delta }(L_2)+q_\mathrm{\Delta }(L_2)^2=\left(q_\mathrm{\Delta }(L_1)+q_\mathrm{\Delta }(L_2)\right)^2\text{,}\end{array}$$
where we made use of the fact that $`L_{1}^{}{}_{}{}^{}L_2+L_{2}^{}{}_{}{}^{}L_1\lambda ^1L_{1}^{}{}_{}{}^{}L_1+\lambda L_{2}^{}{}_{}{}^{}L_2`$ for any $`\lambda >0`$ and $`L_1\text{,}L_2𝔏`$.
* The Hausdorff property of the locally convex spaces $`(𝔏,𝔗_q)`$ and $`(,𝔗_p)`$ can be established using the fact that vectors corresponding to states of bounded energy constitute a dense subspace of $``$. From the very definition of the seminorms $`q_\mathrm{\Delta }`$ and $`p_\mathrm{\Delta }`$ we infer that the conditions $`q_\mathrm{\Delta }(L)=0`$ and $`p_\mathrm{\Delta }(C)=0`$, $`L𝔏`$, $`C`$, imply $`LE(\mathrm{\Delta })=0`$ and $`E(\mathrm{\Delta })CE(\mathrm{\Delta })=0`$ for any bounded Borel set $`\mathrm{\Delta }`$, since the integrands occurring in (2.17c) and (2.17b) vanish identically on $`^s`$, and $`𝔅()_{,1}^+`$ as well as $`𝔅()_{,1}`$ are separating sets of functionals for $`𝔅()`$. By the density property just mentioned, it then follows that $`L=0`$ and $`C=0`$, and the nets of seminorms turn out to separate the elements of the left ideal $`𝔏`$ and the -algebra $``$, respectively.
* The completions $`(\overline{𝔏},\overline{𝔗}_q)`$ and $`(\overline{},\overline{𝔗}_p)`$ as well as $`(𝔄_𝔏,\overline{𝔗}_q^u)`$ and $`(𝔄_{},\overline{𝔗}_p^u)`$ are again locally convex spaces with topologies defined by the unique extensions of the seminorms $`q_\mathrm{\Delta }`$ and $`p_\mathrm{\Delta }`$ and of the norm $`.`$ to $`\overline{𝔏}`$, $`𝔄_𝔏`$ and $`\overline{}`$, $`𝔄_{}`$, respectively \[44, Chapter Four, § 18, 4.\]. Therefore, in the sequel, we shall apply these seminorms to elements of the completions without special mention. Depending on the relation between the underlying uniform structures as being finer or coarser, we infer that $`𝔄_𝔏\overline{𝔏}`$ and $`𝔄_{}\overline{}`$. Furthermore $`𝔄_𝔏`$ and $`𝔄_{}`$ are uniformly closed subspaces of the quasi-local algebra $`𝔄`$.
### 2.3 Characteristics of the Spectral Seminorms
The investigations of the subsequent chapters very much depend on special properties of the seminorms defined above, so these are collected in this section. Interesting in their own right as they may be, we are, in the present context, not aiming at utmost generality of statements, but have future applications in mind.
#### 2.3.1 Basic Properties
###### Proposition 2.9.
The families of seminorms $`q_\mathrm{\Delta }`$ and $`p_\mathrm{\Delta }`$ on $`\overline{𝔏}`$ and $`\overline{}`$, respectively, where the symbols $`\mathrm{\Delta }`$ denote bounded Borel sets, constitute nets with respect to the inclusion relation. For any $`\mathrm{\Delta }`$ and $`\mathrm{\Delta }^{}`$ we have
$`\mathrm{\Delta }\mathrm{\Delta }^{}`$ $`q_\mathrm{\Delta }(L)q_\mathrm{\Delta }^{}(L)\text{,}L\overline{𝔏}\text{,}`$
$`\mathrm{\Delta }\mathrm{\Delta }^{}`$ $`p_\mathrm{\Delta }(C)p_\mathrm{\Delta }^{}(C)\text{,}C\overline{}\text{.}`$
###### Proof.
For the $`q_\mathrm{\Delta }`$-seminorms on $`\overline{𝔏}`$ the assertion follows from the order relation for operators in $`𝔅()^+`$. Let $`L`$ belong to the left ideal $`𝔏`$, then
$$Q_\mathrm{\Delta }^{(L^{}L)}Q_\mathrm{\Delta }^{}^{(L^{}L)}\text{,}$$
which by Definition 2.8 has the consequence
$$q_\mathrm{\Delta }(L)^2=Q_\mathrm{\Delta }^{(L^{}L)}Q_\mathrm{\Delta }^{}^{(L^{}L)}=q_\mathrm{\Delta }^{}(L)^2\text{.}$$
This relation extends by continuity of the seminorms to all of $`\overline{𝔏}`$.
In case of the $`p_\mathrm{\Delta }`$-topologies, note that for any Borel set $`\mathrm{\Delta }`$ the functional $`\varphi ^{E(\mathrm{\Delta })}`$, defined through $`\varphi ^{E(\mathrm{\Delta })}(.)\varphi (E(\mathrm{\Delta }).E(\mathrm{\Delta }))`$, belongs to $`𝔅()_{,1}`$ if $`\varphi `$ does. From this we infer, since moreover $`\mathrm{\Delta }\mathrm{\Delta }^{}`$ implies $`E(\mathrm{\Delta })=E(\mathrm{\Delta })E(\mathrm{\Delta }^{})=E(\mathrm{\Delta }^{})E(\mathrm{\Delta })`$, that
$$\begin{array}{c}\{_^sd^sx\left|\varphi \left(E(\mathrm{\Delta })\alpha x(C)E(\mathrm{\Delta })\right)\right|:\varphi 𝔅()_{,1}\}\hfill \\ \hfill \{_^sd^sx\left|\varphi \left(E(\mathrm{\Delta }^{})\alpha x(C)E(\mathrm{\Delta }^{})\right)\right|:\varphi 𝔅()_{,1}\}\end{array}$$
for any $`C`$ and thus, by (2.17b), that $`p_\mathrm{\Delta }(C)p_\mathrm{\Delta }^{}(C)`$, a relation which by continuity of the seminorms is likewise valid for any operator in the completion $`\overline{}`$. ∎
The continuous extensions of the seminorms $`q_\mathrm{\Delta }`$ and $`p_\mathrm{\Delta }`$ to $`\overline{𝔏}`$ and $`\overline{}`$, respectively, can be explicitly computed on the subspaces $`𝔄_𝔏`$ and $`𝔄_{}`$ of $`𝔄`$.
###### Lemma 2.10.
Let $`\mathrm{\Delta }`$ denote an arbitrary bounded Borel subset of $`^{s+1}`$.
* For any $`L𝔄_𝔏`$ we have
$$q_\mathrm{\Delta }(L)=sup\{_^sd^sx\omega \left(E(\mathrm{\Delta })\alpha x(L^{}L)E(\mathrm{\Delta })\right):\omega 𝔅()_{,1}^+\}^{1/2}\text{.}$$
(2.18a)
* For any $`C𝔄_{}`$ there holds the relation
$$p_\mathrm{\Delta }(C)=sup\{_^sd^sx\left|\varphi \left(E(\mathrm{\Delta })\alpha x(C)E(\mathrm{\Delta })\right)\right|:\varphi 𝔅()_{,1}\}\text{.}$$
(2.18b)
###### Proof.
* Note, that we can define a linear subspace $`𝔄_q^{}`$ of $`𝔄`$ consisting of all those operators $`L^{}`$ which fulfill
$$q_\mathrm{\Delta }^{}(L^{})^2sup\{_^sd^sx\omega \left(E(\mathrm{\Delta })\alpha x(L_{}^{}{}_{}{}^{}L^{})E(\mathrm{\Delta })\right):\omega 𝔅()_{,1}^+\}<\mathrm{}$$
for any bounded Borel set $`\mathrm{\Delta }`$. On this space the mappings $`q_\mathrm{\Delta }^{}`$ act as seminorms whose restrictions to $`𝔏`$ coincide with $`q_\mathrm{\Delta }`$ (cf. the Remark following Definition 2.8). Now let $`L𝔄_𝔏`$ be arbitrary. Given a bounded Borel set $`\mathrm{\Delta }`$ we can then find a *sequence* $`\left\{L_n\right\}_n`$ in $`𝔏`$ satisfying
$$\underset{n\mathrm{}}{lim}q_\mathrm{\Delta }(LL_n)=0\text{and}\underset{n\mathrm{}}{lim}LL_n=0\text{.}$$
The second equation implies
$$\underset{n\mathrm{}}{lim}LE(\mathrm{\Delta })L_nE(\mathrm{\Delta })=0\text{,}$$
so that Lebesgue’s Dominated Convergence Theorem can be applied to get for any functional $`\omega 𝔅()_{,1}^+`$ and any compact $`𝑲^𝒔`$
$$Kd^sx\omega \left(E(\mathrm{\Delta })\alpha x(L^{}L)E(\mathrm{\Delta })\right)=\underset{n\mathrm{}}{lim}Kd^sx\omega \left(E(\mathrm{\Delta })\alpha x(L_{n}^{}{}_{}{}^{}L_n)E(\mathrm{\Delta })\right)\text{.}$$
According to (2.17c) each term in the sequence on the right-hand side is majorized by the corresponding $`q_\mathrm{\Delta }(L_n)^2`$ and this sequence in turn converges to $`q_\mathrm{\Delta }(L)^2`$ by assumption, so that in passing from $`𝑲`$ to $`^s`$ and to the supremum over all $`\omega 𝔅()_{,1}^+`$ we get
$$sup\{_^sd^sx\omega \left(E(\mathrm{\Delta })\alpha x(L^{}L)E(\mathrm{\Delta })\right):\omega 𝔅()_{,1}^+\}q_\mathrm{\Delta }(L)^2\text{.}$$
This final estimate shows, by arbitrariness of $`L𝔄_𝔏`$ and the selected $`\mathrm{\Delta }`$, that $`𝔄_𝔏`$ is a subspace of $`𝔄_q^{}`$ and, from $`q_\mathrm{\Delta }^{}𝔏=q_\mathrm{\Delta }`$, it eventually follows that for all these $`L`$ and $`\mathrm{\Delta }`$
$$q_\mathrm{\Delta }(L)=sup\{_^sd^sx\omega \left(E(\mathrm{\Delta })\alpha x(L^{}L)E(\mathrm{\Delta })\right):\omega 𝔅()_{,1}^+\}^{1/2}\text{.}$$
* The proof of the second part follows the same lines of thought. We introduce the subspace $`𝔄_p^{}𝔄`$ consisting of operators $`C^{}`$ satisfying
$$p_\mathrm{\Delta }^{}(C^{})sup\{_^sd^sx\left|\varphi \left(E(\mathrm{\Delta })\alpha x(C^{})E(\mathrm{\Delta })\right)\right|:\varphi 𝔅()_{,1}\}<\mathrm{}$$
for any bounded Borel set $`\mathrm{\Delta }`$ and furnish it with the locally convex topology defined by the seminorms $`p_\mathrm{\Delta }^{}`$. An arbitrary $`C𝔄_{}`$ is, for given $`\mathrm{\Delta }`$, approximated by a *sequence* $`\left\{C_n\right\}_n`$ with respect to the norm and the $`p_\mathrm{\Delta }`$-topology. As above one has
$$\underset{n\mathrm{}}{lim}E(\mathrm{\Delta })CE(\mathrm{\Delta })E(\mathrm{\Delta })C_nE(\mathrm{\Delta })=0$$
and infers
$$sup\{_^sd^sx\left|\varphi \left(E(\mathrm{\Delta })\alpha x(C^{})E(\mathrm{\Delta })\right)\right|:\varphi 𝔅()_{,1}\}p_\mathrm{\Delta }(C)\text{.}$$
This establishes, by arbitraryness of $`C𝔄_{}`$ and $`\mathrm{\Delta }`$, that $`𝔄_{}𝔄_p^{}`$, and the equation $`p_\mathrm{\Delta }^{}=p_\mathrm{\Delta }`$ implies that for these $`C`$ and $`\mathrm{\Delta }`$
$$p_\mathrm{\Delta }(C)=sup\{_^sd^sx\left|\varphi \left(E(\mathrm{\Delta })\alpha x(C)E(\mathrm{\Delta })\right)\right|:\varphi 𝔅()_{,1}\}\text{.}$$
An immediate consequence of this result is the subsequent lemma, which in some way reverts the arguments given in the concluding remark of the last section in order to establish the Hausdorff property for $`(𝔏,𝔗_q)`$ and $`(,𝔗_p)`$.
###### Lemma 2.11.
Let $`\mathrm{\Delta }`$ be a bounded Borel set.
* For $`L𝔄_𝔏`$ with $`LE(\mathrm{\Delta })=0`$ there holds $`q_\mathrm{\Delta }(L)=0`$.
* If $`C𝔄_{}`$ satisfies $`E(\mathrm{\Delta })CE(\mathrm{\Delta })=0`$ one has $`p_\mathrm{\Delta }(C)=0`$.
Next we deal with an implication of the fact, that $`𝔏`$ is an ideal of the $`C^{}`$-algebra $`𝔄`$, and clarify the relationship between the seminorms $`q_\mathrm{\Delta }`$ and $`p_\mathrm{\Delta }`$.
###### Lemma 2.12.
Let $`\mathrm{\Delta }`$ denote bounded Borel subsets of $`^{s+1}`$.
* $`𝔄_𝔏`$ is a left ideal of the quasi-local algebra $`𝔄`$ and satisfies
$$q_\mathrm{\Delta }(AL)Aq_\mathrm{\Delta }(L)$$
(2.19)
for any $`L𝔄_𝔏`$ and $`A𝔄`$.
* Let $`L_i`$, $`i=1\text{,}2`$, be operators in $`𝔄_𝔏`$ and $`A𝔄`$, then $`L_{1}^{}{}_{}{}^{}AL_2`$ belongs to $`𝔄_{}`$. If in addition the operators $`L_i`$ have energy-momentum transfer in $`\mathrm{\Gamma }_i^{s+1}`$ and $`\mathrm{\Delta }_i`$ are Borel subsets of $`^{s+1}`$ containing $`\mathrm{\Delta }+\mathrm{\Gamma }_i`$, respectively, then
$$p_\mathrm{\Delta }(L_{1}^{}{}_{}{}^{}AL_2)E(\mathrm{\Delta }_1)AE(\mathrm{\Delta }_2)q_\mathrm{\Delta }(L_1)q_\mathrm{\Delta }(L_2)\text{.}$$
(2.20)
###### Proof.
* For any $`L𝔄_𝔏𝔄`$ and arbitrary $`A𝔄`$ the relation $`L^{}A^{}ALA^2L^{}L`$ leads to the estimate
$$_^sd^sx\omega \left(E(\mathrm{\Delta })\alpha x(L^{}A^{}AL)E(\mathrm{\Delta })\right)A^2_^sd^sx\omega \left(E(\mathrm{\Delta })\alpha x(L^{}L)E(\mathrm{\Delta })\right)$$
for any $`\omega 𝔅()_{,1}^+`$ and thus, by (2.18a) and the notation of the proof of Lemma 2.10, to
$$\begin{array}{c}q_\mathrm{\Delta }^{}(AL)=sup\{_^sd^sx\omega \left(E(\mathrm{\Delta })\alpha x(L^{}A^{}AL)E(\mathrm{\Delta })\right):\omega 𝔅()_{,1}^+\}^{1/2}\hfill \\ \hfill Asup\{_^sd^sx\omega \left(E(\mathrm{\Delta })\alpha x(L^{}L)E(\mathrm{\Delta })\right):\omega 𝔅()_{,1}^+\}^{1/2}=Aq_\mathrm{\Delta }(L)\text{.}\end{array}$$
This shows that $`AL`$ belongs to $`𝔄_𝔏`$ and at the same time that the seminorm $`q_\mathrm{\Delta }^{}`$ (on $`𝔄_q^{}`$) can be replaced by $`q_\mathrm{\Delta }`$ to yield (2.19).
* Let $`\varphi `$ be a normal functional on $`𝔅()`$ with $`\varphi 1`$. By polar decomposition there exist a partial isometry $`V`$ and a positive normal functional $`|\varphi |`$ with $`|\varphi |1`$ such that $`\varphi (.)=|\varphi |(.V)`$. Then
$$\begin{array}{c}\left|\varphi \left(E(\mathrm{\Delta })\alpha x(L_{1}^{}{}_{}{}^{}AL_2)E(\mathrm{\Delta })\right)\right|\hfill \\ \hfill =|\varphi |\left(E(\mathrm{\Delta })\alpha x(L_{1}^{}{}_{}{}^{})E(\mathrm{\Delta }_1)\alpha x(A)E(\mathrm{\Delta }_2)\alpha x(L_2)E(\mathrm{\Delta })V\right)\\ \hfill E(\mathrm{\Delta }_1)\alpha x(A)E(\mathrm{\Delta }_2)\sqrt{|\varphi |\left(E(\mathrm{\Delta })\alpha x(L_{1}^{}{}_{}{}^{}L_1)E(\mathrm{\Delta })\right)}\sqrt{|\varphi |\left(V^{}E(\mathrm{\Delta })\alpha x(L_{2}^{}{}_{}{}^{}L_2)E(\mathrm{\Delta })V\right)}\end{array}$$
for any $`𝒙^𝒔`$ and the method used in the proof of Proposition 2.7 can be applied to get, in analogy to (2.16),
$$\begin{array}{c}sup\{_^sd^sx\left|\varphi \left(E(\mathrm{\Delta })\alpha x(L_{1}^{}{}_{}{}^{}AL_2)E(\mathrm{\Delta })\right)\right|:\varphi 𝔅()_{,1}\}\hfill \\ \hfill E(\mathrm{\Delta }_1)AE(\mathrm{\Delta }_2)q_\mathrm{\Delta }(L_1)q_\mathrm{\Delta }(L_2)\text{,}\end{array}$$
where we made use of (2.18a). According to the notation introduced in the proof of Lemma 2.10 this result expressed in terms of the seminorm $`p_\mathrm{\Delta }^{}`$ on $`𝔄_p^{}`$ reads
$$p_\mathrm{\Delta }^{}(L_{1}^{}{}_{}{}^{}AL_2)E(\mathrm{\Delta }_1)AE(\mathrm{\Delta }_2)q_\mathrm{\Delta }(L_1)q_\mathrm{\Delta }(L_2)\text{,}$$
from which we infer, as in the first part of the present proof, not only that $`L_{1}^{}{}_{}{}^{}AL_2`$ is an element of $`𝔄_{}`$ but also that $`p_\mathrm{\Delta }^{}`$ can be substituted by $`p_\mathrm{\Delta }`$ to give (2.20). ∎
The second part of the above lemma means that the product $`L_{1}^{}{}_{}{}^{}L_2`$, defined by two operators $`L_1\text{,}L_2𝔄_𝔏`$, is continuous with respect to the locally convex spaces (cf. \[44, Chapter Four, § 18, 3.(5)\]) $`(𝔄_𝔏,\overline{𝔗}_q^u)\times (𝔄_𝔏,\overline{𝔗}_q^u)`$ and $`(𝔄_{},\overline{𝔗}_p^u)`$.
###### Corollary 2.13.
The sesquilinear mapping on the topological product of the locally convex space $`(𝔄_𝔏,\overline{𝔗}_q^u)`$ with itself, defined by
$$𝔄_𝔏\times 𝔄_𝔏(L_1,L_2)L_{1}^{}{}_{}{}^{}L_2𝔄_{}\text{,}$$
is continuous with respect to the respective locally convex topologies.
In the special case of coincidence of both operators ($`L_1=L_2=L`$) it turns out that $`p_\mathrm{\Delta }(L^{}L)`$ equals the square of $`q_\mathrm{\Delta }(L)`$. Another result involving the operation of adjunction is the fact, that this mapping leaves the $`p_\mathrm{\Delta }`$-seminorms invariant.
###### Lemma 2.14.
Let $`\mathrm{\Delta }`$ denote the bounded Borel sets in $`^{s+1}`$.
* For any operator $`L𝔄_𝔏`$ there hold the relations
$$p_\mathrm{\Delta }(L^{}L)=q_\mathrm{\Delta }(L)^2\text{.}$$
* Let $`C`$ be an element of $`𝔄_{}`$, then $`C^{}`$ lies in $`𝔄_{}`$, too, and satisfies
$$p_\mathrm{\Delta }(C^{})=p_\mathrm{\Delta }(C)\text{.}$$
###### Proof.
* According to Lemma 2.10, we have for any $`L𝔄_𝔏`$
$$\begin{array}{c}q_\mathrm{\Delta }(L)^2=sup\{_^sd^sx\omega \left(E(\mathrm{\Delta })\alpha x(L^{}L)E(\mathrm{\Delta })\right):\omega 𝔅()_{,1}^+\}\hfill \\ \hfill sup\{_^sd^sx\left|\varphi \left(E(\mathrm{\Delta })\alpha x(L^{}L)E(\mathrm{\Delta })\right)\right|:\varphi 𝔅()_{,1}\}=p_\mathrm{\Delta }(L^{}L)\text{,}\end{array}$$
whereas the reverse inequality is a consequence of Lemma 2.12. This proves the assertion.
* Note, that $`𝔅()_{,1}`$ is invariant under the operation of taking adjoints defined by $`\psi \psi ^{}`$ with $`\psi ^{}(A)\overline{\psi (A^{})}`$, $`A𝔅()`$, for any linear functional $`\psi `$ on $`𝔅()`$. Thus
$$\begin{array}{c}p_\mathrm{\Delta }^{}(C^{})=sup\{_^sd^sx\left|\varphi \left(E(\mathrm{\Delta })\alpha x(C^{})E(\mathrm{\Delta })\right)\right|:\varphi 𝔅()_{,1}\}\hfill \\ \hfill =sup\{_^sd^sx\left|\varphi ^{}\left(E(\mathrm{\Delta })\alpha x(C)E(\mathrm{\Delta })\right)\right|:\varphi 𝔅()_{,1}\}=p_\mathrm{\Delta }^{}(C)\end{array}$$
for any $`C𝔄_p^{}`$ (cf. the proof of Lemma 2.10), which is sufficient to establish both of the assertions. ∎
The last statement of this subsection on basic properties of the spectral seminorms establishes their invariance under translations in the $`s+1`$-dimensional configuration space.
###### Lemma 2.15.
The subspaces $`𝔄_𝔏`$ and $`𝔄_{}`$ of the quasi-local algebra $`𝔄`$ are invariant under translations. In particular, let $`\mathrm{\Delta }`$ be a bounded Borel set in $`^{s+1}`$ and let $`x^{s+1}`$ be arbitrary, then
* $`q_\mathrm{\Delta }\left(\alpha _x(L)\right)=q_\mathrm{\Delta }(L)`$, $`L𝔄_𝔏`$;
* $`p_\mathrm{\Delta }\left(\alpha _x(C)\right)=p_\mathrm{\Delta }(C)`$, $`C𝔄_{}`$.
###### Proof.
$`𝔅()_{,1}`$ as well as its intersection $`𝔅()_{,1}^+`$ with the positive cone $`𝔅()_{}^+`$ are invariant under the mapping $`\psi \psi ^U`$ defined by $`\psi ^U(.)\psi (U.U^{})`$ for any unitary operator $`U𝔅()`$ and any linear functional $`\psi `$ on $`𝔅()`$.
* Now, $`\alpha _x(L^{}L)=U_t\alpha x(L^{}L)U_{t}^{}{}_{}{}^{}`$ for any $`x=(t,𝒙)^{𝒔+\mathit{1}}`$. This implies
$$\omega \left(E(\mathrm{\Delta })\alpha y\left(\alpha x(L^{}L)\right)E(\mathrm{\Delta })\right)=\omega \left(U_tE(\mathrm{\Delta })\alpha _{𝒙+𝒚}(L^{}L)E(\mathrm{\Delta })U_{t}^{}{}_{}{}^{}\right)$$
for any $`𝒚^𝒔`$ and any $`\omega 𝔅()_{}`$, henceforth
$$\begin{array}{c}_^sd^sy\omega \left(E(\mathrm{\Delta })\alpha y\left(\alpha _x(L^{}L)\right)E(\mathrm{\Delta })\right)\hfill \\ \hfill =_^sd^sy\omega \left(U_tE(\mathrm{\Delta })\alpha _{𝒙+𝒚}(L^{}L)E(\mathrm{\Delta })U_{t}^{}{}_{}{}^{}\right)=_^sd^sy\omega \left(U_tE(\mathrm{\Delta })\alpha y(L^{}L)E(\mathrm{\Delta })U_{t}^{}{}_{}{}^{}\right)\text{.}\end{array}$$
Therefore the introductory remark in combination with (2.18a) yields for any $`L𝔄_𝔏`$:
$$\begin{array}{c}q_\mathrm{\Delta }^{}\left(\alpha _x(L)\right)^2=sup\{_^sd^sy\omega \left(E(\mathrm{\Delta })\alpha y\left(\alpha _x(L^{}L)\right)E(\mathrm{\Delta })\right):\omega 𝔅()_{,1}^+\}\hfill \\ \hfill =sup\{_^sd^sy\omega \left(U_tE(\mathrm{\Delta })\alpha y(L^{}L)E(\mathrm{\Delta })U_{t}^{}{}_{}{}^{}\right):\omega 𝔅()_{,1}^+\}\\ \hfill =sup\{_^sd^sy\omega \left(E(\mathrm{\Delta })\alpha y(L^{}L)E(\mathrm{\Delta })\right):\omega 𝔅()_{,1}^+\}=q_\mathrm{\Delta }^{}(L)^2\text{,}\end{array}$$
which, as in the proof of Lemma 2.12, establishes the assertions.
* The same argument applies to the seminorm $`p_\mathrm{\Delta }^{}`$, so that for $`C𝔄_{}`$
$$\begin{array}{c}p_\mathrm{\Delta }\left(\alpha _x(C)\right)=sup\{_^sd^sy\left|\varphi \left(E(\mathrm{\Delta })\alpha y\left(\alpha _x(C)\right)E(\mathrm{\Delta })\right)\right|:\varphi 𝔅()_{,1}\}\hfill \\ \hfill =sup\{_^sd^sy\left|\varphi \left(U_tE(\mathrm{\Delta })\alpha y(C)E(\mathrm{\Delta })U_{t}^{}{}_{}{}^{}\right)\right|:\varphi 𝔅()_{,1}\}\\ \hfill =sup\{_^sd^sy\left|\varphi \left(E(\mathrm{\Delta })\alpha y(C)E(\mathrm{\Delta })\right)\right|:\varphi 𝔅()_{,1}\}=p_\mathrm{\Delta }(C)\text{.}\end{array}$$
#### 2.3.2 Continuity and Differentiability
The assumed strong continuity of the automorphism group $`\{\alpha _{(\mathrm{\Lambda },x)}:(\mathrm{\Lambda },x)𝖯_+^{}\}`$ acting on the $`C^{}`$-algebra $`𝔄`$ carries over to the locally convex spaces $`(𝔏,𝔗_q)`$ and $`(,𝔗_p)`$; and even the infinite differentiability of $`(\mathrm{\Lambda },x)\alpha _{(\mathrm{\Lambda },x)}(L_0)`$ for $`L_0𝔏_0`$ is conserved in passing from the uniform topology on $`𝔏_0`$ to that induced by the seminorms $`q_\mathrm{\Delta }`$.
###### Proposition 2.16.
* For fixed $`L𝔏`$ the mapping
$$\mathrm{\Xi }_L:𝖯_+^{}𝔏(\mathrm{\Lambda },x)\mathrm{\Xi }_L(\mathrm{\Lambda },x)\alpha _{(\mathrm{\Lambda },x)}(L)$$
is continuous with respect to the locally convex space $`(𝔏,𝔗_q)`$.
* For given $`C`$ the mapping
$$\mathrm{\Xi }_C:𝖯_+^{}(\mathrm{\Lambda },x)\mathrm{\Xi }_C(\mathrm{\Lambda },x)\alpha _{(\mathrm{\Lambda },x)}(C)$$
is continuous with respect to the locally convex space $`(,𝔗_p)`$.
* Let $`\mathrm{id}_{𝔏_0}`$ denote the identity mapping
$$\mathrm{id}_{𝔏_0}:(𝔏_0,.)(𝔏_0,𝔗_q)L_0\mathrm{id}_{𝔏_0}(L_0)L_0$$
on the space $`𝔏_0`$ , once endowed with the norm topology and once with $`𝔗_q`$. Consider furthermore the family $`𝒳_{𝔏_0}`$ of infinitely often differentiable mappings:
$$𝒳_{𝔏_0}\{\mathrm{\Xi }_{L_0}:L_0𝔏_0\}\text{,}$$
$$\mathrm{\Xi }_{L_0}:𝖯_+^{}𝔏_0(\mathrm{\Lambda },x)\mathrm{\Xi }_{L_0}(\mathrm{\Lambda },x)\alpha _{(\mathrm{\Lambda },x)}(L_0)\text{.}$$
Then the linear operator $`\mathrm{id}_{𝔏_0}`$ is $`𝒳_{𝔏_0}`$-differentiable in the sense of Definition A.16.
###### Remark.
The last assertion means, due to the invariance of $`𝒳_{𝔏_0}`$ under differentiation of arbitrary order, that all the mappings $`\mathrm{\Xi }_{L_0}`$, $`L_0𝔏_0`$, are infinitely often differentiable in the locally convex space $`(𝔏_0,𝔗_q)`$ and, as $`\mathrm{id}_{𝔏_0}`$ and the operator of differentiation $`𝔇`$ commute, inherit the derivatives from the presupposed differentiability of the mappings $`\mathrm{\Xi }_{L_0}`$ with respect to the uniform topology.
###### Proof.
* Note, that continuity of the mapping $`(\mathrm{\Lambda },x)\alpha _{(\mathrm{\Lambda },x)}(L)`$ with respect to the locally convex space $`(𝔏,𝔗_q)`$ is equivalent to its continuity with respect to each of the topologizing seminorms $`q_\mathrm{\Delta }`$.
Let the Borel subset $`\mathrm{\Delta }`$ of $`^{s+1}`$ be arbitrary but fixed. We shall first consider the special point $`(\mathrm{𝟏},0)𝖯_+^{}`$ and restrict attention to an operator $`L^{}𝔏_0`$ having energy-momentum transfer $`\mathrm{\Gamma }`$ which, under transformations from a sufficiently small neighbourhood $`𝒩^{}`$ of the neutral element $`(\mathrm{𝟏},0)`$, stays bounded in a compact and convex subset $`\widehat{\mathrm{\Gamma }}`$ of $`\mathrm{}\overline{V}_+`$. This means that all operators $`\alpha _{(\mathrm{\Lambda },x)}(L^{})𝔏_0`$, $`(\mathrm{\Lambda },x)𝒩^{}`$, have energy-momentum transfer in the common set $`\widehat{\mathrm{\Gamma }}`$ and relation (2.6) of Proposition 2.6 applies to the differences $`\alpha _{(\mathrm{\Lambda },x)}(L^{})L^{}`$ yielding
$$\begin{array}{c}q_\mathrm{\Delta }\left(\alpha _{(\mathrm{\Lambda },x)}(L^{})L^{}\right)^2\hfill \\ \hfill =E(\mathrm{\Delta })_^sd^sy\alpha y\left(\left(\alpha _{(\mathrm{\Lambda },x)}(L^{})L^{}\right)^{}\left(\alpha _{(\mathrm{\Lambda },x)}(L^{})L^{}\right)\right)E(\mathrm{\Delta })\\ \hfill N(\mathrm{\Delta },\widehat{\mathrm{\Gamma }})_^sd^sy[\alpha y\left(\alpha _{(\mathrm{\Lambda },x)}(L^{})L^{}\right),\left(\alpha _{(\mathrm{\Lambda },x)}(L^{})L^{}\right)^{}]\text{.}\end{array}$$
(2.21)
An estimate for the integrand on the right-hand side can be based on relation (2.2c), requiring suitable approximating nets of local operators for $`\alpha _{(\mathrm{\Lambda },x)}(L^{})L^{}`$. Given $`R_0>0`$ there exists a neighbourhood $`𝒩^{\prime \prime }`$ of $`(\mathrm{𝟏},x)`$ such that $`𝒩^{\prime \prime }𝒪_r𝒪_{2r}`$ for $`r>R_0`$, and if $`\{L_r^{}𝔄(𝒪_r):r>0\}`$ is an approximating net of local operators for $`L^{}`$, then $`\alpha _{(\mathrm{\Lambda },x)}(L_r^{})𝔄(𝒪_{2r})`$ for any $`r>R_0`$ and $`(\mathrm{\Lambda },x)𝒩^{\prime \prime }`$. Now
$$r^k\alpha _{(\mathrm{\Lambda },x)}(L^{})\alpha _{(\mathrm{\Lambda },x)}(L_r^{})=r^kL^{}L_r^{}\underset{r\mathrm{}}{\overset{}{}}0$$
(2.22)
holds for any $`k`$, so that the operators $`\alpha _{(\mathrm{\Lambda },x)}(L_r^{})L_r^{}𝔄(𝒪_{2r})`$, $`r>R_0`$, constitute the large radius part of approximating nets for each of $`\alpha _{(\mathrm{\Lambda },x)}(L^{})L^{}`$, $`(\mathrm{\Lambda },x)𝒩^{\prime \prime }`$, subject to the bound
$$\left(\alpha _{(\mathrm{\Lambda },x)}(L^{})L^{}\right)\left(\alpha _{(\mathrm{\Lambda },x)}(L_r^{})L_r^{}\right)L^{}L_r^{}\text{,}$$
which is independent of $`(\mathrm{\Lambda },x)𝒩^{\prime \prime }`$. Then, according to the remark following Definition 2.1, there exist approximating nets $`\{L^{}(\mathrm{\Lambda },x)_r𝔄(𝒪_r):r>0\}`$ for the almost local operators $`\alpha _{(\mathrm{\Lambda },x)}(L^{})L^{}`$ that fulfill the estimates $`L^{}(\mathrm{\Lambda },x)_r\alpha _{(\mathrm{\Lambda },x)}(L^{})L^{}`$ and, for $`r>R_0`$, $`\left(\alpha _{(\mathrm{\Lambda },x)}(L^{})L^{}\right)L^{}(\mathrm{\Lambda },x)_{2r}2L^{}L_r^{}`$, where in view of (2.22) the second inequality amounts to $`\left(\alpha _{(\mathrm{\Lambda },x)}(L^{})L^{}\right)L^{}(\mathrm{\Lambda },x)_{2r}2C_kr^k`$ for suitable $`C_k>0`$. Making use of relation (2.2c) in the same remark we arrive at
$$\begin{array}{c}[\alpha y\left(\alpha _{(\mathrm{\Lambda },x)}(L^{})L^{}\right),\left(\alpha _{(\mathrm{\Lambda },x)}(L^{})L^{}\right)^{}]\hfill \\ \hfill 2\alpha _{(\mathrm{\Lambda },x)}(L^{})L^{}^2\chi _<(𝒚)+\mathit{8}\alpha _{(𝜦,𝒙)}(𝑳^{})𝑳^{}𝑳^{}𝑳_{\mathit{4}^\mathit{1}|𝒚|}^{}\chi _>(𝒚)\\ \hfill 8L^{}^2\chi _<(𝒚)+\mathit{16}𝑳^{}𝑪_𝒌\mathit{\hspace{0.17em}4}^𝒌|𝒚|^𝒌\chi _>(𝒚)\end{array}$$
(2.23)
for any $`𝒚^𝒔`$, where $`\chi _<`$ and $`\chi _>`$ denote the characteristic functions pertaining to the compact ball of radius $`4R_0`$ in $`^s`$ and its complement, respectively. The above relation (2.23) holds for any $`(\mathrm{\Lambda },x)𝒩𝒩^{}𝒩^{\prime \prime }`$, and its right-hand side turns out to be an integrable majorizing function for the mapping
$$𝒚[\alpha 𝒚\left(\alpha _{(𝜦,𝒙)}(𝑳^{})𝑳^{}\right),\left(\alpha _{(𝜦,𝒙)}(𝑳^{})𝑳^{}\right)^{}]\text{,}$$
(2.24)
irrespective of $`(\mathrm{\Lambda },x)𝒩`$, if $`ks+2`$. Another consequence of (2.23) is that the function (2.24) converges pointwise to 0 on $`^s`$ in the limit $`(\mathrm{\Lambda },x)(\mathrm{𝟏},0)`$ due to strong continuity of the automorphism group $`\{\alpha _{(\mathrm{\Lambda },x)}:(\mathrm{\Lambda },x)𝖯_+^{}\}`$. Therefore we can apply Lebesgue’s Dominated Convergence Theorem to the integral on the right-hand side of (2.21), evaluated for any sequence $`\left\{(\mathrm{\Lambda }_n,x_n)\right\}_n𝒩`$ approaching $`(\mathrm{𝟏},0)`$ and infer
$$\underset{n\mathrm{}}{lim}q_\mathrm{\Delta }\left(\alpha _{(\mathrm{\Lambda }_n,x_n)}(L^{})L^{}\right)=0\text{.}$$
Since $`𝖯_+^{}`$ as a topological space satisfies the first axiom of countability, this suffices to establish continuity of the mapping $`(\mathrm{\Lambda },x)\alpha _{(\mathrm{\Lambda },x)}(L^{})`$ in $`(\mathrm{𝟏},0)`$ with respect to the $`q_\mathrm{\Delta }`$-topology.
An arbitrary operator $`L𝔏`$ can be represented as $`L=_{i=1}^NA_iL_i^{}`$ where $`L_i^{}𝔏_0`$ comply with the above assumptions on $`L^{}`$ and the operators $`A_i`$ belong to the quasi-local algebra $`𝔄`$ for any $`i=1\text{,}\mathrm{}\text{,}N`$. According to Lemma 2.12 we have
$$\begin{array}{c}0q_\mathrm{\Delta }\left(\alpha _{(\mathrm{\Lambda },x)}(L)L\right)\hfill \\ \hfill \underset{i=1}{\overset{N}{}}\left(q_\mathrm{\Delta }\left(\alpha _{(\mathrm{\Lambda },x)}(A_i)\left(\alpha _{(\mathrm{\Lambda },x)}(L_i^{})L_i^{}\right)\right)+q_\mathrm{\Delta }\left(\left(\alpha _{(\mathrm{\Lambda },x)}(A_i)A_i\right)L_i^{}\right)\right)\\ \hfill \underset{i=1}{\overset{N}{}}\left(A_iq_\mathrm{\Delta }\left(\alpha _{(\mathrm{\Lambda },x)}(L_i^{})L_i^{}\right)+\alpha _{(\mathrm{\Lambda },x)}(A_i)A_iq_\mathrm{\Delta }(L_i^{})\right)\text{,}\end{array}$$
where the right-hand side vanishes in the limit $`(\mathrm{\Lambda },x)(\mathrm{𝟏},0)`$ due to the preceding result and strong continuity of the group $`\{\alpha _{(\mathrm{\Lambda },x)}:(\mathrm{\Lambda },x)𝖯_+^{}\}`$. Thus the mapping $`(\mathrm{\Lambda },x)\alpha _{(\mathrm{\Lambda },x)}(L)`$ turns out to be continuous in $`(\mathrm{𝟏},0)`$ with respect to $`q_\mathrm{\Delta }`$ for arbitrary $`L𝔏`$. The restriction to the specific point $`(\mathrm{𝟏},0)𝖯_+^{}`$ is inessential in the last step since for arbitrary $`(\mathrm{\Lambda }^{},x^{})\text{,}(\mathrm{\Lambda }_0,x_0)𝖯_+^{}`$ one has
$$q_\mathrm{\Delta }\left(\alpha _{(\mathrm{\Lambda }^{},x^{})}(L)\alpha _{(\mathrm{\Lambda }_0,x_0)}(L)\right)=q_\mathrm{\Delta }\left(\alpha _{(\mathrm{\Lambda }^{},x^{})(\mathrm{\Lambda }_0,x_0)^1}\left(\alpha _{(\mathrm{\Lambda }_0,x_0)}(L)\right)\alpha _{(\mathrm{\Lambda }_0,x_0)}(L)\right)\text{,}$$
explicitly showing that continuity of $`(\mathrm{\Lambda },x)\alpha _{(\mathrm{\Lambda },x)}(L)`$ in $`(\mathrm{\Lambda }_0,x_0)`$ is equivalent to continuity of $`(\mathrm{\Lambda },x)\alpha _{(\mathrm{\Lambda },x)}\left(\alpha _{(\mathrm{\Lambda }_0,x_0)}(L)\right)`$ in $`(\mathrm{𝟏},0)`$ with respect to any of the seminorms $`q_\mathrm{\Delta }`$, where $`\alpha _{(\mathrm{\Lambda }_0,x_0)}(L)`$ belongs to $`𝔏`$.
* Continuity of a mapping with values in the locally convex space $`(,𝔗_p)`$ is equivalent to its continuity with respect to all seminorms $`p_\mathrm{\Delta }`$. The problem at hand thus reduces to the one already solved in the first part, if one takes into account the shape of general elements of $``$ according to Definition 2.5 and Corollary 2.13.
* According to Definition A.16 we have to show that for any vacuum annihilation operator $`L_0𝔏_0`$ the mapping $`(\mathrm{\Lambda },x)\mathrm{\Xi }_{L_0}(\mathrm{\Lambda },x)=\alpha _{(\mathrm{\Lambda },x)}(L_0)`$ is differentiable in the locally convex space $`(𝔏_0,𝔗_q)`$ and has derivatives coinciding with those existing in the uniform topology by assumption.
Let $`L_0𝔏_0`$ be given and consider the local chart $`(𝖴,\varphi )`$ around $`(\mathrm{\Lambda }_0,x_0)=\varphi ^1(𝒕_\mathit{0})`$. Due to the presupposed differentiability of the mapping $`\mathrm{\Xi }_{L_0}`$ with respect to the uniform topology the corresponding residual term at $`(\mathrm{\Lambda }_0,x_0)`$ with respect to $`(𝖴,\varphi )`$ is given by
$$R[\mathrm{\Xi }_{L_0}\varphi ^1,𝒕_\mathit{0}](𝒉)=\alpha _{(𝜦𝒉,\mathrm{𝒙𝒉})}(𝑳_\mathit{0})\alpha _{(𝜦_\mathit{0},𝒙_\mathit{0})}(𝑳_\mathit{0})𝔇_\varphi 𝜩_{𝑳_\mathit{0}}(𝜦_\mathit{0},𝒙_\mathit{0})𝒉\text{,}$$
(2.25a)
using the notation $`(\mathrm{\Lambda }_𝒉^{},x_𝒉^{})\varphi ^1(𝒕_\mathit{0}+𝒉^{})`$ for elements of $`𝖴`$, and satisfies the limit condition
$$\underset{𝒉\mathit{0}}{lim}|𝒉|^\mathit{1}𝑹[𝜩_{𝑳_\mathit{0}}\varphi ^\mathit{1},𝒕_\mathit{0}](𝒉)=\mathit{0}\text{,}$$
(2.25b)
To prove the assertion it has to be shown that (2.25b) stays true when the norm is replaced by any of the seminorms $`q_\mathrm{\Delta }`$. Now, according to the Mean Value Theorem A.7, we have for small $`𝒉`$
$$\alpha _{(\mathrm{\Lambda }h,xh)}(L_0)\alpha _{(\mathrm{\Lambda }_0,x_0)}(L_0)=_0^1𝑑\vartheta 𝔇_\varphi \mathrm{\Xi }_{L_0}(\mathrm{\Lambda }_{\vartheta 𝒉},x_{\vartheta 𝒉})𝒉\text{,}$$
where the integral is to be understood with respect to the norm topology of $`𝔄`$. Thus the residual term (2.25a) can be re-written as
$$\begin{array}{c}R[\mathrm{\Xi }_{L_0}\varphi ^1,𝒕_\mathit{0}](𝒉)=_\mathit{0}^\mathit{1}𝒅\vartheta \left(𝔇_\varphi 𝜩_{𝑳_\mathit{0}}(𝜦_{\vartheta 𝒉},𝒙_{\vartheta 𝒉})𝔇_\varphi 𝜩_{𝑳_\mathit{0}}(𝜦_\mathit{0},𝒙_\mathit{0})\right)𝒉\hfill \\ \hfill =\underset{i,j=1}{\overset{d_𝖯}{}}_0^1𝑑\vartheta h_i\left(C_{ij}(\mathrm{\Lambda }_{\vartheta 𝒉},x_{\vartheta 𝒉})\alpha _{(\mathrm{\Lambda }_{\vartheta 𝒉},x_{\vartheta 𝒉})}\left(\delta ^j(L_0)\right)C_{ij}(\mathrm{\Lambda }_0,x_0)\alpha _{(\mathrm{\Lambda }_0,x_0)}\left(\delta ^j(L_0)\right)\right)\text{,}\end{array}$$
where in the last equation (A.20a) is used to represent the linear operator $`𝔇_\varphi \mathrm{\Xi }_{L_0}`$ in terms of partial derivatives of $`\mathrm{\Xi }_{L_0}`$ which can be expressed by means of analytic functions $`C_{ij}`$ on $`𝖴`$ and Poincaré transformed derivatives $`\delta ^j(L_0)`$ of $`L_0`$ ($`d_𝖯`$ is the dimension of the Poincaré group). As a consequence of the first statement of this proposition, the integrand on the right-hand side is continuous with respect to all seminorms $`q_\mathrm{\Delta }`$, so that the integral exists in the complete locally convex space $`(\overline{𝔏},\overline{𝔗}_q)`$. By \[26, II.6.2 and 5.4\] this leads to the following estimate for the residual term
$$\begin{array}{c}|𝒉|^\mathit{1}𝒒_𝜟\left(𝑹[𝜩_{𝑳_\mathit{0}}\varphi ^\mathit{1},𝒕_\mathit{0}](𝒉)\right)\hfill \\ \hfill \underset{i,j=1}{\overset{d_𝖯}{}}_0^1𝑑\vartheta \frac{|h_i|}{|𝒉|}q_\mathrm{\Delta }\left(C_{ij}(\mathrm{\Lambda }_{\vartheta 𝒉},x_{\vartheta 𝒉})\alpha _{(\mathrm{\Lambda }_{\vartheta 𝒉},x_{\vartheta 𝒉})}\left(\delta ^j(L_0)\right)C_{ij}(\mathrm{\Lambda }_0,x_0)\alpha _{(\mathrm{\Lambda }_0,x_0)}\left(\delta ^j(L_0)\right)\right)\\ \hfill \underset{i,j=1}{\overset{d_𝖯}{}}\underset{0\vartheta 1}{\mathrm{max}}q_\mathrm{\Delta }\left(C_{ij}(\mathrm{\Lambda }_{\vartheta 𝒉},x_{\vartheta 𝒉})\alpha _{(\mathrm{\Lambda }_{\vartheta 𝒉},x_{\vartheta 𝒉})}\left(\delta ^j(L_0)\right)C_{ij}(\mathrm{\Lambda }_0,x_0)\alpha _{(\mathrm{\Lambda }_0,x_0)}\left(\delta ^j(L_0)\right)\right)\text{,}\end{array}$$
where evidently the right-hand side vanishes in the limit $`𝒉\mathit{0}`$. Thus condition (A.1b) for differentiability of mappings with values in a locally convex space is fulfilled, and according to the counterpart (2.25a) of (A.1a) the derivatives of $`\mathrm{\Xi }_{L_0}`$ with respect to both the uniform and locally convex topologies on $`𝔏_0`$ coincide. ∎
#### 2.3.3 Integrability
Having established Proposition 2.16 on continuity of the mappings $`(\mathrm{\Lambda },x)\alpha _{(\mathrm{\Lambda },x)}(L)`$ and $`(\mathrm{\Lambda },x)\alpha _{(\mathrm{\Lambda },x)}(C)`$ for given $`L𝔏`$ and $`C`$, it turns out to be possible to construct new elements of $`(𝔏_0,𝔗_q)`$, $`(𝔄_𝔏,\overline{𝔗}_q^u)`$ and $`(𝔄_{},\overline{𝔗}_p^u)`$ through integration with respect to the Haar measure on $`𝖯_+^{}`$.
###### Lemma 2.17.
Let the function $`FL^1(𝖯_+^{},d\mu (\mathrm{\Lambda },x))`$ have compact support $`𝖲`$.
* For any $`L_0𝔏_0`$ the operator
$$\alpha _F(L_0)𝑑\mu (\mathrm{\Lambda },x)F(\mathrm{\Lambda },x)\alpha _{(\mathrm{\Lambda },x)}(L_0)$$
(2.26)
belongs to $`𝔏_0`$, too.
* If $`L𝔏`$ and $`C`$, then
$`\alpha _F(L)`$ $`{\displaystyle 𝑑\mu (\mathrm{\Lambda },x)F(\mathrm{\Lambda },x)\alpha _{(\mathrm{\Lambda },x)}(L)\text{,}}`$ (2.27a)
$`\alpha _F(C)`$ $`{\displaystyle 𝑑\mu (\mathrm{\Lambda },x)F(\mathrm{\Lambda },x)\alpha _{(\mathrm{\Lambda },x)}(C)}`$ (2.27b)
exist as integrals in the complete locally convex spaces $`(𝔄_𝔏,\overline{𝔗}_q^u)`$ and $`(𝔄_{},\overline{𝔗}_p^u)`$, respectively, and for any bounded Borel set $`\mathrm{\Delta }`$ there hold the estimates
$`q_\mathrm{\Delta }\left(\alpha _F(L)\right)`$ $`F_1\underset{(\mathrm{\Lambda },x)𝖲}{sup}q_\mathrm{\Delta }\left(\alpha _{(\mathrm{\Lambda },x)}(L)\right)\text{,}`$ (2.28a)
$`p_\mathrm{\Delta }\left(\alpha _F(C)\right)`$ $`F_1\underset{(\mathrm{\Lambda },x)𝖲}{sup}p_\mathrm{\Delta }\left(\alpha _{(\mathrm{\Lambda },x)}(C)\right)\text{.}`$ (2.28b)
###### Proof.
* By assumption $`(\mathrm{\Lambda },x)|F(\mathrm{\Lambda },x)|\alpha _{(\mathrm{\Lambda },x)}(L_0)=|F(\mathrm{\Lambda },x)|L_0`$ \[54, Corollary I.5.4\] is an integrable majorizing function for the integrand of (2.26), so $`\alpha _F(L_0)`$ exists as a Bochner integral in $`𝔄`$. The same holds true for the integrals constructed by use of an approximating net $`\{L_{0,r}𝔄(𝒪_r):r>0\}`$ for the almost local operator $`L_0𝔏_0`$:
$$\alpha _F(L_{0,r})_𝖲𝑑\mu (\mathrm{\Lambda },x)F(\mathrm{\Lambda },x)\alpha _{(\mathrm{\Lambda },x)}(L_{0,r})\text{.}$$
Due to compactness of $`𝖲`$, these operators belong to the local algebras $`𝔄\left(𝒪_{r(𝖲)}\right)`$ pertaining to standard diamonds in $`^{s+1}`$ which have each an $`s`$-dimensional basis of radius $`r(𝖲)a(𝖲)r+b(𝖲)`$ where $`a(𝖲)`$ and $`b(𝖲)`$ are suitable positive constants. Now,
$$\alpha _F(L_0)\alpha _F(L_{0,r})=_𝖲𝑑\mu (\mathrm{\Lambda },x)F(\mathrm{\Lambda },x)\left(\alpha _{(\mathrm{\Lambda },x)}(L_0)\alpha _{(\mathrm{\Lambda },x)}(L_{0,r})\right)\text{,}$$
so that we arrive at the estimate ($`\mu (𝖲)`$ is the measure of the compact set $`𝖲`$)
$$r(𝖲)^k\alpha _F(L_0)\alpha _F(L_{0,r})\mu (𝖲)F_1\left(a(𝖲)r+b(𝖲)\right)^kL_0L_{0,r}$$
which holds for any $`k`$. Due to almost locality of $`L_0`$, the right-hand side vanishes in the limit of large $`r`$, so that the operator $`\alpha _F(L_0)`$ itself turns out to be almost local: $`\alpha _F(L_0)𝔄_𝒮`$ with approximating net $`\{\alpha _F(L_{0,r})𝔄\left(𝒪_{r(𝖲)}\right):r>0\}`$.
Let $`\mathrm{\Gamma }\mathrm{}\overline{V}_+`$ denote the energy-momentum transfer of the vacuum annihilation operator $`L_0`$, then, by the Fubini Theorem \[26, II.16.3\], the following equation is valid for any $`gL^1(^{s+1},d^{s+1}y)`$
$$_{^{s+1}}d^{s+1}yg(y)\alpha _y\left(\alpha _F(L_0)\right)=_𝖲𝑑\mu (\mathrm{\Lambda },x)F(\mathrm{\Lambda },x)_{^{s+1}}d^{s+1}yg(y)\alpha _y\left(\alpha _{(\mathrm{\Lambda },x)}(L_0)\right)\text{.}$$
In the special case $`\mathrm{supp}\stackrel{~}{g}_{(\mathrm{\Lambda },x)𝖲}\mathrm{}(\mathrm{\Lambda }\mathrm{\Gamma })`$, $`\stackrel{~}{g}`$ the Fourier transform of $`g`$, the inner integrals on the right-hand side vanish for any $`(\mathrm{\Lambda },x)𝖲`$ so that we infer
$$_{^{s+1}}d^{s+1}yg(y)\alpha _y\left(\alpha _F(L_0)\right)=0\text{,}$$
which shows that the energy-momentum transfer of $`\alpha _F(L_0)`$ is contained in the compact subset $`_{(\mathrm{\Lambda },x)𝖲}\mathrm{\Lambda }\mathrm{\Gamma }`$ of $`\mathrm{}\overline{V}_+`$. Therefore $`\alpha _F(L_0)`$ is indeed a vacuum annihilation operator from $`𝔄_{\text{ann}}`$.
Finally, infinite differentiability with respect to the uniform topology of the mapping
$$(\mathrm{\Lambda },x)\mathrm{\Xi }_{\alpha _F(L_0)}(\mathrm{\Lambda },x)=\alpha _{(\mathrm{\Lambda },x)}\left(\alpha _F(L_0)\right)$$
has to be established. By assumption $`L_0`$ is infinitely often differentiable with respect to the Poincaré group, which implies that likewise all the operators $`\alpha _{(\mathrm{\Lambda },x)}(L_0)`$ belong to $`𝒟^{(\mathrm{})}(𝔄)`$ for any $`(\mathrm{\Lambda },x)𝖯_+^{}`$. Their residual terms at $`(\mathrm{𝟏},0)=\varphi _0^1(\mathit{0})`$ with respect to the canonical coordinates $`(𝖴_0,\varphi _0)`$ of the first kind, as introduced in \[55, Section 2.10\], can, using the notation $`(\mathrm{\Lambda }_𝒉^{},x_𝒉^{})\varphi _0^1(𝒉^{})`$ for the transformations in $`𝖴_0`$, be expressed by
$$\begin{array}{c}R[\mathrm{\Xi }_{\alpha _{(\mathrm{\Lambda },x)}(L_0)}\varphi _0^1,\mathit{0}](𝒉)\hfill \\ \hfill =\alpha _{(\mathrm{\Lambda }h,xh)}\left(\alpha _{(\mathrm{\Lambda },x)}(L_0)\right)\alpha _{(\mathrm{\Lambda },x)}(L_0)𝔇_{\varphi _0}\mathrm{\Xi }_{\alpha _{(\mathrm{\Lambda },x)}(L_0)}(\mathrm{𝟏},0)𝒉\\ \hfill =_0^1𝑑\vartheta \left(𝔇_{\varphi _0}\mathrm{\Xi }_{\alpha _{(\mathrm{\Lambda },x)}(L_0)}(\mathrm{\Lambda }_{\vartheta 𝒉},x_{\vartheta 𝒉})𝔇_{\varphi _0}\mathrm{\Xi }_{\alpha _{(\mathrm{\Lambda },x)}(L_0)}(\mathrm{𝟏},0)\right)𝒉\text{,}\end{array}$$
(2.29)
where the last equation stems from an application of the Mean Value Theorem A.7, which holds true for small $`𝒉`$. By Proposition A.11 the term $`𝔇_{\varphi _0}\mathrm{\Xi }_{\alpha _{(\mathrm{\Lambda },x)}(L_0)}(\mathrm{𝟏},0)𝒉`$ on the second line is continuous in $`(\mathrm{\Lambda },x)`$, so that it is possible to multiply (2.29) with the function $`F(\mathrm{\Lambda },x)`$ and subsequently integrate over its compact support $`𝖲`$. Taking into account that each of the automorphisms $`\alpha _{(\mathrm{\Lambda }h,xh)}`$ is uniformly continuous, thus commuting with Bochner integrals, this yields
$$\begin{array}{c}_𝖲𝑑\mu (\mathrm{\Lambda },x)F(\mathrm{\Lambda },x)R[\mathrm{\Xi }_{\alpha _{(\mathrm{\Lambda },x)}(L_0)}\varphi _0^1,\mathit{0}](𝒉)\hfill \\ \hfill =\alpha _{(\mathrm{\Lambda }h,xh)}\left(\alpha _F(L_0)\right)\alpha _F(L_0)_𝖲𝑑\mu (\mathrm{\Lambda },x)F(\mathrm{\Lambda },x)𝔇_{\varphi _0}\mathrm{\Xi }_{\alpha _{(\mathrm{\Lambda },x)}(L_0)}(\mathrm{𝟏},0)𝒉\text{,}\end{array}$$
(2.30)
which has the shape of a residual term for $`\mathrm{\Xi }_{\alpha _F(L_0)}`$ at $`(\mathrm{𝟏},0)`$. Now, the operator-norm of $`𝔇_{\varphi _0}\mathrm{\Xi }_{\alpha _{(\mathrm{\Lambda },x)}(L_0)}(\mathrm{\Lambda }_{\vartheta 𝒉},x_{\vartheta 𝒉})`$ can be estimated according to (A.17) by
$$𝔇_{\varphi _0}\mathrm{\Xi }_{\alpha _{(\mathrm{\Lambda },x)}(L_0)}(\mathrm{\Lambda }_{\vartheta 𝒉},x_{\vartheta 𝒉})𝔇_{\varphi _0}\mathrm{\Xi }_{L_0}(1\mathrm{}𝐍\mathrm{}\mathrm{\Lambda }\mathrm{}\mathrm{x}\mathrm{}𝐌^{\varphi _0}\mathrm{}\mathrm{\Lambda }_{\vartheta 𝒉}\mathrm{}\mathrm{x}_{\vartheta 𝒉}\mathrm{}\text{,}$$
which, due to continuity of $`(\mathrm{\Lambda },x)𝐍(\mathrm{\Lambda },x)`$ and $`\vartheta 𝐌^{\varphi _0}(\mathrm{\Lambda }_{\vartheta 𝒉},x_{\vartheta 𝒉})`$ with respect to the operator-norm topology, is majorized on the compact set $`𝖲\times [0,1]`$ by a constant $`K(𝖲)`$. As a consequence of the last equation in (2.29) we then get for any $`(\mathrm{\Lambda },x)𝖲`$ and small $`𝒉`$ the bound
$$|𝒉|^\mathit{1}𝑭(𝜦,𝒙)𝑹[𝜩_{\alpha _{(𝜦,𝒙)}(𝑳_\mathit{0})}\varphi _\mathit{0}^\mathit{1},\mathit{0}](𝒉)\mathit{2}𝑲(𝖲)|𝑭(𝜦,𝒙)|\text{,}$$
(2.31)
which is integrable over $`𝖲`$ by assumption; restricting furthermore attention to sequences $`\left\{𝒉_𝒏\right\}_𝒏`$ converging to $`\mathit{0}`$, we see that the left-hand side of (2.31) converges pointwise to $`0`$. With this information at hand it is possible to apply Lebesgue’s Dominated Convergence Theorem \[26, II.5.6\] to the left-hand side of (2.30) to get
$$\underset{n\mathrm{}}{lim}|𝒉_𝒏|^\mathit{1}_𝖲𝒅\mu (𝜦,𝒙)𝑭(𝜦,𝒙)𝑹[𝜩_{\alpha _{(𝜦,𝒙)}(𝑳_\mathit{0})}\varphi _\mathit{0}^\mathit{1},\mathit{0}](𝒉_𝒏)=\mathit{0}\text{,}$$
(2.32)
which is sufficient to establish condition (A.1b) for differentiability of the mapping $`\mathrm{\Xi }_{\alpha _F(L_0)}`$ at $`(\mathrm{𝟏},0)`$. The linear operator defining the corresponding derivative is according to the right-hand side of (2.30) in connection with (A.17) given by
$$\begin{array}{c}𝔇_{\varphi _0}\mathrm{\Xi }_{\alpha _F(L_0)}(\mathrm{𝟏},0)𝒉=_𝖲𝒅\mu (𝜦,𝒙)𝑭(𝜦,𝒙)𝔇_{\varphi _\mathit{0}}𝜩_{\alpha _{(𝜦,𝒙)}(𝑳_\mathit{0})}(\mathrm{𝟏},\mathit{0})𝒉\hfill \\ \hfill =_𝖲𝑑\mu (\mathrm{\Lambda },x)F(\mathrm{\Lambda },x)\alpha _{(\mathrm{\Lambda },x)}𝔇_{\varphi _0}\mathrm{\Xi }_{L_0}(\mathrm{𝟏},0)𝐍(\mathrm{\Lambda },x)𝒉\\ \hfill =\underset{i,j=1}{\overset{d_𝖯}{}}h_i_𝖲𝑑\mu (\mathrm{\Lambda },x)F_{ji}(\mathrm{\Lambda },x)\alpha _{(\mathrm{\Lambda },x)}\left(\delta ^j(L_0)\right)=\underset{i,j=1}{\overset{d_𝖯}{}}h_i\alpha _{F_{ji}}\left(\delta ^j(L_0)\right)\text{,}\end{array}$$
where $`F_{ji}(\mathrm{\Lambda },x)F(\mathrm{\Lambda },x)𝐍_{ji}(\mathrm{\Lambda },x)`$ are functions from $`L^1(𝖯_+^{},d\mu (\mathrm{\Lambda },x))`$ with compact support $`𝖲`$. Since $`𝔏_0`$ is invariant under differentiation we conclude from the first two paragraphs of the present proof and the above considerations that the partial derivatives
$$\delta ^i\left(\alpha _F(L_0)\right)=\underset{j=1}{\overset{d_𝖯}{}}\alpha _{F_{ji}}\left(\delta ^j(L_0)\right)$$
are again almost local vacuum annihilation operators which belong to $`𝒟^{(1)}(𝔄)`$. Thus by induction, repeatedly using these methods, $`\alpha _F(L_0)`$ is seen to be an element of $`𝒟^{(\mathrm{})}(𝔄)`$ with almost local derivatives of any order, i. e. $`\alpha _F(L_0)𝔏_0`$.
* By Proposition 2.16 the mappings $`(\mathrm{\Lambda },x)\alpha _{(\mathrm{\Lambda },x)}(L)`$ and $`(\mathrm{\Lambda },x)\alpha _{(\mathrm{\Lambda },x)}(C)`$ are continuous with respect to the uniform topology and all the $`q_\mathrm{\Delta }`$\- and $`p_\mathrm{\Delta }`$-topologies, respectively, staying bounded on the compact set $`𝖲`$. This implies their measurability in the locally convex spaces $`(𝔄_𝔏,\overline{𝔗}_q^u)`$ and $`(𝔄_{},\overline{𝔗}_p^u)`$ together with the fact that their product with the integrable function $`f`$ is majorized in each of the norm and seminorm topologies by a multiple of $`|F|`$. As a consequence the integrals $`\alpha _F(L)`$ and $`\alpha _F(C)`$ exist in the complete locally convex spaces $`(𝔄_𝔏,\overline{𝔗}_q^u)`$ and $`(𝔄_{},\overline{𝔗}_p^u)`$, respectively, and (2.28) is an immediate upshot \[26, II.6.2 and 5.4\]. ∎
There exists a version of the second part of the above lemma for functions on $`^{s+1}`$ that are Lebesgue-integrable but no longer have to be compactly supported.
###### Lemma 2.18.
Let $`L𝔏`$ and let $`gL^1(^{s+1},d^{s+1}x)`$. Then
$$\alpha _g(L)_{^{s+1}}d^{s+1}xg(x)\alpha _x(L)$$
(2.33)
is an operator in $`(𝔄_𝔏,\overline{𝔗}_q^u)`$, satisfying the estimates
$$q_\mathrm{\Delta }\left(\alpha _g(L)\right)g_1q_\mathrm{\Delta }(L)$$
(2.34)
for any bounded Borel set $`\mathrm{\Delta }`$. The energy-momentum transfer of $`\alpha _g(L)`$ is contained in $`\mathrm{supp}\stackrel{~}{g}`$, the support of the Fourier transform $`\stackrel{~}{g}`$ of $`g`$.
###### Proof.
By translation invariance of the norm $`.`$ as well as of the seminorms $`q_\mathrm{\Delta }`$ (cf. Lemma 2.15) the (measurable) integrand on the right-hand side of (2.33) is majorized by the functions $`x|g(x)|L`$ and $`x|g(x)|q_\mathrm{\Delta }(L)`$ for any bounded Borel set $`\mathrm{\Delta }`$. These are Lebesgue-integrable and therefore $`\alpha _g(L)`$ exists as a unique element of $`(𝔄_𝔏,\overline{𝔗}_q^u)`$, satisfying the claimed estimates (2.34).
Next, we consider an arbitrary function $`hL^1(^{s+1},d^{s+1}x)`$. By Fubini’s Theorem \[26, II.16.3\] and translation invariance of Lebesgue measure
$$\begin{array}{c}_{^{s+1}}d^{s+1}yh(y)\alpha _y\left(\alpha _g(L)\right)=_{^{s+1}}d^{s+1}yh(y)\alpha _y\left(_{^{s+1}}d^{s+1}xg(x)\alpha _x(L)\right)\hfill \\ \hfill =_{^{s+1}}d^{s+1}y_{^{s+1}}d^{s+1}xh(y)g(x)\alpha _{x+y}(L)\\ \hfill =_{^{s+1}}d^{s+1}x\left(_{^{s+1}}d^{s+1}yh(y)g(xy)\right)\alpha _x(L)\text{,}\end{array}$$
where the term in brackets on the right-hand side of the last equation is the convolution product $`hg`$ of $`h`$ and $`g`$. Its Fourier transform $`\stackrel{~}{hg}`$ is given by $`\stackrel{~}{hg}(p)=(2\pi )^{(s+1)/2}\stackrel{~}{h}(p)\stackrel{~}{g}(p)`$ (cf. \[39, Theorem VI.(21.41)\]), so that this function vanishes if $`\stackrel{~}{h}`$ and $`\stackrel{~}{g}`$ have disjoint supports. Therefore $`\mathrm{supp}\stackrel{~}{h}\mathrm{supp}\stackrel{~}{g}=\mathrm{}`$ entails
$$_{^{s+1}}d^{s+1}yh(y)\alpha _y\left(\alpha _g(L)\right)=0\text{,}$$
and this shows that the Fourier transform of $`y\alpha _y\left(\alpha _g(L)\right)`$ has support in $`\mathrm{supp}\stackrel{~}{g}`$, which henceforth contains the energy-momentum transfer of $`\alpha _g(L)`$. ∎
#### 2.3.4 Decay Property
Eventually we are able to establish a property of rapid decay with respect to the seminorms $`q_\mathrm{\Delta }`$ for commutators of elements of $`𝔏`$ which are almost local.
###### Lemma 2.19.
Let $`L_1`$ and $`L_2`$ belong to $`𝔏_0`$ and let $`A_1\text{,}A_2𝔄`$ be almost local. Then for any bounded Borel subset $`\mathrm{\Delta }`$ of $`^{s+1}`$
$$^s𝒙𝒒_𝜟\left([\alpha 𝒙(𝑨_\mathit{1}𝑳_\mathit{1}),𝑨_\mathit{2}𝑳_\mathit{2}]\right)$$
decreases with $`|𝐱|\mathrm{}`$ faster than any power of $`|𝐱|^\mathrm{𝟏}`$.
###### Proof.
First we consider the special case of two elements $`L_a`$ and $`L_b`$ in $`𝔏_0`$ having energy-momentum transfer in compact and convex subsets $`\mathrm{\Gamma }_a`$ and $`\mathrm{\Gamma }_b`$ of $`\mathrm{}\overline{V}_+`$, respectively, with the additional property that $`\mathrm{\Gamma }_{a,b}(\mathrm{\Gamma }_a+\mathrm{\Gamma }_b)\mathrm{\Gamma }_a`$ and $`\mathrm{\Gamma }_{b,a}(\mathrm{\Gamma }_a+\mathrm{\Gamma }_b)\mathrm{\Gamma }_b`$ lie in the complement of $`\overline{V}_+`$, too. According to the Lemmas 2.14 and 2.12
$$\begin{array}{c}q_\mathrm{\Delta }\left([\alpha x(L_a),L_b]\right)^2=p_\mathrm{\Delta }\left([\alpha x(L_a),L_b]^{}[\alpha x(L_a),L_b]\right)\hfill \\ \hfill q_\mathrm{\Delta }(L_b)q_\mathrm{\Delta }\left(\alpha x(L_a)^{}[\alpha x(L_a),L_b]\right)+q_\mathrm{\Delta }(L_a)q_\mathrm{\Delta }\left(L_{b}^{}{}_{}{}^{}[\alpha x(L_a),L_b]\right)\text{,}\end{array}$$
(2.35)
and we are left with the task to investigate for large $`|𝒙|`$ the behaviour of the functions $`q_\mathrm{\Delta }\left(\alpha x(L_a)^{}[\alpha x(L_a),L_b]\right)`$ and $`q_\mathrm{\Delta }\left(L_{b}^{}{}_{}{}^{}[\alpha x(L_a),L_b]\right)`$. Since the arguments of both terms belong to $`𝔏_0`$, having energy-momentum transfer in the compact and convex subsets $`\mathrm{\Gamma }_{a,b}`$ and $`\mathrm{\Gamma }_{b,a}`$ of $`\mathrm{}\overline{V}_+`$, we can apply (2.6) of Proposition 2.6 in connection with (2.17a) to get the estimate (for the second term)
$$\begin{array}{c}|𝒙|^{\mathit{2}𝒌}𝒒_𝜟\left(𝑳_{𝒃}^{}{}_{}{}^{}[\alpha 𝒙(𝑳_𝒂),𝑳_𝒃]\right)^\mathit{2}\hfill \\ \hfill N(\mathrm{\Delta },\mathrm{\Gamma }_{b,a})_^sd^sy|𝒙|^{\mathit{2}𝒌}[\alpha 𝒚\left(𝑳_{𝒃}^{}{}_{}{}^{}[\alpha 𝒙(𝑳_𝒂),𝑳_𝒃]\right),\left(𝑳_{𝒃}^{}{}_{}{}^{}[\alpha 𝒙(𝑳_𝒂),𝑳_𝒃]\right)^{}]\text{.}\end{array}$$
(2.36)
Let $`\{L_{a,r}𝔄(𝒪_r):r>0\}`$ and $`\{L_{b,r}𝔄(𝒪_r):r>0\}`$ be approximating nets for $`L_a`$ and $`L_b`$, respectively, satisfying $`L_{a,r}L_a`$ and $`L_{b,r}L_b`$. Then the elements
$$L_{b,r}^{}{}_{}{}^{}[\alpha x(L_{a,r}),L_{b,r}]𝔄(𝒪_r+𝒙)𝔄(𝒪_{𝒓+|𝒙|})$$
constitute the large radius part of approximating nets for the almost local operators $`L_{b}^{}{}_{}{}^{}[\alpha x(L_a),L_b]`$, $`𝒙^𝒔`$, subject to the estimate
$$\begin{array}{c}L_{b}^{}{}_{}{}^{}[\alpha x(L_a),L_b]L_{b,r}^{}{}_{}{}^{}[\alpha x(L_{a,r}),L_{b,r}]\hfill \\ \hfill 4L_aL_bL_bL_{b,r}+2L_b^2L_aL_{a,r}C_lr^l\end{array}$$
(2.37)
for any $`l`$ with suitable $`C_l>0`$. Now, as suggested by the remark following Definition 2.1, there exist approximating nets $`\{L(a,b;𝒙)_𝒓𝔄(𝒪_𝒓):𝒓>\mathit{0}\}`$, $`𝒙^𝒔`$, with $`L(a,b;𝒙)_𝒓𝑳_{𝒃}^{}{}_{}{}^{}[\alpha 𝒙(𝑳_𝒂),𝑳_𝒃]`$ and $`L_{b}^{}{}_{}{}^{}[\alpha x(L_a),L_b]L(a,b;𝒙)_{𝒓+|𝒙|}\mathit{2}𝑪_𝒍𝒓^𝒍`$, so that, according to (2.2c), the integrand of (2.36) is bounded by
$$\begin{array}{c}|𝒙|^{\mathit{2}𝒌}[\alpha 𝒚\left(𝑳_{𝒃}^{}{}_{}{}^{}[\alpha 𝒙(𝑳_𝒂),𝑳_𝒃]\right),\left(𝑳_{𝒃}^{}{}_{}{}^{}[\alpha 𝒙(𝑳_𝒂),𝑳_𝒃]\right)^{}]\hfill \\ \hfill |𝒙|^{\mathit{2}𝒌}\mathit{4}𝑳_{𝒃}^{}{}_{}{}^{}[\alpha 𝒙(𝑳_𝒂),𝑳_𝒃]𝑳_{𝒃}^{}{}_{}{}^{}[\alpha 𝒙(𝑳_𝒂),𝑳_𝒃]𝑳(𝒂,𝒃;𝒙)_{\mathit{2}^\mathit{1}|𝒚|}\\ \hfill \{\begin{array}{cc}8|𝒙|^{\mathit{2}𝒌}𝑳_𝒃^\mathit{2}[\alpha 𝒙(𝑳_𝒂),𝑳_𝒃]^\mathit{2}\hfill & \text{,}|𝒚|\mathit{2}(|𝒙|+\mathit{1})\text{,}\hfill \\ 8L_b|𝒙|^{\mathit{2}𝒌}[\alpha 𝒙(𝑳_𝒂),𝑳_𝒃]𝑪_𝒍(\mathit{2}^\mathit{1}|𝒚||𝒙|)^𝒍\hfill & \text{,}|𝒚|>\mathit{2}(|𝒙|+\mathit{1})\text{,}\hfill \end{array}\end{array}$$
(2.38)
which implies
$$\begin{array}{c}|𝒙|^{\mathit{2}𝒌}𝒒_𝜟\left(𝑳_{𝒃}^{}{}_{}{}^{}[\alpha 𝒙(𝑳_𝒂),𝑳_𝒃]\right)^\mathit{2}\hfill \\ \hfill N(\mathrm{\Delta },\mathrm{\Gamma }_{b,a})[8L_b^2|𝒙|^{\mathit{2}𝒌}[\alpha 𝒙(𝑳_𝒂),𝑳_𝒃]^\mathit{2}\underset{|𝒚|\mathit{2}(|𝒙|+\mathit{1})}{}𝒅^𝒔𝒚\\ \hfill +8C_lL_b|𝒙|^{\mathit{2}𝒌}[\alpha 𝒙(𝑳_𝒂),𝑳_𝒃]\underset{|𝒚|>\mathit{2}(|𝒙|+\mathit{1})}{}𝒅^𝒔𝒚(\mathit{2}^\mathit{1}|𝒚||𝒙|)^𝒍]\text{.}\end{array}$$
(2.39)
Evaluation of the integrals on the right-hand side yields (for $`ls+2`$) polynomials of degree $`s`$ in $`|𝒙|`$, so that, due to the decay properties of $`𝒙[\alpha 𝒙(𝑳_𝒂),𝑳_𝒃]`$, there exists a uniform bound
$$|𝒙|^𝒌𝒒_𝜟\left(𝑳_{𝒃}^{}{}_{}{}^{}[\alpha 𝒙(𝑳_𝒂),𝑳_𝒃]\right)^\mathit{2}𝑴\text{,}𝒙^𝒔\text{.}$$
(2.40)
The same reasoning applies to the term $`q_\mathrm{\Delta }\left(\alpha x(L_a)^{}[\alpha x(L_a),L_b]\right)`$, thus establishing the asserted rapid decrease for the mapping $`𝒙𝒒_𝜟\left([\alpha 𝒙(𝑳_𝒂),𝑳_𝒃]\right)`$, according to relation (2.35).
In the general case of almost local elements $`A_1\text{,}A_2𝔄`$ and $`L_1\text{,}L_2𝔏_0`$ one has, by Lemma 2.12
$$\begin{array}{c}q_\mathrm{\Delta }\left([\alpha x(A_1L_1),A_2L_2]\right)\hfill \\ \hfill A_1[\alpha x(L_1),A_2]q_\mathrm{\Delta }(L_2)+A_1A_2q_\mathrm{\Delta }\left([\alpha x(L_1),L_2]\right)\\ \hfill +[\alpha x(A_1),A_2]L_2q_\mathrm{\Delta }(L_1)+A_2[\alpha x(A_1),L_2]q_\mathrm{\Delta }(L_1)\text{,}\end{array}$$
and rapid decay is an immediate consequence of almost locality for all terms but the second one on the right-hand side of this inequality. Using suitable decompositions of $`L_1`$ and $`L_2`$ in terms of elements of $`𝔏_0`$ complying pairwise with the special properties exploited in the previous paragraph, the remaining problem of decrease of the mapping $`𝒙𝒒_𝜟\left([\alpha 𝒙(𝑳_\mathit{1}),𝑳_\mathit{2}]\right)`$ reduces to the case that has already been solved above, thus completing the proof. ∎
## Chapter 3 Particle Weights as Asymptotic Plane Waves
Having analysed in great detail the nets of seminorms $`q_\mathrm{\Delta }`$ and $`p_\mathrm{\Delta }`$, indexed by the bounded Borel sets $`\mathrm{\Delta }^{s+1}`$, on $`𝔏`$ and $``$, respectively, we now turn to the investigation of the topological dual spaces:
###### Definition 3.1.
* The linear functionals on $``$ which are continuous with respect to the seminorm $`p_\mathrm{\Delta }`$ constitute a vector space $`_{\mathrm{\Delta }}^{}{}_{}{}^{}`$, which is a normed space via
$$\varsigma _\mathrm{\Delta }sup\{|\varsigma (C)|:C,p_\mathrm{\Delta }(C)1\}\text{,}\varsigma _{\mathrm{\Delta }}^{}{}_{}{}^{}\text{.}$$
* The topological duals of the locally convex spaces $`(𝔏_0,𝔗_q)`$, $`(𝔏,𝔗_q)`$ and $`(,𝔗_p)`$ are denoted $`𝔏_{0}^{}{}_{}{}^{}`$, $`𝔏^{}`$ and $`^{}`$, respectively.
###### Remark.
Due to the net property (Proposition 2.9) of the family of seminorms $`p_\mathrm{\Delta }`$, a linear functional belongs to the topological dual $`^{}`$ of $`(,𝔗_p)`$ if and only if it is continuous with respect to one specific seminorm $`p_\mathrm{\Delta }^{}`$, $`\mathrm{\Delta }^{}`$ a bounded Borel subset of $`^{s+1}`$ \[42, Proposition 1.2.8\]. Hence
$$^{}=\{_{\mathrm{\Delta }}^{}{}_{}{}^{}:\mathrm{\Delta }^{s+1}\text{ a bounded Borel set}\}\text{.}$$
(3.1)
By continuous linear extension \[44, Chapter One, § 5, 4.(4)\], the functionals from $`^{}`$ are moreover in one-to-one correspondence with the elements of the topological dual $`\overline{}^{}`$ of the complete locally convex space $`(\overline{},\overline{𝔗}_p)`$. By the same argument, they are furthermore embedded in the topological dual $`𝔄_{}^{}{}_{}{}^{}`$ of $`(𝔄_{},\overline{𝔗}_p^u)`$. We shall make use of these properties without special mention.
### 3.1 General Properties
Before proceeding to extract certain elements from $`^{}`$ to be interpreted, on the grounds of their specific properties, as representing asymptotic mixtures of particle-like quantities, we are first going to collect a number of important properties common to *all* functionals from the topological dual of $``$ whose proof does not depend on special assumptions. First of all, continuity as established in Proposition 2.16 directly carries over to functionals in $`^{}`$.
###### Lemma 3.2.
Continuous linear functionals $`\varsigma ^{}`$ have the following properties.
* The mapping $`𝖯_+^{}(\mathrm{\Lambda },x)\varsigma \left(L_{1}^{}{}_{}{}^{}\alpha _{(\mathrm{\Lambda },x)}(L_2)\right)`$ is continuous for arbitrary but fixed $`L_1\text{,}L_2𝔏`$.
* The mapping $`𝖯_+^{}(\mathrm{\Lambda },x)\varsigma \left(\alpha _{(\mathrm{\Lambda },x)}(C)\right)`$ is continuous for given $`C`$.
###### Proof.
Due to the assumed continuity of $`\varsigma `$, the assertions follow from Proposition 2.16 in connection with Corollary 2.13. ∎
Every *positive* functional $`\varsigma `$ on the -algebra $`=𝔏^{}𝔏`$ defines a non-negative sesquilinear form on $`𝔏`$ through
$$.|._\varsigma :𝔏\times 𝔏(L_1,L_2)L_1|L_2_\varsigma \varsigma (L_{1}^{}{}_{}{}^{}L_2)\text{,}$$
(3.2a)
and thus induces a seminorm $`q_\varsigma `$ on $`𝔏`$ via
$$q_\varsigma :𝔏_+Lq_\varsigma (L)L|L_\varsigma ^{1/2}\text{.}$$
(3.2b)
Denoting by $`𝔑_\varsigma `$ the null space of $`q_\varsigma `$, one can construct the quotient $`𝔏_\varsigma 𝔏/𝔑_\varsigma `$, which is a normed space through the definition
$$._\varsigma :𝔏/𝔑_\varsigma _+[L]_\varsigma [L]_\varsigma _\varsigma q_\varsigma (L)\text{,}$$
(3.2c)
where we used square brackets to designate the cosets in $`𝔏/𝔑_\varsigma `$. These concepts can be applied to formulate, parallel to Proposition 2.16, differentiability of the Poincaré automorphisms with respect to continuous positive functionals on $``$.
###### Lemma 3.3.
Let $`\varsigma `$ be a continuous positive functional on the -algebra $``$, i. e. $`\varsigma _+^{}`$. Then the restriction of the canonical homomorphism
$$𝒬_\varsigma :𝔏𝔏/𝔑_\varsigma L𝒬_\varsigma (L)[L]_\varsigma $$
to the subspace $`𝔏_0`$ is $`𝒳_{𝔏_0}`$-differentiable in the sense of Definition A.16, where
$$𝒳_{𝔏_0}=\{\mathrm{\Xi }_{L_0}:L_0𝔏_0\}$$
is the family of infinitely often differentiable mappings defined in Proposition 2.16.
###### Proof.
Due to the assumed continuity of the functional $`\varsigma `$, there exists a bounded Borel set $`\mathrm{\Delta }`$ such that, according to (3.2) in connection with Definition 3.1 and Lemma 2.14, for any $`L𝔏`$ there holds the inequality
$$[L]_\varsigma _\varsigma ^2=q_\varsigma (L)^2=\varsigma (L^{}L)\varsigma _\mathrm{\Delta }p_\mathrm{\Delta }(L^{}L)=\varsigma _\mathrm{\Delta }q_\mathrm{\Delta }(L)^2\text{.}$$
Therefore the linear operator
$$𝒬_\varsigma 𝔏_0:(𝔏_0,𝔗_q)(𝔏/𝔑_\varsigma ,._\varsigma )$$
turns out to be continuous, so that the assertion follows by an application of Corollary A.15 from the result of Proposition 2.16, stating that the mappings
$$\mathrm{\Xi }_{L_0}:𝖯_+^{}(𝔏_0,𝔗_q)(\mathrm{\Lambda },x)\mathrm{\Xi }_{L_0}(\mathrm{\Lambda },x)\alpha _{(\mathrm{\Lambda },x)}(L_0)$$
are differentiable for any $`L_0𝔏_0`$ (cf. the remark of that place). ∎
The next lemmas are concerned with integrability properties of functionals $`\varsigma ^{}`$, parallel to those established in Subsection 2.3.3. The first one, Lemma 3.4, is an immediate consequence of Lemmas 2.17 and 2.18, whereas the second one, Lemma 3.5, prepares the proof of a kind of Cluster Property for *positive* functionals in $`^{}`$, formulated in the subsequent Proposition 3.6.
###### Lemma 3.4.
Let $`\varsigma ^{}`$, $`L_1\text{,}L_2𝔏`$ and $`C`$.
* Let $`FL^1(𝖯_+^{},d\mu (\mathrm{\Lambda },x))`$ have compact support $`𝖲`$, then
$`\varsigma \left(L_{1}^{}{}_{}{}^{}\alpha _F(L_2)\right)`$ $`={\displaystyle 𝑑\mu (\mathrm{\Lambda },x)F(\mathrm{\Lambda },x)\varsigma \left(L_{1}^{}{}_{}{}^{}\alpha _{(\mathrm{\Lambda },x)}(L_2)\right)\text{,}}`$ (3.3a)
$`\varsigma \left(\alpha _F(C)\right)`$ $`={\displaystyle 𝑑\mu (\mathrm{\Lambda },x)F(\mathrm{\Lambda },x)\varsigma \left(\alpha _{(\mathrm{\Lambda },x)}(C)\right)\text{,}}`$ (3.3b)
and there hold the estimates
$`\left|\varsigma \left(L_{1}^{}{}_{}{}^{}\alpha _F(L_2)\right)\right|`$ $`F_1\varsigma _\mathrm{\Delta }q_\mathrm{\Delta }(L_1)\underset{(\mathrm{\Lambda },x)𝖲}{sup}q_\mathrm{\Delta }\left(\alpha _{(\mathrm{\Lambda },x)}(L_2)\right)\text{,}`$ (3.4a)
$`\left|\varsigma \left(\alpha _F(C)\right)\right|`$ $`F_1\varsigma _\mathrm{\Delta }\underset{(\mathrm{\Lambda },x)𝖲}{sup}p_\mathrm{\Delta }\left(\alpha _{(\mathrm{\Lambda },x)}(C)\right)`$ (3.4b)
for any $`\mathrm{\Delta }`$ such that $`\varsigma _{\mathrm{\Delta }}^{}{}_{}{}^{}`$.
* For any function $`gL^1(^{s+1},d^{s+1}x)`$
$$\varsigma \left(L_{1}^{}{}_{}{}^{}\alpha _g(L_2)\right)=_{^{s+1}}d^{s+1}xg(x)\varsigma \left(L_{1}^{}{}_{}{}^{}\alpha _x(L_2)\right)\text{,}$$
(3.5)
and a bound is given by
$$\left|\varsigma \left(L_{1}^{}{}_{}{}^{}\alpha _g(L_2)\right)\right|g_1\varsigma _\mathrm{\Delta }q_\mathrm{\Delta }(L_1)q_\mathrm{\Delta }(L_2)$$
(3.6)
for any $`\mathrm{\Delta }`$ satisfying $`\varsigma _{\mathrm{\Delta }}^{}{}_{}{}^{}`$.
###### Proof.
Lemmas 2.17 and 2.18 state that
$`\alpha _F(L_2)`$ $`={\displaystyle 𝑑\mu (\mathrm{\Lambda },x)F(\mathrm{\Lambda },x)\alpha _{(\mathrm{\Lambda },x)}(L_2)\text{,}}`$
$`\alpha _F(C)`$ $`={\displaystyle 𝑑\mu (\mathrm{\Lambda },x)F(\mathrm{\Lambda },x)\alpha _{(\mathrm{\Lambda },x)}(C)\text{,}}`$
$`\alpha _g(L_2)`$ $`={\displaystyle _{^{s+1}}}d^{s+1}xg(x)\alpha _x(L_2)`$
exist in the complete locally convex spaces $`(𝔄_𝔏,\overline{𝔗}_q^u)`$ and $`(𝔄_{},\overline{𝔗}_p^u)`$, respectively. Now, the functional $`\varsigma `$, which lies in $`𝔄_{}^{}{}_{}{}^{}`$ according to the remark following Definition 3.1, is linear and continuous with respect to $`\alpha _F(C)𝔄_{}`$ and, by Corollary 2.13, also with respect to both $`\alpha _F(L_2)\text{,}\alpha _g(L_2)𝔄_𝔏`$. Therefore it commutes with the locally convex integrals \[26, Proposition II.5.7 adapted to integrals in locally convex spaces\], which proves the assertion. The annexed estimates are a further simple application of the results contained in Lemmas 2.17 and 2.18. ∎
###### Lemma 3.5.
Let $`L^{}𝔏`$ and let $`L𝔏(\mathrm{\Gamma })=𝔏\stackrel{~}{𝔄}(\mathrm{\Gamma })`$, $`\mathrm{\Gamma }^{s+1}`$ compact, i. e. $`L`$ has energy-momentum transfer in $`\mathrm{\Gamma }`$. If $`\varsigma ^+`$ is a positive functional which belongs to $`_{\mathrm{\Delta }}^{}{}_{}{}^{}`$ and $`\mathrm{\Delta }^{}`$ denotes any bounded Borel set containing $`\mathrm{\Delta }+\mathrm{\Gamma }`$, then
$$_^sd^sx\varsigma \left(L^{}\alpha x(L_{}^{}{}_{}{}^{}L^{})L\right)\varsigma _\mathrm{\Delta }q_\mathrm{\Delta }(L)^2q_\mathrm{\Delta }^{}(L^{})^2\text{.}$$
(3.7)
###### Proof.
Let $`𝑲`$ be an arbitrary compact subset of $`^s`$ and note that
$$Kd^sx\alpha x(L_{}^{}{}_{}{}^{}L^{})𝔄\text{.}$$
Thus, according to the construction of $``$,
$$Kd^sxL^{}\alpha x(L_{}^{}{}_{}{}^{}L^{})L=L^{}Kd^sx\alpha x(L_{}^{}{}_{}{}^{}L^{})L$$
belongs to the algebra of counters and exists furthermore as an integral in the locally convex space $`(𝔄_{},\overline{𝔗}_p^u)`$. Therefore the functional $`\varsigma 𝔄_{}^{}{}_{}{}^{}`$ can be interchanged with the integral \[26, Proposition II.5.7\] to give
$$Kd^sx\varsigma \left(L^{}\alpha x(L_{}^{}{}_{}{}^{}L^{})L\right)=\varsigma \left(L^{}Kd^sx\alpha x(L_{}^{}{}_{}{}^{}L^{})L\right)\text{.}$$
Application of Lemma 2.12 then leads to the estimate
$$\begin{array}{c}0Kd^sx\varsigma \left(L^{}\alpha x(L_{}^{}{}_{}{}^{}L^{})L\right)\varsigma _\mathrm{\Delta }p_\mathrm{\Delta }\left(L^{}Kd^sx\alpha x(L_{}^{}{}_{}{}^{}L^{})L\right)\hfill \\ \hfill \varsigma _\mathrm{\Delta }q_\mathrm{\Delta }(L)^2E(\mathrm{\Delta }^{})Kd^sx\alpha x(L_{}^{}{}_{}{}^{}L^{})E(\mathrm{\Delta }^{})=\varsigma _\mathrm{\Delta }q_\mathrm{\Delta }(L)^2Q_{\mathrm{\Delta }^{},𝑲}^{(L_{}^{}{}_{}{}^{}L^{})}\text{,}\end{array}$$
where we made use of the positivity of $`\varsigma `$. The above inequality survives in the limit $`𝑲^𝒔`$ and the convergence of the right-hand side to a finite real number establishes the integrability of the function
$$^s𝒙\varsigma \left(𝑳^{}\alpha 𝒙(𝑳_{}^{}{}_{}{}^{}𝑳^{})𝑳\right)$$
as a consequence of the Monotone Convergence Theorem \[26, II.2.7\]. In view of (2.17a), one finally arrives at the asserted bound
$$_^sd^sx\varsigma \left(L^{}\alpha x(L_{}^{}{}_{}{}^{}L^{})L\right)\varsigma _\mathrm{\Delta }q_\mathrm{\Delta }(L)^2Q_\mathrm{\Delta }^{}^{(L_{}^{}{}_{}{}^{}L^{})}=\varsigma _\mathrm{\Delta }q_\mathrm{\Delta }(L)^2q_\mathrm{\Delta }^{}(L^{})^2\text{.}$$
After these preparations we are in a position to prove the announced Cluster Property for positive functionals in $`^{}`$.
###### Proposition 3.6 (Cluster Property).
Let $`L_i`$ and $`L_i^{}`$ be elements of $`𝔏_0`$ and let $`A_i𝔄`$, $`i=1\text{,}2`$, be almost local operators, then the function
$$^s𝒙\varsigma \left((𝑳_{\mathit{1}}^{}{}_{}{}^{}𝑨_\mathit{1}𝑳_\mathit{1}^{})\alpha 𝒙(𝑳_{\mathit{2}}^{}{}_{}{}^{}𝑨_\mathit{2}𝑳_\mathit{2}^{})\right)$$
is an element of $`L^1(^s,d^sx)`$ for any $`\varsigma ^+`$ and satisfies
$$_^sd^sx\left|\varsigma \left((L_{1}^{}{}_{}{}^{}A_1L_1^{})\alpha x(L_{2}^{}{}_{}{}^{}A_2L_2^{})\right)\right|\varsigma _\mathrm{\Delta }M_\mathrm{\Delta }$$
(3.8)
for any bounded Borel set $`\mathrm{\Delta }`$ for which $`\varsigma `$ belongs to $`_{\mathrm{\Delta }}^{}{}_{}{}^{}`$, where the constant $`M_\mathrm{\Delta }`$ depends on $`\mathrm{\Delta }`$ and the operators involved.
###### Proof.
First, we re-write the argument $`(L_{1}^{}{}_{}{}^{}A_1L_1^{})\alpha x(L_{2}^{}{}_{}{}^{}A_2L_2^{})`$, commuting the operators $`A_1L_1^{}`$ and $`\alpha x(L_{2}^{}{}_{}{}^{}A_2)`$, to get
$$\begin{array}{c}(L_{1}^{}{}_{}{}^{}A_1L_1^{})\alpha x(L_{2}^{}{}_{}{}^{}A_2L_2^{})\hfill \\ \hfill =L_{1}^{}{}_{}{}^{}[A_1L_1^{},\alpha x(L_{2}^{}{}_{}{}^{}A_2)]\alpha x(L_2^{})+L_{1}^{}{}_{}{}^{}\alpha x(L_{2}^{}{}_{}{}^{}A_2)A_1L_1^{}\alpha x(L_2^{})\text{.}\end{array}$$
(3.9)
This implies
$$\begin{array}{c}\left|\varsigma \left((L_{1}^{}{}_{}{}^{}A_1L_1^{})\alpha x(L_{2}^{}{}_{}{}^{}A_2L_2^{})\right)\right|\hfill \\ \hfill \left|\varsigma \left(L_{1}^{}{}_{}{}^{}[A_1L_1^{},\alpha x(L_{2}^{}{}_{}{}^{}A_2)]\alpha x(L_2^{})\right)\right|+\left|\varsigma \left(L_{1}^{}{}_{}{}^{}\alpha x(L_{2}^{}{}_{}{}^{}A_2)A_1L_1^{}\alpha x(L_2^{})\right)\right|\text{,}\end{array}$$
(3.10)
where the first term on the right-hand side is evidently integrable over $`^s`$, due to almost locality of the operators encompassed by the commutator. For $`\varsigma _{\mathrm{\Delta }}^{}{}_{}{}^{}`$ we have the estimate
$$\begin{array}{c}_^sd^sx\left|\varsigma \left(L_{1}^{}{}_{}{}^{}[A_1L_1^{},\alpha x(L_{2}^{}{}_{}{}^{}A_2)]\alpha x(L_2^{})\right)\right|\hfill \\ \hfill \varsigma _\mathrm{\Delta }q_\mathrm{\Delta }(L_1)q_\mathrm{\Delta }(L_2^{})_^sd^sx[A_1L_1^{},\alpha x(L_{2}^{}{}_{}{}^{}A_2)]\text{.}\end{array}$$
(3.11)
The second term can be estimated by use of the Cauchy-Schwarz inequality applied to the positive functional $`\varsigma `$:
$$\begin{array}{c}2\left|\varsigma \left(L_{1}^{}{}_{}{}^{}\alpha x(L_{2}^{}{}_{}{}^{}A_2)A_1L_1^{}\alpha x(L_2^{})\right)\right|\hfill \\ \hfill 2\varsigma \left(L_{1}^{}{}_{}{}^{}\alpha x(L_{2}^{}{}_{}{}^{}A_2A_{2}^{}{}_{}{}^{}L_2)L_1\right)^{1/2}\varsigma \left(\alpha x(L_{2}^{}{}_{}{}^{})L_{1}^{}{}_{}{}^{}A_{1}^{}{}_{}{}^{}A_1L_1^{}\alpha x(L_2^{})\right)^{1/2}\\ \hfill =\underset{\lambda >0}{inf}\left(\lambda ^1\varsigma \left(L_{1}^{}{}_{}{}^{}\alpha x(L_{2}^{}{}_{}{}^{}A_2A_{2}^{}{}_{}{}^{}L_2)L_1\right)+\lambda \varsigma \left(\alpha x(L_{2}^{}{}_{}{}^{})L_{1}^{}{}_{}{}^{}A_{1}^{}{}_{}{}^{}A_1L_1^{}\alpha x(L_2^{})\right)\right)\text{.}\end{array}$$
(3.12)
Integration of the first term on the right-hand side is possible according to the previous Lemma 3.5 and gives
$$_^sd^sx\varsigma \left(L_{1}^{}{}_{}{}^{}\alpha x(L_{2}^{}{}_{}{}^{}A_2A_{2}^{}{}_{}{}^{}L_2)L_1\right)\varsigma _\mathrm{\Delta }q_\mathrm{\Delta }(L_1)^2q_{\mathrm{\Delta }_1}(A_{2}^{}{}_{}{}^{}L_2)^2\text{,}$$
(3.13)
where $`\mathrm{\Delta }_1`$ is any bounded Borel set containing the sum of $`\mathrm{\Delta }`$ and the energy-momentum transfer $`\mathrm{\Gamma }_1`$ of $`L_1`$. Concerning the second term on the right of (3.12), we get, upon commuting $`\alpha x(L_{2}^{}{}_{}{}^{})`$ and $`\alpha x(L_2^{})`$ to the interior,
$$\begin{array}{c}\varsigma \left(\alpha x(L_{2}^{}{}_{}{}^{})L_{1}^{}{}_{}{}^{}A_{1}^{}{}_{}{}^{}A_1L_1^{}\alpha x(L_2^{})\right)\hfill \\ \hfill \left|\varsigma \left([\alpha x(L_{2}^{}{}_{}{}^{}),L_{1}^{}{}_{}{}^{}]A_{1}^{}{}_{}{}^{}A_1L_1^{}\alpha x(L_2^{})\right)\right|+\left|\varsigma \left(L_{1}^{}{}_{}{}^{}\alpha x(L_{2}^{}{}_{}{}^{})A_{1}^{}{}_{}{}^{}A_1[L_1^{},\alpha x(L_2^{})]\right)\right|\\ \hfill +A_1^2\left|\varsigma \left(L_{1}^{}{}_{}{}^{}\alpha x(L_{2}^{}{}_{}{}^{}L_2^{})L_1^{}\right)\right|\text{,}\end{array}$$
(3.14)
where again use was made of the positivity of $`\varsigma `$. The rapid decay of commutators of almost local operators with respect to the $`q_\mathrm{\Delta }`$-seminorm established in Lemma 2.19 of Subsection 2.3.4 can be combined with Lemma 3.5 to show integrability over $`^s`$:
$$\begin{array}{c}_^sd^sx\varsigma \left(\alpha x(L_{2}^{}{}_{}{}^{})L_{1}^{}{}_{}{}^{}A_{1}^{}{}_{}{}^{}A_1L_1^{}\alpha x(L_2^{})\right)\hfill \\ \hfill \varsigma _\mathrm{\Delta }A_1^2(q_\mathrm{\Delta }(L_1^{})^2q_{\mathrm{\Delta }_1^{}}(L_2^{})^2+(L_1^{}q_\mathrm{\Delta }(L_2^{})+L_2^{}q_\mathrm{\Delta }(L_1^{}))\\ \hfill _^sd^sxq_\mathrm{\Delta }\left([L_1^{},\alpha x(L_2^{})]\right))\text{,}\end{array}$$
(3.15)
which holds for any bounded Borel set $`\mathrm{\Delta }_1^{}\mathrm{\Delta }+\mathrm{\Gamma }_1^{}`$, where $`\mathrm{\Gamma }_1^{}`$ denotes the energy-momentum transfer of $`L_1^{}`$. By (3.14) and (3.15), the left-hand side of (3.12) turns out to be integrable, and a bound for this integral is proportional to $`\varsigma _\mathrm{\Delta }`$. In connection with (3.11) this establishes the assertion for a suitable constant $`M_\mathrm{\Delta }`$ that can be deduced from relations (3.11), (3.12), (3.14) and (3.15). ∎
The Cluster Property has been proved above under the fairly general assumption of almost locality of the operators involved. If for given $`L_1\text{,}L_2𝔏`$ the mapping
$$^s𝒙𝒑_𝜟\left(𝑳_{\mathit{1}}^{}{}_{}{}^{}\alpha 𝒙(𝑳_\mathit{2})\right)$$
happens to belong to the space $`L^1(^s,d^sx)`$ for the bounded Borel set $`\mathrm{\Delta }`$, (3.8) is obviously fulfilled in case that $`\varsigma ^{}`$ belongs to $`_{\mathrm{\Delta }}^{}{}_{}{}^{}`$. As an example consider almost local operators $`L_1^{}\text{,}L_2^{}𝔏`$ having energy-momentum transfer $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$, respectively, such that $`(\mathrm{\Delta }+\mathrm{\Gamma }_1+\mathrm{\Gamma }_2)\overline{V}_+=\mathrm{}`$. This implies $`L_1^{}\alpha x(L_2^{})E(\mathrm{\Delta })=0`$ for any $`𝒙^𝒔`$ and, by Lemmas 2.11 and 2.12, $`p_\mathrm{\Delta }\left(L_{1}^{}{}_{}{}^{}\alpha x(L_{2}^{}{}_{}{}^{})L_1^{}\alpha x(L_2^{})\right)=0`$. An application of Lemma 2.12 in connection with translation invariance of $`q_\mathrm{\Delta }`$ (Lemma 2.15) then yields for the counters $`C_i^{}L_{i}^{}{}_{}{}^{}L_i^{}`$, $`i=1\text{,}2`$,
$$\begin{array}{c}p_\mathrm{\Delta }\left(C_{1}^{}{}_{}{}^{}\alpha x(C_2^{})\right)=p_\mathrm{\Delta }\left(L_{1}^{}{}_{}{}^{}L_1^{}\alpha x(L_{2}^{}{}_{}{}^{}L_2^{})\right)\hfill \\ \hfill =p_\mathrm{\Delta }\left(L_{1}^{}{}_{}{}^{}[L_1^{},\alpha x(L_{2}^{}{}_{}{}^{})]\alpha x(L_2^{})\right)[L_1^{},\alpha x(L_{2}^{}{}_{}{}^{})]q_\mathrm{\Delta }(L_1^{})q_\mathrm{\Delta }(L_2^{})\text{,}\end{array}$$
where, due to the assumed almost locality of $`L_1^{}`$ and $`L_2^{}`$, the right-hand side is seen to belong to $`L^1(^s,d^sx)`$. The integrability of a mapping $`𝒙𝒑_𝜟\left(𝑳_{\mathit{1}}^{}{}_{}{}^{}\alpha 𝒙(𝑳_\mathit{2})\right)`$, $`L_1\text{,}L_2𝔏`$, has another consequence concerning weakly convergent nets $`\{\varsigma _\iota :\iota J\}`$ of functionals from $`^{}`$, which are contained in bounded subsets of $`_{\mathrm{\Delta }}^{}{}_{}{}^{}`$ with respect to the norm $`._\mathrm{\Delta }`$: a kind of Dominated Convergence Theorem.
###### Lemma 3.7.
Let $`L_1\text{,}L_2𝔏`$ be such that $`𝐱𝐩_𝚫\left(𝐋_{\mathrm{𝟏}}^{}{}_{}{}^{}\alpha 𝐱(𝐋_\mathrm{𝟐})\right)`$ is integrable and consider the weakly convergent net $`\{\varsigma _\iota :\iota J\}`$ in the $`D`$-ball of $`_{\mathrm{\Delta }}^{}{}_{}{}^{}`$ with limit $`\varsigma `$.This means that for any $`C`$
$$\underset{\iota }{lim}\varsigma _\iota (C)=\varsigma (C)$$
and for any $`\iota J`$
$`|\varsigma _\iota (C)|`$ $`Dp_\mathrm{\Delta }(C)\text{,}`$ (3.16a)
$`|\varsigma (C)|`$ $`Dp_\mathrm{\Delta }(C)\text{,}`$ (3.16b)
the latter relation being implied by the former. Then
$$_^sd^sx\varsigma \left(L_{1}^{}{}_{}{}^{}\alpha x(L_2)\right)=\underset{\iota }{lim}_^sd^sx\varsigma _\iota \left(L_{1}^{}{}_{}{}^{}\alpha x(L_2)\right)\text{.}$$
(3.17)
###### Proof.
As implied by Proposition 2.16 and Corollary 2.13, $`𝒙𝑳_{\mathit{1}}^{}{}_{}{}^{}\alpha 𝒙(𝑳_\mathit{2})`$ is a continuous mapping on $`^s`$ with respect to the $`p_\mathrm{\Delta }`$-topology, hence it is uniformly continuous on any compact set $`𝑲`$. This means that to $`ϵ>0`$ there exists $`\delta >0`$ such that $`𝒙\text{,}𝒙^{}𝑲`$ and $`|𝒙𝒙^{}|<\delta `$ imply
$$p_\mathrm{\Delta }\left(L_{1}^{}{}_{}{}^{}\alpha x(L_2)L_{1}^{}{}_{}{}^{}\alpha _𝒙^{}(L_2)\right)<\frac{ϵ}{6D|𝑲|}\text{,}$$
where $`|𝑲|`$ denotes the $`s`$-dimensional volume of $`𝑲`$. Consequently, under the above assumption on $`𝒙`$ and $`𝒙^{}`$, we infer from (3.16)
$`\left|\varsigma _\iota \left(L_{1}^{}{}_{}{}^{}\alpha x(L_2)\right)\varsigma _\iota \left(L_{1}^{}{}_{}{}^{}\alpha _𝒙^{}(L_2)\right)\right|`$ $`=\left|\varsigma _\iota \left(L_{1}^{}{}_{}{}^{}\alpha x(L_2)L_{1}^{}{}_{}{}^{}\alpha _𝒙^{}(L_2)\right)\right|<{\displaystyle \frac{ϵ}{6|𝑲|}}\text{,}`$
$`\left|\varsigma \left(L_{1}^{}{}_{}{}^{}\alpha x(L_2)\right)\varsigma \left(L_{1}^{}{}_{}{}^{}\alpha _𝒙^{}(L_2)\right)\right|`$ $`=\left|\varsigma \left(L_{1}^{}{}_{}{}^{}\alpha x(L_2)L_{1}^{}{}_{}{}^{}\alpha _𝒙^{}(L_2)\right)\right|<{\displaystyle \frac{ϵ}{6|𝑲|}}\text{.}`$
By compactness of $`𝑲`$, there exist finitely many elements $`𝒙_\mathit{1}\text{,}\mathrm{}\text{,}𝒙_𝑵𝑲`$ such that the $`\delta `$-balls around these points cover all of $`𝑲`$; moreover, since $`\varsigma `$ is the weak limit of the net $`\{\varsigma _\iota :\iota J\}`$, we can find $`\iota _0J`$ such that $`\iota \iota _0`$ implies
$$\left|\varsigma \left(L_{1}^{}{}_{}{}^{}\alpha _{𝒙_𝒊}(L_2)\right)\varsigma _\iota \left(L_{1}^{}{}_{}{}^{}\alpha _{𝒙_𝒊}(L_2)\right)\right|<\frac{ϵ}{6|𝑲|}$$
for any $`i=1\text{,}\mathrm{}\text{,}N`$. Now, for $`\iota J`$ and $`k\{1,\mathrm{},N\}`$,
$$\begin{array}{c}\left|\varsigma \left(L_{1}^{}{}_{}{}^{}\alpha x(L_2)\right)\varsigma _\iota \left(L_{1}^{}{}_{}{}^{}\alpha x(L_2)\right)\right|\hfill \\ \hfill \left|\varsigma \left(L_{1}^{}{}_{}{}^{}\alpha x(L_2)\right)\varsigma \left(L_{1}^{}{}_{}{}^{}\alpha _{𝒙_𝒊}(L_2)\right)\right|+\left|\varsigma \left(L_{1}^{}{}_{}{}^{}\alpha _{𝒙_𝒊}(L_2)\right)\varsigma _\iota \left(L_{1}^{}{}_{}{}^{}\alpha _{𝒙_𝒊}(L_2)\right)\right|\\ \hfill +\left|\varsigma _\iota \left(L_{1}^{}{}_{}{}^{}\alpha _{𝒙_𝒊}(L_2)\right)\varsigma _\iota \left(L_{1}^{}{}_{}{}^{}\alpha x(L_2)\right)\right|\text{,}\end{array}$$
and, selecting for $`𝒙𝑲`$ an appropriate $`𝒙_𝒌`$ in a distance less than $`\delta `$, we can put the above results together to get the estimate
$$\left|\varsigma \left(L_{1}^{}{}_{}{}^{}\alpha x(L_2)\right)\varsigma _\iota \left(L_{1}^{}{}_{}{}^{}\alpha x(L_2)\right)\right|<\frac{ϵ}{2|𝑲|}\text{,}$$
which holds for any $`𝒙𝑲`$ and $`\iota \iota _0`$. Thus weak (i. e. pointwise) convergence of the net $`\{\varsigma _\iota :\iota J\}`$ is indeed uniform convergence on compact subsets of $`^s`$. Upon integration over $`𝑲`$ we arrive at
$$\begin{array}{c}\left|Kd^sx\left(\varsigma \left(L_{1}^{}{}_{}{}^{}\alpha x(L_2)\right)\varsigma _\iota \left(L_{1}^{}{}_{}{}^{}\alpha x(L_2)\right)\right)\right|\hfill \\ \hfill Kd^sx\left|\varsigma \left(L_{1}^{}{}_{}{}^{}\alpha x(L_2)\right)\varsigma _\iota \left(L_{1}^{}{}_{}{}^{}\alpha x(L_2)\right)\right|<\frac{ϵ}{2}\text{.}\end{array}$$
(3.18)
Now, by assumption
$$_^sd^sxp_\mathrm{\Delta }\left(L_{1}^{}{}_{}{}^{}\alpha x(L_2)\right)<\mathrm{}\text{,}$$
so that to $`ϵ>0`$ there exists a compact subset $`𝑲_ϵ`$ satisfying
$$_{\mathrm{}𝑲_ϵ}d^sxp_\mathrm{\Delta }\left(L_{1}^{}{}_{}{}^{}\alpha x(L_2)\right)<\frac{ϵ}{4D}\text{.}$$
Then, as a consequence of (3.16a) and (3.16b), we get for any $`\iota J`$
$`{\displaystyle _{\mathrm{}𝑲_ϵ}}d^sx\left|\varsigma _\iota \left(L_{1}^{}{}_{}{}^{}\alpha x(L_2)\right)\right|`$ $`<{\displaystyle \frac{ϵ}{4}}\text{,}`$ (3.19a)
$`{\displaystyle _{\mathrm{}𝑲_ϵ}}d^sx\left|\varsigma \left(L_{1}^{}{}_{}{}^{}\alpha x(L_2)\right)\right|`$ $`<{\displaystyle \frac{ϵ}{4}}\text{.}`$ (3.19b)
Combining (3.19) with (3.18) for the compact set $`𝑲_ϵ`$ yields for $`\iota \iota _0`$ (note, that $`\iota _0`$ only depends on $`ϵ`$)
$$\begin{array}{c}\left|_^sd^sx\varsigma \left(L_{1}^{}{}_{}{}^{}\alpha x(L_2)\right)_^sd^sx\varsigma _\iota \left(L_{1}^{}{}_{}{}^{}\alpha x(L_2)\right)\right|\hfill \\ \hfill _{\mathrm{}𝑲_ϵ}d^sx\left|\varsigma \left(L_{1}^{}{}_{}{}^{}\alpha x(L_2)\right)\right|+Kd^sx\left|\varsigma \left(L_{1}^{}{}_{}{}^{}\alpha x(L_2)\right)\varsigma _\iota \left(L_{1}^{}{}_{}{}^{}\alpha x(L_2)\right)\right|\\ \hfill +_{\mathrm{}𝑲_ϵ}d^sx\left|\varsigma _\iota \left(L_{1}^{}{}_{}{}^{}\alpha x(L_2)\right)\right|<\frac{ϵ}{4}+\frac{ϵ}{2}+\frac{ϵ}{4}=ϵ\text{.}\end{array}$$
By arbitrariness of $`ϵ`$ this proves the possibility to interchange integration and the limit with respect to $`\iota `$ as asserted in (3.17). ∎
The spectral support of not necessarily positive functionals $`\varsigma ^{}`$ (considered as distributions) depends, as expressed in the subsequent proposition, on the bounded Borel sets $`\mathrm{\Delta }`$ for which $`\varsigma _{\mathrm{\Delta }}^{}{}_{}{}^{}`$. This property will prove to be of importance when it comes to defining the energy-momentum of particle weights.
###### Proposition 3.8 (Spectral Property).
Let $`L_1\text{,}L_2𝔏`$ and $`\varsigma ^{}`$. Then the support of the Fourier transform of the distribution
$$^{s+1}x\varsigma \left(L_{1}^{}{}_{}{}^{}\alpha _x(L_2)\right)$$
is contained in the shifted light cone $`\overline{V}_+q`$ for some $`q\overline{V}_+`$. More specifically, $`q`$ is such that a bounded Borel set $`\mathrm{\Delta }`$, satisfying $`\varsigma _{\mathrm{\Delta }}^{}{}_{}{}^{}`$, is contained in $`q\overline{V}_+`$.
###### Proof.
If a function $`g`$ belongs to the space $`L^1(^{s+1},d^{s+1}x)`$, then the operator
$$\alpha _g(L_2)=_{^{s+1}}d^{s+1}xg(x)\alpha _x(L_2)$$
lies in $`𝔄_𝔏`$, according to Lemma 2.18, and has energy-momentum transfer in $`\mathrm{supp}\stackrel{~}{g}`$, the support of the Fourier transform of $`g`$. If this happens to satisfy $`\mathrm{supp}\stackrel{~}{g}\mathrm{}(\overline{V}_+\mathrm{\Delta })`$, we infer $`\alpha _g(L_2)E(\mathrm{\Delta })=0`$ and henceforth, by Lemma 2.11, $`q_\mathrm{\Delta }\left(\alpha _g(L_2)\right)=0`$. Since $`\varsigma `$ is assumed to belong to $`_{\mathrm{\Delta }}^{}{}_{}{}^{}`$, Lemma 3.4 results in
$$\left|_{^{s+1}}d^{s+1}xg(x)\varsigma \left(L_{1}^{}{}_{}{}^{}\alpha _x(L_2)\right)\right|=\left|\varsigma \left(L_{1}^{}{}_{}{}^{}\alpha _g(L_2)\right)\right|\varsigma _\mathrm{\Delta }q_\mathrm{\Delta }(L_1)q_\mathrm{\Delta }\left(\alpha _g(L_2)\right)\text{,}$$
which, according to the preceding considerations, entails
$$_{^{s+1}}d^{s+1}xg(x)\varsigma \left(L_{1}^{}{}_{}{}^{}\alpha _x(L_2)\right)=0\text{.}$$
(3.20)
Now, let $`g^{}`$ be an arbitrary function from $`L^1(^{s+1},d^{s+1}x)`$ with $`\mathrm{supp}\stackrel{~}{g^{}}\mathrm{}(\overline{V}_+q)`$, $`\mathrm{\Delta }q\overline{V}_+`$, then $`\mathrm{supp}\stackrel{~}{g^{}}\mathrm{}(\overline{V}_+\mathrm{\Delta })`$, so that (3.20) is fulfilled for any function of this kind, proving the assertion. ∎
### 3.2 Asymptotic Functionals
Now we turn to functionals in $`^{}`$ that carry additional properties, reflecting the fact that the present investigation is concerned with the structure of the totality of physical states at asymptotic times (scattering states). The temporal development of such a state of bounded energy, $`\omega 𝒮(\mathrm{\Delta })`$, $`\mathrm{\Delta }`$ a bounded Borel set, can be explored by considering an integral of the following shape:
$$_^sd^svh(𝒗)\omega \left(\alpha _{(\tau ,\tau 𝒗)}(𝑪)\right)\text{,}$$
(3.21)
where $`h`$ denotes a bounded measurable function on the unit ball of $`^s`$, where the elements $`𝒗`$ represent velocities. Apart from this function, (3.21) coincides with the integral (2.5) encountered on page 2.5 in the heuristic considerations of Chapter 2. The investigations carried through in that part (cf. Proposition 2.7) imply that (3.21) takes on a finite value for any counter $`C`$ at any time $`\tau `$ and, according to Lemma 2.10, the integral (3.21) even exists for all $`C𝔄_{}`$.
The physical interpretation is as follows: Consider a function $`h`$ of bounded support $`𝖵^s\{0\}`$ in velocity space, then the integral (3.21) corresponds to summing up, for given time $`\tau `$, the expectation values of measurements of $`C`$ in the state $`\omega `$, where these measurements extend over the bounded section $`\tau 𝖵`$ of configuration space. For growing $`\tau `$ the distance of this portion from the origin increases together with its total extension. More exactly, the measurements take place in a cone with apex at the point $`0`$ of space-time, its direction is determined by the support of $`h`$, and for different times $`\tau `$ the counter $`C`$ is set up in specific parts of that cone, their extension growing as $`|\tau |^s`$ (compensating for the quantum mechanical spreading of wave packets) while their distance from the origin increases proportional to $`|\tau |`$. If the physical state $`\omega `$ has, in the limit of large (positive or negative) times, evolved into a configuration containing a particle (incoming or outgoing) travelling with velocity $`𝒗_\mathit{0}𝖵`$, then a counter $`C_0`$, sensitive for that specific particle, is expected to asymptotically produce a stable signal under the above experimental conditions.
The mathematical equivalent of this situation is the existence of limits of the above integral at asymptotic times, evaluated for the counter $`C_0`$ and a function $`h_0`$ with support containing $`𝒗_\mathit{0}`$. Thus the problem has to be settled in which (topological) sense such limits can be established, if they happen to exist at all. To tackle this assignment we turn to a slightly modified version of (3.21) in Definition 3.9, involving, for technical reasons, a certain time average.
###### Definition 3.9.
Let $`\mathrm{\Delta }`$ be a bounded Borel subset of $`^{s+1}`$, let $`\omega 𝒮(\mathrm{\Delta })`$ denote a physical state of bounded energy and let $`𝒗𝒉(𝒗)`$ be a bounded measurable function on the unit ball of $`^s`$. Furthermore suppose that $`tT(t)`$ is a continuous real-valued function, approaching $`+\mathrm{}`$ or $`\mathrm{}`$ for asymptotic positive or negative times, respectively, not as fast as $`|t|`$. Then we define a net $`\{\rho _{h,t}:t\}`$ of linear functionals on $``$ by setting
$$\begin{array}{cc}\hfill \rho _{h,t}(C)& T(t)^1_t^{t+T(t)}𝑑\tau \tau ^s_^sd^svh(𝒗)\omega \left(\alpha _{(\tau ,\tau 𝒗)}(𝑪)\right)\hfill \\ & =T(t)^1_t^{t+T(t)}𝑑\tau _^sd^sxh(\tau ^1𝒙)\omega \left(\alpha _{(\tau ,𝒙)}(𝑪)\right)\text{,}𝑪\text{.}\hfill \end{array}$$
(3.22)
Under the above assumptions the functionals $`\rho _{h,t}`$ turn out to be continuous with respect to the seminorm $`p_\mathrm{\Delta }`$ pertaining to the energy-momentum support of the physical state $`\omega (.)=\omega (E(\mathrm{\Delta }).E(\mathrm{\Delta }))`$, i. e. $`\rho _{h,t}_{\mathrm{\Delta }}^{}{}_{}{}^{}`$. This can be seen as follows: First, note that the operators $`U(\tau )`$ implementing time translations commute with $`E(\mathrm{\Delta })`$, so that
$$\omega \left(E(\mathrm{\Delta })\alpha _{(\tau ,𝒙)}(C)E(\mathrm{\Delta })\right)=\omega \left(U(\tau )E(\mathrm{\Delta })\alpha x(C)E(\mathrm{\Delta })U(\tau )^{}\right)\text{,}$$
which allows (3.22) to be re-written as
$$\rho _{h,t}(C)=T(t)^1_t^{t+T(t)}𝑑\tau _^sd^sxh(\tau ^1𝒙)\omega \left(𝑼(\tau )𝑬(𝜟)\alpha 𝒙(𝑪)𝑬(𝜟)𝑼(\tau )^{}\right)\text{.}$$
(3.23)
Now, all the functionals $`\omega (U(\tau ).U(\tau )^{})`$, $`\tau `$, belong to $`𝔅()_{,1}`$, so that the absolute value of $`\rho _{h,t}(C)`$ can be estimated, making use of $`p_\mathrm{\Delta }`$ as defined in (2.17b). Abbreviating the interval of $`\tau `$-integration depending on $`t`$ as $`I_t`$, this gives
$$\begin{array}{c}\left|\rho _{h,t}(C)\right|\underset{\tau I_t}{sup}\left|_^sd^sxh(\tau ^1𝒙)\omega \left(𝑼(\tau )𝑬(𝜟)\alpha 𝒙(𝑪)𝑬(𝜟)𝑼(\tau )^{}\right)\right|\hfill \\ \hfill h_{\mathrm{}}\underset{\varphi 𝔅()_{,1}}{sup}_^sd^sx\left|\varphi \left(E(\mathrm{\Delta })\alpha x(C)E(\mathrm{\Delta })\right)\right|=h_{\mathrm{}}p_\mathrm{\Delta }(C)\text{.}\end{array}$$
(3.24)
The above inequality implies that the functionals $`\rho _{h,t}`$ belong to the dual space $`^{}`$ of $`(,𝔗_p)`$. Moreover, the estimate (3.24) is uniform in $`t`$, so that the net $`\{\rho _{h,t}:t\}`$ is even an equicontinuous subset of $`^{}`$. The Theorem of Alaoğlu-Bourbaki \[41, Theorem 8.5.2\] then tells us, that this net is relatively compact with respect to the weak topology, leading to the following fundamental result.
###### Theorem 3.10 (Existence of Limits).
Under the assumptions of Definition 3.9 the net $`\{\rho _{h,t}:t\}_{\mathrm{\Delta }}^{}{}_{}{}^{}`$ possesses weak limit points in $`^{}`$ at asymptotic times. This means that there exist functionals $`\sigma _{h,\omega }^{(+)}`$ and $`\sigma _{h,\omega }^{()}`$ on $``$ together with corresponding subnets $`\{\rho _{h,t_\iota }:\iota J\}`$ and $`\{\rho _{h,t_\kappa }:\kappa K\}`$, i. e. $`lim_\iota t_\iota =+\mathrm{}`$ and $`lim_\kappa t_\kappa =\mathrm{}`$, such that for arbitrary $`C`$
$`\rho _{h,t_\iota }(C)`$ $`\underset{𝜄}{}\sigma _{h,\omega }^{(+)}(C)\text{,}`$ (3.25a)
$`\rho _{h,t_\kappa }(C)`$ $`\underset{𝜅}{}\sigma _{h,\omega }^{()}(C)\text{.}`$ (3.25b)
The heuristic picture laid open above suggests, that in theories which are reasonable from a physicist’s point of view the net $`\{\rho _{h,t}:t\}`$ actually converges, but as yet we have not been able to give rigorously formulated conditions under which to prove this conjecture. This question seems to be connected with the problem of asymptotic completeness of quantum field theoretic models; one has to assure that in the limit of large times multiple scattering does no longer withhold the measurement results $`\rho _{h,t}(C)`$ from growing stable. Another possibility is the disappearance of the limit functionals $`\sigma _{h,\omega }^{(+)}`$ and $`\sigma _{h,\omega }^{()}`$ on all of the algebra of counters $``$, a phenomenon that we anticipate to encounter in theories without a particle interpretation (e. g. generalized free field). The denomination of the asymptotic functionals ‘$`\sigma `$’ is chosen to reflect their *singular* nature: the values that the functionals $`\rho _{h,t}`$ return for finite times $`t`$ when applied to the identity operator $`\mathrm{𝟏}`$ (which is not contained in $``$) are divergent as $`|t|^s`$ at asymptotic times.
The convergence problem as yet only partially solved in the sense of Theorem 3.10, one can nevertheless establish a number of distinctive properties of the limit functionals $`\sigma `$ (from now on we will skip sub- and superscripts not to overburden the notation), that allow for their interpretation in terms of asymptotic configurations of particles. An immediate first consequence of the above construction is the following proposition.
###### Proposition 3.11 (Positivity and Continuity of Limits).
Suppose that $`\mathrm{\Delta }`$ is a bounded Borel subset of $`^{s+1}`$, $`\omega 𝒮(\mathrm{\Delta })`$ a physical state of bounded energy and $`hL^{\mathrm{}}(^s,d^sx)`$ a non-negative function. Then the limit functionals $`\sigma `$ for the net $`\{\rho _{h,t}:t\}`$ are positive elements of $`_{\mathrm{\Delta }}^{}{}_{}{}^{}`$:
$`\left|\sigma (C)\right|`$ $`h_{\mathrm{}}p_\mathrm{\Delta }(C)\text{,}C\text{;}`$ (3.26a)
$`0`$ $`\sigma (C)\text{,}C^+\text{.}`$ (3.26b)
###### Remark.
Due to the continuity of $`\rho _{h,t}`$ and $`\sigma `$ with respect to the $`p_\mathrm{\Delta }`$-topology, these functionals can be continuously extended to $`\overline{}`$ as well as $`𝔄_{}`$, where $`\rho _{h,t}`$ are explicitly given on $`𝔄_{}`$ by the formula (3.22) with $`C𝔄_{}`$. It is then easily established, by use of elements $`C^{}`$ from $``$ lying in suitable $`p_\mathrm{\Delta }`$-neighbourhoods of $`C`$, that the relations (3.25) remain valid on this larger subspace of the quasi-local algebra $`𝔄`$.
The next result deals with the effect that space-time translations exert on these limit functionals. A further assumption on the velocity implementation $`hL^{\mathrm{}}(^s,d^sv)`$ turns out to be indispensible in their proof: $`h`$ has to be continuous, approximating a constant value in the limit $`|𝒗|\mathrm{}`$, i. e. $`hM_hC_0(^s)`$ for a suitable constant $`M_h`$; these functions constitute a subspace of $`C(^s)`$ that will be denoted $`C_{0,c}(^s)`$ in the sequel.
###### Proposition 3.12 (Translation Invariance).
Let $`\mathrm{\Delta }^{s+1}`$ be a bounded Borel set, let $`\omega 𝒮(\mathrm{\Delta })`$ and $`hC_{0,c}(^s)`$. Then the limit functionals $`\sigma `$ of $`\{\rho _{h,t_\iota }:\iota J\}`$ are invariant under space-time translations:
$$\sigma \left(\alpha _x(C)\right)=\sigma (C)$$
(3.27)
for any $`C𝔄_{}`$ and any $`x^{s+1}`$.
###### Proof.
Taking into account the fact that the Lebesgue measure on $`^{s+1}`$ is invariant under translations, one can express $`\rho _{h,t}\left(\alpha _{(x^0,𝒙)}(C)\right)`$ for any finite time $`t`$ and any given $`x=(x^0,𝒙)^{𝒔+\mathit{1}}`$ by the following integral
$$\rho _{h,t}\left(\alpha _{(x^0,𝒙)}(C)\right)=T(t)^1_{t+x^0}^{t+x^0+T(t)}𝑑\tau _^sd^syh\left((\tau x^0)^1(𝒚𝒙)\right)\omega \left(\alpha _{(\tau ,𝒚)}(𝑪)\right)\text{.}$$
Next, we want to evaluate $`\left|\rho _{h,t}(C)\rho _{h,t}\left(\alpha _{(x^0,𝒙)}(C)\right)\right|`$ which, according to the respective limits of $`\tau `$-integration, can be split into a sum of three integrals to be estimated separately:
$`\left|T(t)^1{\displaystyle _t^{t+x^0}}𝑑\tau {\displaystyle _^s}d^syh(\tau ^1𝒚)\omega \left(\alpha _{(\tau ,𝒚)}(𝑪)\right)\right|`$ $`|T(t)|^1|x^0|h_{\mathrm{}}p_\mathrm{\Delta }(C)\text{,}`$
$`\left|T(t)^1{\displaystyle _{t+x^0+T(t)}^{t+T(t)}}𝑑\tau {\displaystyle _^s}d^syh(\tau ^1𝒚)\omega \left(\alpha _{(\tau ,𝒚)}(𝑪)\right)\right|`$ $`|T(t)|^1|x^0|h_{\mathrm{}}p_\mathrm{\Delta }(C)\text{;}`$
both $`\rho _{h,t}(C)`$ and $`\rho _{h,t}\left(\alpha _{(x^0,𝒙)}(C)\right)`$ contribute to the third integral
$$\begin{array}{c}\left|T(t)^1_{t+x^0}^{t+x^0+T(t)}𝑑\tau _^sd^sy\left[h(\tau ^1𝒚)𝒉\left((\tau 𝒙^\mathit{0})^\mathit{1}(𝒚𝒙)\right)\right]\omega \left(\alpha _{(\tau ,𝒚)}(𝑪)\right)\right|\hfill \\ \hfill \underset{\tau I_{t,x^0}}{sup}\underset{𝒚^𝒔}{sup}\left|h(\tau ^1𝒚)𝒉\left((\tau 𝒙^\mathit{0})^\mathit{1}(𝒚𝒙)\right)\right|𝒑_𝜟(𝑪)\text{,}\end{array}$$
where we used the abbreviation $`I_{t,x^0}`$ for the interval of $`\tau `$-integration. Setting (for $`|\tau |`$ large enough)
$$𝒛_\tau 𝒛+(\tau 𝒙^\mathit{0})^\mathit{1}(𝒙^\mathit{0}𝒛𝒙)$$
we finally arrive at the estimate
$$\left|\rho _{h,t}(C)\rho _{h,t}\left(\alpha _{(x^0,𝒙)}(C)\right)\right|\left(2|T(t)|^1|x^0|h_{\mathrm{}}+\underset{\tau I_{t,x^0}}{sup}\underset{𝒛^𝒔}{sup}|h(𝒛)𝒉(𝒛_\tau )|\right)𝒑_𝜟(𝑪)\text{.}$$
(3.28)
The net $`\{𝒛_\tau :\tau \}`$ approximates $`𝒛`$ uniformly on compact subsets of $`^s`$ in the limit of large $`|\tau |`$, i. e. given $`ϵ>0`$ and $`R^{}>0`$ there exists a positive number $`T^{}`$ such that $`|\tau |>T^{}`$ implies $`|𝒛𝒛_\tau |<ϵ`$ for any $`𝒛^𝒔`$ with $`|𝒛|𝑹^{}`$. On the other hand, given $`R^{\prime \prime }>0`$ there exists $`T^{\prime \prime }>0`$ such that $`|𝒛_\tau |>\frac{\mathit{1}}{\mathit{2}}𝑹^{\prime \prime }`$ for any $`|𝒛|>𝑹^{\prime \prime }`$ and any $`|\tau |>T^{\prime \prime }`$. Combining these results with the special properties of $`hC_{0,c}(^s)`$, i. e. uniform continuity on compact balls in $`^s`$ and approximate constancy at infinity, we infer that for large $`|\tau |`$ the term $`sup_{𝒛^𝒔}\left|h(𝒛)𝒉(𝒛_\tau )\right|`$ falls below any given positive bound. Therefore the right-hand side of (3.28) vanishes with $`|t|\mathrm{}`$ since $`|T(t)|`$ exceeds any positive value in this limit.
Now, let $`\sigma `$ be the weak limit of the subnet $`\{\rho _{h,t_\iota }:\iota J\}`$, i. e.
$$\rho _{h,t_\iota }(C)\underset{𝜄}{}\sigma (C)\text{,}$$
then there holds for any $`\iota J`$, any $`C𝔄_{}`$ and any $`x^{s+1}`$ the subsequent inequality
$$\begin{array}{c}0\left|\sigma \left(\alpha _x(C)\right)\sigma (C)\right|\hfill \\ \hfill \left|\sigma \left(\alpha _x(C)\right)\rho _{h,t_\iota }\left(\alpha _x(C)\right)\right|+\left|\rho _{h,t_\iota }\left(\alpha _x(C)\right)\rho _{h,t_\iota }(C)\right|+\left|\rho _{h,t_\iota }(C)\sigma (C)\right|\text{.}\end{array}$$
By the reasoning of the preceding paragraph and the above condition for subnet convergence, all three terms on the right-hand side vanish with respect to the directed set $`J`$, since in this limit $`|t_\iota |\mathrm{}`$. As a result the intermediate term has to be equal to $`0`$, thereby establishing translation invariance of $`\sigma `$. ∎
The last property that we are going to demonstrate in this section for those special elements $`\sigma _{\mathrm{\Delta }}^{}{}_{}{}^{+}`$, that arise as limits of nets of functionals $`\{\rho _{h,t_\iota }:\iota J\}`$, complements the Cluster Property 3.6. It asserts, given certain specific operators $`C`$, the existence of *lower* bounds for integrals of the functions $`𝒙\sigma \left(𝑪^{}\alpha 𝒙(𝑪)\right)`$.
###### Proposition 3.13 (Existence of Lower Bounds).
Let $`C`$ be a counter which has the property that the function $`𝐱𝐩_𝚫\left(𝐂^{}\alpha 𝐱(𝐂)\right)`$ is integrable (cf. Lemma 3.7). Let furthermore $`\sigma _{\mathrm{\Delta }}^{}{}_{}{}^{+}`$ be the limit of a net of functionals $`\{\rho _{h,t_\iota }:\iota J\}`$, each defined by (3.22), where the velocity function $`h`$ is non-negative and belongs to $`C_{0,c}(^s)`$. Under these assumptions
$$\left|\sigma (C)\right|^2h_{\mathrm{}}_^sd^sx\sigma \left(C^{}\alpha x(C)\right)\text{.}$$
(3.29)
###### Proof.
Consider the functional $`\rho _{h,t}`$ at finite time $`t`$. Applying to the absolute value of its defining equation (3.22) the Cauchy-Schwarz inequality with respect to the inner product ($`|t|`$ large enough)
$$(f,g)_tT(t)^1_t^{t+T(t)}𝑑\tau \overline{f(\tau )}g(\tau )$$
of square-integrable functions $`f`$ and $`g`$ depending on the time variable $`\tau I_t`$, one gets in the special case of
$$f(\tau )1\text{and}g(\tau )=_^sd^sxh(\tau ^1𝒙)\omega \left(\alpha _{(\tau ,𝒙)}(𝑪)\right)$$
the estimate
$$\begin{array}{c}\left|\rho _{h,t}(C)\right|^2=\left|T(t)^1_t^{t+T(t)}𝑑\tau _^sd^sxh(\tau ^1𝒙)\omega \left(\alpha _{(\tau ,𝒙)}(𝑪)\right)\right|^\mathit{2}\hfill \\ \hfill T(t)^1_t^{t+T(t)}𝑑\tau \left|_^sd^sxh(\tau ^1𝒙)\omega \left(\alpha _\tau \left(\alpha 𝒙(𝑪)\right)\right)\right|^\mathit{2}\text{.}\end{array}$$
(3.30)
Now, let $`𝑲`$ be a compact subset of $`^s`$; then, by positivity of the functional $`\omega 𝒮(\mathrm{\Delta })`$, \[11, Proposition 2.3.11(b)\] together with the Fubini Theorem \[26, II.16.3\] leads for arbitrary $`\tau `$ to
$$\begin{array}{c}\left|\omega \left(\alpha _\tau \left(Kd^sxh(\tau ^1𝒙)\alpha 𝒙(𝑪)\right)\right)\right|^\mathit{2}\hfill \\ \hfill \omega \left(\alpha _\tau \left(Kd^sxKd^syh(\tau ^1𝒚)𝒉(\tau ^\mathit{1}𝒙)\alpha 𝒚(𝑪^{})\alpha 𝒙(𝑪)\right)\right)\text{,}\end{array}$$
which is preserved in the limit $`𝑲^𝒔`$, which exists on account of the assumed integrability of the mapping $`𝒙𝒑_𝜟\left(𝑪^{}\alpha 𝒙(𝑪)\right)`$. On commuting $`\omega \alpha _\tau `$ and the integrals one arrives at
$$\begin{array}{c}\left|_^sd^sxh(\tau ^1𝒙)\omega \left(\alpha _\tau \left(\alpha 𝒙(𝑪)\right)\right)\right|^\mathit{2}\hfill \\ \hfill _^sd^sx_^sd^syh(\tau ^1𝒚)𝒉(\tau ^\mathit{1}𝒙)\omega \left(\alpha _\tau \left(\alpha 𝒚(𝑪^{})\alpha 𝒙(𝑪)\right)\right)𝒉_{\mathrm{}}^\mathit{2}_^𝒔𝒅^𝒔𝒙𝒑_𝜟\left(𝑪^{}\alpha 𝒙(𝑪)\right)\end{array}$$
(3.31)
and the combination of (3.30) and (3.31) gives
$$\left|\rho _{h,t}(C)\right|^2T(t)^1_t^{t+T(t)}𝑑\tau _^sd^sx_^sd^syh(\tau ^1𝒚)𝒉(\tau ^\mathit{1}𝒙)\omega \left(\alpha _\tau \left(\alpha 𝒚(𝑪^{})\alpha 𝒙(𝑪)\right)\right)\text{.}$$
(3.32)
We want to replace the term $`h(\tau ^1𝒙)`$ by the norm $`h_{\mathrm{}}`$ and, to do so, define the function $`h_+(h_{\mathrm{}}hh^2)^{1/2}`$, which is a non-negative element of $`C_{0,c}(^s)`$ as is $`h`$ itself. Then for any $`𝒛\text{,}𝒛^{}^𝒔`$ there holds the equation
$$h_{\mathrm{}}h(𝒛)=𝒉(𝒛)𝒉(𝒛^{})+𝒉_+(𝒛)𝒉_+(𝒛^{})+𝒉_+(𝒛)\left(𝒉_+(𝒛)𝒉_+(𝒛^{})\right)+𝒉(𝒛)\left(𝒉(𝒛)𝒉(𝒛^{})\right)\text{.}$$
(3.33)
Next, consider for an arbitrary function $`gC_{0,c}(^s)`$ the following inequality, based on an application of Fubini’s Theorem and the reasoning of (3.24),
$$\begin{array}{c}\left|T(t)^1_t^{t+T(t)}𝑑\tau _^sd^sx_^sd^syg(\tau ^1𝒚)\left(𝒈(\tau ^\mathit{1}𝒚)𝒈(\tau ^\mathit{1}𝒙)\right)\omega \left(\alpha _\tau \left(\alpha 𝒚(𝑪^{})\alpha 𝒙(𝑪)\right)\right)\right|\hfill \\ \hfill =\left|_^sd^sxT(t)^1_t^{t+T(t)}𝑑\tau _^sd^sz\tau ^sg(𝒛)\left(𝒈(𝒛)𝒈\left(𝒛_\tau (𝒙)\right)\right)\omega \left(\alpha _{(\tau ,\tau 𝒛)}\left(𝑪^{}\alpha 𝒙(𝑪)\right)\right)\right|\\ \hfill g_{\mathrm{}}_^sd^sx\underset{\tau I_t}{sup}\underset{𝒛^𝒔}{sup}\left|g(𝒛)𝒈\left(𝒛_\tau (𝒙)\right)\right|𝒑_𝜟\left(𝑪^{}\alpha 𝒙(𝑪)\right)\text{,}\end{array}$$
(3.34)
where we made use of the coordinate transformation $`𝒙𝒙+𝒚`$ followed by the transformation $`𝒚𝒛\tau ^\mathit{1}𝒚`$ and introduced the abbreviations $`𝒛_\tau (𝒙)\tau ^\mathit{1}𝒙+𝒛`$ as well as $`I_t`$ for the interval of $`\tau `$-integration. Similar to the proof of Proposition 3.12, the expression $`sup_{\tau I_t}sup_{𝒛^𝒔}|g(𝒛)𝒈\left(𝒛_\tau (𝒙)\right)|`$ is seen to vanish for all $`𝒙^𝒔`$ in the limit of large $`|t|`$, so that by Lebesgue’s Dominated Convergence Theorem the left-hand side of (3.34) converges to 0. This reasoning in particular applies to the functions $`h`$ as well as $`h_+`$ and thus to the third and fourth term on the right of equation (3.33). On the other hand, substitution of $`h`$ by $`h_+`$ in the integral of (3.32) likewise gives a non-negative result for all times $`t`$. Combining all these informations and specializing to a subnet $`\{t_\iota :\iota J\}`$ approximating $`+\mathrm{}`$ or $`\mathrm{}`$, one arrives at the following version of (3.32), valid for asymptotic times:
$$\begin{array}{c}\underset{\iota }{lim}\left|\rho _{h,t_\iota }(C)\right|^2\hfill \\ \hfill \underset{\iota }{lim}h_{\mathrm{}}T(t_\iota )^1_{t_\iota }^{t_\iota +T(t_\iota )}𝑑\tau _^sd^sx_^sd^syh(\tau ^1𝒚)\omega \left(\alpha _\tau \left(\alpha 𝒚(𝑪^{})\alpha 𝒙(𝑪)\right)\right)\\ \hfill h_{\mathrm{}}\underset{\iota }{lim}_^sd^sxT(t_\iota )^1_{t_\iota }^{t_\iota +T(t_\iota )}𝑑\tau _^sd^syh(\tau ^1𝒚)\omega \left(\alpha _{(\tau ,𝒚)}\left(𝑪^{}\alpha 𝒙(𝑪)\right)\right)\\ \hfill =h_{\mathrm{}}\underset{\iota }{lim}_^sd^sx\rho _{h,t_\iota }\left(C^{}\alpha x(C)\right)\text{.}\end{array}$$
Making use of Lemma 3.7, this result can be expressed in terms of the functional $`\sigma =lim_\iota \rho _{h,t_\iota }`$ to yield
$$\left|\sigma (C)\right|^2h_{\mathrm{}}_^sd^sx\sigma \left(C^{}\alpha x(C)\right)\text{.}$$
### 3.3 Particle Weights
The features of limit functionals $`\sigma _{\mathrm{\Delta }}^{}{}_{}{}^{+}`$ collected thus far, point to their interpretation as representatives of mixtures of particle-like quantities with sharp energy-momentum: being translationally invariant according to Proposition 3.12, they appear as plane waves, i. e. energy-momentum eigenstates, on the other hand they are singly localized at all times by Proposition 3.6, thereby exhibiting properties of particle-like systems, their energy-momentum spectrum being determined by Proposition 3.8. We shall summarize systems of the above kind under the concept of *particle weights*, a term chosen to reflect the connection to the notion of ‘weights’ or ‘extended positive functionals’ in the theory of $`C^{}`$-algebras, going back to Dixmier \[24, Section I.4.2\] (cf. also \[48, Section 5.1\] and ). These designate functions on the positive cone $`𝔄^+`$ of a $`C^{}`$-algebra $`𝔄`$ which can attain infinite values, a property they share with the singular functionals constructed in Theorem 3.10: it was seen to be of importance that their domain $``$ does not comprise the element $`\mathrm{𝟏}`$ of the quasi-local algebra, for the defining approximation would then lead to the value $`\sigma (\mathrm{𝟏})=+\mathrm{}`$.
As already mentioned in Section 3.1, every positive functional $`\sigma `$ on $`=𝔏^{}𝔏`$ defines a non-negative sesquilinear form $`.|._\sigma `$ on $`𝔏\times 𝔏`$ via
$$L_1|L_2_\sigma \sigma (L_{1}^{}{}_{}{}^{}L_2)$$
(3.35)
for any $`L_1\text{,}L_2𝔏`$, which induces a seminorm $`q_\sigma `$ on $`𝔏`$ and a norm $`._\sigma `$ on the corresponding quotient of $`𝔏`$ by the null space $`𝔑_\sigma `$ of $`q_\sigma `$. Taking advantage of these constructions, we shall depart from functionals and proceed to sesquilinear forms, a step which is necessitated by the special demands of the subsequent analysis. The following definition consists of a résumé of the essence of our knowledge on asymptotic functionals acquired in the above sequence of propositions.
###### Definition 3.14.
A particle weight is a non-trivial, non-negative sesquilinear form on $`𝔏`$, written $`.|.`$, which induces by (3.2) on the ideal $`𝔏`$ a seminorm $`q_w`$ with null space $`𝔑_w`$ as well as a norm $`._w`$ on the quotient $`𝔏/𝔑_w`$, and which complies with the following assumptions:
* for any $`L_1\text{,}L_2𝔏`$ and $`A𝔄`$ there holds the relation
$$L_1|AL_2=A^{}L_1|L_2\text{;}$$
* for given $`L𝔏`$ the following mapping is continuous with respect to $`q_w`$:
$$\mathrm{\Xi }_L:𝖯_+^{}𝔏(\mathrm{\Lambda },x)\mathrm{\Xi }_L(\mathrm{\Lambda },x)=\alpha _{(\mathrm{\Lambda },x)}(L)\text{;}$$
* the restriction to the subspace $`𝔏_0`$ of the canonical homomorphism
$$𝒬_w:𝔏𝔏/𝔑_wL𝒬_w(L)[L]_w$$
is $`𝒳_{𝔏_0}`$-differentiable in the sense of Definition A.16;
* the sesquilinear form is invariant with respect to space-time translations $`x^{s+1}`$, i. e.
$$\alpha _x(L_1)|\alpha _x(L_2)=L_1|L_2\text{,}L_1\text{,}L_2𝔏\text{,}$$
and the $`(s+1)`$-dimensional Fourier transforms of the distributions
$$xL_1|\alpha _x(L_2)$$
have support in a shifted forward light cone $`\overline{V}_+q`$, where $`q\overline{V}_+`$.
###### Remark.
* Note, that we did not impose on $`.|.`$ any restrictions concerning continuity with respect to the $`q_\mathrm{\Delta }`$-topology of $`𝔏`$, for in general such conditions will get lost in the disintegration of particle weights to be expounded in Chapter 4. The continuity property, which actually depends on the topology of $`𝔏`$, is formulated in terms of the seminorm $`q_w`$ induced by the sesquilinear form under consideration. The constituent properties of the above definition are preserved under the operations of addition and of multiplication by positive numbers, so that the totality of particle weights supplemented by the trivial form proves to be a positive (proper convex) cone (cf. ), denoted $`\mathrm{W}`$, in the linear space of all sesquilinear forms on $`𝔏`$. This ascertainment is the foundation for the constructions of Chapter 6.
* One could be tempted to go the way back from a sesquilinear form of the above type to a positive linear functional on $``$, but this is by no means self-evident. It is only possible under restrictive assumptions on the structure of the algebra $``$ to make the definition of the associated functional unambiguous.
A completely equivalent characterization of particle weights can be given in terms of representations $`(\pi _w,_w)`$ of the quasi-local algebra $`𝔄`$, obtained by means of a GNS-construction (cf. \[47, Theorem 3.2\] and \[48, Proposition 5.1.3\]).
###### Theorem 3.15.
* To any particle weight $`.|.`$ there corresponds a non-zero, non-degenerate representation $`(\pi _w,_w)`$ of the quasi-local $`C^{}`$-algebra $`𝔄`$ with the following properties:
+ there exists a linear mapping $`|.`$ from $`𝔏`$ onto a dense subspace of $`_w`$
$$|.:𝔏_wL|L\text{,}$$
such that the representation $`\pi _w`$ is given by
$$\pi _w(A)|L=|AL\text{,}A𝔄\text{,}L𝔏\text{;}$$
+ the following mapping is continuous for given $`L𝔏`$:
$$|\mathrm{\Xi }_L(.):𝖯_+^{}_w(\mathrm{\Lambda },x)|\mathrm{\Xi }_L(\mathrm{\Lambda },x)=|\alpha _{(\mathrm{\Lambda },x)}(L)\text{;}$$
+ the restriction of the linear mapping $`|.`$ to $`𝔏_0`$ with range in the subspace of $`_w`$ spanned by all vectors $`|L_0`$, $`L_0𝔏_0`$, is $`𝒳_{𝔏_0}`$-differentiable;
+ there exists a strongly continuous unitary representation $`xU_w(x)`$ of space-time translations $`x^{s+1}`$ on $`_w`$ defined by
$$U_w(x)|L|\alpha _x(L)\text{,}L𝔏\text{,}$$
with spectrum in a displaced forward light cone $`\overline{V}_+q`$, $`q\overline{V}_+`$.
* Any representation $`(\pi _w,_w)`$ which has the above characteristics defines a particle weight through the scalar product on $`_w`$.
###### Remark.
By their very definition, the unitaries $`U_w(x)`$ implement the automorphism group $`\{\alpha _x:x^{s+1}\}\mathrm{Aut}𝔄`$ through
$$U_w(x)\pi _w(A)U_w(x)^{}=\pi _w\left(\alpha _x(A)\right)\text{,}A𝔄\text{}x^{s+1}\text{,}$$
(3.36)
in the representation $`(_w,\pi _w)`$.
###### Proof.
* The proof of the various properties stated in the Theorem is readily carried out, once the GNS-construction has been realized.
+ Since a particle weight satisfies the Cauchy-Schwarz inequality its null space
$$𝔑_w\{N𝔏:N|N=0\}$$
turns out to be a left ideal in $`𝔏`$ (and hence in $`𝔄`$). The defining sesquilinear form endows the quotient space of $`𝔏`$ by $`𝔑_w`$ with a pre-Hilbert space structure; its completion $`_w`$ contains by construction the range of the canonical homomorphism
$$|.:𝔏𝔏/𝔑_wL|L[L]_w$$
as a dense subspace. $`𝔏`$ and $`𝔑_w`$ being left ideals in $`𝔄`$, the definition
$$\pi _w(A)|L|AL\text{,}A𝔄\text{,}$$
makes sense on the range of $`|.`$ and can be extended to all of $`_w`$ due to the estimate
$$\pi _w(A)|L^2=AL|AL=L|A^{}ALA^2L|L=A^2|L^2\text{,}$$
(3.37)
which is founded on the fact that the particle weight is a non-negative sesquilinear form and the operator $`A^2\mathbf{\hspace{0.17em}1}A^{}A`$ is positive. Since $`𝔄`$ is unital, this yields a non-zero, non-degenerate representation of the quasi-local algebra on the Hilbert space $`_w`$.
+ The norm on $`_w`$ induces a seminorm on $`𝔏`$ via the linear mapping $`|.`$ and this coincides with $`q_w`$ as defined for particle weights. Therefore the asserted continuity of the mapping $`(\mathrm{\Lambda },x)|\alpha _{(\mathrm{\Lambda },x)}(L)`$ is an immediate consequence of the respective property in Definition 3.14.
+ By construction, the canonical homomorphisms $`|.`$ and $`𝒬_w`$ coincide and furthermore $`|L=[L]_w_w`$, so that the assumption of $`𝒳_{𝔏_0}`$-differentiability is self-evident.
+ The existence of a strongly continuous unitary representation of space-time translations in $`(\pi _w,_w)`$ is a direct consequence of translation invariance of the particle weight $`.|.`$ and its continuity under Poincaré transformations with respect to $`q_w`$. Stone’s Theorem (cf. \[6, Chapter 6, § 2\] and \[38, Theorem VIII.(33.8)\]) connects the spectrum of its generator $`P_w=(P_w^\mu )`$ to the support of the Fourier transform of $`xL_1|\alpha _x(L_2)`$ in Definition 3.14 by virtue of the relation
$$\begin{array}{c}_{^{s+1}}d^{s+1}xg(x)L_1|\alpha _x(L_2)=_{^{s+1}}d^{s+1}xg(x)L_1\left|U_w(x)\right|L_2\hfill \\ \hfill =(2\pi )^{(s+1)/2}L_1\left|\stackrel{~}{g}(P_w)\right|L_2\text{,}\end{array}$$
(3.38)
which holds for any $`L_1\text{,}L_2𝔏`$ and any $`gL^1(^{s+1},d^{s+1}x)`$. To clarify this fact, note, that the projection-valued measure $`E_w(.)`$ corresponding to $`P_w`$ is regular, i. e. $`E_w(\mathrm{\Delta }^{})`$ is for any Borel set $`\mathrm{\Delta }^{}`$ the strong limit of the net $`\{E_w(\mathrm{\Gamma }^{}):\mathrm{\Gamma }^{}\mathrm{\Delta }^{}\text{compact}\}`$. For each compact $`\mathrm{\Gamma }\mathrm{}(\overline{V}_+q)`$ consider an infinitely often differentiable function $`\stackrel{~}{g}_\mathrm{\Gamma }`$ with support in $`\mathrm{}(\overline{V}_+q)`$ that envelops the characteristic function for $`\mathrm{\Gamma }`$ (cf. \[40, Satz 7.7\]): $`0\chi _\mathrm{\Gamma }\stackrel{~}{g}_\mathrm{\Gamma }`$. According to the assumption of Definition 3.14 the left-hand side of (3.38) vanishes for any $`g_\mathrm{\Gamma }`$ of the above kind, and this means that all the bounded operators $`\stackrel{~}{g}_\mathrm{\Gamma }(P_w)`$ equal $`0`$ not only on the dense subspace spanned by vectors $`|L`$, $`L𝔏`$, but on all of $`_w`$. Due to the fact that $`\stackrel{~}{g}_\mathrm{\Gamma }`$ majorizes $`\chi _\mathrm{\Gamma }`$, this in turn implies $`\chi _\mathrm{\Gamma }(P_w)=E_w(\mathrm{\Gamma })=0`$ and thus, by arbitrariness of $`\mathrm{\Gamma }\mathrm{}(\overline{V}_+q)`$ in connection with regularity, the desired relation $`E_w\left(\mathrm{}(\overline{V}_+q)\right)=0`$.
* The reversion of the above arguments in order to establish that the scalar product on $`_w`$ possesses the characteristics of a particle weight is self-evident. ∎
The following analogue of Lemmas 2.17 and 2.18 in terms of the $`q_w`$-topology induced on $`𝔏`$ by a particle weight is of importance not only for the remaining results of this chapter, but plays an important role in the constructions that underlie the theory of disintegration to be expounded in Chapter 4.
###### Lemma 3.16.
Let $`L𝔏`$ and let $`.|.`$ be a particle weight.
* Let $`FL^1(𝖯_+^{},d\mu (\mathrm{\Lambda },x))`$ have compact support $`𝖲`$, then the Bochner integral
$$\alpha _F(L)=𝑑\mu (\mathrm{\Lambda },x)F(\mathrm{\Lambda },x)\alpha _{(\mathrm{\Lambda },x)}(L)$$
(3.39a)
lies in the completion of $`𝔏`$ with respect to the locally convex topology induced on it by the initial norm $`.`$ and the $`q_w`$-seminorm defined by the particle weight. Moreover $`|\alpha _F(L)`$ is a vector in the corresponding Hilbert space $`_w`$ and can be written
$$|\alpha _F(L)=𝑑\mu (\mathrm{\Lambda },x)F(\mathrm{\Lambda },x)|\alpha _{(\mathrm{\Lambda },x)}(L)\text{,}$$
(3.39b)
satisfying the inequality
$$|\alpha _F(L)F_1\underset{(\mathrm{\Lambda },x)𝖲}{sup}|\alpha _{(\mathrm{\Lambda },x)}(L)\text{.}$$
(3.39c)
* For any function $`gL^1(^{s+1},d^{s+1}x)`$ the Bochner integral
$$\alpha _g(L)=_{^{s+1}}d^{s+1}xg(x)\alpha _x(L)$$
(3.40a)
likewise lies in the completion of $`𝔏`$ with respect to the locally convex topology mentioned above. $`|\alpha _g(L)`$ is a vector in the Hilbert space $`_w`$ subject to the relation
$$|\alpha _g(L)=_{^{s+1}}d^{s+1}xg(x)|\alpha _x(L)=(2\pi )^{(s+1)/2}\stackrel{~}{g}(P_w)|L\text{,}$$
(3.40b)
so that
$$|\alpha _g(L)g_1|L\text{.}$$
(3.40c)
###### Proof.
* Due to continuity of the particle weight $`.|.`$ with respect to Poincaré transformations as claimed in Definition 3.14, the integrand of (3.39a) can be estimated with respect to the seminorm $`q_w`$ induced on $`𝔏`$, which gives the Lebesgue-integrable function $`(\mathrm{\Lambda },x)|F(\mathrm{\Lambda },x)|sup_{(\mathrm{\Lambda },x)𝖲}q_w(\alpha _{(\mathrm{\Lambda },x)}(L))`$. Therefore the integral in question indeed exists in the completion of the locally convex space $`𝔏`$ not only with respect to the norm topology but also with respect to the seminorm $`q_w`$. Furthermore the corresponding GNS-construction of $`(\pi _w,_w)`$ implies that $`|L`$ coincides with $`q_w(L)`$ for any $`L𝔏`$, a relation which extends to the respective completions (cf. \[44, Chapter One, § 5 4.(4)\]) thus resulting in (3.39b). (3.39c) is then an immediate consequence, again on grounds of continuity under Poincaré transformations.
* According to Definition 3.14, the particle weight $`.|.`$ is invariant under space-time translations and so is the seminorm $`q_w`$. Therefore the integrand of (3.40a) is majorized by the Lebesgue-integrable function $`x|g(x)|q_w(L)`$, so that the respective integral exists in the completion of $`𝔏`$. The first equation of (3.40b) arises from the same arguments that were already applied above, whereas the second one is then a consequence of Stone’s Theorem (cf. (3.38)). Again on the ground of translation invariance, the estimate (3.40c) is an immediate conclusion from (3.40b). ∎
Having this preparatory result at our disposal, we are in the position to prove a statement on spectral subspaces of $`_w`$, that will be significant in the next chapter as well as for the subsequent proof of the Cluster Property for particle weights.
###### Proposition 3.17 (Spectral Subspaces).
Let $`L`$ be an element of $`𝔏(\mathrm{\Delta }^{})=𝔏\stackrel{~}{𝔄}(\mathrm{\Delta }^{})`$, which means that $`L𝔏`$ has energy-momentum transfer in the Borel subset $`\mathrm{\Delta }^{}`$ of $`^{s+1}`$. Then, in the representation $`(\pi _w,_w)`$ corresponding to the particle weight $`.|.`$, the vector $`|L`$ belongs to the spectral subspace which pertains to $`\mathrm{\Delta }^{}`$ with respect to the intrinsic unitary representation $`xU_w(x)`$ of space-time translations:
$$|L=E_w(\mathrm{\Delta }^{})|L\text{.}$$
(3.41)
###### Proof.
The energy-momentum transfer of an operator $`A𝔄`$ can be stated in terms of the support properties of the Fourier transform of the mapping $`x\alpha _x(A)`$ considered as an operator-valued distribution (cf. the remark following Definition 2.2). For the operator $`L𝔏(\mathrm{\Delta }^{})`$ this has the consequence that $`\alpha _g(L)=0`$ if $`g`$ is any Lebesgue-integrable function with $`\mathrm{supp}\stackrel{~}{g}\mathrm{\Delta }^{}=\mathrm{}`$. In this case we have, by an application of Lemma 3.16,
$$_{^{s+1}}d^{s+1}xg(x)|\alpha _x(L)=|\alpha _g(L)=0\text{.}$$
(3.42)
Upon insertion of (3.42) into the formulation (3.38) of Stone’s Theorem, the reasoning applied in the proof of Theorem 3.15 yields the assertion. ∎
The particle weights enjoy a Cluster Property parallel to that established in Proposition 3.6 for functionals in $`_{\mathrm{\Delta }}^{}{}_{}{}^{+}`$. This characteristic, shared by the asymptotic functionals $`\sigma `$, could have been included in Definition 3.14, but it turns out, that it is already enforced by the other features.
###### Proposition 3.18 (Cluster Property for Particle Weights).
Let $`L_i`$ and $`L_i^{}`$ be elements of $`𝔏_0`$ with energy-momentum transfer $`\mathrm{\Gamma }_i`$ respectively $`\mathrm{\Gamma }_i^{}`$, and let $`A_i𝔄`$, $`i=1\text{,}2`$, be almost local operators. Suppose furthermore that $`.|.`$ is a particle weight with associated GNS-representation $`(\pi _w,_w)`$, then
$$^s𝒙𝑳_{\mathit{1}}^{}{}_{}{}^{}𝑨_\mathit{1}𝑳_\mathit{1}^{}|\alpha 𝒙(𝑳_{\mathit{2}}^{}{}_{}{}^{}𝑨_\mathit{2}𝑳_\mathit{2}^{})=𝑳_{\mathit{1}}^{}{}_{}{}^{}𝑨_\mathit{1}𝑳_\mathit{1}^{}\left|𝑼_𝒘(𝒙)\right|𝑳_{\mathit{2}}^{}{}_{}{}^{}𝑨_\mathit{2}𝑳_\mathit{2}^{}$$
is a function in $`L^1(^s,d^sx)`$.
###### Proof.
To establish this result we follow in the main the strategy of the proof of Proposition 3.6. Applied to the problem at hand in terms of $`(\pi _w,_w)`$, this yields initially the estimate
$$\begin{array}{c}\left|L_{1}^{}{}_{}{}^{}A_1L_1^{}\left|U_w(𝒙)\right|𝑳_{\mathit{2}}^{}{}_{}{}^{}𝑨_\mathit{2}𝑳_\mathit{2}^{}\right|\hfill \\ \hfill \left|L_1^{}\left|\pi _w\left([A_{1}^{}{}_{}{}^{}L_1,\alpha x(L_{2}^{}{}_{}{}^{}A_2)]\right)U_w(𝒙)\right|𝑳_\mathit{2}^{}\right|+\left|𝑳_\mathit{1}^{}\left|\pi _𝒘\left(\alpha 𝒙(𝑳_{\mathit{2}}^{}{}_{}{}^{}𝑨_\mathit{2})𝑨_{\mathit{1}}^{}{}_{}{}^{}𝑳_\mathit{1}\right)𝑼_𝒘(𝒙)\right|𝑳_\mathit{2}^{}\right|\end{array}$$
(3.43)
for any $`𝒙^𝒔`$. The first term on the right-hand side turns out to be majorized by $`[A_{1}^{}{}_{}{}^{}L_1,\alpha x(L_{2}^{}{}_{}{}^{}A_2)]|L_1^{}|L_2^{}`$ in view of the fact that the particle weight is invariant under translations and that the representation $`\pi _w`$ is continuous. As the operators involved are almost local without exception, the norm of the commutator taking part in this expression decreases rapidly, thus rendering it integrable. The second term requires a closer inspection. One has
$$\begin{array}{c}2\left|L_1^{}\left|\pi _w\left(\alpha x(L_{2}^{}{}_{}{}^{}A_2)A_{1}^{}{}_{}{}^{}L_1\right)U_w(𝒙)\right|𝑳_\mathit{2}^{}\right|\hfill \\ \hfill 2\pi _w\left(\alpha x(A_{2}^{}{}_{}{}^{}L_2)\right)|L_1^{}\pi _w\left(A_{1}^{}{}_{}{}^{}L_1\right)U_w(𝒙)|𝑳_\mathit{2}^{}\\ \hfill \pi _w\left(\alpha x(A_{2}^{}{}_{}{}^{}L_2)\right)|L_1^{}^2+\pi _w\left(\alpha _{(𝒙)}(A_{1}^{}{}_{}{}^{}L_1)\right)|L_2^{}^2\text{,}\end{array}$$
(3.44)
again by translation invariance of the particle weight in the last estimate. Now, a substitute of Lemma 3.5 has to be sought for, which was applied in the proof of Proposition 3.6 to get an estimate for (3.12), corresponding to the right-hand side of (3.44). Note, that $`\pi _w(A^{})`$ has the same energy-momentum transfer with respect to the unitary group $`\{U_w(x):x^{s+1}\}`$ as the operator $`A^{}𝔄`$ has regarding the underlying positive energy representation, and that, according to Proposition 3.17, $`|L_1^{}=E_w(\mathrm{\Gamma }_1^{})|L_1^{}`$ and $`|L_2^{}=E_w(\mathrm{\Gamma }_2^{})|L_2^{}`$ belong to the spectral subspaces pertaining to the compact sets $`\mathrm{\Gamma }_1^{}`$ and $`\mathrm{\Gamma }_2^{}`$. As in addition the spectrum of $`\{U_w(x):x^{s+1}\}`$ is restricted to a displaced forward light cone, all of the arguments given in the proofs of Propositions 2.6 and 2.7 also apply to the representation $`(\pi _w,_w)`$, so that e. g.
$$_^sd^sxE_w(\mathrm{\Gamma }_1^{})\pi _w\left(\alpha x(L_{2}^{}{}_{}{}^{}A_2A_{2}^{}{}_{}{}^{}L_2)\right)E_w(\mathrm{\Gamma }_1^{})$$
is seen to exist in the $`\sigma \mathrm{weak}`$-topology on $`𝔅(_w)`$. For this term we thus have
$$\begin{array}{c}_^sd^sx\pi _w\left(\alpha x(A_{2}^{}{}_{}{}^{}L_2)\right)E_w(\mathrm{\Gamma }_1^{})|L_1^{}^2\hfill \\ \hfill =_^sd^sxL_1^{}\left|E_w(\mathrm{\Gamma }_1^{})\pi _w\left(\alpha x(L_{2}^{}{}_{}{}^{}A_2A_{2}^{}{}_{}{}^{}L_2)\right)E_w(\mathrm{\Gamma }_1^{})\right|L_1^{}<\mathrm{}\text{.}\end{array}$$
(3.45)
The same holds true for the other expression on the right-hand side of (3.44), which shows that $`𝒙\left|𝑳_\mathit{1}^{}\left|\pi _𝒘\left(\alpha 𝒙(𝑳_{\mathit{2}}^{}{}_{}{}^{}𝑨_\mathit{2})𝑨_{\mathit{1}}^{}{}_{}{}^{}𝑳_\mathit{1}\right)𝑼_𝒘(𝒙)\right|𝑳_\mathit{2}^{}\right|`$ is an integrable function, too. Altogether, we have thus established the Cluster Property for particle weights. ∎
###### Remark.
Note, that the above result is independent of the differentiability properties of a particle weight (item (iii) in both Definition 3.14 and Theorem 3.15), since these did not enter into its proof.
At this point a brief comment on the notation chosen seems appropriate (cf. ). We deliberately utilize the typographical token $`|.`$ introduced by Dirac \[23, § 23\] for ket vectors describing improper momentum eigenstates $`|𝒑`$, $`𝒑^𝒔`$. These act as distributions on the space of momentum wave functions with values in the physical Hilbert space $``$, thereby presupposing a superposition principle to hold without limitations. This assumption collapses in an infraparticle situation as described in the Introduction. In contrast to this, the *pure* particle weights, that will shortly have their appearance in connection with elementary physical systems, are seen to be associated with sharp momentum and yet capable of describing infraparticles. Here the operators $`L𝔏`$ take on the role of the previously mentioned momentum space wave functions in that they localize the particle weight in order to produce a normalizable vector $`|L`$ in the pertaining Hilbert space $`_w`$. This in turn substantiates the terminology introduced in Definition 2.4. As they describe elementary physical systems, pure particle weights should give rise to irreducible representations of the quasi-local algebra, thus motivating the subsequent definition. It is supplemented by a certain specific regularity condition of technical importance, which we anticipate to hold in physically relevant situations, and by a notion of boundedness which is in particular shared by the positive asymptotic functionals $`\sigma `$, as shown in Lemma 3.20.
###### Definition 3.19.
A particle weight is said to be
* *pure*, if the corresponding representation $`(\pi _w,_w)`$ is irreducible;
* *regular*, if for any $`L𝔏`$ the following implication is valid:
$$L^{}L|L^{}L=0L|L=0\text{;}$$
* *$`\mathrm{\Delta }`$-bounded*, if to any bounded Borel subset $`\mathrm{\Delta }^{}`$ of $`^{s+1}`$ there exists another such set $`\overline{\mathrm{\Delta }}\mathrm{\Delta }+\mathrm{\Delta }^{}`$, such that the GNS-representation $`(\pi _w,_w)`$ of the particle weight and the defining representation are connected by the inequality
$$E_w(\mathrm{\Delta }^{})\pi _w(A)E_w(\mathrm{\Delta }^{})cE(\overline{\mathrm{\Delta }})AE(\overline{\mathrm{\Delta }})$$
(3.46)
for any $`A𝔄`$ with a suitable positive constant $`c`$ (independent of the Borel sets). Evidently, $`\mathrm{\Delta }`$ ought to be a bounded Borel set as well.
###### Lemma 3.20.
Any positive asymptotic functional $`\sigma _{\mathrm{\Delta }}^{}{}_{}{}^{+}`$, constructed according to Theorem 3.10 under the assumptions of Proposition 3.11, gives rise to a $`\mathrm{\Delta }`$-bounded particle weight $`.|._\sigma `$.
###### Proof.
Let $`(\pi _\sigma ,_\sigma )`$ denote the GNS-representation of the particle weight $`\sigma `$ with associated spectral measure $`E_\sigma (.)`$ for the generator $`P_\sigma =(P_\sigma ^\mu )`$ of the intrinsic space-time translations. For the time being, suppose that $`\mathrm{\Delta }^{}`$ is an *open* bounded Borel set in $`^{s+1}`$. Let furthermore $`L`$ be an arbitrary element of $`𝔏`$ and $`A𝔄`$. We are interested in an estimate of the term $`L|E_\sigma (\mathrm{\Delta }^{})\pi _\sigma (A)E_\sigma (\mathrm{\Delta }^{})|L_\sigma `$. Note, that the spectral measure is regular, so that $`E_\sigma (\mathrm{\Delta }^{})`$ is the strong limit of the net $`\{E_\sigma (\mathrm{\Gamma }):\mathrm{\Gamma }\mathrm{\Delta }^{}\text{compact}\}`$. As $`\mathrm{\Delta }^{}`$ is assumed to be open, there exists for each compact subset $`\mathrm{\Gamma }`$ of $`\mathrm{\Delta }^{}`$ an infinitely often differentiable function $`\stackrel{~}{g}_\mathrm{\Gamma }`$ with $`\mathrm{supp}\stackrel{~}{g}_\mathrm{\Gamma }\mathrm{\Delta }^{}`$ that fits between the corresponding characteristic functions \[40, Satz 7.7\]: $`\chi _\mathrm{\Gamma }\stackrel{~}{g}_\mathrm{\Gamma }\chi _\mathrm{\Delta }^{}`$. Thus the respective operators are subject to the relation
$$0\left(E_\sigma (\mathrm{\Delta }^{})\stackrel{~}{g}_\mathrm{\Gamma }(P_\sigma )\right)^2\left(E_\sigma (\mathrm{\Delta }^{})E_\sigma (\mathrm{\Gamma })\right)^2\text{,}$$
from which we infer that for arbitrary $`L^{}𝔏`$
$$0\left(E_\sigma (\mathrm{\Delta }^{})\stackrel{~}{g}_\mathrm{\Gamma }(P_\sigma )\right)|L^{}^2\left(E_\sigma (\mathrm{\Delta }^{})E_\sigma (\mathrm{\Gamma })\right)|L^{}^2\underset{\mathrm{\Gamma }\mathrm{\Delta }^{}}{\overset{}{}}0\text{.}$$
(3.47)
By density of all the vectors $`|L^{}`$ in $`_\sigma `$, it is thereby established that
$$E_\sigma (\mathrm{\Delta }^{})=\mathrm{strong}\underset{\mathrm{\Gamma }\mathrm{\Delta }^{}}{lim}\stackrel{~}{g}_\mathrm{\Gamma }(P_\sigma )\text{,}$$
(3.48)
which implies for the scalar product in to be considered here
$$L|E_\sigma (\mathrm{\Delta }^{})\pi _\sigma (A)E_\sigma (\mathrm{\Delta }^{})|L_\sigma =\underset{\mathrm{\Gamma }\mathrm{\Delta }^{}}{lim}L|\stackrel{~}{g}_\mathrm{\Gamma }(P_\sigma )\pi _\sigma (A)\stackrel{~}{g}_\mathrm{\Gamma }(P_\sigma )|L\text{.}$$
(3.49)
Since $`\stackrel{~}{g}_\mathrm{\Gamma }`$ is the Fourier transform of a rapidly decreasing function $`g_\mathrm{\Gamma }`$, which therefore belongs to the space $`L^1(^{s+1},d^{s+1}x)`$, Lemma 3.16 can be applied to yield for the right-hand side of (3.49)
$$\begin{array}{c}L|\stackrel{~}{g}_\mathrm{\Gamma }(P_\sigma )\pi _\sigma (A)\stackrel{~}{g}_\mathrm{\Gamma }(P_\sigma )|L=(2\pi )^{(s+1)}\alpha _{g_\mathrm{\Gamma }}(L)|\pi _\sigma (A)|\alpha _{g_\mathrm{\Gamma }}(L)_\sigma \hfill \\ \hfill =(2\pi )^{(s+1)}\sigma \left(\alpha _{g_\mathrm{\Gamma }}(L)^{}A\alpha _{g_\mathrm{\Gamma }}(L)\right)\text{,}\end{array}$$
(3.50)
where, following the remark pertaining to Proposition 3.11, the ultimate expression is based on the fact that $`\alpha _{g_\mathrm{\Gamma }}(L)^{}A\alpha _{g_\mathrm{\Gamma }}(L)𝔄_{}`$ as a consequence of Lemmas 2.18 and 2.12 in connection with Corollary 2.13. The approximating functionals $`\rho _{h,t}`$ for $`\sigma `$ in the form (3.23) with a non-negative function $`hL^{\mathrm{}}(^s,d^sx)`$ allow, through an application of \[11, Proposition 2.3.11\], for the following estimate of their integrand:
$$\begin{array}{c}\left|h(\tau ^1𝒙)\omega \left(𝑼(\tau )𝑬(𝜟)\alpha 𝒙(\alpha _{𝒈_𝜞}(𝑳)^{}𝑨\alpha _{𝒈_𝜞}(𝑳))𝑬(𝜟)𝑼(\tau )^{}\right)\right|\hfill \\ \hfill =h(\tau ^1𝒙)\left|\omega \left(𝑼(\tau )𝑬(𝜟)\alpha 𝒙(\alpha _{𝒈_𝜞}(𝑳)^{})𝑬(\overline{𝜟})\alpha 𝒙(𝑨)𝑬(\overline{𝜟})\alpha 𝒙(\alpha _{𝒈_𝜞}(𝑳))𝑬(𝜟)𝑼(\tau )^{}\right)\right|\\ \hfill E(\overline{\mathrm{\Delta }})AE(\overline{\mathrm{\Delta }})h(\tau ^1𝒙)\omega \left(𝑼(\tau )𝑬(𝜟)\alpha 𝒙(\alpha _{𝒈_𝜞}(𝑳)^{}\alpha _{𝒈_𝜞}(𝑳))𝑬(𝜟)𝑼(\tau )^{}\right)\text{.}\end{array}$$
Here the spectral projections $`E(\overline{\mathrm{\Delta }})`$ pertaining to the Borel set $`\overline{\mathrm{\Delta }}=\mathrm{\Delta }+\mathrm{\Delta }^{}`$, which is both bounded and open, could be introduced, since, according to Lemma 2.18, the energy-momentum transfer of $`\alpha _{g_\mathrm{\Gamma }}(L)`$ is contained in $`\mathrm{\Delta }^{}`$ by construction. An immediate consequence of the above relation is
$$\left|\rho _{h,t}\left(\alpha _{g_\mathrm{\Gamma }}(L)^{}A\alpha _{g_\mathrm{\Gamma }}(L)\right)\right|E(\overline{\mathrm{\Delta }})AE(\overline{\mathrm{\Delta }})\rho _{h,t}\left(\alpha _{g_\mathrm{\Gamma }}(L)^{}\alpha _{g_\mathrm{\Gamma }}(L)\right)\text{,}$$
which extends to the limit functional $`\sigma `$:
$$\left|\sigma \left(\alpha _{g_\mathrm{\Gamma }}(L)^{}A\alpha _{g_\mathrm{\Gamma }}(L)\right)\right|E(\overline{\mathrm{\Delta }})AE(\overline{\mathrm{\Delta }})\sigma \left(\alpha _{g_\mathrm{\Gamma }}(L)^{}\alpha _{g_\mathrm{\Gamma }}(L)\right)\text{.}$$
(3.51)
Insertion of this result into (3.50) yields
$$\left|L|\stackrel{~}{g}_\mathrm{\Gamma }(P_\sigma )\pi _\sigma (A)\stackrel{~}{g}_\mathrm{\Gamma }(P_\sigma )|L\right|E(\overline{\mathrm{\Delta }})AE(\overline{\mathrm{\Delta }})L|\stackrel{~}{g}_\mathrm{\Gamma }(P_\sigma )^2|L$$
(3.52)
and in the limit $`\mathrm{\Gamma }\mathrm{\Delta }^{}`$, in compliance with (3.49),
$$\left|L|E_\sigma (\mathrm{\Delta }^{})\pi _\sigma (A)E_\sigma (\mathrm{\Delta }^{})|L_\sigma \right|E(\overline{\mathrm{\Delta }})AE(\overline{\mathrm{\Delta }})L|E_\sigma (\mathrm{\Delta }^{})|L_\sigma E(\overline{\mathrm{\Delta }})AE(\overline{\mathrm{\Delta }})L|L_\sigma \text{.}$$
(3.53)
Passing to the supremum with respect to all $`L𝔏`$ such that $`|L_\sigma 1`$ (these constitute a dense subset of the unit ball in $`_\sigma `$), we get through an application of \[56, Satz 4.4\]
$$E_\sigma (\mathrm{\Delta }^{})\pi _\sigma (A)E_\sigma (\mathrm{\Delta }^{})2E(\overline{\mathrm{\Delta }})AE(\overline{\mathrm{\Delta }})\text{.}$$
(3.54)
This establishes the defining condition (3.46) for $`\mathrm{\Delta }`$-boundedness with $`c=2`$ in the case of an *open* bounded Borel set $`\mathrm{\Delta }^{}`$. But this is not an essential restriction, since an arbitrary bounded Borel set $`\mathrm{\Delta }^{}`$ is contained in the open set $`\mathrm{\Delta }_\eta ^{}`$, $`\eta >0`$, consisting of all those points $`p^{s+1}`$ for which $`inf_{p^{}\mathrm{\Delta }^{}}|pp^{}|<\eta `$. Since $`\mathrm{\Delta }_\eta ^{}`$ is likewise a bounded Borel set, we get
$$E_\sigma (\mathrm{\Delta }^{})\pi _\sigma (A)E_\sigma (\mathrm{\Delta }^{})E_\sigma (\mathrm{\Delta }_\eta ^{})\pi _\sigma (A)E_\sigma (\mathrm{\Delta }_\eta ^{})2E(\overline{\mathrm{\Delta }}_\eta )AE(\overline{\mathrm{\Delta }}_\eta )$$
(3.55)
as an immediate consequence of (3.54), where $`\overline{\mathrm{\Delta }}_\eta \mathrm{\Delta }+\mathrm{\Delta }_\eta ^{}`$. This covers the general case and thereby proves $`\mathrm{\Delta }`$-boundedness for the asymptotic functionals $`\sigma _{\mathrm{\Delta }}^{}{}_{}{}^{+}`$. ∎
## Chapter 4 Disintegration of Particle Weights
<sup>1</sup><sup>1</sup>footnotetext: A german translation can be found on page Acknowledgements.
In Section III of their treatment of collision cross sections for massive theories within the framework of local quantum physics, Araki and Haag got to the following asymptotic relation which holds true for the counters $`C`$ they had selected, for arbitrary vectors $`\mathrm{\Phi }`$ and certain specific vectors $`\mathrm{\Psi }`$ representing outgoing particle configurations \[3, Theorem 4\]:
$$\underset{t\mathrm{}}{lim}\mathrm{\Phi }\left|t^3C(h,t)\right|\mathrm{\Psi }=\underset{i,j}{^{}}d^3p\mathrm{\Gamma }_{ij}(𝒑)𝜱\left|𝒂_𝒋^{\text{out}}(𝒑)𝒂_𝒊^{\text{out}}(𝒑)\right|𝜳𝒉(𝒗_𝒊)\text{,}$$
(4.1)
where
$`\mathrm{\Gamma }_{ij}(𝒑)`$ $`8\pi ^3\mathrm{𝒑𝒋}\left|𝑪(\mathit{0})\right|𝒑𝒊\text{,}`$
$`𝒗_𝒊`$ $`(𝒑^\mathit{2}+𝒎_𝒊^\mathit{2})^{\mathit{1}/\mathit{2}}𝒑\text{.}`$
The indices $`i`$ and $`j`$ in the above formula denote the particle types including spin, and summation runs over pairs of particles with equal mass: $`m_i=m_j`$. The structure of the right-hand side of this equation is based on the *a priori* knowledge of the particle content of the theory they considered. Comparing this result with the concepts developed in the preceding chapter (cf. Theorem 3.10), one has an asymptotic functional $`\sigma _h^{(+)}`$ standing on the left-hand side of equation (4.1) that is decomposed with respect to momentum eigenstates $`|𝒑𝒊`$, hidden in the definition of $`\mathrm{\Gamma }_{ij}`$. If we accept such an interpretation of this theorem of Araki and Haag, it is possible to re-write it in the form
$$\sigma _h^{(+)}(C)=\underset{i,j}{^{}}𝑑\mu _{i,j}(𝒑)\mathrm{𝒑𝒋}\left|𝑪(\mathit{0})\right|𝒑𝒊\text{,}$$
where all expressions occurring in (4.1) apart from $`\mathrm{\Gamma }_{ij}`$ are absorbed into the measures $`\mu _{i,j}`$. This presents the asymptotic functional as a mixture of linear forms on $``$ (an algebra which is part of that selected in ) defined by Dirac kets representing improper momentum eigenstates; thus we happen to meet exactly those constructs that we already hinted at in the remarks concerning our notation that led to Definition 3.19. The aim of the present chapter is to establish a corresponding formula in the general setting, i. e. without any previous knowledge of the particle content.
As indicated by (4.1), representations resulting from the construction of asymptotic functionals as expounded in Chapter 3 will be highly reducible, whereas elementary physical systems are expected to be connected with pure particle weights, giving rise to irreducible representations of the quasi-local $`C^{}`$-algebra $`𝔄`$. In view of the preceding paragraph the obvious problem to be tackled now is to develop a theory for the decomposition or rather *disintegration* of generic particle weights into pure ones. Two approaches to this problem will be presented in this work:
* Decomposition of the GNS-representation pertaining to a particle weight into a direct integral of representations (spatial disintegration):
$$(\pi _w,_w)X^{}𝑑\nu (\xi )(\pi _\xi ,_\xi )\text{.}$$
* Barycentric decomposition of a given particle weight with respect to a base $`\mathrm{BW}`$ of the positive cone $`\mathrm{W}`$ of all particle weights in the space of sesquilinear forms on $`𝔏`$ (Choquet theory):
$$.|.=_{\mathrm{BW}}d\upsilon (\zeta ).|._\zeta \text{.}$$
Although the technical problems to come to grips with in these two constructions are quite different, we anticipate equivalence of their results: the separability assumptions essential in the first one are substituted by compactness conditions in the second. So evidently both of them require certain restrictions in the number of degrees of freedom, which seem to be complementary in one way or another. While the partial results achieved so far in connection with the barycentric decomposition will be discussed in Chapter 6, the spatial disintegration of the GNS-representation of a particle weight is the subject we will elaborate on first.
### 4.1 Separable Reformulation of Local Quantum Physics and its Associated Algebra of Detectors
The theory of spatial disintegration of representations $`(\overline{\pi },\overline{})`$ of a $`C^{}`$-algebra $`\overline{𝔄}`$ is a common theme of the pertinent textbooks (cf. ), an indispensable presupposition being that of separability of the algebra $`\overline{𝔄}`$ as well as of the Hilbert space $`\overline{}`$ in their respective uniform topologies. Note, that in this way the statements of \[11, Section 4.4\] are incorrect (cf. also \[12, Corrigenda\]). These separability assumptions are too restrictive to be encountered in physically reasonable theories from the outset, so first of all a *countable* version of the fundamental assumptions of local quantum field theory in terms of the net $`𝒪𝔄(𝒪)`$ and of the symmetry group $`𝖯_+^{}`$ has to be formulated before one can benefit from the extensive theory made available in the literature. This construction will be accomplished in a sequence of steps:
* With respect to its initial topology, the Poincaré group $`𝖯_+^{}`$ contains a numerable dense subgroup that we signify by $`𝖯^c`$. It is itself the semi-direct product of countable dense subgroups of Lorentz transformations $`𝖫^c`$ in $`𝖫_+^{}`$ and of space-time translations $`𝖳^c`$ in $`^{s+1}`$: $`𝖯^c=𝖫^c𝖳^c`$.
* Consider the standard diamonds with *rational* radii, centred around the origin. Subjecting these regions to all of the transformations in $`𝖯^c`$ yields a countable family $`^c`$ of open bounded regions, which is invariant with respect to the selected Poincaré transformations and constitutes a covering of $`^{s+1}`$. Note, that arbitrarily small regions belong to $`^c`$ in the sense, that any region in Minkowski space contains an element of this numerable collection as a subset.
* As shown in Appendix B, any unital $`C^{}`$-algebra of operators on a separable Hilbert space $``$ contains a strongly dense (i. e. dense with respect to the strong-operator topology), norm-separable $`C^{}`$-subalgebra, that includes the identity. Applied to the local $`C^{}`$-algebras $`𝔄(𝒪)`$ of the defining positive-energy representation, this result has the consequence that to each open bounded region $`𝒪`$ in Minkowski space one can associate a norm-separable, unital $`C^{}`$-algebra $`𝔄_{}(𝒪)`$, that lies strongly dense in $`𝔄(𝒪)`$. This means, that the algebra $`𝔄_{}(𝒪)`$ in turn contains a countable -subalgebra $`𝔄^c(𝒪)`$ over the field $`+i`$, which is uniformly dense in $`𝔄_{}(𝒪)`$, strongly dense in $`𝔄(𝒪)`$ and can likewise be chosen to comprise the unit.
Let $`𝒪_k`$, $`k`$, be a denumeration of the countable family $`^c`$ of open bounded regions in Minkowski space constructed above. We define $`𝔄^{}(𝒪_k)`$ as the $`C^{}`$-algebra (over $``$) which is generated by the union of all $`\alpha _{(\mathrm{\Lambda },x)}\left(𝔄^c(𝒪_i)\right)`$, where $`(\mathrm{\Lambda },x)𝖯^c`$ and $`𝒪_i^c`$ run through all combinations for which $`\mathrm{\Lambda }𝒪_i+x𝒪_k`$. By construction this algebra is norm-separable and satisfies
$$𝔄^c(𝒪_k)𝔄^{}(𝒪_k)𝔄(𝒪_k)\text{,}$$
(4.2a)
so that $`𝔄^{}(𝒪_k)`$ turns out to be strongly dense in $`𝔄(𝒪_k)`$.
The net of local $`C^{}`$-algebras $`\{𝔄^{}(𝒪_k):k\}`$ fulfills the conditions of isotony, locality and covariance with respect to $`^c`$ and $`𝖯^c`$. Isotony is an immediate consequence of the construction whereas locality follows from (4.2a) in connection with locality of the defining net $`𝒪𝔄(𝒪)`$. To establish covariance one has to observe that, given any $`(\mathrm{\Lambda },x)𝖯^c`$, the algebra $`\alpha _{(\mathrm{\Lambda },x)}\left(𝔄^{}(𝒪_k)\right)`$ is generated by all $`\alpha _{(\mathrm{\Lambda },x)}\left(\alpha _{(\mathrm{\Lambda }^{},x^{})}\left(𝔄^c(𝒪_i)\right)\right)`$, where $`(\mathrm{\Lambda }^{},x^{})𝖯^c`$ and $`𝒪_i^c`$ run through those combinations which satisfy the relation $`\mathrm{\Lambda }^{}𝒪_i+x^{}𝒪_k`$. This can equivalently be expressed by saying that the algebra in question is generated by all $`\alpha _{(\mathrm{\Lambda }^{\prime \prime },x^{\prime \prime })}\left(𝔄^c(𝒪_i)\right)`$, for which $`(\mathrm{\Lambda }^{\prime \prime },x^{\prime \prime })𝖯^c`$ and $`𝒪_i^c`$ have the property $`\mathrm{\Lambda }^{\prime \prime }𝒪_i+x^{\prime \prime }\mathrm{\Lambda }𝒪_k+x`$. In this formulation $`\alpha _{(\mathrm{\Lambda },x)}\left(𝔄^{}(𝒪_k)\right)`$ turns out to be equal to the algebra $`𝔄^{}(\mathrm{\Lambda }𝒪_k+x)`$. The somewhat intricate construction of $`𝔄^{}(𝒪_k)`$ is necessitated by the requisite to have the standard properties of a net of local algebras at our disposition.
By construction, the countable -algebra $`𝔄^c`$ over $`+i`$, which is generated by the union of all the algebras $`𝔄^c(𝒪_k)`$, $`𝒪_k^c`$, and thus invariant under transformations from $`𝖯^c`$, lies uniformly dense in the $`C^{}`$-inductive limit $`𝔄^{}`$ of the net $`𝒪_k𝔄^{}(𝒪_k)`$, and is, on account of (4.2a), even strongly dense in the quasi-local algebra $`𝔄`$ itself. We thus have the inclusions
$$𝔄^c𝔄^{}𝔄\text{,}$$
(4.2b)
with a norm-separable $`C^{}`$-algebra $`𝔄^{}`$, which lies strongly dense in $`𝔄`$ and contains $`𝔄^c`$ as a numerable uniformly dense subalgebra (over $`+i`$).
Into this restricted setting of Local Quantum Physics defined above, we now introduce countable counterparts of the left ideal of localizing operators $`𝔏`$, of the algebra of detectors $``$ and, most important of all, of the subspace $`𝔏_0𝔏`$ of almost local vacuum annihilation operators.
First of all note, that it is possible to select a *numerable* subspace over $`+i`$ in $`𝔏_0`$, which consists of almost local vacuum annihilation operators with energy-momentum transfer in arbitrarily small regions. E. g. let $`\{\mathrm{\Gamma }_n\}_n`$ be a countable cover of $`\mathrm{}\overline{V}_+`$, constituted by compact and convex subsets of the complement of the forward light cone, with the additional property that any bounded region in $`\mathrm{}\overline{V}_+`$ contains one of these compacta. Let, for instance, $`\{p_i\}_i`$ be a dense sequence in $`\mathrm{}\overline{V}_+`$ and associate to each $`p_i`$ the compact balls of rational radius $`r`$ that satisfy $`\overline{_r(p_i)}\mathrm{}\overline{V}_+`$ in addition. The Lorentz group $`𝖫_+^{}`$, being locally compact, can be covered by a countable family of arbitrarily small compact sets $`\{\mathrm{\Theta }_m\}_m`$ as well. Now, the spaces $`𝒟_{\mathrm{\Gamma }_n}`$ and $`𝒟_{\mathrm{\Theta }_m}`$ of test functions with support in $`\mathrm{\Gamma }_n`$ or else $`\mathrm{\Theta }_m`$ (cf. \[40, § 12\]) are separable as subspaces of the respective Banach spaces $`L^p(^{s+1},d^{s+1}x)`$ and $`L^p(^{d_𝖫},d^{d_𝖫}t)`$ ($`d_𝖫2^1s(s+1)`$ is the dimension of $`𝖫_+^{}`$), which in turn are separable due to an application of \[39, Theorem IV.(13.20)\] using elements of the numerable set of simple functions with rational values on intervals with rational end points. Thus there exist dense sequences $`\stackrel{~}{g}_n^l`$ and $`h_m^k`$ in the spaces $`𝒟_{\mathrm{\Gamma }_n}`$ and $`𝒟_{\mathrm{\Theta }_m}`$, respectively. Consider the countable family of operators in $`𝔏_0`$, which are defined through
$$\alpha _{h_m^kg_n^l}(A_j)_{𝖯_+^{}}𝑑\mu (\mathrm{\Lambda },x)h_m^k(\mathrm{\Lambda })g_n^l(x)\alpha _{(\mathrm{\Lambda },x)}(A_j)\text{,}$$
(4.3)
for any $`A_j𝔄^c`$ in the uniform topology of $`𝔄`$, and supplement this selection by all orders of partial derivatives with respect to the canonical coordinates around $`(\mathrm{𝟏},0)`$ (cf. Appendix A):
$$\delta ^{\iota _M}\left(\alpha _{h_m^kg_n^l}(A_j)\right)=\delta ^{i_M}\mathrm{}\delta ^{i_1}\left(\alpha _{h_m^kg_n^l}(A_j)\right)𝔏_0$$
for any $`M`$-tuple $`\iota _M=(i_1,\mathrm{},i_M)`$ with integer entries from the set $`\{1,\mathrm{},d_𝖯\}`$, where $`d_𝖯=d_𝖫+(s+1)`$. Upon application of all transformations from $`𝖯^c`$ to these constructs, we get a sequence of vacuum annihilation operators, comprising elements with energy-momentum transfer in arbitrarily small regions, which generates a countable subspace $`𝔏_0^c`$ over the field $`+i`$ in $`𝔏_0`$, invariant under transformations from $`𝖯^c`$ and under arbitrary partial derivations. When this construct is to be used in connection with a given particle weight $`.|.`$ that is non-negative by definition, it does not cause any problems to supplement the set of operators defined in (4.3) by a countable number of other elements from $`𝔏_0`$, on which the particle weight attains non-vanishing values. In this way the imminent restriction of $`.|.`$ to a subset of $`𝔏`$ can be protected from getting trivial.
The above selection of vacuum annihilation operators does not yet meet the requirements for the disintegration. For it to be feasible we compactly regularize these operators: Take a *countable* set of compactly supported test functions $`F`$ on $`𝖯_+^{}`$ with a support $`𝖲_F\mathrm{supp}F`$ which contains the unit $`(\mathrm{𝟏},0)`$ of $`𝖯_+^{}`$. Then all the Bochner integrals
$$\alpha _F(L_0)=_{𝖲_F}𝑑\mu (\mathrm{\Lambda },x)F(\mathrm{\Lambda },x)\alpha _{(\mathrm{\Lambda },x)}(L_0)\text{,}L_0𝔏_0^c\text{,}$$
(4.4)
are elements of the $`C^{}`$-algebra $`𝔄^{}`$ and of $`𝔏_0`$ according to Lemma 2.17 with energy-momentum transfer contained in $`_{(\mathrm{\Lambda },x)𝖲_F}\mathrm{\Lambda }\mathrm{\Gamma }`$ given $`L_0𝔏_0(\mathrm{\Gamma })`$ (cf. the proof of the quoted Lemma). The specific property of operators of type (4.4) in contrast to those from $`𝔏_0^c`$ is, that their differentiability with respect to the Poincaré group can be expressed in terms of derivatives of the infinitely differentiable test function $`F𝒟_{𝖲_F}`$, a feature that will be of great significance later on. By choosing the support of the functions $`F`$ small enough, one can impose an energy-momentum transfer in arbitrarily small regions on the operators $`\alpha _F(L_0)`$ as was the case for the elements of $`𝔏_0^c`$ itself. Furthermore, a particle weight that did not vanish on the set $`𝔏_0^c`$ is also non-zero when restricted to all of the operators $`\alpha _F(L_0)`$ constructed in (4.4). This fact is easily established with relation (3.39b) of Lemma 3.16 and the continuity of the particle weight under Poincaré transformations in mind. The numerable set of vacuum annihilation operators that consists of those explicitly presented in (4.4) together with all their partial derivations of arbitrary order (that share this specific style of construction) will be denoted $`\overline{𝔏_0^c}`$ in the sequel. It might happen that two of these elements of $`𝔏_0`$ are connected by a Poincaré transformation not yet included in $`𝖯^c`$. For technical reasons, which are motivated by the exigencies for the proof of the central Theorem 4.4 of this chapter, we supplement $`𝖯^c`$ by all of the (countably many) transformations arising in this way and consider henceforth the countable subgroup $`\overline{𝖯}^c=\overline{𝖫}^c\overline{𝖳}^c𝖯_+^{}`$ generated by them. The set $`\overline{𝔏_0^c}`$ is then invariant under the operation of taking derivatives as well as under all transformations from the numerable dense subgroup $`\overline{𝖯}^c`$.
Here is a list of the countable substitutes for the algebraic concepts used thus far:
* We have defined an isotonous, local and $`𝖯^c`$-covariant net $`𝒪_k𝔄^{}(𝒪_k)`$, $`𝒪_k^c`$, which has $`𝔄^{}`$ as $`C^{}`$-inductive limit. This is a norm-separable $`C^{}`$-algebra (over the field $``$) with unit $`\mathrm{𝟏}`$, containing $`𝔄^c`$, which is generated by the countable local algebras $`𝔄^c(𝒪_k)𝔄^{}(𝒪_k)`$, as a likewise unital, numerable, uniformly dense $``$-subalgebra over $`+i`$. $`𝔄^{}`$ itself lies strongly dense in the quasi-local algebra $`𝔄`$ and, due to uniform continuity of the mappings $`(\mathrm{\Lambda },x)\alpha _{(\mathrm{\Lambda },x)}(A)`$, $`A𝔄`$, it is invariant with respect to the whole Poincaré group. In contrast to this, note, that the invariance property for $`𝔄^c`$ is restricted to $`𝖯^c`$.
* $`\overline{𝔏_0^c}𝔏_0𝔄^{}`$ is a countable set of vacuum annihilation operators of the special construction (4.4), which is invariant under transformations from $`\overline{𝖯}^c`$ and under the operation of taking partial derivations. Depending on a given particle weight, it can be chosen in such a way, that the particle weight restricted to $`\overline{𝔏_0^c}`$ remains non-trivial.
* The image of $`\overline{𝔏_0^c}`$ under all Poincaré transformations is denoted $`\overline{𝔏_0}`$:
$$\overline{𝔏_0}\{\alpha _{(\mathrm{\Lambda },x)}(L_0):L_0\overline{𝔏_0^c},(\mathrm{\Lambda },x)𝖯_+^{}\}\text{.}$$
(4.5a)
* $`\overline{𝔄^c}𝔄^{}`$ in turn denotes the numerable, unital -algebra over $`+i`$ which is generated by $`𝔄^c\overline{𝔏_0^c}`$. It is thus stable with respect to $`𝖯^c`$ and uniformly dense in $`𝔄^{}`$.
* The countable counterpart $`𝔏^c`$ of the left ideal $`𝔏`$ in $`𝔄`$ is defined as the linear span with respect to the field $`+i`$ of operators of the form $`L=AL_0`$ with $`A\overline{𝔄^c}`$ and $`L_0\overline{𝔏_0^c}`$:
$$𝔏^c\overline{𝔄^c}\overline{𝔏_0^c}=\mathrm{span}_{+i}\left\{AL_0:A\overline{𝔄^c},L_0\overline{𝔏_0^c}\right\}\text{.}$$
(4.5b)
This constitutes a left ideal of the algebra $`\overline{𝔄^c}`$, likewise invariant under transformations from $`𝖯^c`$.
* Finally, one can introduce the countable -subalgebra $`^c`$ via
$$^c𝔏_{}^{c}{}_{}{}^{}𝔏^c=\mathrm{span}_{+i}\left\{L_{1}^{}{}_{}{}^{}L_2:L_1,L_2𝔏^c\right\}\text{.}$$
(4.5c)
### 4.2 Restricted $`𝕶_\mathrm{𝟎}^𝒄`$-Particle Weights
The subsequent developments in this chapter have to be founded on a mitigated version for the concept of particle weights as it was introduced in Definition 3.14. The reason is that the sesquilinear forms occurring in the decomposition theory of Section 4.3 do not share all the desired properties. Therefore we insert the present section which deals with the necessary restrictions that have to be imposed on the concepts of Chapter 3. The essential cuts are indicated by the work previously accomplished.
###### Definition 4.1 (Restricted $`𝕶_\mathrm{𝟎}^𝒄`$-Particle Weights).
Suppose that we are given a sextuple $`(\underset{¯}{\pi },\underset{¯}{},\underset{¯}{𝔄},\underset{¯}{\alpha },\underset{¯}{𝖯}^c,𝔎_0^c)`$ with entries of the following sense:
* $`\underset{¯}{𝔄}`$ is a norm separable $`C^{}`$-subalgebra of the quasi-local algebra $`𝔄`$, which arises as the $`C^{}`$-inductive limit of a countable net of local $`C^{}`$-algebras.
* $`\underset{¯}{𝖯}^c`$ is a numerable dense subgroup of the Poincaré group. $`𝖯_+^{}`$ as a whole is implemented in $`\underset{¯}{𝔄}`$ by the strongly continuous group of automorphisms
$$\{\underset{¯}{\alpha }_{(\mathrm{\Lambda },x)}:(\mathrm{\Lambda },x)𝖯_+^{}\}\mathrm{Aut}\underset{¯}{𝔄}\text{.}$$
* $`𝔎_0^c`$ designates a countable set of almost local vacuum annihilation operators in $`\underset{¯}{𝔄}`$, stable with respect to transformations from $`\underset{¯}{𝖯}^c`$. The image of $`𝔎_0^c`$ under all Poincaré transformations is denoted $`𝔎_0`$:
$$𝔎_0\{\underset{¯}{\alpha }_{(\mathrm{\Lambda },x)}(K_0):K_0𝔎_0^c,(\mathrm{\Lambda },x)𝖯_+^{}\}\text{.}$$
(4.6a)
* Together with the numerable uniformly dense -subalgebra of $`\underset{¯}{𝔄}`$, which exists by construction, $`𝔎_0^c`$ generates a countable -algebra over $`+i`$, denoted $`\underset{¯}{𝔄}^c`$ and likewise invariant under $`\underset{¯}{𝖯}^c`$.
* A countable left ideal in $`\underset{¯}{𝔄}^c`$ is then defined by
$$𝔎^c\underset{¯}{𝔄}^c𝔎_0^c=\mathrm{span}_{+i}\left\{AK_0:A\underset{¯}{𝔄}^c,K_0𝔎_0^c\right\}\text{.}$$
(4.6b)
It is invariant under the automorphism group $`\{\underset{¯}{\alpha }_{(\mathrm{\Lambda },x)}:(\mathrm{\Lambda },x)\underset{¯}{𝖯}^c\}`$ as well.
* Finally, $`(\underset{¯}{\pi },\underset{¯}{})`$ is a non-zero, non-degenerate representation of the $`C^{}`$-algebra $`\underset{¯}{𝔄}`$.
The sextuple $`(\underset{¯}{\pi },\underset{¯}{},\underset{¯}{𝔄},\underset{¯}{\alpha },\underset{¯}{𝖯}^c,𝔎_0^c)`$ is called a restricted $`𝔎_0^c`$-particle weight, in case that it complies with the following list of features:
* There exists a $`(+i)`$-linear mapping $`|.`$ from $`𝔎^c`$ onto a dense subset $`\underset{¯}{}^c\underset{¯}{}`$:
$$|.:𝔎^c\underset{¯}{}^cK|K\text{,}$$
(4.7a)
such that the representation $`\underset{¯}{\pi }`$ acts on this space according to
$$\underset{¯}{\pi }(A)|K=|AK\text{,}A\underset{¯}{𝔄}^c\text{,}K𝔎^c\text{.}$$
(4.7b)
* The above linear mapping allows for an extension to any operator in $`𝔎_0`$, such that (in the notation of Theorem 3.15)
$$|\mathrm{\Xi }_K^{}(.):𝖯_+^{}\underset{¯}{}(\mathrm{\Lambda },x)|\mathrm{\Xi }_K^{}(\mathrm{\Lambda },x)|\underset{¯}{\alpha }_{(\mathrm{\Lambda },x)}(K^{})\text{,}K^{}𝔎_0\text{,}$$
(4.8)
is a continuous mapping.
* There exists a strongly continuous unitary representation $`x\underset{¯}{U}(x)`$ of space-time translations $`x^{s+1}`$ with spectral measure $`\mathrm{\Delta }\underset{¯}{E}(\mathrm{\Delta })`$, supported by a displaced forward light cone $`\overline{V}_+q`$, $`q\overline{V}_+`$, which implements these transformations in the representation $`(\underset{¯}{\pi },\underset{¯}{})`$ via
$$\underset{¯}{U}(x)\underset{¯}{\pi }(A)\underset{¯}{U}(x)^{}=\underset{¯}{\pi }(\underset{¯}{\alpha }_x(A))\text{,}A\underset{¯}{𝔄}\text{,}x^{s+1}\text{.}$$
(4.9a)
On the subset $`\left\{|K^{}:K^{}𝔎_0\right\}`$ of $`\underset{¯}{}`$ this unitary group acts according to
$$\underset{¯}{U}(x)|K^{}=|\underset{¯}{\alpha }_x(K^{})\text{,}K^{}𝔎_0\text{,}$$
(4.9b)
and there holds the relation
$$\underset{¯}{E}(\mathrm{\Delta }^{})|K=|K\text{,}K𝔎(\mathrm{\Delta }^{})\text{,}$$
(4.9c)
where $`𝔎(\mathrm{\Delta }^{})`$ denotes the set of operators from $`𝔎^c𝔎_0`$ with energy-momentum transfer in the Borel set $`\mathrm{\Delta }^{}^{s+1}`$.
Through (4.9c) we have explicitly installed into the definition of restricted $`𝔎_0^c`$-particle weights the result of Proposition 3.17 for generic particle weights. A spectral assumption of this kind is of great importance since it constitutes the basis for the proof of the Cluster Property of Proposition 3.18, and the arguments presented there can be adopted literally, on condition that the obvious substitutions are observed, to implement it in the present reduced setting as well.
###### Proposition 4.2.
A restricted $`𝔎_0^c`$-particle weight $`(\underset{¯}{\pi },\underset{¯}{},\underset{¯}{𝔄},\underset{¯}{\alpha },\underset{¯}{𝖯}^c,𝔎_0^c)`$ has the Cluster Property presented in Proposition 3.18, with the reservation that the replacements $`𝔏_0𝔎_0^c`$ and $`𝔄\underset{¯}{𝔄}^c`$ have to be carried out.
###### Remark.
The rather intricate Definition 4.1 will find its justification in the subsequent section, where it turns out, that the characteristics listed above are exactly those which survive in the process of spatial disintegration—at least, it did regrettably not lie within our reach to establish a more complete list of features to be preserved. Nevertheless, it should be noted, that those characteristics motivating the interpretation of particle weights as asymptotic plane waves are perpetuated (cf. the first paragraph of Section 3.3).
Now, it does not come as a surprise that, with respect to the countable and separable notions introduced in Section 4.1, a particle weight of the general type gives rise to a restricted $`\overline{𝔏_0^c}`$-particle weight.
###### Theorem 4.3.
Let $`(\pi _w,_w)`$ be the GNS-representation corresponding to a given particle weight $`.|.`$ according to Theorem 3.15. Then $`(\pi ^{},^{},𝔄^{},\alpha ^{},𝖯^c,\overline{𝔏_0^c})`$ is a restricted $`\overline{𝔏_0^c}`$-particle weight, where the individual entries (if not already fixed by Section 4.1) are defined as follows:
* $`^{}`$ designates the Hilbert subspace of $`_w`$, which is the closed $``$-linear span of the assortment of vectors $`\left\{|L_w:L𝔏^c=\overline{𝔄^c}\overline{𝔏_0^c}\right\}`$ and thus separable;
* $`\pi ^{}\pi _w𝔄^{}`$ denotes the restriction of the initial representation to the algebra $`𝔄^{}`$, where the representatives have their limited domain as well as range on $`^{}`$;
* $`\{\alpha _{(\mathrm{\Lambda },x)}^{}\alpha _{(\mathrm{\Lambda },x)}𝔄^{}:(\mathrm{\Lambda },x)𝖯_+^{}\}`$ is the restriction of the initial automorphism group to $`𝔄^{}`$.
###### Proof.
With the definitions $`|.^{}|.𝔏^c`$ and $`U^{}(x)U_w(x)^{}`$, $`x^{s+1}`$, where the latter obviously leaves invariant $`^{}`$ and is such that the corresponding spectral measure turns out to be $`E^{}(\mathrm{\Delta })E_w(\mathrm{\Delta })^{}`$ for any Borel set $`\mathrm{\Delta }`$, all features of the restricted $`\overline{𝔏_0^c}`$-particle weight are readily checked on the grounds of Theorem 3.15 and Proposition 3.17. ∎
### 4.3 Spatial Disintegration of Particle Weights
We now get to the central result of this chapter: the construction of the spatial disintegration of a particle weight in terms of pure ones, or rather of the corresponding restricted $`\overline{𝔏_0^c}`$-particle weight into a direct integral of pure representations, which again are associated with restricted $`\overline{𝔏_0^c}`$-particle weights. In Theorem 4.3 the representation $`(\pi ^{},^{})`$ of the norm-separable $`C^{}`$-algebra $`𝔄^{}`$ on the separable Hilbert space $`^{}`$ was derived from the given particle weight $`.|.`$. This places the method of spatial disintegration expounded in the relevant literature at our disposal to apply it to the problem at hand. In order to express $`\pi ^{}`$ in terms of an integral of irreducible representations, a last preparatory step has to be taken: a *maximal abelian* von Neumann algebra $`𝔐`$ in the commutant of $`\pi ^{}(𝔄^{})`$ has to be selected in view of \[25, Theorem 8.5.2\]. The choice of such an algebra is restricted by our further objective to arrive at a disintegration in terms of restricted $`\overline{𝔏_0^c}`$-particle weights, which means that one has to provide for the possibility to establish the relations (4.9).
The unitary group $`\{U^{}(x):x^{s+1}\}`$ has generators with joint spectrum in a displaced forward light cone. Through multiplication by suitably chosen exponential factors $`\mathrm{exp}(iqx)`$ with fixed $`q\overline{V}_+`$ we can pass to another unitary group which likewise implements the space-time translations but has spectrum contained in $`\overline{V}_+`$. This places \[9, Theorem IV.5\] at our disposal, implying that one can find a strongly continuous unitary group of this kind with elements belonging to $`\pi ^{}(𝔄^{})^{\prime \prime }`$, the weak closure of $`\pi ^{}(𝔄^{})`$ (cf. \[11, Corollary 2.4.15\]). This result can again be tightened up by use of \[10, Theorem 3.3\] in the sense that among all the unitary groups complying with the above features there exists exactly one which is characterized by the further requirement that the lower boundary of the joint spectrum of its generators be Lorentz invariant. It is denoted as
$$\{U_c^{}(x)\pi ^{}(𝔄^{})^{\prime \prime }:x^{s+1}\}\text{.}$$
(4.10a)
At this point it turns out to be significant that the $`C^{}`$-algebra $`𝔄^{}`$ has been constructed in Section 4.1 by using local operators so that the reasoning given in applies to the present situation. Another unitary group can be defined through
$$\{V^{}(x)U_c^{}(x)U^{}(x)^1:x^{s+1}\}\text{.}$$
(4.10b)
By their very construction, all the operators $`V^{}(x)`$, $`x^{s+1}`$, are elements of $`\pi ^{}(𝔄^{})^{}`$. The maximal commutative von Neumann algebra $`𝔐`$ that we are going to work with in the sequel is now selected in compliance with the condition
$$\{V^{}(x):x^{s+1}\}^{\prime \prime }𝔐\left(\pi ^{}(𝔄^{})\{U^{}(x):x^{s+1}\}\right)^{}\text{.}$$
(4.11)
The main result to be acquired in the present chapter can then be summarized in the subsequent theorem.
###### Theorem 4.4.
Let $`.|.`$ be a generic particle weight with representation $`(\pi _w,_w)`$ inducing, by Theorem 4.3, the restricted $`\overline{𝔏_0^c}`$-particle weight $`(\pi ^{},^{},𝔄^{},\alpha ^{},𝖯^c,\overline{𝔏_0^c})`$. With respect to the representation $`(\pi ^{},^{})`$ of the separable $`C^{}`$-algebra $`𝔄^{}`$ on the separable Hilbert space $`^{}`$, we select a maximal abelian von Neumann algebra $`𝔐`$ such that (4.11) is fulfilled. Then there exist a standard Borel space $`\mathrm{X}`$, a bounded positive measure $`\nu `$ on $`\mathrm{X}`$, and a field of restricted $`\overline{𝔏_0^c}`$-particle weights
$$\mathrm{X}\xi (\pi _\xi ,_\xi ,𝔄^{},\alpha ^{},𝖯^\mathrm{c},\overline{𝔏_0^\mathrm{c}})\text{,}$$
(4.12)
such that the following assertions hold true:
* The field $`\xi (\pi _\xi ,_\xi )`$, as part of (4.12), is a $`\nu `$-measurable field of irreducible representations of $`𝔄^{}`$.
* The non-zero representation $`(\pi ^{},^{})`$ is unitarily equivalent to the direct integral of this field of irreducible representations:
$$(\pi ^{},^{})X^{}𝑑\nu (\xi )(\pi _\xi ,_\xi )\text{,}$$
(4.13a)
and, when $`W`$ denotes the unitary operator connecting both sides of (4.13a), the vectors in both spaces are linked up by the relation
$$W|L^{}=\{|L_\xi :\xi \mathrm{X}\}\mathrm{X}^{}\mathrm{d}\nu (\xi )|\mathrm{L}_\xi \text{,}\mathrm{L}𝔏^\mathrm{c}\overline{𝔏_0}\text{,}$$
(4.13b)
where $`|._\xi `$ denotes the linear mapping characteristic for the restricted $`\overline{𝔏_0^c}`$-particle weight $`(\pi _\xi ,_\xi ,𝔄^{},\alpha ^{},𝖯^c,\overline{𝔏_0^c})`$, according to (4.7a) in Definition 4.1.
* The von Neumann algebra $`𝔐`$ coincides with the algebra of those operators which are diagonalisable with respect to the above disintegration of $`(\pi ^{},^{})`$: any operator $`T𝔐`$ corresponds to an essentially bounded measurable complex-valued function $`g_T`$ according to
$$WTW^{}=X^{}𝑑\nu (\xi )g_T(\xi )\mathbf{\hspace{0.17em}1}_\xi \text{,}$$
(4.13c)
where $`\mathrm{𝟏}_\xi `$, $`\xi \mathrm{X}`$, are the unit operators of the algebras $`𝔅(_\xi )`$, respectively.
* Let $`\{U_\xi (x):x^{s+1}\}𝔅(_\xi )`$ denote the unitary group, which implements the space-time translations in the restricted $`\overline{𝔏_0^c}`$-particle weight pertaining to $`\xi \mathrm{X}`$ according to (4.9a), and let $`E_\xi (\mathrm{\Delta })𝔅(_\xi )`$ designate the corresponding spectral measure belonging to the Borel set $`\mathrm{\Delta }^{s+1}`$. Then the fields of operators
$$\xi U_\xi (x)\text{and}\xi E_\xi (\mathrm{\Delta })$$
are measurable and satisfy for any $`x`$ and any Borel set $`\mathrm{\Delta }`$ the following equations:
$`WU^{}(x)W^{}`$ $`={\displaystyle X^{}𝑑\nu (\xi )U_\xi (x)\text{,}}`$ (4.13d)
$`WE^{}(\mathrm{\Delta })W^{}`$ $`={\displaystyle X^{}𝑑\nu (\xi )E_\xi (\mathrm{\Delta })\text{.}}`$ (4.13e)
* There exists a canonical choice of a strongly continuous unitary group in each Hilbert space $`_\xi `$
$$\{U_\xi ^c(x)\pi _\xi (𝔄^{})^{\prime \prime }=𝔅(_\xi ):x^{s+1}\}\text{,}$$
(4.13f)
which is measurable with respect to $`\xi `$, implements the space-time translations in the representation $`(\pi _\xi ,_\xi )`$ and has generators $`P_\xi ^c`$ whose joint spectrum lies in the closed forward light cone $`\overline{V}_+`$. It is defined by
$$U_\xi ^c(x)\mathrm{exp}(ip_\xi x)U_\xi (x)\text{,}x^{s+1}\text{,}$$
(4.13g)
where $`p_\xi `$ is the unequivocal vector in $`^{s+1}`$ that is to be interpreted as the sharp energy-momentum corresponding to the respective particle weight.
###### Remark.
The concepts occurring in the theory of direct integrals of Hilbert spaces (standard Borel space, decomposable and diagonalisable operators, and the like) are expounded in \[4, Chapter 3\], \[24, Part II\] and likewise \[54, Section IV.8 and Appendix\].
###### Proof.
The presuppositions of this theorem meet the requirements for an application of \[25, Theorem 8.5.2\]. This supplies us with
* a standard Borel space $`\overline{\mathrm{X}}`$,
* a bounded positive measure $`\overline{\nu }`$ on $`\overline{\mathrm{X}}`$,
* a $`\overline{\nu }`$-measurable field $`\xi (\pi _\xi ,_\xi )`$ on $`\overline{\mathrm{X}}`$ consisting of irreducible representations $`\pi _\xi `$ of the $`C^{}`$-algebra $`𝔄^{}`$ on the Hilbert spaces $`_\xi `$,
* and an isomorphism (a linear isometry) $`\overline{W}`$ from $`^{}`$ onto the direct integral of these Hilbert spaces, such that
$$\overline{W}:^{}_{\overline{\mathrm{X}}}^{}𝑑\overline{\nu }(\xi )_\xi \text{,}$$
(4.14a)
transforms $`\pi ^{}`$ into the direct integral of the representations $`\pi _\xi `$ according to
$$\overline{W}\pi ^{}(A)\overline{W}^{}=_{\overline{\mathrm{X}}}^{}𝑑\overline{\nu }(\xi )\pi _\xi (A)\text{,}A𝔄^{}\text{,}$$
(4.14b)
and the maximal abelian von Neumann algebra $`𝔐`$ can be identified with the algebra of diagonalisable operators via
$$\overline{W}T\overline{W}^{}=_{\overline{\mathrm{X}}}^{}𝑑\overline{\nu }(\xi )g_T(\xi )\mathbf{\hspace{0.17em}1}_\xi \text{,}T𝔐\text{,}$$
(4.14c)
with an appropriate function $`g_TL^{\mathrm{}}(\overline{\mathrm{X}},d\overline{\nu }(\xi ))`$.
At first sight, the different statements of \[25, Theorem 8.5.2\] listed above seem to cover almost all of the assertions of the present Theorem 4.4, but one must not forget that the disintegration is to be expressed in terms of a field of restricted $`\overline{𝔏_0^c}`$-particle weights. So we are left with the task to establish their defining properties in the representations $`(\pi _\xi ,_\xi )`$ supplied by the standard disintegration theory. In accomplishing this assignment, one has to see to it that simultaneously relation (4.13b) is to be satisfied, which means that one is faced with the following problem: In general the isomorphism $`\overline{W}`$ connects a given vector $`\mathrm{\Psi }^{}`$ not with a unique vector field $`\{\mathrm{\Psi }_\xi :\xi \mathrm{X}\}`$ but rather with an equivalence class of such fields, characterized by the fact that its elements differ pairwise at most on $`\overline{\nu }`$-null sets. In contrast to this, (4.13b) associates the vector field $`\left\{|L_\xi :\xi \mathrm{X}\right\}`$ with $`|L^{}`$ for any $`L𝔏^c\overline{𝔏_0}`$, leaving no room for any ambiguity. In particular, the algebraic relations prevailing in the set $`𝔏^c\overline{𝔏_0}`$ which carry over to $`|.`$ have to be observed in defining each of the mappings $`|._\xi `$ which are characteristic of a restricted $`\overline{𝔏_0^c}`$-particle weight. The contents of the theorem quoted above, important as they are, can therefore only serve as the starting point for the constructions carried out below, in the course of which again and again $`\overline{\nu }`$-null sets have to be removed from $`\overline{\mathrm{X}}`$ to secure definiteness of the remaining components in the disintegration of a given vector. In doing so, one has to be cautious not to apply this procedure uncountably many times; for, otherwise, by accident the standard Borel space $`\mathrm{X}\overline{\mathrm{X}}`$ arising in the end could happen to be itself a $`\overline{\nu }`$-null set. Then, if $`\nu `$ denotes the restriction of $`\overline{\nu }`$ to this set, one would have $`\overline{\nu }(\mathrm{X})=\nu (\mathrm{X})=0`$, in contradiction to the disintegration (4.13a) of the *non-zero* representation $`(\pi ^{},^{})`$.
* As indicated above, our first task in view of (4.7a) and (4.7b) of Definition 4.1 will be to establish the existence of $`(+i)`$-linear mappings
$$|._\xi :𝔏^c_\xi ^cL|L_\xi \text{,}$$
(4.15a)
from $`𝔏^c`$ onto a dense subset $`_\xi ^c`$ of each of the component Hilbert spaces supplied by \[25, Theorem 8.5.2\] with the property
$$\pi _\xi (A)|L_\xi =|AL_\xi \text{,}A\overline{𝔄^c}\text{,}L𝔏^c\text{.}$$
(4.15b)
Now, by relation (4.14a), there exists to each $`L𝔏^c`$ an equivalence class of vector fields on $`\mathrm{X}`$ which corresponds to the element $`|L^{}`$ in $`^{}`$. The assumed $`(+i)`$-linearity of the mapping $`|.^{}:𝔏^c^{}`$ carries first of all over to these equivalence classes, but, upon selection of a single representative from each class, it turns out that every algebraic relation in question is fulfilled in all components of the representatives involved, possibly apart from those pertaining to a $`\overline{\nu }`$-null set. So, if we pick out one representative of the vector $`|L^{}`$ for every $`L`$ in the numerable set $`𝔏^c`$ and designate it as $`\left\{|L_\xi :\xi \mathrm{X}\right\}`$, all of the countably many relations that constitute $`(+i)`$-linearity are satisfied for $`\overline{\nu }`$-almost all of the components of these representatives. They can thus be taken to define the linear mappings of the form (4.15a) for all $`\xi `$ in a Borel subset $`\overline{\mathrm{X}}_1`$ of $`\overline{\mathrm{X}}`$, which is left by the procedure of dismissing an appropriate $`\overline{\nu }`$-null set for each algebraic relation to be satisfied.
The same reasoning can be applied to the disintegration of vectors of the form $`|AL^{}=\pi ^{}(A)|L^{}`$ with $`A\overline{𝔄^c}`$ and $`L𝔏^c`$. Again with (4.13b) in mind, the number of relations (4.15b) to be satisfied is countable, so that in view of relation (4.14b) the mere removal of an appropriate $`\overline{\nu }`$-null set from $`\overline{\mathrm{X}}_1`$ leaves only those indices $`\xi `$ behind, for which the mappings $`|._\xi `$ indeed have the desired property (4.15b).
In this way we have implemented by hand the first defining property of restricted $`\overline{𝔏_0^c}`$-particle weights in the representations $`(\pi _\xi ,_\xi )`$ for $`\overline{\nu }`$-almost all indices $`\xi `$. The only thing that remains to be done in this connection is to show that $`\left\{|L_\xi :L𝔏^c\right\}`$ is a dense subset $`_\xi ^c`$ in $`_\xi `$. But, according to \[24, Section II.1.6, Proposition 8\], the fact that the set $`\left\{|L^{}:L𝔏^c\right\}`$ is total in $`^{}`$ by assumption implies that the corresponding property holds for $`\overline{\nu }`$-almost all $`\xi `$ in the disintegration. Thus there exists a non-null Borel set $`\overline{\mathrm{X}}_2\overline{\mathrm{X}}_1`$, such that the corresponding mappings $`|._\xi `$, $`\xi \overline{\mathrm{X}}_2`$, have this property, too. In this way all of the characteristics presented in the first item of Definition 4.1 are fulfilled for $`\xi \overline{\mathrm{X}}_2`$ by the mappings (4.15a) constructed above, and additionally we have
$$\overline{W}|L^{}=_{\overline{\mathrm{X}}_2}^{}𝑑\overline{\nu }(\xi )|L_\xi \text{,}L𝔏^c\text{.}$$
(4.16)
* In the next step, the mappings $`|._\xi `$ have to be extended to the set $`\overline{𝔏_0}`$ of all Poincaré transforms of operators from $`\overline{𝔏_0^c}`$ in such a way that the counterpart of (4.8) in Definition 4.1 is continuous. In the present notation this is the mapping
$$𝖯_+^{}(\mathrm{\Lambda },x)|\alpha _{(\mathrm{\Lambda },x)}^{}(L^{})_\xi _\xi \text{,}L^{}\overline{𝔏_0}\text{.}$$
(4.17)
At this point the special selection of $`\overline{𝔏_0^c}`$ as consisting of compactly regularized vacuum annihilation operators comes into play, and also the invariance of this set under transformations $`(\mathrm{\Lambda },x)\overline{𝖯}^c`$ will be of importance. Great care has to be taken in these investigations based on the differentiability properties of the operators in question, that not uncountably many conditions are imposed on the mappings $`|._\xi `$, since anew not all of them will share the claimed extension property, but only a $`\overline{\nu }`$-null subset of $`\overline{\mathrm{X}}_2`$ shall get lost on the way.
To start with, note that the Poincaré group $`𝖯_+^{}`$ can be covered by a sequence of open sets $`𝖵_i`$ with compact closures $`𝖢_i`$, $`i`$, contained in corresponding open charts $`(𝖴_i,\varphi _i)`$ with the additional property that the sets $`\varphi _i(𝖢_i)^{d_𝖯}`$ are convex (e. g. consider the translates of the canonical coordinates $`(𝖴_0,\varphi _0)`$ around $`(\mathrm{𝟏},0)`$ to all elements of $`𝖯^c`$ and take suitable open subsets thereof). Select one of these compacta, say $`𝖢_k`$, and fix an element $`\widehat{L}_0\overline{𝔏_0^c}`$, which by assumption is given as a compactly supported regularization of an element $`L_0𝔏_0^c`$:
$$\widehat{L}_0=\alpha _F(L_0)_{𝖲_F}𝑑\mu (\mathrm{\Lambda },x)F(\mathrm{\Lambda },x)\alpha _{(\mathrm{\Lambda },x)}(L_0)\text{,}$$
(4.18a)
where $`F`$ is an infinitely often differentiable function on $`𝖯_+^{}`$ with compact support $`𝖲_F`$ in the Poincaré group $`𝖯_+^{}`$. According to Lemma 3.16 the mapping $`|.`$ commutes with this integral so that the vector $`|\widehat{L}_0`$ in $`_w`$ takes on the shape
$$|\widehat{L}_0=_{𝖲_F}𝑑\mu (\mathrm{\Lambda },x)F(\mathrm{\Lambda },x)|\alpha _{(\mathrm{\Lambda },x)}(L_0)\text{.}$$
(4.18b)
The same equation holds for the Poincaré transforms of the operator $`\widehat{L}_0`$ as well, so that invariance of the Haar measure on $`𝖯_+^{}`$ implies for any $`(\mathrm{\Lambda }_0,x_0)𝖢_k`$ the equations
$$\begin{array}{c}|\alpha _{(\mathrm{\Lambda }_0,x_0)}(\widehat{L}_0)=_{𝖲_F}𝑑\mu (\mathrm{\Lambda },x)F(\mathrm{\Lambda },x)|\alpha _{(\mathrm{\Lambda }_0,x_0)(\mathrm{\Lambda },x)}(L_0)\hfill \\ \hfill =_{(\mathrm{\Lambda }_0,x_0)𝖲_F}𝑑\mu (\mathrm{\Lambda },x)F\left((\mathrm{\Lambda }_0,x_0)^1(\mathrm{\Lambda },x)\right)|\alpha _{(\mathrm{\Lambda },x)}(L_0)\\ \hfill =_{𝖢_k𝖲_F}𝑑\mu (\mathrm{\Lambda },x)F\left((\mathrm{\Lambda }_0,x_0)^1(\mathrm{\Lambda },x)\right)|\alpha _{(\mathrm{\Lambda },x)}(L_0)\text{.}\end{array}$$
(4.18c)
The derivatives of the mapping $`(\mathrm{\Lambda }_0,x_0)|\alpha _{(\mathrm{\Lambda }_0,x_0)}(\widehat{L}_0)`$, the domain of $`(\mathrm{\Lambda }_0,x_0)`$ restricted to the neighbourhood $`𝖵_k`$ in $`𝖢_k`$, are thus explicitly seen to be expressible in terms of derivatives of the functions
$$F^{(\mathrm{\Lambda },x)}:𝖵_k(\mathrm{\Lambda }_0,x_0)F^{(\mathrm{\Lambda },x)}(\mathrm{\Lambda }_0,x_0)F\left((\mathrm{\Lambda }_0,x_0)^1(\mathrm{\Lambda },x)\right)\text{.}$$
So, let $`(\mathrm{\Lambda }_1,x_1)`$ and $`(\mathrm{\Lambda }_2,x_2)`$ be a pair of Poincaré transformations lying in the common neighbourhood $`𝖵_k`$; then the following equation results from an application of the Mean Value Theorem A.7 to the $`𝒳_{𝔏_0}`$-differentiable mapping $`|.`$ (cf. Theorem 3.15):
$$\begin{array}{c}|\alpha _{(\mathrm{\Lambda }_1,x_1)}(\widehat{L}_0)\alpha _{(\mathrm{\Lambda }_2,x_2)}(\widehat{L}_0)\hfill \\ \hfill =|\alpha _{\varphi _k^1(𝒔)}(\widehat{L}_0)\alpha _{\varphi _k^1(𝒕)}(\widehat{L}_0)=_0^1𝑑\vartheta |𝔇(\mathrm{\Xi }_{\widehat{L}_0}\varphi _k^1)(𝒕+\vartheta (𝒔𝒕))(𝒔𝒕)\\ \hfill =_0^1𝑑\vartheta _{𝖢_k𝖲_F}𝑑\mu (\mathrm{\Lambda },x)𝔇(F^{(\mathrm{\Lambda },x)}\varphi _k^1)(𝒕+\vartheta (𝒔𝒕))(𝒔𝒕)|\alpha _{(𝜦,𝒙)}(𝑳_\mathit{0})\text{,}\end{array}$$
(4.18d)
where $`𝒔\varphi _𝒌(𝜦_\mathit{1},𝒙_\mathit{1})`$ and $`𝒕\varphi _𝒌(𝜦_\mathit{2},𝒙_\mathit{2})`$ belong to the compact and *convex* set $`\varphi _k(𝖢_k)`$.
Now, the vector $`|\alpha _{(\mathrm{\Lambda }_1,x_1)}(\widehat{L}_0)\alpha _{(\mathrm{\Lambda }_2,x_2)}(\widehat{L}_0)`$ defines a positive functional on the algebra $`𝔅(_w)`$, and we want to show that this vector functional can be majorized by a positive normal functional in $`𝔅()_{}`$. To establish this fact, note, that the integrals in (4.18d) exist in the uniform topology of $`_w`$, so that they commute with every bounded linear operator $`B𝔅(_w)`$. Hence
$$\begin{array}{c}\alpha _{(\mathrm{\Lambda }_1,x_1)}(\widehat{L}_0)\alpha _{(\mathrm{\Lambda }_2,x_2)}(\widehat{L}_0)\left|B\right|\alpha _{(\mathrm{\Lambda }_1,x_1)}(\widehat{L}_0)\alpha _{(\mathrm{\Lambda }_2,x_2)}(\widehat{L}_0)\hfill \\ \hfill =\underset{[0,1]^2}{}𝑑\vartheta 𝑑\vartheta ^{}𝔇(\mathrm{\Xi }_{\widehat{L}_0}\varphi _k^1)(𝒕+\vartheta ^{}(𝒔𝒕))(𝒔𝒕)\left|𝑩\right|𝔇(𝜩_{\widehat{𝑳}_\mathit{0}}\varphi _𝒌^\mathit{1})(𝒕+\vartheta (𝒔𝒕))(𝒔𝒕)\text{.}\end{array}$$
(4.19a)
This equation is invariant with respect to an exchange of $`\vartheta `$ and $`\vartheta ^{}`$. In the case of a *positive* operator $`B`$ the following relation holds for arbitrary vectors $`\mathrm{\Psi }`$ and $`\mathrm{\Phi }`$ in $`_w`$:
$$\mathrm{\Psi }|B|\mathrm{\Phi }+\mathrm{\Phi }|B|\mathrm{\Psi }\mathrm{\Psi }|B|\mathrm{\Psi }+\mathrm{\Phi }|B|\mathrm{\Phi }\text{,}$$
which, applied to the integrand of (4.19a) and to that resulting from an interchange of $`\vartheta `$ and $`\vartheta ^{}`$, yields
$$\begin{array}{c}\alpha _{(\mathrm{\Lambda }_1,x_1)}(\widehat{L}_0)\alpha _{(\mathrm{\Lambda }_2,x_2)}(\widehat{L}_0)\left|B\right|\alpha _{(\mathrm{\Lambda }_1,x_1)}(\widehat{L}_0)\alpha _{(\mathrm{\Lambda }_2,x_2)}(\widehat{L}_0)\hfill \\ \hfill _0^1𝑑\vartheta 𝔇(\mathrm{\Xi }_{\widehat{L}_0}\varphi _k^1)(𝒕+\vartheta (𝒔𝒕))(𝒔𝒕)\left|𝑩\right|𝔇(𝜩_{\widehat{𝑳}_\mathit{0}}\varphi _𝒌^\mathit{1})(𝒕+\vartheta (𝒔𝒕))(𝒔𝒕)\end{array}$$
(4.19b)
upon execution of a trivial integration over $`\vartheta `$ and $`\vartheta ^{}`$, respectively. As in (4.18d) we can pass to the following representation for the integrand on the right-hand side of (4.19b):
$$\begin{array}{c}𝔇(\mathrm{\Xi }_{\widehat{L}_0}\varphi _k^1)(𝒕+\vartheta (𝒔𝒕))(𝒔𝒕)\left|𝑩\right|𝔇(𝜩_{\widehat{𝑳}_\mathit{0}}\varphi _𝒌^\mathit{1})(𝒕+\vartheta (𝒔𝒕))(𝒔𝒕)\hfill \\ \hfill =_{𝖢_k𝖲_F}d\mu (\mathrm{\Lambda },x)_{𝖢_k𝖲_F}d\mu (\mathrm{\Lambda }^{},x^{})\overline{𝔇(F^{(\mathrm{\Lambda }^{},x^{})}\varphi _k^1)(𝒕+\vartheta (𝒔𝒕))(𝒔𝒕)}\\ \hfill 𝔇(F^{(\mathrm{\Lambda },x)}\varphi _k^1)(𝒕+\vartheta (𝒔𝒕))(𝒔𝒕)\alpha _{(𝜦^{},𝒙^{})}(𝑳_\mathit{0})\left|𝑩\right|\alpha _{(𝜦,𝒙)}(𝑳_\mathit{0})\text{.}\end{array}$$
(4.19c)
The derivatives which show up in (4.19c) depend by construction continuously on the parameters $`𝒔`$ and $`𝒕`$, $`\vartheta `$ and $`\vartheta ^{}`$ as well as $`(\mathrm{\Lambda },x)`$ and $`(\mathrm{\Lambda }^{},x^{})`$, so that their absolute values, taken on the compact domains $`\varphi _k(𝖢_k)`$, $`[0,1]`$ and $`𝖢_k𝖲_F`$, respectively, are bounded by
$$\left|𝔇(F^{(\mathrm{\Lambda },x)}\varphi _k^1)(𝒕+\vartheta (𝒔𝒕))(𝒔𝒕)\right|𝑫(𝑭;𝖢_𝒌)|𝒔𝒕|<\mathrm{}$$
for all $`(\mathrm{\Lambda },x)𝖢_k𝖲_F`$ with a suitable non-negative constant $`D(F;𝖢_k)`$. Hence the non-negative matrix element in (4.19c) can be estimated by
$$\begin{array}{c}𝔇(\mathrm{\Xi }_{\widehat{L}_0}\varphi _k^1)(𝒕+\vartheta (𝒔𝒕))(𝒔𝒕)\left|𝑩\right|𝔇(𝜩_{\widehat{𝑳}_\mathit{0}}\varphi _𝒌^\mathit{1})(𝒕+\vartheta (𝒔𝒕))(𝒔𝒕)\hfill \\ \hfill D(F;𝖢_k)^2|𝒔𝒕|^\mathit{2}_{𝖢_𝒌𝖲_𝑭}𝒅\mu (𝜦,𝒙)_{𝖢_𝒌𝖲_𝑭}𝒅\mu (𝜦^{},𝒙^{})\left|\alpha _{(𝜦^{},𝒙^{})}(𝑳_\mathit{0})\left|𝑩\right|\alpha _{(𝜦,𝒙)}(𝑳_\mathit{0})\right|\text{,}\end{array}$$
(4.19d)
which is independent of $`\vartheta `$, so that insertion into (4.19b) yields
$$\begin{array}{c}\alpha _{(\mathrm{\Lambda }_1,x_1)}(\widehat{L}_0)\alpha _{(\mathrm{\Lambda }_2,x_2)}(\widehat{L}_0)\left|B\right|\alpha _{(\mathrm{\Lambda }_1,x_1)}(\widehat{L}_0)\alpha _{(\mathrm{\Lambda }_2,x_2)}(\widehat{L}_0)\hfill \\ \hfill D(F;𝖢_k)^2|𝒔𝒕|^\mathit{2}_{𝖢_𝒌𝖲_𝑭}𝒅\mu (𝜦,𝒙)_{𝖢_𝒌𝖲_𝑭}𝒅\mu (𝜦^{},𝒙^{})\left|\alpha _{(𝜦^{},𝒙^{})}(𝑳_\mathit{0})\left|𝑩\right|\alpha _{(𝜦,𝒙)}(𝑳_\mathit{0})\right|\text{.}\end{array}$$
(4.19e)
Since the positive operator $`B`$ can be written as $`C^{}C`$ for suitable $`C𝔅(_w)`$, the integrand on the right-hand side allows for the following estimate, making use of the relation between the geometric and the arithmetic mean of two non-negative numbers:
$$\begin{array}{c}\left|\alpha _{(\mathrm{\Lambda }^{},x^{})}(L_0)\left|B\right|\alpha _{(\mathrm{\Lambda },x)}(L_0)\right|C|\alpha _{(\mathrm{\Lambda }^{},x^{})}(L_0)C|\alpha _{(\mathrm{\Lambda },x)}(L_0)\hfill \\ \hfill 2^1\left(\alpha _{(\mathrm{\Lambda }^{},x^{})}(L_0)\left|B\right|\alpha _{(\mathrm{\Lambda }^{},x^{})}(L_0)+\alpha _{(\mathrm{\Lambda },x)}(L_0)\left|B\right|\alpha _{(\mathrm{\Lambda },x)}(L_0)\right)\text{.}\end{array}$$
As a consequence of this inequality entered into (4.19e), one integration over $`𝖢_k𝖲_F`$ can be carried out on its right-hand side for each resulting term of the sum, so that finally
$$\begin{array}{c}\alpha _{(\mathrm{\Lambda }_1,x_1)}(\widehat{L}_0)\alpha _{(\mathrm{\Lambda }_2,x_2)}(\widehat{L}_0)\left|B\right|\alpha _{(\mathrm{\Lambda }_1,x_1)}(\widehat{L}_0)\alpha _{(\mathrm{\Lambda }_2,x_2)}(\widehat{L}_0)\hfill \\ \hfill D(F;𝖢_k)^2|𝒔𝒕|^\mathit{2}\mu (𝖢_𝒌𝖲_𝑭)_{𝖢_𝒌𝖲_𝑭}𝒅\mu (𝜦,𝒙)\alpha _{(𝜦,𝒙)}(𝑳_\mathit{0})\left|𝑩\right|\alpha _{(𝜦,𝒙)}(𝑳_\mathit{0})\text{,}\end{array}$$
(4.19f)
where the last integral can be viewed as a positive normal functional on $`𝔅(_w)`$ in the variable $`B`$, as announced at the beginning of this paragraph.
Now, let $`\mathrm{M}`$ be a measurable subset of $`\overline{\mathrm{X}}`$ then, according to (4.14c), it corresponds via the associated characteristic function $`\chi M`$ to a projection $`PM`$ in the selected maximal abelian von Neumann algebra $`𝔐`$. If $`P^{}`$ in turn denotes the orthogonal projection from $`_w`$ onto the Hilbert space $`^{}`$, we can define $`BMP^{}PMP^{}`$ as a positive operator in $`𝔅(_w)`$, which is therefore subject to (4.19f). This relation can then be re-written for $`B=BM`$ in terms of the restricted $`\overline{𝔏_0^c}`$-particle weight $`(\pi ^{},^{},𝔄^{},\alpha ^{},𝖯^c,\overline{𝔏_0^c})`$:
$$\begin{array}{c}PM|\alpha _{(\mathrm{\Lambda }_1,x_1)}^{}(\widehat{L}_0)\alpha _{(\mathrm{\Lambda }_2,x_2)}^{}(\widehat{L}_0)^{}^2=PMP^{}|\alpha _{(\mathrm{\Lambda }_1,x_1)}(\widehat{L}_0)\alpha _{(\mathrm{\Lambda }_2,x_2)}(\widehat{L}_0)^2\hfill \\ \hfill D(F;𝖢_k)^2|𝒔𝒕|^\mathit{2}\mu (𝖢_𝒌𝖲_𝑭)_{𝖢_𝒌𝖲_𝑭}𝒅\mu (𝜦,𝒙)\alpha _{(𝜦,𝒙)}(𝑳_\mathit{0})\left|𝑷^{}\mathrm{𝑷𝑴𝑷}^{}\right|\alpha _{(𝜦,𝒙)}(𝑳_\mathit{0})\text{,}\end{array}$$
(4.20a)
where now the integral on the right-hand side defines a positive normal functional on the von Neumann algebra $`𝔐`$ through
$$\phi [\widehat{L}_0;𝖢_k](T)_{𝖢_k𝖲_F}𝑑\mu (\mathrm{\Lambda },x)\alpha _{(\mathrm{\Lambda },x)}(L_0)\left|P^{}TP^{}\right|\alpha _{(\mathrm{\Lambda },x)}(L_0)\text{,}T𝔐\text{.}$$
(4.20b)
Specializing to Poincaré transformations $`(\mathrm{\Lambda }_1,x_1)`$ and $`(\mathrm{\Lambda }_2,x_2)`$ from the countable subgroup $`\overline{𝖯}^c`$, the unique disintegration of the vector $`|\alpha _{(\mathrm{\Lambda }_1,x_1)}^{}(\widehat{L}_0)\alpha _{(\mathrm{\Lambda }_2,x_2)}^{}(\widehat{L}_0)^{}`$ occurring on the left-hand side of (4.20a) is already explicitly given by (4.15a) for all $`\xi \overline{\mathrm{X}}_2`$ so that
$$\overline{W}|\alpha _{(\mathrm{\Lambda }_1,x_1)}^{}(\widehat{L}_0)\alpha _{(\mathrm{\Lambda }_2,x_2)}^{}(\widehat{L}_0)^{}=_{\overline{\mathrm{X}}_2}^{}𝑑\nu (\xi )|\alpha _{(\mathrm{\Lambda }_1,x_1)}^{}(\widehat{L}_0)\alpha _{(\mathrm{\Lambda }_2,x_2)}^{}(\widehat{L}_0)_\xi \text{.}$$
(4.20c)
On the other hand, the positive normal functional $`\phi [\widehat{L}_0,𝖢_k]𝔐_{}`$ of (4.20b) is easily seen by \[54, Proposition IV.8.34\] in connection with (4.14c) to correspond to a unique integrable field $`\{\phi [\widehat{L}_0,𝖢_k]_\xi :\xi \overline{\mathrm{X}}\}`$ of positive normal functionals on the von Neumann algebras $`\mathrm{𝟏}_\xi `$ in the direct integral decomposition of $`𝔐`$. Explicitly,
$$\phi [\widehat{L}_0;𝖢_k](T)=_{\overline{\mathrm{X}}}𝑑\overline{\nu }(\xi )g_T(\xi )\phi [\widehat{L}_0,𝖢_k]_\xi (\mathrm{𝟏}_\xi )$$
(4.20d)
for any $`T𝔐`$ with an appropriate function $`g_TL^{\mathrm{}}(\overline{\mathrm{X}},d\overline{\nu }(\xi ))`$. The above relation stays true, if we replace $`\overline{\mathrm{X}}`$ by $`\overline{\mathrm{X}}_2`$, since both differ at most by a $`\overline{\nu }`$-null set. So, in view of relations (4.20b) through (4.20d), (4.20a) can for any measurable subset $`\mathrm{M}`$ of $`\overline{\mathrm{X}}_2`$ corresponding to the orthogonal projection $`PM𝔐`$ be expressed in terms of integrals according to
$$\begin{array}{c}_\mathrm{M}𝑑\overline{\nu }(\xi )|\alpha _{(\mathrm{\Lambda }_1,x_1)}^{}(\widehat{L}_0)\alpha _{(\mathrm{\Lambda }_2,x_2)}^{}(\widehat{L}_0)_\xi ^2\hfill \\ \hfill D(F;𝖢_k)^2|𝒔𝒕|^\mathit{2}\mu (𝖢_𝒌𝖲_𝑭)_\mathrm{M}𝒅\overline{\nu }(\xi )\phi [\widehat{𝑳}_\mathit{0},𝖢_𝒌]_\xi (\mathrm{𝟏}_\xi )\text{.}\end{array}$$
(4.20e)
Due to arbitrariness of $`\mathrm{M}\overline{\mathrm{X}}_2`$, we then infer, making use of elementary results of integration theory \[35, Chapter V, viz. § 25, Theorem D\], that for $`\overline{\nu }`$-almost all $`\xi \overline{\mathrm{X}}_2`$ there holds the estimate
$$\begin{array}{c}|\alpha _{(\mathrm{\Lambda }_1,x_1)}^{}(\widehat{L}_0)\alpha _{(\mathrm{\Lambda }_2,x_2)}^{}(\widehat{L}_0)_\xi \hfill \\ \hfill \left|\varphi _k(\mathrm{\Lambda }_1,x_1)\varphi _k(\mathrm{\Lambda }_2,x_2)\right|D(F;𝖢_k)\mu (𝖢_k𝖲_F)\phi [\widehat{L}_0,𝖢_k]_\xi (\mathrm{𝟏}_\xi )\text{,}\end{array}$$
(4.20f)
where we replaced the points $`𝒔`$ and $`𝒕`$ from the space $`^{d_𝖯}`$ of coordinates for $`𝖯_+^{}`$ by their pre-images $`(\mathrm{\Lambda }_1,x_1)`$ and $`(\mathrm{\Lambda }_2,x_2)`$ from $`𝖵_k\overline{𝖯}^c`$. The important thing to notice at this point is that, apart from the factor $`\left|\varphi _k(\mathrm{\Lambda }_1,x_1)\varphi _k(\mathrm{\Lambda }_2,x_2)\right|`$, the terms on the right-hand side of (4.20f) hinge upon the operator $`\widehat{L}_0`$ (determining the function $`F`$ as well as its support $`𝖲_F`$) and on the neighbourhood $`𝖵_k`$ with compact closure $`𝖢_k`$ containing $`(\mathrm{\Lambda }_1,x_1)\text{,}(\mathrm{\Lambda }_2,x_2)\overline{𝖯}^c`$. Therefore this estimate also holds for any other pair of Lorentz transformations in $`𝖵_k\overline{𝖯}^c`$ with the same $`(\widehat{L}_0,𝖵_k)`$-dependent factor; of course, in each of the resulting countably many relations one possibly loses a further $`\overline{\nu }`$-null subset of $`\overline{\mathrm{X}}_2`$. The reasoning leading up to this point can then be applied to any combination of an operator in the numerable selection $`\overline{𝔏_0^c}`$ with an open set from the countable cover of $`𝖯_+^{}`$ to produce in each case a relation of the form of (4.20f) for the respective Poincaré transformations in $`\overline{𝖯}^c`$. Simultaneously, the domain of indices $`\xi `$, for which *all* of these inequalities are valid, shrinks to an appropriate $`\overline{\nu }`$-measurable non-null subset $`\overline{\mathrm{X}}_3`$ of $`\overline{\mathrm{X}}_2`$.
Consider now an arbitrary element $`(\mathrm{\Lambda }_0,x_0)𝖯_+^{}`$, which belongs to at least one of the open sets $`𝖵_j`$ from the covering of the Poincaré group already used above. By density of $`\overline{𝖯}^c`$ in $`𝖯_+^{}`$, the transformation $`(\mathrm{\Lambda }_0,x_0)`$ can be approximated by a sequence $`\left\{(\mathrm{\Lambda }_n,x_n)\right\}_n\overline{𝖯}^c𝖵_j`$. This is a Cauchy sequence in the initial topology of $`𝖯_+^{}`$, so that relation (4.20f) implies that for each $`\xi \overline{\mathrm{X}}_3`$ the corresponding sequences
$$\left\{|\alpha _{(\mathrm{\Lambda }_n,x_n)}^{}(\widehat{L}_0)_\xi \right\}_n_\xi \text{,}\widehat{L}_0\overline{𝔏_0^c}\text{,}$$
(4.21a)
likewise have the Cauchy property with respect to the Hilbert space norms. Their limits in each of the spaces $`_\xi `$, $`\xi \overline{\mathrm{X}}_3`$, thus exist and are obviously independent of the approximating sequence of Lorentz transformations from $`\overline{𝖯}^c`$. Therefore, we can write
$$|\widehat{L}_0;(\mathrm{\Lambda }_0,x_0)_\xi \underset{n\mathrm{}}{lim}|\alpha _{(\mathrm{\Lambda }_n,x_n)}^{}(\widehat{L}_0)_\xi \text{,}$$
(4.21b)
a result that holds for arbitrary $`\widehat{L}_0\overline{𝔏_0^c}`$ as long as $`\xi `$ is taken from the non-null set $`\overline{\mathrm{X}}_3`$. According to \[26, Definition II.4.1\], which lays down the notion of measurability for vector fields, the mapping
$$\overline{\mathrm{X}}_3\xi |\widehat{L}_0;(\mathrm{\Lambda }_0,x_0)_\xi _\xi \text{,}$$
(4.21c)
that arises as the pointwise limit of measurable vector fields on $`\overline{\mathrm{X}}_3`$, is itself measurable with respect to the restriction of $`\overline{\nu }`$ to this subset of $`\overline{\mathrm{X}}`$ and turns out to be a representative of the vector $`|\alpha _{(\mathrm{\Lambda }_0,x_0)}^{}(\widehat{L}_0)^{}^{}`$ (cf. \[24, Section II.1.5, Proof of Proposition 5(ii)\], and note that we can neglect the null set missing in $`\overline{\mathrm{X}}_3`$ compared to $`\overline{\mathrm{X}}`$).
The question now is, if the limits $`|\widehat{L}_0;(\mathrm{\Lambda }_0,x_0)_\xi `$, constructed by the above method for arbitrary operators $`\widehat{L}_0\overline{𝔏_0^c}`$ and any transformation $`(\mathrm{\Lambda }_0,x_0)𝖯_+^{}`$, can *unambiguously* be identified for all $`\xi `$ in $`\overline{\mathrm{X}}_3`$ with vectors $`|\alpha _{(\mathrm{\Lambda }_0,x_0)}^{}(\widehat{L}_0)_\xi _\xi `$, which satisfy a relation of the form (4.13b). One of the situations, in which an inconsistency possibly arises, is the appearance of two different representations for a single element $`L^{}\overline{𝔏_0}`$:
$$L^{}=\alpha _{(\mathrm{\Lambda }_1,x_1)}^{}(\widehat{L}_1)=\alpha _{(\mathrm{\Lambda }_2,x_2)}^{}(\widehat{L}_2)\text{,}$$
(4.21d)
where $`\widehat{L}_1\text{,}\widehat{L}_2\overline{𝔏_0^c}`$, and $`(\mathrm{\Lambda }_1,x_1)\text{,}(\mathrm{\Lambda }_2,x_2)𝖯_+^{}`$. In this case the pair of operators is connected by the Poincaré transformation $`(\mathrm{\Lambda }_1,x_1)^1(\mathrm{\Lambda }_2,x_2)`$, which belongs to the subgroup $`\overline{𝖯}^c`$ of $`𝖯_+^{}`$ according to the constructions of Section 4.1. Therefore
$$\widehat{L}_1=\alpha _{(\mathrm{\Lambda }_1,x_1)^1(\mathrm{\Lambda }_2,x_2)}^{}(\widehat{L}_2)\text{,}$$
which implies that
$$\alpha _{(\mathrm{\Lambda }_{1,n},x_{1,n})}^{}(\widehat{L}_1)=\alpha _{(\mathrm{\Lambda }_{1,n},x_{1,n})(\mathrm{\Lambda }_1,x_1)^1(\mathrm{\Lambda }_2,x_2)}^{}(\widehat{L}_2)$$
for any sequence $`\left\{(\mathrm{\Lambda }_{1,n},x_{1,n})\right\}_n\overline{𝖯}^c`$ approximating $`(\mathrm{\Lambda }_1,x_1)`$. But then the transformations on the right-hand side of the last equation constitute another sequence in $`\overline{𝖯}^c`$, which in this case tends to $`(\mathrm{\Lambda }_2,x_2)`$ in the limit $`n\mathrm{}`$. As a consequence of the independence of the limits (4.21b) from the selected sequence in $`\overline{𝖯}^c`$, we could define
$$|L^{}_\xi |\widehat{L}_1;(\mathrm{\Lambda }_1,x_1)_\xi =|\widehat{L}_2;(\mathrm{\Lambda }_2,x_2)_\xi \text{.}$$
(4.21e)
The only problem that is still left open with respect to an unequivocal definition of vectors of the form $`|L^{}_\xi `$, $`L^{}\overline{𝔏_0}`$, occurs when the vacuum annihilation operator $`L^{}`$ happens to be an element of $`𝔏^c`$, so that its components in the Hilbert spaces $`_\xi `$ have already been fixed in the initial step. But, as $`𝔏^c`$ is a numerable set, such a coincidence will be encountered at most countably often, so that relation (4.21e) indeed turns out to be the unique definition of $`|L^{}_\xi `$ for all $`\xi \overline{\mathrm{X}}_4`$, such that the relation
$$\overline{W}|L^{}^{}=_{\overline{\mathrm{X}}_4}^{}𝑑\overline{\nu }(\xi )|L^{}_\xi \text{,}L^{}\overline{𝔏_0}\text{,}$$
(4.21f)
is satisfied, where again $`\overline{\mathrm{X}}_4`$ is a $`\overline{\nu }`$-measurable subset which differs from $`\overline{\mathrm{X}}_3`$ only by a null set.
The $`_\xi `$-vectors corresponding to elements of $`\overline{𝔏_0}`$ that arise as Poincaré transforms of $`L^{}=\alpha _{(\mathrm{\Lambda }_0,x_0)}^{}(\widehat{L}_0)\overline{𝔏_0}`$ are defined according to (4.21), in particular by the relations (4.21e) and (4.21b). As a result, when $`(\mathrm{\Lambda }_1,x_1)`$ and $`(\mathrm{\Lambda }_2,x_2)`$ are closely neighbouring elements of $`𝖯_+^{}`$ so that their products with $`(\mathrm{\Lambda }_0,x_0)`$ lie in the common open neighbourhood $`𝖵_k`$, we get the following estimate, which is a direct consequence of the above constructions inserted into relation (4.20f) and which holds for any $`\xi \overline{\mathrm{X}}_4`$:
$$\begin{array}{c}|\alpha _{(\mathrm{\Lambda }_1,x_1)}^{}(L^{})_\xi |\alpha _{(\mathrm{\Lambda }_2,x_2)}^{}(L^{})_\xi \hfill \\ \hfill \left|\varphi _k\left((\mathrm{\Lambda }_0,x_0)(\mathrm{\Lambda }_1,x_1)\right)\varphi _k\left((\mathrm{\Lambda }_0,x_0)(\mathrm{\Lambda }_2,x_2)\right)\right|D(F;𝖢_k)\mu (𝖢_k𝖲_F)\phi [\widehat{L}_0,𝖢_k]_\xi (\mathrm{𝟏}_\xi )\text{.}\end{array}$$
(4.22)
This shows that the continuity property with respect to generic Poincaré transformations as expressed in (4.8) of Definition 4.1 is fulfilled by all the extended mappings $`|._\xi `$ introduced above for arbitrary $`L^{}\overline{𝔏_0}`$.
* The last property of restricted $`\overline{𝔏_0^c}`$-particle weights to be established is the existence of unitary groups $`\{U_\xi (x):x^{s+1}\}`$ in the representations $`(\pi _\xi ,_\xi )`$ which satisfy the relations (4.9). To construct them we first consider one element $`L`$ of the countable space $`𝔏^c`$ together with a single space-time translation $`y`$ in the numerable dense subgroup $`𝖳^c`$ of $`^{s+1}`$. By assumption (4.11), the von Neumann algebra $`𝔐`$ is contained in the commutant of $`\{U^{}(x):x^{s+1}\}`$, which means that for any measurable subset $`\mathrm{M}`$ of $`\overline{\mathrm{X}}_4`$ with associated orthogonal projection $`PM𝔐`$ there holds the equation
$$\begin{array}{c}_\mathrm{M}𝑑\overline{\nu }(\xi )|\alpha _y^{}(L)_\xi ^2=PM|\alpha _y^{}(L)^{}^2=PMU^{}(y)|L^{}^2\hfill \\ \hfill =U^{}(y)PM|L^{}^2=PM|L^{}^2=_\mathrm{M}𝑑\overline{\nu }(\xi )|L_\xi ^2\text{.}\end{array}$$
(4.23a)
Since this result is valid for arbitrary measurable sets $`\mathrm{M}`$, we infer by \[35, Chapter V, § 25, Theorem E\] that for $`\overline{\nu }`$-almost all $`\xi `$ the vectors are subject to the relation
$$|\alpha _y^{}(L)_\xi =|L_\xi \text{.}$$
(4.23b)
A corresponding equation can be derived for any other of the countable number of combinations of elements in $`𝔏^c`$ and $`𝖳^c`$, so that (4.23b) is true in all of these cases when the domain of $`\xi `$ is restricted to the $`\overline{\nu }`$-measurable subset $`\overline{\mathrm{X}}_5`$, which again differs from $`\overline{\mathrm{X}}_4`$ only by a null set. On $`\overline{\mathrm{X}}_5`$ we can then define for arbitrary $`y𝖳^c`$ the mappings
$$\overline{U}_\xi (y):_\xi ^c_\xi ^c\overline{U}_\xi (y)|L_\xi |\alpha _y^{}(L)_\xi $$
(4.23c)
which are indeed determined unambiguously according to (4.23b). By the same relation they are norm-preserving and, moreover, turn out to be $`(+i)`$-linear operators on the countable spaces $`_\xi ^c_\xi `$.
We want to extend the definition given by (4.23c) in two respects: All space-time translations $`y^{s+1}`$ should be permissible, and all vectors of $`_\xi `$ are to belong to the domain of the resulting operators. Now, let $`L`$ be an arbitrary element of $`𝔏^c`$, i. e.
$$L=\underset{i=1}{\overset{N}{}}A_iL_i\text{with}A_i\overline{𝔄^c}\text{and}L_i\overline{𝔏_0^c}\text{,}$$
(4.24a)
and consider $`x^{s+1}`$ approximated by the sequence $`\left\{x_n\right\}_n𝖳^c`$. Then, by definition (4.23c) in connection with property (4.15b), the translates by $`x_k`$ and $`x_l`$ of the vectors $`|L_\xi `$ are for $`\xi \overline{\mathrm{X}}_5`$ subject to the following relation:
$$\begin{array}{c}\overline{U}_\xi (x_k)|L_\xi \overline{U}_\xi (x_l)|L_\xi =|\alpha _{x_k}^{}\left(\underset{i=1}{\overset{N}{}}A_iL_i\right)_\xi |\alpha _{x_l}^{}\left(\underset{i=1}{\overset{N}{}}A_iL_i\right)_\xi \hfill \\ \hfill =\underset{i=1}{\overset{N}{}}\pi _\xi \left(\alpha _{x_k}^{}(A_i)\right)|\alpha _{x_k}^{}(L_i)_\xi \underset{i=1}{\overset{N}{}}\pi _\xi \left(\alpha _{x_l}^{}(A_i)\right)|\alpha _{x_l}^{}(L_i)_\xi \\ \hfill =\underset{i=1}{\overset{N}{}}\pi _\xi \left(\alpha _{x_k}^{}(A_i)\alpha _{x_l}^{}(A_i)\right)|\alpha _{x_k}^{}(L_i)_\xi +\underset{i=1}{\overset{N}{}}\pi _\xi \left(\alpha _{x_l}^{}(A_i)\right)\left(|\alpha _{x_k}^{}(L_i)_\xi |\alpha _{x_l}^{}(L_i)_\xi \right)\text{.}\end{array}$$
(4.24b)
Since $`\overline{\mathrm{X}}_3`$ is a subset of $`\overline{\mathrm{X}}_5`$, we have relation (4.20f) at our disposal; moreover, by assumption, the group of automorphisms $`\{\alpha _{(\mathrm{\Lambda },x)}:(\mathrm{\Lambda },x)𝖯_+^{}\}`$ is strongly continuous. As a result, the sequences
$$\left\{|\alpha _{x_k}^{}(L)_\xi \right\}_k\text{,}L\overline{𝔏_0^c}\text{,}\left\{\pi _\xi \left(\alpha _{x_k}^{}(A)\right)\right\}_k\text{,}A\overline{𝔄^c}\text{,}$$
(4.24c)
both have the Cauchy property and are thus convergent as well as bounded in the respective norm topologies. Applied to the elements of $`\overline{𝔏_0^c}`$ and $`\overline{𝔄^c}`$ appearing in the representation (4.24a) of $`L𝔏^c`$ this has the consequence, that the right-hand side of (4.24b) can be made arbitrarily small for all pairs $`k\text{,}l`$ exceeding a certain number. The terms on the left-hand side of this inequality thus turn out to be part of Cauchy sequences $`\left\{\overline{U}_\xi (x_k)|L_\xi \right\}_k`$ which converge in the Hilbert spaces $`_\xi `$. Since a renewed application of the above arguments shows that the arising limits are independent of the approximating sequence in $`𝖳^c`$, the following relation unambiguously defines the mappings $`\overline{U}_\xi (x)`$ for arbitrary $`x^{s+1}`$ and $`L𝔏^c`$:
$$\overline{U}_\xi (x)|L_\xi \underset{k\mathrm{}}{lim}\overline{U}_\xi (x_k)|L_\xi =\underset{k\mathrm{}}{lim}|\alpha _{x_k}^{}(L)_\xi \text{,}$$
(4.24d)
$$\text{where}𝖳^cx_k\underset{k\mathrm{}}{\overset{}{}}x^{s+1}\text{.}$$
Again these mappings act as $`(+i)`$-linear operators on the spaces $`_\xi ^c`$ and preserve the Hilbert space norm. As a consequence they can, by the standard procedure used for completions of uniform spaces, be continuously extended in a unique fashion to all of the Hilbert space on condition that their countable domain constitutes a dense subset of $`_\xi `$; but this is the case as $`\overline{\mathrm{X}}_5`$ is contained in $`\overline{\mathrm{X}}_2`$. Changing the notation from $`\overline{U}_\xi `$ to $`U_\xi `$ for these extensions, their definition on arbitrary vectors $`\mathrm{\Psi }_\xi _\xi `$ then reads for any $`x^{s+1}`$
$$U_\xi (x)\mathrm{\Psi }_\xi \underset{l\mathrm{}}{lim}\overline{U}_\xi (x)|L^{(l)}_\xi \text{,}$$
(4.24e)
$$\text{where}_\xi ^c|L^{(l)}_\xi \underset{l\mathrm{}}{\overset{}{}}\mathrm{\Psi }_\xi _\xi \text{,}$$
and this definition is again independent of the selected sequence. For any $`L𝔏^c`$ and any $`x^{s+1}`$ the vector field $`\left\{U_\xi (x)|L_\xi :\xi \overline{\mathrm{X}}_5\right\}`$, which is the pointwise limit of a sequence of measurable vector fields by (4.24d) and hence itself measurable \[26, Definition II.4.1\], corresponds to $`|\alpha _x^{}(L)^{}`$, the equivalent limit in $`^{}`$ (where we neglect the difference between $`\overline{\mathrm{X}}`$ and $`\overline{\mathrm{X}}_5`$ which is of measure $`0`$):
$$\overline{W}U^{}(x)|L^{}=\overline{W}|\alpha _x^{}(L)^{}=_{\overline{\mathrm{X}}_5}^{}𝑑\overline{\nu }(\xi )U_\xi (x)|L_\xi \text{.}$$
(4.24f)
Having defined the family of mappings $`\{U_\xi (x):x^{s+1}\}𝔅(_\xi )`$ for $`\xi \overline{\mathrm{X}}_5`$, we now have to check that they obey (4.9). First of all, note that, as an immediate consequence of the way in which they were introduced, these mappings are $``$-linear and norm-preserving. Another property that is readily checked by use of the relations (4.24e) and (4.24d) in connection with the estimates arising from (4.24b) with $`L^{(l)}`$ replacing $`L`$ is the fact that for arbitrary $`x\text{,}y^{s+1}`$
$$U_\xi (x)U_\xi (y)=U_\xi (x+y)\text{.}$$
(4.25)
From this we infer that, as $`U_\xi (0)=\mathrm{𝟏}_\xi `$, each operator $`U_\xi (x)`$ has the inverse $`U_\xi (x)`$ and thus proves to be an isometric isomorphism of $`_\xi `$. Hence, in accordance with (4.25), the set $`\{U_\xi (x):x^{s+1}\}`$ indeed turns out to be a unitary group in $`𝔅(_\xi )`$.
The strong continuity of this group is easily seen: Consider the operator $`L𝔏^c`$ as defined in (4.24a) and two sequences $`\{x_k\}_k`$, $`\{y_l\}_l`$ in $`𝖳^c`$ converging to $`x`$ and $`y`$, respectively. Then (4.24b) stays valid if we replace the translations $`x_l`$ by $`y_l`$ in each case. In compliance with (4.24d) it is then possible to pass to the limit, which results in the obvious estimate
$$\begin{array}{c}\overline{U}_\xi (x)|L_\xi \overline{U}_\xi (y)|L_\xi \hfill \\ \hfill \underset{i=1}{\overset{N}{}}\alpha _x^{}(A_i)\alpha _y^{}(A_i)|\alpha _x^{}(L_i)_\xi +\underset{i=1}{\overset{N}{}}A_i|\alpha _x^{}(L_i)_\xi |\alpha _y^{}(L_i)_\xi \text{.}\end{array}$$
(4.26a)
This explicit inequality shows that the right-hand side can be made arbitrarily small for all $`y`$ in an appropriate neighbourhood of $`x`$; for the first term this is brought about by the strong continuity of the automorphism group $`\{\alpha _{(\mathrm{\Lambda },x)}:(\mathrm{\Lambda },x)𝖯_+^{}\}`$, whereas for the second term it is a consequence of relation (4.22). The defining condition for strong continuity is therefore satisfied for vectors in the dense subset $`_\xi ^c`$. If now an arbitrary vector $`\mathrm{\Psi }_\xi _\xi `$ is considered, we can expand the difference $`U_\xi (x)\mathrm{\Psi }_\xi U_\xi (y)\mathrm{\Psi }_\xi `$ by introducing the corresponding translates of any element $`|L_\xi _\xi ^c`$ and, making use of the property of norm-preservation of the unitaries, arrive at
$$\begin{array}{c}U_\xi (x)\mathrm{\Psi }_\xi U_\xi (y)\mathrm{\Psi }_\xi \hfill \\ \hfill \mathrm{\Psi }_\xi |L_\xi +\overline{U}_\xi (x)|L_\xi \overline{U}_\xi (y)|L_\xi +|L_\xi \mathrm{\Psi }_\xi \text{.}\end{array}$$
(4.26b)
The right-hand side of this inequality can again be made smaller than any given bound by first choosing a suitable element $`|L_\xi _\xi ^c`$ from a small neighbourhood of $`\mathrm{\Psi }_\xi `$ and then, in dependence on this selected vector $`|L_\xi `$ but irrelevant for the statement, selecting an appropriate neighbourhood of translations $`y`$ around $`x`$ as implied by (4.26a). Thereby we have established strong continuity of the unitary group $`\{U_\xi (x):x^{s+1}\}`$.
Before considering the support of the spectral measure $`E_\xi (.)`$ associated with this strongly continuous unitary group, we mention a result on the interchange of integrations with respect to the Lebesgue measure on $`^{s+1}`$ and the bounded positive measure $`\overline{\nu }`$ on $`\overline{\mathrm{X}}_5`$, which proves to be necessary as Fubini’s Theorem does not apply to the situation in question. Let $`g`$ be a continuous bounded function in $`L^1(^{s+1},d^{s+1}x)`$, then
$$^{s+1}xg(x)L_1\left|U_\xi (x)\right|L_2_\xi $$
is an integrable mapping for any $`L_1\text{,}L_2𝔏^c`$ and $`\xi \overline{\mathrm{X}}_5`$. Moreover, it is Riemann integrable over any compact $`(s+1)`$-dimensional interval $`𝑲`$, and this integral is the limit of a Riemann sequence (cf. \[37, Kapitel XXIII, Abschnitt 197 and Lebesguesches Integrabilitätskriterium 199.3\]):
$$Kd^{s+1}xg(x)L_1\left|U_\xi (x)\right|L_2_\xi =\underset{i\mathrm{}}{lim}\underset{m=1}{\overset{n_i}{}}\left|𝒵_m^{(i)}\right|g\left(x_m^{(i)}\right)L_1\left|U_\xi \left(x_m^{(i)}\right)\right|L_2_\xi \text{,}$$
(4.27a)
where $`\{𝒵_m^{(i)}:m=1,\mathrm{},n_i\}`$ denotes the $`i`$-th subdivision of $`𝑲`$, $`\left|𝒵_m^{(i)}\right|`$ are the Lebesgue measures of these sets, and $`x_m^{(i)}𝒵_m^{(i)}`$ are corresponding intermediate points. The sums on the right-hand side of this equation turn out to be $`\overline{\nu }`$-measurable when their dependence on $`\xi `$ is taken into account, and so is the limit on the left-hand side. Moreover this property is preserved in passing to the limit $`𝑲^{𝒔+\mathit{1}}`$, so that
$$\overline{\mathrm{X}}_5\xi _{^{s+1}}d^{s+1}xg(x)L_1\left|U_\xi (x)\right|L_2_\xi $$
is $`\overline{\nu }`$-measurable and, in addition, integrable since
$$\begin{array}{c}_{\overline{\mathrm{X}}_5}𝑑\overline{\nu }(\xi )\left|_{^{s+1}}d^{s+1}xg(x)L_1\left|U_\xi (x)\right|L_2_\xi \right|\hfill \\ \hfill _{\overline{\mathrm{X}}_5}𝑑\overline{\nu }(\xi )_{^{s+1}}d^{s+1}x|g(x)||L_1_\xi |L_2_\xi \\ \hfill =g_1_{\overline{\mathrm{X}}_5}𝑑\overline{\nu }(\xi )|L_1_\xi |L_2_\xi g_1|L_1^{}|L_2^{}\text{.}\end{array}$$
(4.27b)
The counterpart of (4.27a) is valid in $`^{}`$, too, and, if $`\mathrm{M}`$ denotes a measurable subset of $`\overline{\mathrm{X}}_5`$ with associated orthogonal projection $`PM𝔐`$, this equation reads
$$Kd^{s+1}xg(x)L_1\left|PMU^{}(x)\right|L_2=\underset{i\mathrm{}}{lim}\underset{m=1}{\overset{n_i}{}}\left|𝒵_m^{(i)}\right|g\left(x_m^{(i)}\right)L_1\left|PMU^{}\left(x_m^{(i)}\right)\right|L_2\text{.}$$
(4.27c)
Then (4.16) and (4.24f) in connection with (4.14c) imply
$$\begin{array}{c}Kd^{s+1}xg(x)L_1\left|PMU^{}(x)\right|L_2=\underset{i\mathrm{}}{lim}_\mathrm{M}𝑑\overline{\nu }(\xi )\underset{m=1}{\overset{n_i}{}}\left|𝒵_m^{(i)}\right|g\left(x_m^{(i)}\right)L_1\left|U_\xi \left(x_m^{(i)}\right)\right|L_2_\xi \hfill \\ \hfill =_\mathrm{M}𝑑\overline{\nu }(\xi )Kd^{s+1}xg(x)L_1\left|U_\xi (x)\right|L_2_\xi \text{,}\end{array}$$
(4.27d)
where, in the second equation, use was made of Lebesgue’s Dominated Convergence Theorem in view of the fact that the integrable function
$$\overline{\mathrm{X}}_5\xi g_1|L_1_\xi |L_2_\xi $$
majorizes both sides of (4.27a). Relation (4.27d) again stays true in passing to the limit $`𝑲^{𝒔+\mathit{1}}`$:
$$_{^{s+1}}d^{s+1}xg(x)L_1\left|PMU^{}(x)\right|L_2=_\mathrm{M}𝑑\overline{\nu }(\xi )_{^{s+1}}d^{s+1}xg(x)L_1\left|U_\xi (x)\right|L_2_\xi \text{,}$$
(4.27e)
which is the announced result on the commutability of integrations in the present context.
The support of the spectral measure $`E_\xi (.)`$ associated with the generators $`P_\xi `$ of the unitary group $`\{U_\xi (x):x^{s+1}\}`$ can now be investigated by use of the method applied in the proof of the fourth item of the first part of Theorem 3.15. Note, that the complement of the closed set $`\overline{V}_+q^{s+1}`$ can be covered by an increasing sequence $`\left\{\mathrm{\Gamma }_N\right\}_N`$ of compact subsets (take e. g. the intersection of the compact ball of radius $`N`$ with the complement of the open $`N^1`$-neighbourhood of $`\overline{V}_+q`$). To each of these sets one can find an infinitely often differentiable function $`\stackrel{~}{g}_N`$ with support in $`\mathrm{}(\overline{V}_+q)`$ that has the property $`0\chi _{\mathrm{\Gamma }_N}\stackrel{~}{g}_N`$. As before, let $`\mathrm{M}`$ be a measurable subset of $`\overline{\mathrm{X}}_5`$ with associated orthogonal projection $`PM𝔐`$, then, by assumption on the spectral support of the unitary group implementing space-time translations in the underlying particle weight, we have
$$_{^{s+1}}d^{s+1}xg_N(x)L_1\left|PMU^{}(x)\right|L_2=0$$
(4.28a)
for any $`N`$ and any pair of vectors $`|L_1`$ and $`|L_2`$ with $`L_1\text{,}L_2𝔏^c`$, and this, according to (4.27e), implies
$$M𝑑\overline{\nu }(\xi )_{^{s+1}}d^{s+1}xg_N(x)L_1\left|U_\xi (x)\right|L_2_\xi =0\text{.}$$
(4.28b)
By arbitrariness of $`\mathrm{M}`$ in the last expression, we conclude once more that for $`\overline{\nu }`$-almost all $`\xi \overline{\mathrm{X}}_5`$
$$_{^{s+1}}d^{s+1}xg_N(x)L_1\left|U_\xi (x)\right|L_2_\xi =0\text{.}$$
(4.28c)
Eventually, if we want this equation to hold for any element of the countable set of triples $`(g_N,|L_1_\xi ,|L_2_\xi )`$, a non-null set $`\overline{\mathrm{X}}_6\overline{\mathrm{X}}_5`$ is left, and (4.28c) stays valid for the remaining $`\xi \overline{\mathrm{X}}_6`$ even if the special vectors $`|L_1_\xi `$ and $`|L_2_\xi `$ are replaced by arbitrary ones. Stone’s Theorem then implies (cf. (3.38)) that $`\stackrel{~}{g}_N(P_\xi )=0`$ and therefore, by the order relation inherent in the definition of $`\stackrel{~}{g}_N`$, we have $`E_\xi (\mathrm{\Gamma }_N)=\chi _{\mathrm{\Gamma }_N}(P_\xi )=0`$ for any $`N`$. As the spectral measure $`E_\xi (.)`$ is regular, one can pass to the limit $`N\mathrm{}`$ and thereby arrives at the desired result
$$E_\xi \left(\mathrm{}(\overline{V}_+q)\right)=0\text{,}\xi \overline{\mathrm{X}}_6\text{.}$$
(4.28d)
The defining equation (4.23c) in connection with (4.15b) implies that for arbitrary operators $`A^{}\overline{𝔄^c}`$ and $`L𝔏^c`$ and for any translation $`x^{}𝖳^c`$ one can write
$$\begin{array}{c}\pi _\xi \left(\alpha _x^{}^{}(A^{})\right)|L_\xi =|\alpha _x^{}^{}(A^{})L_\xi =|\alpha _x^{}^{}\left(A^{}\alpha _x^{}^{}(L)\right)_\xi =\overline{U}_\xi (x^{})|A^{}\alpha _x^{}^{}(L)_\xi \hfill \\ \hfill =\overline{U}_\xi (x^{})\pi _\xi (A^{})\overline{U}_\xi (x^{})|L_\xi =\overline{U}_\xi (x^{})\pi _\xi (A^{})\overline{U}_\xi (x^{})^{}|L_\xi \text{.}\end{array}$$
(4.29a)
Since the vectors $`|L_\xi _\xi ^c`$, $`L𝔏^c`$, constitute a dense subset of $`_\xi `$ for $`\xi \overline{\mathrm{X}}_2\overline{\mathrm{X}}_5`$ we infer from this equation that
$$\pi _\xi \left(\alpha _x^{}^{}(A^{})\right)=\overline{U}_\xi (x^{})\pi _\xi (A^{})\overline{U}_\xi (x^{})^{}\text{,}$$
(4.29b)
an equation that readily extends to all translations $`x`$ in $`^{s+1}`$ and, by uniform density of $`\overline{𝔄^c}`$ in $`𝔄^{}`$, also to any operator $`A`$ in the $`C^{}`$-algebra $`𝔄^{}`$:
$$\pi _\xi \left(\alpha _x^{}(A)\right)=U_\xi (x)\pi _\xi (A)U_\xi (x)^{}\text{,}A𝔄^{}\text{,}x^{s+1}\text{.}$$
(4.29c)
This proves the counterpart of (4.9a). The action of the group $`\{U_\xi (x):x^{s+1}\}`$ on $`\left\{|L^{}_\xi :L^{}\overline{𝔏_0}\right\}`$ according to (4.9b) is an immediate consequence of the defining relations (4.24d) and (4.24e) in connection with (4.21e) and the continuity property as expressed by (4.22). In the present setting we thus have
$$U_\xi (x)|L^{}_\xi |\alpha _x^{}(L^{})_\xi \text{,}L^{}\overline{𝔏_0}\text{.}$$
(4.30)
Now, let $`L`$ be an arbitrary element of $`𝔏^c`$ having energy-momentum transfer $`\mathrm{\Gamma }_L`$. Defined as the support of the Fourier transform of an operator-valued distribution (cf. the Remark following Definition 2.2), $`\mathrm{\Gamma }_L`$ is a closed Borel set, so that the arguments given in the preceding paragraph can again be applied when $`L`$ replaces $`L_1`$ and $`L_2`$ and the functions $`\stackrel{~}{g}_N`$ now correspond to an increasing sequence of compact sets $`\mathrm{\Gamma }_N^{}`$ which constitute a cover of $`\mathrm{}\mathrm{\Gamma }_L`$. As a result we arrive at the equivalent of (4.28c), so that
$$_{^{s+1}}d^{s+1}xg_N(x)L\left|U_\xi (x)\right|L_\xi =0$$
(4.31a)
holds for $`\overline{\nu }`$-almost all $`\xi \overline{\mathrm{X}}_6`$ even if the index $`N`$ is allowed to run through all natural numbers. As in the preceding paragraph we then conclude that for all of these $`\xi `$ and all $`N`$ one has $`L\left|\stackrel{~}{g}_N(P_\xi )\right|L_\xi =0`$ and hence $`L\left|E_\xi (\mathrm{\Gamma }_N^{})\right|L_\xi =0`$. According to the regularity of the spectral measure $`E_\xi (.)`$, passage to the limit with respect to $`N`$ yields the equation $`E_\xi \left(\mathrm{}\mathrm{\Gamma }_L\right)|L_\xi =0`$. By countability, this last result is valid for arbitrary $`L𝔏^c`$ if a $`\overline{\nu }`$-measurable non-null set $`\overline{\mathrm{X}}_7\overline{\mathrm{X}}_6`$ is appropriately selected, from which the indices $`\xi `$ are to be taken. The complementary statement thus constitutes a restricted version of the counterpart of (4.9c):
$$E_\xi (\mathrm{\Gamma }_L)|L_\xi =|L_\xi \text{,}L𝔏^c\text{,}\xi \overline{\mathrm{X}}_7\text{.}$$
(4.31b)
Now, let $`\widehat{L}_0`$ be an arbitrary element of $`\overline{𝔏_0^c}`$, then the energy-momentum transfer of its Poincaré transform by $`(\mathrm{\Lambda },x)𝖯^c`$, i. e. of the operator $`\alpha _{(\mathrm{\Lambda },x)}^{}(\widehat{L}_0)\overline{𝔏_0^c}𝔏^c`$, is given by $`\mathrm{\Lambda }\mathrm{\Gamma }_{\widehat{L}_0}`$, so that, according to (4.31b),
$$E_\xi (\mathrm{\Lambda }\mathrm{\Gamma }_{\widehat{L}_0})|\alpha _{(\mathrm{\Lambda },x)}^{}(\widehat{L}_0)_\xi =|\alpha _{(\mathrm{\Lambda },x)}^{}(\widehat{L}_0)_\xi \text{,}\xi \overline{\mathrm{X}}_7\text{.}$$
(4.31c)
This result can be applied to investigate the case of a generic element from $`\overline{𝔏_0}`$. For arbitrary $`(\mathrm{\Lambda }_0,x_0)𝖯_+^{}`$ approximated by the sequence $`\left\{(\mathrm{\Lambda }_n,x_n)\right\}_n𝖯^c`$ we have, by virtue of the relevant parts of (4.21),
$$|\alpha _{(\mathrm{\Lambda }_0,x_0)}^{}(\widehat{L}_0)_\xi =\underset{n\mathrm{}}{lim}|\alpha _{(\mathrm{\Lambda }_n,x_n)}^{}(\widehat{L}_0)_\xi \text{,}$$
and Lebesgue’s Dominated Convergence Theorem in connection with Stone’s Theorem yields for any function $`gL^1(^{s+1},d^{s+1}x)`$ and any index $`\xi \overline{\mathrm{X}}_7`$
$$\begin{array}{c}_{^{s+1}}d^{s+1}xg(x)\alpha _{(\mathrm{\Lambda }_0,x_0)}^{}(\widehat{L}_0)\left|U_\xi (x)\right|\alpha _{(\mathrm{\Lambda }_0,x_0)}^{}(\widehat{L}_0)_\xi \hfill \\ \hfill =\underset{n\mathrm{}}{lim}_{^{s+1}}d^{s+1}xg(x)\alpha _{(\mathrm{\Lambda }_n,x_n)}^{}(\widehat{L}_0)\left|U_\xi (x)\right|\alpha _{(\mathrm{\Lambda }_n,x_n)}^{}(\widehat{L}_0)_\xi \\ \hfill =(2\pi )^{(s+1)/2}\underset{n\mathrm{}}{lim}\alpha _{(\mathrm{\Lambda }_n,x_n)}^{}(\widehat{L}_0)\left|\stackrel{~}{g}(P_\xi )\right|\alpha _{(\mathrm{\Lambda }_n,x_n)}^{}(\widehat{L}_0)_\xi \text{.}\end{array}$$
(4.31d)
In the limit of large $`n`$ one finds the energy-momentum transfer $`\mathrm{\Lambda }_n\mathrm{\Gamma }_{\widehat{L}_0}`$ of $`\alpha _{(\mathrm{\Lambda }_n,x_n)}^{}(\widehat{L}_0)`$ in a small $`\epsilon `$-neighbourhood of $`\mathrm{\Lambda }_0\mathrm{\Gamma }_{\widehat{L}_0}`$. Therefore, in view of (4.31c), the right-hand side of (4.31d) vanishes for all $`n`$ exceeding a certain bound $`N`$ if $`g`$ is chosen in such a way that $`\mathrm{supp}\stackrel{~}{g}\mathrm{}(\mathrm{\Lambda }_0\mathrm{\Gamma }_{\widehat{L}_0})`$. The Fourier transform of the distribution
$$^{s+1}x\alpha _{(\mathrm{\Lambda }_0,x_0)}^{}(\widehat{L}_0)\left|U_\xi (x)\right|\alpha _{(\mathrm{\Lambda }_0,x_0)}^{}(\widehat{L}_0)_\xi $$
is thus seen to be supported by $`\mathrm{\Lambda }_0\mathrm{\Gamma }_{\widehat{L}_0}`$ from which we infer
$$E_\xi (\mathrm{\Lambda }_0\mathrm{\Gamma }_{\widehat{L}_0})|\alpha _{(\mathrm{\Lambda }_0,x_0)}^{}(\widehat{L}_0)_\xi =|\alpha _{(\mathrm{\Lambda }_0,x_0)}^{}(\widehat{L}_0)_\xi \text{,}\xi \overline{\mathrm{X}}_7\text{,}$$
(4.31e)
which is the equivalent of (4.31c) for arbitrary operators in $`\overline{𝔏_0}`$. Equations (4.31b) and (4.31e) are readily generalized, making use of the order structure of spectral projections reflecting the inclusion relation of Borel subsets of $`^{s+1}`$. If $`\overline{𝔏}(\mathrm{\Delta }^{})`$ denotes the set of operators from $`𝔏^c\overline{𝔏_0^c}`$ having energy-momentum transfer in the Borel set $`\mathrm{\Delta }^{}`$, then
$$E_\xi (\mathrm{\Delta }^{})|L_\xi =|L_\xi \text{,}L\overline{𝔏}(\mathrm{\Delta }^{})\text{,}$$
(4.31f)
and thus the counterpart of (4.9c) is established for the remaining indices $`\xi `$ from the non-null subset $`\overline{\mathrm{X}}_7`$ of $`\overline{\mathrm{X}}`$.
The above construction has supplied us with a measurable subset $`\mathrm{X}\overline{\mathrm{X}}_7`$ of the standard Borel space $`\overline{\mathrm{X}}`$, that was introduced at the outset, emerging from an application of \[25, Theorem 8.5.2\]. $`\mathrm{X}`$ is a non-null set, differing from $`\overline{\mathrm{X}}`$ only by a $`\overline{\nu }`$-null set. Moreover it is itself a standard Borel space (cf. the definition in \[4, Section 3.3\]), and we shall denote the restriction of the measure $`\overline{\nu }`$ to it by $`\nu `$; $`\nu \overline{\nu }\mathrm{X}`$ is again a bounded positive measure. Moreover, and this has been the central aim of the previous investigations, the field
$$\mathrm{X}\xi (\pi _\xi ,_\xi ,𝔄^{},\alpha ^{},𝖯^\mathrm{c},\overline{𝔏_0^\mathrm{c}})\text{,}$$
indeed consists of $`\overline{𝔏_0^c}`$-particle weights. What remains to be done now is a verification of the properties listed in (4.13).
* Arising as the restriction to a measurable subset in $`\overline{\mathrm{X}}`$ of a field of irreducible representations, the family of representations
$$\mathrm{X}\xi (\pi _\xi ,_\xi )$$
is obviously $`\nu `$-measurable and its components inherit the feature of irreducibility.
* As $`\mathrm{X}`$ and $`\overline{\mathrm{X}}`$ only differ by a $`\overline{\nu }`$-null set, one has
$$_{\overline{\mathrm{X}}}^{}𝑑\overline{\nu }(\xi )_\xi X^{}𝑑\nu _\xi \text{,}$$
(4.32)
and the relations (4.14) can be reformulated, using the right-hand side of (4.32) and an isomorphism $`W`$, which is the composition of $`\overline{W}`$ with the isometry that implements the above unitary equivalence. As an immediate consequence of (4.14a) and (4.14b) we then get the equivalence assertion of (4.13a). Moreover, by (4.16) and (4.21f), the operator $`W`$ connects the vector fields $`\left\{|L_\xi :\xi \mathrm{X}\right\}`$ with vectors $`|L^{}`$ for $`L𝔏^c\overline{𝔏_0}`$ in the desired way as expressed in (4.13b).
* (4.13c) is the mere reformulation of (4.14c) in terms of $`\mathrm{X}`$ and $`W`$.
* According to the argument preceding (4.24f), the mappings $`\xi L_1|U_\xi (x)|L_2_\xi `$, with $`\xi `$ restricted to $`\mathrm{X}`$ and $`L_1`$ as well as $`L_2`$ taken from $`𝔏^c`$, are measurable for all vectors $`|L_1_\xi `$ and $`|L_2_\xi `$ in the dense subsets $`_\xi ^c`$, and this suffices, by \[24, Section II.2.1, Proposition 1\], to establish measurability of the field $`\xi U_\xi (x)`$ for arbitrary $`x^{s+1}`$. Moreover, this is a bounded field of operators, so that it defines a bounded operator on $`^{}`$ which is given by (4.13d) as an immediate consequence of (4.24f), bearing in mind that $`\mathrm{X}`$ and $`\overline{\mathrm{X}}_5`$ only differ by a $`\overline{\nu }`$-null set. The demonstration of (4.13e) on the other hand is less straightforward. Assume first of all that the Borel set $`\mathrm{\Delta }`$ in question is open. Then we can make use of the regularity of spectral measures and construct, according to \[26, Definition II.8.2\], a sequence of compact subsets $`\left\{\mathrm{\Gamma }_N\right\}_N`$ as well as of infinitely often differentiable functions $`\left\{\stackrel{~}{g}_N\right\}_N`$ with support in $`\mathrm{\Delta }`$ such that $`0\chi _{\mathrm{\Gamma }_N}\stackrel{~}{g}_N\chi _\mathrm{\Delta }`$ and furthermore
$`L\left|E_\xi (\mathrm{\Delta })\right|L_\xi `$ $`=\underset{N\mathrm{}}{lim}L\left|E_\xi (\mathrm{\Gamma }_N)\right|L_\xi =\underset{N\mathrm{}}{lim}L\left|\stackrel{~}{g}_N(P_\xi )\right|L_\xi \text{,}`$ (4.33a)
$`L\left|E^{}(\mathrm{\Delta })\right|L`$ $`=\underset{N\mathrm{}}{lim}L\left|E^{}(\mathrm{\Gamma }_N)\right|L=\underset{N\mathrm{}}{lim}L\left|\stackrel{~}{g}_N(P^{})\right|L`$ (4.33b)
for any $`L𝔏^c`$. The discussion on page 3 f.—with $`\mathrm{X}`$ replacing $`\overline{\mathrm{X}}_5`$ and $`\nu `$ instead of $`\overline{\nu }`$—demonstrates, making use of Stone’s Theorem, that the sequence appearing on the right-hand side consists of $`\nu `$-measurable functions of $`\xi `$, so that we infer that its limit
$$\mathrm{X}\xi \mathrm{L}\left|\mathrm{E}_\xi (\mathrm{\Delta })\right|\mathrm{L}_\xi $$
is $`\nu `$-measurable, too. Another application of Stone’s Theorem in connection with (4.27e) in terms of $`\mathrm{X}`$ then shows that
$$\begin{array}{c}(2\pi )^{(s+1)/2}L\left|\stackrel{~}{g}_N(P^{})\right|L=_{^{s+1}}d^{s+1}xg_N(x)L\left|U^{}(x)\right|L\hfill \\ \hfill =_\mathrm{X}𝑑\overline{\nu }(\xi )_{^{s+1}}d^{s+1}xg_N(x)L\left|U_\xi (x)\right|L_\xi =(2\pi )^{(s+1)/2}_\mathrm{X}𝑑\overline{\nu }(\xi )L\left|\stackrel{~}{g}_N(P_\xi )\right|L_\xi \text{,}\end{array}$$
(4.33c)
and, as Lebesgue’s Dominated Convergence Theorem allows for the passage to the limit function under the last integral, we get according to (4.33a) and (4.33b)
$$L\left|E^{}(\mathrm{\Delta })\right|L=_\mathrm{X}𝑑\overline{\nu }(\xi )L\left|E_\xi (\mathrm{\Delta })\right|L_\xi \text{.}$$
(4.33d)
This formula, as yet valid only for open Borel sets $`\mathrm{\Delta }`$, is readily generalized to closed Borel sets and then, since by regularity the spectral measure of an arbitrary Borel set is approximated by a sequence in terms of compact subsets of it, to any Borel set. By polarization and the fact that the ket vectors with entries from $`𝔏^c`$ are dense in $`^{}`$ and $`_\xi `$, respectively, we first conclude with \[24, Section II.2.1, Proposition 1\] that all the fields $`\xi E_\xi (\mathrm{\Delta })`$ are measurable for arbitrary Borel sets $`\mathrm{\Delta }`$ and then pass to the aspired formula (4.13e) from (4.33d).
* The unitary operators $`V^{}(x)`$, $`x^{s+1}`$, defined in (4.10b) belong to the von Neumann algebra $`𝔐`$, according to (4.11), and are thus diagonalisable in the form
$$WV^{}(x)W^{}=X^{}𝑑\nu (\xi )\mathrm{exp}(ip_\xi x)\mathbf{\hspace{0.17em}1}_\xi \text{.}$$
(4.34a)
According to the construction of these operators, we can re-express this result in terms of the canonical unitary group of (4.10a):
$$WU_c^{}(x)W^{}=X^{}𝑑\nu (\xi )\mathrm{exp}(ip_\xi x)U_\xi (x)\text{.}$$
(4.34b)
The definition
$$U_\xi ^c(x)\mathrm{exp}(ip_\xi x)U_\xi (x)\text{,}x^{s+1}\text{,}\xi \mathrm{X}\text{,}$$
(4.35)
then provides the asserted canonical choice of a strongly continuous unitary group on each Hilbert space $`_\xi `$. Its spectral properties are derived from those of the canonical group $`\{U_c^{}(x):x^{s+1}\}`$ by the methods that have already been used repeatedly above. Possibly a further $`\nu `$-null subset of $`\mathrm{X}`$ gets lost by this procedure.
This finishes the proof of the assertions of Theorem 4.4. ∎
###### Remark.
Theorem 4.4 includes the existence of a spatial disintegration of the strongly continuous unitary group implementing space-time translations in the representation $`(\pi ^{},^{})`$ as well as of the spectral measure associated with it. The method used in the demonstration of this fact can be generalized to other symmetry groups; however obvious a problem of this kind may seem in the present context, it has, to our knowledge, not been treated in the literature. Nevertheless, the disintegration of unbounded closed operators in Hilbert spaces (the self-adjoint generators of strongly continuous unitary groups being an example) is the topic of and also presented in \[52, Chapter 12\].
At present we have no control over the range of energy-momenta $`p_\xi `$ which enter into the above disintegration theory. It still has to be investigated if, starting from a physical state of bounded energy in the constructions of Chapter 3 and passing to the asymptotic limit with respect to a function $`h`$ that has support on a small part of the velocity domain, the occurring momenta are correlated with those defined by the geometric momenta involved in the limiting procedure. Even tachyonic states cannot be ruled out to date. These problems might be tackled by introducing a certain property of ‘closability’ for particle weights, stating that, in case that a sequence of operators $`\{L_n\}_n`$ approaches $`0`$ in a suitable topology and at the same time $`\left\{|L_n\right\}_n`$ is convergent, the limit of the sequence of vectors likewise vanishes.
Moreover, the spatial disintegration presented above is subject to arbitrariness in two respects. There exist different constructions of the type expounded in Section 4.1 and therefore, according to Theorem 4.3, one has to deal with a number of different restricted $`\overline{𝔏_0^c}`$-particle weights $`(\pi ^{},^{},𝔄^{},\alpha ^{},𝖯^c,\overline{𝔏_0^c})`$ derived from the GNS-representation $`(\pi _w,_w)`$ pertaining to a given particle weight. As a result, the object to be disintegrated according to Theorem 4.4 is by no means uniquely fixed. And even if one has decided to select a system complying with the requirements of this theorem, there still remains an ambiguity as to the choice of a maximal abelian von Neumann algebra, with respect to which the disintegration is to be performed. The same problem is encountered in the framework of Choquet disintegration theory in its present status (cf. the end of Chapter 6). There a suitable base in the positive cone of particle weights has to be chosen with respect to which the disintegration is to be carried through.
But it should be stressed that these open questions only arise on the basis of the fact that a disintegration of general particle weights into pure ones, representing elementary systems, has successfully been constructed.
## Chapter 5 Phase Space Restrictions and Local Normality
A number of criteria have been introduced into the analysis of generic quantum field theories in order to dismiss those which are not reasonable form a physical point of view in that they do not allow for an interpretation in terms of particles. These attempts can be traced back to the year 1965 when Haag and Swieca proposed a compactness condition in , imposing an effective restriction to the size of phase space. Subsequently, the notion of nuclearity has entered the stage, determining maximum values for the number of local degrees of freedom for physical states of bounded energy (cf. the discussions of , and in addition for a treatment of the interdependence of these various concepts). In the present context we want to make use of the compactness condition proposed by Fredenhagen and Hertel to show that, under this physically motivated presupposition, the arbitrariness in the choice of a separable $`C^{}`$-subalgebra $`𝔄^{}`$ of the quasi-local algebra $`𝔄`$ in Chapter 4 can be removed.
###### Compactness Criterion (Fredenhagen–Hertel).
A local quantum field theory, as introduced in Chapter 1, qualifies the Fredenhagen–Hertel Compactness Condition if the mappings $`T_\mathrm{\Delta }^{}^𝒪`$, which are defined for any bounded Borel set $`\mathrm{\Delta }^{}^{s+1}`$ and any bounded region $`𝒪`$ of Minkowski space through
$$T_\mathrm{\Delta }^{}^𝒪:𝔄(𝒪)𝔅()AT_\mathrm{\Delta }^{}^𝒪(A)E(\mathrm{\Delta }^{})AE(\mathrm{\Delta }^{})\text{,}$$
send bounded subsets of $`𝔄(𝒪)`$ onto precompact subsets of $`𝔅()`$ with respect to its uniform topology. Precompactness is synonymous with totally boundedness and, in the present situation, equivalent to relative compactness \[44, Chapter One, § 4, 5.\].
To be able to demonstrate the main result of this chapter, Theorem 5.3, we have to fall back upon $`\mathrm{\Delta }`$-bounded particle weights as introduced in Definition 3.19. This restriction can be motivated on physical grounds, as opposed to mere technical needs, since, according to Lemma 3.20, the asymptotic functionals in $`_{\mathrm{\Delta }}^{}{}_{}{}^{+}`$, constructed by use of physical states of bounded energy in Chapter 3, give rise to particle weights of this special kind. The corresponding GNS-representations $`(\pi _w,_w)`$ then meet the Fredenhagen–Hertel Compactness Condition if the underlying local quantum field theory does, and the same holds true for the restricted $`\overline{𝔏_0^c}`$-particle weights which can be derived from them as expounded in Chapter 4.
###### Proposition 5.1.
Suppose that the given local quantum field theory satisfies the Compactness Criterion of Fredenhagen and Hertel.
* If $`.|.`$ is a $`\mathrm{\Delta }`$-bounded particle weight on $`𝔏\times 𝔏`$, then the associated GNS-representation $`(\pi _w,_w)`$ of the quasi-local algebra $`𝔄`$ is subject to the compactness condition as well.
* The restricted $`\overline{𝔏_0^c}`$-particle weight $`(\pi ^{},^{},𝔄^{},\alpha ^{},𝖯^c,\overline{𝔏_0^c})`$ derived from the above GNS-representation by virtue of Theorem 4.3 likewise inherits the compactness property in question.
###### Proof.
* $`\mathrm{\Delta }`$-boundedness of the particle weight $`.|.`$ means, according to Definition 3.19, that to any bounded Borel set $`\mathrm{\Delta }^{}^{s+1}`$ there exist another such set $`\overline{\mathrm{\Delta }}`$ containing $`\mathrm{\Delta }+\mathrm{\Delta }^{}`$ and an appropriate positive constant $`c`$, so that the estimate
$$E_w(\mathrm{\Delta }^{})\pi _w(A)E_w(\mathrm{\Delta }^{})cE(\overline{\mathrm{\Delta }})AE(\overline{\mathrm{\Delta }})$$
(3.46)
holds for any $`A𝔄`$. Then a finite cover of $`T_{\overline{\mathrm{\Delta }}}^𝒪\left(𝔄_r(𝒪)\right)=E(\overline{\mathrm{\Delta }})𝔄_r(𝒪)E(\overline{\mathrm{\Delta }})`$ by sets of diameter less than a given $`\delta >0`$ (which exists on account of the hypothesis of precompactness) induces a corresponding cover of $`E_w(\mathrm{\Delta }^{})\pi _w\left(𝔄_r(𝒪)\right)E_w(\mathrm{\Delta }^{})`$, which is composed of sets with a diameter smaller than $`c\delta `$ as (3.46) shows. This establishes totally boundedness of the set
$$E_w(\mathrm{\Delta }^{})\pi _w\left(𝔄_r(𝒪)\right)E_w(\mathrm{\Delta }^{})𝔅(_w)\text{.}$$
By arbitrariness of $`\mathrm{\Delta }^{}`$ as well as of the bounded region $`𝒪`$, the representation $`(\pi _w,_w)`$ is thus seen to satisfy the Compactness Criterion of Fredenhagen and Hertel in the sense that the mappings
$$T_{w,\mathrm{\Delta }^{}}^𝒪:𝔄(𝒪)𝔅(_w)AT_{w,\mathrm{\Delta }^{}}^𝒪(A)E_w(\mathrm{\Delta }^{})\pi _w(A)E_w(\mathrm{\Delta }^{})$$
are altogether precompact.
* According to the construction of $`(\pi ^{},^{},𝔄^{},\alpha ^{},𝖯^c,\overline{𝔏_0^c})`$ from $`(\pi _w,_w)`$ explained in the proof of Theorem 4.3, both of these representations are related by the inequality
$$E^{}(\mathrm{\Delta }^{})\pi ^{}(A)E^{}(\mathrm{\Delta }^{})E_w(\mathrm{\Delta }^{})\pi _w(A)E_w(\mathrm{\Delta }^{})\text{,}$$
(5.1a)
which holds for any $`A𝔄^{}`$. Then $`\mathrm{\Delta }`$-boundedness of the underlying particle weight again implies the existence of a bounded Borel set $`\overline{\mathrm{\Delta }}\mathrm{\Delta }+\mathrm{\Delta }^{}`$ such that
$$E^{}(\mathrm{\Delta }^{})\pi ^{}(A)E^{}(\mathrm{\Delta }^{})cE(\overline{\mathrm{\Delta }})AE(\overline{\mathrm{\Delta }})\text{.}$$
(5.1b)
This relation replaces (3.46) in the proof of the first part, so that we conclude that indeed $`(\pi ^{},^{})`$ inherits the precompactness properties of the underlying quantum field theory in the sense that all the sets
$$E^{}(\mathrm{\Delta }^{})\pi ^{}\left(𝔄_r^{}(𝒪_k)\right)E^{}(\mathrm{\Delta }^{})𝔅(^{})$$
are totally bounded for any $`r>0`$ whenever $`\mathrm{\Delta }^{}`$ is an arbitrary bounded Borel set and $`𝒪_k`$ is one of the countably many localization regions from $`^c`$. Again that is sufficient to establish the fact that the Fredenhagen–Hertel Compactness Condition is satisfied in the restricted setting for Local Quantum Physics introduced in Section 4.1. ∎
Under the presupposition of the Compactness Criterion, a result corresponding to Proposition 5.1 can be proved for the irreducible representations $`(\pi _\xi ,_\xi )`$ arising in the spatial disintegration of the restricted $`\overline{𝔏_0^c}`$-particle weight $`(\pi ^{},^{},𝔄^{},\alpha ^{},𝖯^c,\overline{𝔏_0^c})`$ by virtue of Theorem 4.4 if the domain of $`\xi `$ is further astricted to a $`\nu `$-measurable non-null subset $`\mathrm{X}_0`$ of $`\mathrm{X}`$.
###### Proposition 5.2.
Let $`(\pi _w,_w)`$ be the GNS-representation of the quasi-local algebra $`𝔄`$ corresponding to the $`\mathrm{\Delta }`$-bounded particle weight $`.|.`$, and let $`(\pi ^{},^{},𝔄^{},\alpha ^{},𝖯^c,\overline{𝔏_0^c})`$ denote the restricted $`\overline{𝔏_0^c}`$-particle weight associated with it according to Theorem 4.3. Under the hypothesis that the Compactness Criterion of Fredenhagen and Hertel is in force in the underlying quantum field theory, $`\nu `$-almost all of the irreducible representations $`(\pi _\xi ,_\xi )`$ occurring in the spatial disintegration (4.13a) of $`(\pi ^{},^{})`$ by course of Theorem 4.4 comply with this condition as well.
###### Proof.
Select a dense sequence $`\{A_k\}_k`$ in the norm-separable $`C^{}`$-algebra $`𝔄^{}`$ and consider the countable set of compact balls $`\mathrm{\Gamma }_N`$ of radius $`N`$ in $`^{s+1}`$. The corresponding operators $`E^{}(\mathrm{\Gamma }_N)\pi ^{}(A_k)E^{}(\mathrm{\Gamma }_N)𝔅(^{})`$ are decomposable according to Theorem 4.4:
$$WE^{}(\mathrm{\Gamma }_N)\pi ^{}(A_k)E^{}(\mathrm{\Gamma }_N)W^{}=X^{}𝑑\nu (\xi )E_\xi (\mathrm{\Gamma }_N)\pi _\xi (A_k)E_\xi (\mathrm{\Gamma }_N)\text{,}$$
(5.2a)
and \[24, Section II.2.3, Proposition 2\] tells us that the respective norms are related in compliance with the equation
$$E^{}(\mathrm{\Gamma }_N)\pi ^{}(A_k)E^{}(\mathrm{\Gamma }_N)=\nu \mathrm{ess}\mathrm{sup}\left\{E_\xi (\mathrm{\Gamma }_N)\pi _\xi (A_k)E_\xi (\mathrm{\Gamma }_N):\xi \mathrm{X}\right\}\text{.}$$
(5.2b)
With regard to all possible combinations of operators $`A_k`$ and compact balls $`\mathrm{\Gamma }_N`$ we thus infer that there exists a measurable non-null subset $`\mathrm{X}_0`$ of $`\mathrm{X}`$ such that for all $`\xi \mathrm{X}_0`$ and all indices $`k`$ and $`N`$ the estimate
$$E_\xi (\mathrm{\Gamma }_N)\pi _\xi (A_k)E_\xi (\mathrm{\Gamma }_N)E^{}(\mathrm{\Gamma }_N)\pi ^{}(A_k)E^{}(\mathrm{\Gamma }_N)$$
(5.3)
holds. Now, let $`\mathrm{\Delta }^{}`$ be an arbitrary bounded Borel set which is thus contained in a compact ball $`\mathrm{\Gamma }_{N_0}`$ and note that, by continuity of the representations $`\pi _\xi `$ and $`\pi ^{}`$, the inequality (5.3) extends to arbitrary operators $`A𝔄^{}`$. Therefore
$$E_\xi (\mathrm{\Delta }^{})\pi _\xi (A)E_\xi (\mathrm{\Delta }^{})E_\xi (\mathrm{\Gamma }_{N_0})\pi _\xi (A)E_\xi (\mathrm{\Gamma }_{N_0})E^{}(\mathrm{\Gamma }_{N_0})\pi ^{}(A)E^{}(\mathrm{\Gamma }_{N_0})\text{,}$$
(5.4a)
and this implies, according to (5.1b), the existence of a bounded Borel set $`\overline{\mathrm{\Delta }}\mathrm{\Delta }+\mathrm{\Delta }^{}`$ such that
$$E_\xi (\mathrm{\Delta }^{})\pi _\xi (A)E_\xi (\mathrm{\Delta }^{})cE(\overline{\mathrm{\Delta }})AE(\overline{\mathrm{\Delta }})\text{.}$$
(5.4b)
The arguments given in the proof of Proposition 5.1 can then again be applied to the present situation to show that for $`\xi \mathrm{X}_0`$ the irreducible representations $`(\pi _\xi ,_\xi )`$ altogether meet the requirements of the Fredenhagen–Hertel Compactness Condition. ∎
The central result of the present chapter is the perception that, under the above assumptions on the structure of phase space, the representations of the quasi-local $`C^{}`$-algebras $`𝔄`$ and $`𝔄^{}`$ which we have come across, i. e. $`(\pi _w,_w)`$ and $`(\pi ^{},^{})`$, respectively, as well as $`\nu `$-almost all of the irreducible representations $`(\pi _\xi ,_\xi )`$ occurring in the direct integral decomposition of the latter, are locally normal. The representations of $`𝔄^{}`$ can thus be continuously extended to all of $`𝔄`$ in such a way that the formula (4.13a) describing the disintegration stays valid for the extended representations when $`\mathrm{X}`$ is replaced by the non-null set $`\mathrm{X}_0`$ occurring in Proposition 5.2.
###### Theorem 5.3 (Local Normality).
Under the presumptions of Proposition 5.2 the following assertions are valid:
* The GNS-representation $`(\pi _w,_w)`$ of the quasi-local algebra $`𝔄`$ is locally normal, i. e. continuous with respect to the relative $`\sigma `$-weak topologies of both $`𝔄(𝒪)𝔅()`$ and $`\pi _w\left(𝔄(𝒪)\right)𝔅(_w)`$ for arbitrary bounded regions $`𝒪`$.
* The representation $`(\pi ^{},^{})`$ of the quasi-local algebra $`𝔄^{}`$ is locally normal (continuous with respect to the relative $`\sigma `$-weak topologies of both $`𝔄^{}(𝒪_k)𝔅()`$ and $`\pi ^{}\left(𝔄^{}(𝒪_k)\right)𝔅(^{})`$ for arbitrary bounded regions $`𝒪_k^c`$). The same holds true for the irreducible representations $`(\pi _\xi ,_\xi )`$ occurring in the spatial disintegration of $`(\pi ^{},^{})`$ when the indices $`\xi `$ are astricted to $`\mathrm{X}_0`$.
* The representation $`(\pi ^{},^{})`$ as well as the irreducible ones $`(\pi _\xi ,_\xi )`$ with $`\xi \mathrm{X}_0`$ allow for unique locally normal extensions to the whole of the original quasi-local algebra $`𝔄`$ designated $`(\overline{\pi }^{},^{})`$ and $`(\overline{\pi }_\xi ,_\xi )`$, respectively, which are related by
$$(\overline{\pi }^{},^{})_{\mathrm{X}_0}^{}𝑑\nu (\xi )(\overline{\pi }_\xi ,_\xi )\text{,}$$
(5.5)
where the representations $`(\overline{\pi }_\xi ,_\xi )`$ are again irreducible.
###### Proof.
* Let $`\overline{\mathrm{\Delta }}`$ be a bounded Borel set and suppose that $`\rho `$ is a normal functional on $`𝔅()`$. Then the same applies to the functional $`\rho _{\overline{\mathrm{\Delta }}}(.)\rho (E(\overline{\mathrm{\Delta }}).E(\overline{\mathrm{\Delta }}))`$, and therefore the mapping
$$T_{\overline{\mathrm{\Delta }}}:𝔄𝔅()AT_{\overline{\mathrm{\Delta }}}(A)E(\overline{\mathrm{\Delta }})AE(\overline{\mathrm{\Delta }})$$
is continuous with respect to the relative $`\sigma `$-weak topology of $`𝔄`$. Now, according to the Compactness Condition, $`T_{\overline{\mathrm{\Delta }}}𝔄(𝒪)=T_{\overline{\mathrm{\Delta }}}^𝒪`$ maps the unit ball $`𝔄_1(𝒪)`$ of the local $`C^{}`$-algebra $`𝔄(𝒪)`$ onto the relatively compact set $`E(\overline{\mathrm{\Delta }})𝔄_1(𝒪)E(\overline{\mathrm{\Delta }})`$. The restriction of $`T_{\overline{\mathrm{\Delta }}}^𝒪`$ to $`𝔄_1(𝒪)`$ is now obviously continuous with respect to the relative $`\sigma `$-weak topologies, but this result can be tightened up in the following sense: The relative $`\sigma `$-weak topology, being Hausdorff and coarser than the relative norm topology, and the relative uniform topology itself coincide on the compact norm closure of $`E(\overline{\mathrm{\Delta }})𝔄_1(𝒪)E(\overline{\mathrm{\Delta }})`$ due to a conclusion of general topology \[44, Chapter One, § 3, 2.(6)\]. Therefore $`T_{\overline{\mathrm{\Delta }}}^𝒪`$ is still continuous on $`𝔄_1(𝒪)`$ when its image is furnished with the norm topology instead.
Now, suppose that $`\mathrm{\Delta }^{}`$ is an arbitrary bounded Borel set and let $`\overline{\mathrm{\Delta }}\mathrm{\Delta }+\mathrm{\Delta }^{}`$ be another bounded Borel set with the property that (3.46) is satisfied. Then the linear mapping
$$E(\overline{\mathrm{\Delta }})AE(\overline{\mathrm{\Delta }})E_w(\mathrm{\Delta }^{})\pi _w(A)E_w(\mathrm{\Delta }^{})$$
(5.6)
is well-defined and continuous with respect to the uniform topologies of both domain and image. Therefore, as a consequence of the previous paragraph, we infer that the composition of this map with the restriction of $`T_{\overline{\mathrm{\Delta }}}`$ to $`𝔄_1(𝒪)`$ is continuous, when $`𝔄_1(𝒪)`$ is endowed with the $`\sigma `$-weak topology whereas the range carries the relative norm topology. The resulting map is explicitly determined as the restriction to $`𝔄_1(𝒪)`$ of
$$\pi _{w,\mathrm{\Delta }^{}}:𝔄𝔅(_w)A\pi _{w,\mathrm{\Delta }^{}}(A)E_w(\mathrm{\Delta }^{})\pi _w(A)E_w(\mathrm{\Delta }^{})\text{.}$$
(5.7)
If $`\eta `$ denotes a $`\sigma `$-weakly continuous functional on $`𝔅(_w)`$, the same is true regarding $`\eta _\mathrm{\Delta }^{}(.)\eta (E_w(\mathrm{\Delta }^{}).E_w(\mathrm{\Delta }^{}))`$ for any bounded Borel set $`\mathrm{\Delta }^{}^{s+1}`$, and moreover, due to strong continuity of the spectral measure, $`\eta `$ is the uniform limit of the net of functionals $`\eta _\mathrm{\Delta }^{}`$ for $`\mathrm{\Delta }^{}^{s+1}`$. Given a $`\sigma `$-weakly convergent net $`\{A_\iota :\iota J\}𝔄_1(𝒪)`$ with limit $`A𝔄_1(𝒪)`$, we conclude from the discussion in the preceding paragraph that
$$\eta _\mathrm{\Delta }^{}\left(\pi _w(A_\iota A)\right)=\eta \left(\pi _{w,\mathrm{\Delta }^{}}(A_\iota A)\right)\underset{\iota J}{\overset{}{}}0\text{.}$$
(5.8)
Therefore, by means of the estimate
$$\begin{array}{c}\left|\eta \pi _w(A_\iota A)\right|\left|\eta \left(\pi _w(A_\iota A)\right)\eta _\mathrm{\Delta }^{}\left(\pi _w(A_\iota A)\right)\right|+\left|\eta _\mathrm{\Delta }^{}\left(\pi _w(A_\iota A)\right)\right|\hfill \\ \hfill \eta \eta _\mathrm{\Delta }^{}\pi _w(A_\iota A)+\left|\eta _\mathrm{\Delta }^{}\left(\pi _w(A_\iota A)\right)\right|2\eta \eta _\mathrm{\Delta }^{}+\left|\eta _\mathrm{\Delta }^{}\left(\pi _w(A_\iota A)\right)\right|\text{,}\end{array}$$
(5.9)
it is easily seen that, upon selection of a suitable bounded Borel set $`\mathrm{\Delta }^{}`$, the right-hand side can be made smaller than any given bound for $`\iota \iota _0`$ with an appropriate index $`\iota _0`$. This being true for any $`\sigma `$-weakly continuous functional $`\eta `$ on $`𝔅(_w)`$ and arbitrary nets $`\{A_\iota :\iota J\}`$ in $`𝔄_1(𝒪)`$ converging to $`A𝔄_1(𝒪)`$ with respect to the $`\sigma `$-weak topology of $`𝔅()`$, we have thus established that the restrictions of the representation $`\pi _w`$ to each of the unit balls $`𝔄_1(𝒪)`$ are $`\sigma `$-weakly continuous. According to \[43, Lemma 10.1.10\] this assertion extends to the entire local $`C^{}`$-algebra $`𝔄(𝒪)`$, so that $`\pi _w`$ indeed turns out to be locally normal.
* The arguments given above in the case of $`\pi _w`$ can be transferred literally to the representations $`\pi ^{}`$ and $`\pi _\xi `$, $`\xi \mathrm{X}_0`$, in view of the relations (5.1b) and (5.4b) established in the proofs of Propositions 5.1 and 5.2, which substitute (3.46) used in the first part. The evident modifications to be applied include the restriction to local algebras $`𝔄^{}(𝒪_k)`$ where $`𝒪_k`$ is a member of the countable family $`^c`$.
* Complementary to the statements of the second part, \[43, Lemma 10.1.10\] exhibits that the representations $`\pi ^{}`$ and $`\pi _\xi `$, $`\xi \mathrm{X}_0`$, allow for unique $`\sigma `$-weakly continuous extensions $`\overline{\pi }^{}`$ and $`\overline{\pi }_\xi `$ onto the weak closures $`𝔄^{}(𝒪_k)^{\prime \prime }`$ \[11, Corollary 2.4.15\] of the local algebras, which, due to the Bicommutant Theorem \[11, Theorem 2.4.11\], coincide with the strong closures and thus, by the very construction of $`𝔄^{}(𝒪_k)`$, $`𝒪_k^c`$, expounded in Section 4.1, contain the corresponding local $`C^{}`$-algebras $`𝔄(𝒪_k)`$ of the underlying quantum field theory. Now, due to the net structure of $`𝒪_k𝔄(𝒪_k)`$, the quasi-local $`C^{}`$-algebra $`𝔄`$ is its $`C^{}`$-inductive limit, i. e. the norm closure of the -algebra $`_{𝒪_k^c}𝔄(𝒪_k)`$. As the representations $`\overline{\pi }^{}`$ and $`\overline{\pi }_\xi `$, $`\xi \mathrm{X}_0`$, are altogether uniformly continuous on this -algebra \[48, Theorem 1.5.7\], they can in a unique way be continuously extended to its completion $`𝔄`$ \[44, Chapter One, § 5, 4.(4)\], and these extensions, again denoted $`\overline{\pi }^{}`$ and $`\overline{\pi }_\xi `$, respectively, are easily seen to be compatible with the algebraic structure of $`𝔄`$. $`(\overline{\pi }^{},^{})`$ and $`(\overline{\pi }_\xi ,_\xi )`$ are thus representations of this quasi-local algebra, evidently irreducible in the case of $`\overline{\pi }_\xi `$, and moreover locally normal, since, by construction, they are $`\sigma `$-weakly continuous when restricted to local algebras $`𝔄(𝒪_k)`$ pertaining to the countable subclass of regions $`𝒪_k^c`$, and an arbitrary local algebra $`𝔄(𝒪)`$ is contained in at least one of these. The statement on uniqueness of the extensions is then an immediate consequence of the fact that they are uniquely determined by the property of being $`\sigma `$-weakly continuous on the local $`C^{}`$-algebras $`𝔄(𝒪_k)`$.
Regarding the disintegration of operators $`\overline{\pi }^{}(A)`$ for arbitrary $`A𝔄`$, note that any operator $`B𝔄(𝒪_k)`$ is the $`\sigma `$-weak limit of a *sequence* $`\{B_n\}_n`$ in $`𝔄_r^{}(𝒪_k)`$ with $`r=B`$. For nets in $`𝔄_r^{}(𝒪_k)`$ this statement is a consequence of Kaplansky’s Density Theorem \[54, Theorem II.4.8\] in connection with \[54, Lemma II.2.5\] and the various relations between the different locally convex topologies on $`𝔅()`$. The specialization to sequences is justified by \[54, Proposition II.2.7\] in view of the separability of $``$. The operators $`L𝔏^c`$ define fundamental sequences of measurable vector fields $`\left\{|L_\xi :\xi \mathrm{X}_0\right\}`$ (cf. \[24, Section II.1.3, Definition 1\]) and, as the operators $`\pi ^{}(B_n)`$ are decomposable, all the functions
$$h_n:\mathrm{X}_0\xi \mathrm{h}_\mathrm{n}(\xi )\mathrm{L}_1\left|\pi _\xi (\mathrm{B}_\mathrm{n})\right|\mathrm{L}_2_\xi $$
are measurable for arbitrary $`L_1\text{,}L_2𝔏^c`$. The same is valid for the pointwise limit of this sequence \[26, II.1.10\]
$$h:\mathrm{X}_0\xi \mathrm{h}(\xi )\mathrm{L}_1\left|\overline{\pi }_\xi (\mathrm{B})\right|\mathrm{L}_2_\xi \text{,}$$
and that suffices, according to \[24, Section II.2.1, Proposition 1\], to demonstrate that $`\{\overline{\pi }_\xi (B):\xi \mathrm{X}_0\}`$ is a measurable field of operators. As the sequence $`\left\{\pi ^{}(B_n)\right\}_n`$ converges $`\sigma `$-weakly to $`\overline{\pi }^{}(B)`$ by assumption and since, moreover, $`\nu (\mathrm{X}_0)`$ is finite and the family of operators $`\{\pi _\xi (B_n):\xi \mathrm{X}_0\}`$ is bounded by $`B`$ for any $`n`$, we conclude with Lebesgue’s Dominated Convergence Theorem applied to the decompositions of $`\pi ^{}(B_n)`$ with respect to $`\mathrm{X}_0`$ (which differs from $`\mathrm{X}`$ only by a null set), that
$$\begin{array}{c}L_1\left|\pi ^{}(B_n)\right|L_2=_{\mathrm{X}_0}𝑑\nu (\xi )L_1\left|\pi _\xi (B_n)\right|L_2_\xi \hfill \\ \hfill \underset{n\mathrm{}}{\overset{}{}}_{\mathrm{X}_0}𝑑\nu (\xi )L_1\left|\overline{\pi }_\xi (B)\right|L_2_\xi =L_1\left|\overline{\pi }^{}(B)\right|L_2\text{.}\end{array}$$
(5.10)
If $`W_0`$ denotes the isometry which implements the unitary equivalence
$$(\pi ^{},^{})_{\mathrm{X}_0}^{}𝑑\nu (\xi )(\pi _\xi ,_\xi )$$
and has all the properties of the operator $`W`$ introduced in Theorem 4.4, then, by density of the set $`\left\{|L^{}:L𝔏^c\right\}`$ in $`^{}`$, we infer from (5.10) that
$$W_0\overline{\pi }^{}(B)W_0^{}=_{\mathrm{X}_0}^{}𝑑\nu (\xi )\overline{\pi }_\xi (B)\text{.}$$
(5.11a)
This relation has been established under the presupposition that $`B`$ belongs to some local $`C^{}`$-algebra $`𝔄(𝒪_k)`$. Now, it is possible to reapply the above reasoning in the case of an arbitrary element $`A`$ of the quasi-local algebra $`𝔄`$, which can be approximated uniformly by a sequence $`\{A_n\}_n`$ from $`_{𝒪_k^c}𝔄(𝒪_k)`$. In this way, (5.11a) is extended to all of $`𝔄`$ so that we end up with the final equation
$$W_0\overline{\pi }^{}(A)W_0^{}=_{\mathrm{X}_0}^{}𝑑\nu (\xi )\overline{\pi }_\xi (A)\text{,}A𝔄\text{,}$$
(5.11b)
demonstrating that indeed
$$(\overline{\pi }^{},^{})_{\mathrm{X}_0}^{}𝑑\nu (\xi )(\overline{\pi }_\xi ,_\xi )\text{.}$$
Theorem 5.3 shows that, under the assumption of sensible phase space restrictions, no information on a physical system described by a normal state of bounded energy $`\omega 𝒮(\mathrm{\Delta })`$ gets lost in the entirety of constructions presented in Chapters 3 and 4. These have led us from $`\omega `$ via an associated particle weight with representation $`(\pi _w,_w)`$ of the quasi-local algebra $`𝔄`$ to the induced restricted $`\overline{𝔏_0^c}`$-particle weight $`(\pi ^{},^{},𝔄^{},\alpha ^{},𝖯^c,\overline{𝔏_0^c})`$, which comprises a representation $`(\pi ^{},^{})`$ of the algebra $`𝔄^{}`$ allowing for a disintegration in terms of irreducible representations $`\{(\pi _\xi ,_\xi ):\xi \mathrm{X}_0\}`$. Then, according to the preceding theorem, this disintegration is extendable in a unique fashion to one in terms of locally normal representations of the original algebra $`𝔄`$ as expressed by (5.5). Now, due to the explicit construction in Theorem 4.3 of $`(\pi ^{},^{})`$ from $`(\pi _w,_w)`$, the local normality of both of these representations implies that, actually, $`\overline{\pi }^{}`$ coincides with the restriction of $`\pi _w`$ to the subspace $`^{}`$ of $`_w`$. Thus we arrive at a partial reconstruction of the GNS-representation $`(\pi _w,_w)`$, which only depends on the initial choice of a subspace of the Hilbert space $`_w`$. Moreover, by Theorem 5.3, this entails a spatial disintegration of $`\mathrm{\Delta }`$-bounded particle weights $`.|.`$ according to the following reformulation of (5.5):
$$(\pi _w,^{})_{\mathrm{X}_0}^{}𝑑\nu (\xi )(\overline{\pi }_\xi ,_\xi )\text{.}$$
(5.12)
## Chapter 6 Disintegration Revisited: Choquet Theory
The spatial disintegration as expounded in Chapter 4 suffered from a couple of awkward drawbacks, which, in our belief, are inessential concomitants of this special approach to a decomposition theory for particle weights and have no bearing on the physical significance of the concept proper. It should be noted in this connection that, to be able to apply the standard disintegration theory for representations made available in the literature on $`C^{}`$-algebras, we had to fall back upon separable constructs and countable dense subsets thereof. As a consequence it had to be accepted, that a theory of disintegration could only be formulated in terms of restricted $`\overline{𝔏_0^c}`$-particle weights. But these technical difficulties seem to be accidental, and the question obtrudes on us if it is possible to carry through a disintegration, in the course of which no need arises to leave the class of particle weights proper, which means that the disintegration can indeed be formulated in terms of *pure* particle weights. This is the topic of the present chapter, presenting the partial results we were able to produce in this direction to date.
As already noted in the remark following Definition 3.14, the totality of particle weights constitutes a positive proper convex cone when supplemented by the trivial form. This observation opens the way to an application of another concept of disintegration: the barycentric decomposition in the special form of a generalization of the well-known Theorem of Krein–Milman \[42, Theorem 1.4.3\]. This approach, initiated by Choquet and further developed by Bishop and de Leeuw , is especially well-adapted to the study of convex sets in infinite-dimensional spaces. An introduction to this theory can be found in and also in \[11, Section 4.1.2\] where it is applied to get a decomposition of states on the quasi-local $`C^{}`$-algebra $`𝔄`$ in terms of pure ones. The mathematical structure in this case is easily accessible from the point of view of Choquet theory:
* The positive linear functionals on $`𝔄`$ constitute a positive convex cone $`\mathrm{K}`$ in its topological dual $`𝔄^{}`$.
* The quasi-local algebra contains a unit $`\mathrm{𝟏}`$, which defines a base of this cone when it is considered as a continuous linear functional on $`𝔄^{}`$, thus introducing a convex function which intersects $`\mathrm{K}`$.
* This convex base of $`\mathrm{K}`$ coincides with the set of states on $`𝔄`$ and thus turns out to be compact with respect to the weak topology \[11, Propositon 2.3.11 and Theorem 2.3.15\].
The situation is much more complicated when particle weights are considered. These do constitute a positive cone $`\mathrm{W}`$ in the space of sesquilinear forms on the left ideal $`𝔏`$ of localizing operators; but one of the key features in the construction of an algebra of detectors was the absence of a unit element, the existence of which would produce infinite values of the asymptotic functionals $`\sigma `$ arising in Chapter 3. The obvious questions to be answered are:
* What is an appropriate (metrizable) topology $`𝒯`$ to be introduced on $`\mathrm{W}`$ to render relevant subsets compact?
* How can a convex base $`\mathrm{BW}`$ be fitted into the cone $`\mathrm{W}`$ in a physically meaningful way and such that this base is compact with respect to the beforementioned topology?
In our approach the class with respect to which the disintegration is to be performed will be further restricted (cf. Definition 6.6 below).
The answer to the first of the above problems which we are going to present in this chapter is based on an effective control of the dislocalization of almost local operators combined with the Fredenhagen–Hertel compactness condition. To be more specific, a norm will be introduced on the space $`𝔏_0`$ of almost local vacuum annihilation operators which in a way measures their deviation from being contained in a local algebra. Making use of this norm on the set $`𝔏_0(\mathrm{\Gamma })`$ of those operators with energy-momentum transfer in a compact and convex subset $`\mathrm{\Gamma }`$ of $`\mathrm{}\overline{V}_+`$, it can be shown that $`𝔏_0(\mathrm{\Gamma })`$ is small in the $`q_\mathrm{\Delta }`$-topology under the assumption of the Compactness Criterion of Fredenhagen and Hertel. By its very definition, the notion of almost locality as introduced in Definition 2.1 imposes a condition of rapid decrease on the norm difference between almost local operators and strictly local ones according to the growing extension of the localization regions. This contrives to introduce the following norms on $`𝔏_0`$.
###### Definition 6.1.
Let $`m`$ be an arbitrary natural number, then the equation
$$𝒬^m(L_0)\underset{r>0}{sup}inf\{r^mL_0L_r:L_r𝔄(𝒪_r)\}\text{,}L_0𝔏_0\text{,}$$
(6.1)
defines a norm on the vector space $`𝔏_0`$.
###### Remark.
The seminorm properties of the mapping $`𝒬^m`$ are self-evident, they even hold for arbitrary almost local operators as arguments. To infer $`L_0=0`$ from the condition $`𝒬^m(L_0)=0`$ one has to restrict attention to vacuum annihilation operators. The latter equation means that the operator $`L_0`$ is contained in the norm closure of any local algebra $`𝔄(𝒪_r)`$, $`r>0`$. As these are $`C^{}`$-algebras, hence uniformly closed, $`L_0`$ itself turns out to be a local operator, a property which can be reconciled with it being a vacuum annihilation operator only for $`L_0=0`$.
The information concerning the localization of the operator $`L_0𝔏_0`$ embodied by the value of $`𝒬^m(L_0)`$ is highly dependent on the norm which $`L_0`$ carries as an element of the $`C^{}`$-algebra $`𝔄`$. Therefore we combine both topologies in the subsequent definition.
###### Definition 6.2.
For any natural number $`m`$ a norm on $`𝔏_0`$ is defined by
$$|||L_0|||_mL_0+𝒬^m(L_0)\text{,}L_0𝔏_0\text{.}$$
(6.2)
As announced above this topology is now to be related to the $`q_\mathrm{\Delta }`$-seminorms on the subspace $`𝔏_0(\mathrm{\Gamma })`$, where $`\mathrm{\Gamma }`$ denotes a compact, convex subset of the complement of the forward light cone. Although we have the inequality (2.10) at our disposition, we want to reformulate it here in order to make explicit the dependence of the integrand on its right-hand side upon the bounded Borel set $`\mathrm{\Delta }`$ and upon the energy-momentum transfer $`\mathrm{\Gamma }`$. To this end one has to reapply the arguments given in the Appendix of .
###### Proposition 6.3.
Let $`\mathrm{\Delta }`$ be a bounded Borel set and $`\mathrm{\Gamma }`$ a compact and convex subset of $`\mathrm{}\overline{V}_+`$. There exists a bounded Borel set $`\mathrm{\Delta }^{}(\mathrm{\Delta },\mathrm{\Gamma })^{s+1}`$, depending on $`\mathrm{\Delta }`$ and $`\mathrm{\Gamma }`$ only, such that for any $`L_0𝔏_0(\mathrm{\Gamma })`$ there holds the estimate
$$q_\mathrm{\Delta }(L_0)^2N^{}(\mathrm{\Delta },\mathrm{\Gamma })_^sd^sxE\left(\mathrm{\Delta }^{}(\mathrm{\Delta },\mathrm{\Gamma })\right)[\alpha x(L_0),L_{0}^{}{}_{}{}^{}]E\left(\mathrm{\Delta }^{}(\mathrm{\Delta },\mathrm{\Gamma })\right)$$
(6.3)
with a suitable constant $`N^{}(\mathrm{\Delta },\mathrm{\Gamma })`$, which is again specified by the sets $`\mathrm{\Delta }`$ and $`\mathrm{\Gamma }`$.
###### Proof.
In a first step it will be shown that, setting
$$QKKd^sx\alpha x(L_{0}^{}{}_{}{}^{}L_0)$$
for any compact subset $`𝑲`$ of $`^s`$, the following estimate is in force for arbitrary bounded Borel sets $`\mathrm{\Delta }_0`$:
$$E(\mathrm{\Delta }_0)QKE(\mathrm{\Delta }_0)N^{\prime \prime }_{𝑲𝑲}d^sxE(\mathrm{\Delta }^{\prime \prime })[\alpha x(L_0),L_{0}^{}{}_{}{}^{}]E(\mathrm{\Delta }^{\prime \prime })$$
(6.4)
with a suitable constant $`N^{\prime \prime }`$ and an appropriate bounded Borel set $`\mathrm{\Delta }^{\prime \prime }`$. If $`\omega _\mathrm{\Psi }`$ denotes a state on $`𝔅()`$ which is induced by a vector $`\mathrm{\Psi }E(\mathrm{\Delta }_0)`$ we can immediately adopt the inequalities of \[15, p. 640\] to get
$$\begin{array}{c}\omega _\mathrm{\Psi }(QK)^2\omega _\mathrm{\Psi }(QKQK)\hfill \\ \hfill \omega _\mathrm{\Psi }(QK)\underset{𝒚^𝒔}{sup}_{𝑲𝑲}d^sxL_0U(𝒚)𝜳^\mathit{1}[\alpha 𝒙(𝑳_\mathit{0}),𝑳_{\mathit{0}}^{}{}_{}{}^{}]𝑳_\mathit{0}𝑼(𝒚)𝜳\\ \hfill +Kd^sx\omega _\mathrm{\Psi }\left(\alpha x(L_{0}^{}{}_{}{}^{})QK\alpha x(L_0)\right)\text{.}\end{array}$$
(6.5a)
The integrand of the second term on the right-hand side is subject to the relation
$$\omega _\mathrm{\Psi }\left(\alpha x(L_{0}^{}{}_{}{}^{})QK\alpha x(L_0)\right)E(\overline{\mathrm{\Delta }}_0+\mathrm{\Gamma })QKE(\overline{\mathrm{\Delta }}_0+\mathrm{\Gamma })\omega _\mathrm{\Psi }\left(\alpha x(L_{0}^{}{}_{}{}^{}L_0)\right)$$
(6.5b)
with $`\overline{\mathrm{\Delta }}_0`$ denoting the closure of $`\mathrm{\Delta }_0`$. Upon insertion into (6.5a), removal of the resulting common factor $`\omega _\mathrm{\Psi }(QK)`$ on both sides and passing to the supremum with respect to all unit vectors $`\mathrm{\Psi }E(\mathrm{\Delta }_0)`$, we get
$$E(\mathrm{\Delta }_0)QKE(\mathrm{\Delta }_0)_{𝑲𝑲}d^sx[\alpha x(L_0),L_{0}^{}{}_{}{}^{}]E(\overline{\mathrm{\Delta }}_0+\mathrm{\Gamma })+E(\overline{\mathrm{\Delta }}_0+\mathrm{\Gamma })QKE(\overline{\mathrm{\Delta }}_0+\mathrm{\Gamma })\text{,}$$
(6.5c)
where use is made of the fact that all the vectors $`L_0U(𝒚)𝜳`$ belong to the subspace $`E(\overline{\mathrm{\Delta }}_0+\mathrm{\Gamma })`$ for arbitrary $`𝒚^𝒔`$ and $`\mathrm{\Psi }E(\mathrm{\Delta }_0)`$. The preparatory estimate (6.4) is now established by complete induction on $`n`$, where this natural number is defined in dependence on the sets $`\mathrm{\Delta }_0`$ and $`\mathrm{\Gamma }`$ through the condition $`(\overline{\mathrm{\Delta }}_0+\mathrm{\Gamma }_n)\overline{V}_+=\mathrm{}`$ (cf. the proof of Proposition 2.6 on page 2.2).
For $`n=1`$ we have, according to the spectrum condition, $`E(\overline{\mathrm{\Delta }}_0+\mathrm{\Gamma })=0`$ so that (6.4) is trivially fulfilled since its left-hand side vanishes. Now assume that the condition $`(\overline{\mathrm{\Delta }}_0+\mathrm{\Gamma }_{n+1})\overline{V}_+=\mathrm{}`$ is valid, which, stated another way, means that the intersection of $`(\overline{\mathrm{\Delta }}_0+\mathrm{\Gamma })+\mathrm{\Gamma }_n`$ with the complement of $`\overline{V}_+`$ is empty. As $`\overline{\mathrm{\Delta }}_0+\mathrm{\Gamma }`$ is a bounded Borel set we can apply the induction hypothesis for $`n`$, i. e. (6.4) with $`\mathrm{\Delta }_0`$ replaced by $`\overline{\mathrm{\Delta }}_0+\mathrm{\Gamma }`$, to infer that there exists a bounded Borel set $`\mathrm{\Delta }_0^{\prime \prime }`$ which satisfies
$$E(\overline{\mathrm{\Delta }}_0+\mathrm{\Gamma })QKE(\overline{\mathrm{\Delta }}_0+\mathrm{\Gamma })N_0^{\prime \prime }_{𝑲𝑲}d^sxE(\mathrm{\Delta }_0^{\prime \prime })[\alpha x(L_0),L_{0}^{}{}_{}{}^{}]E(\mathrm{\Delta }_0^{\prime \prime })$$
(6.6a)
for an appropriate constant $`N_0^{\prime \prime }`$. This estimate inserted into (6.5c) amounts to
$$\begin{array}{c}E(\mathrm{\Delta }_0)QKE(\mathrm{\Delta }_0)\hfill \\ \hfill _{𝑲𝑲}d^sx[\alpha x(L_0),L_{0}^{}{}_{}{}^{}]E(\overline{\mathrm{\Delta }}_0+\mathrm{\Gamma })+N_0^{\prime \prime }_{𝑲𝑲}d^sxE(\mathrm{\Delta }_0^{\prime \prime })[\alpha x(L_0),L_{0}^{}{}_{}{}^{}]E(\mathrm{\Delta }_0^{\prime \prime })\text{,}\end{array}$$
(6.6b)
from which to conclude the validity of (6.4) with suitable constant $`N^{\prime \prime }=N_0^{\prime \prime }+1`$ and proper bounded Borel set $`\mathrm{\Delta }^{\prime \prime }=\mathrm{\Delta }_0^{\prime \prime }\left((\overline{\mathrm{\Delta }}_0+\mathrm{\Gamma })+\mathrm{\Gamma }\mathrm{\Gamma }\right)`$ is an obvious task.
Now, having established (6.4), we can specialize it to $`\mathrm{\Delta }_0\mathrm{\Delta }`$ and pass to the limit $`𝑲^𝒔`$ as in the proof of Proposition 2.6, noting that
$$Q_{\mathrm{\Delta },𝑲}^{(L_{0}^{}{}_{}{}^{}L_0)}=E(\mathrm{\Delta })QKE(\mathrm{\Delta })$$
and that, due to almost locality of $`L_0`$, the integral on the right-hand side can be extended over all of $`^s`$. As a result, in view of Definition 2.8, one arrives at the desired inequality (6.3), where the formulation chosen makes explicit its dependence on $`\mathrm{\Delta }`$ and $`\mathrm{\Gamma }`$. ∎
The formula (6.3) just established is to be applied in the sequel to produce an estimate of the seminorm $`q_\mathrm{\Delta }(L_0)`$ for operators $`L_0𝔏_0(\mathrm{\Gamma })`$ with compact and convex $`\mathrm{\Gamma }\mathrm{}\overline{V}_+`$ in terms of the initial operator norm $`.`$ and of the norm $`𝒬^m(.)`$ introduced in Definition 6.1. In order to get a manageable result we specialize to the case $`m=2s`$.
###### Lemma 6.4.
Let $`\mathrm{\Delta }`$ be a bounded Borel set and $`\mathrm{\Gamma }`$ a compact and convex subset of the complement of $`\overline{V}_+`$. Then there exists a bounded Borel set $`\mathrm{\Delta }^{\prime \prime }(\mathrm{\Delta },\mathrm{\Gamma })`$, depending on $`\mathrm{\Delta }`$ and $`\mathrm{\Gamma }`$, such that for any $`L_0𝔏_0(\mathrm{\Gamma })`$ the estimate
$$\begin{array}{c}q_\mathrm{\Delta }(L_0)N^{}(\mathrm{\Delta },\mathrm{\Gamma })^{1/2}(a(s)L_0^2+b(s)L_0+c(s))^{1/2}\hfill \\ \hfill 𝒬^{2s}(L_0)^{1/4}E\left(\mathrm{\Delta }^{\prime \prime }(\mathrm{\Delta },\mathrm{\Gamma })\right)L_0E\left(\mathrm{\Delta }^{\prime \prime }(\mathrm{\Delta },\mathrm{\Gamma })\right)^{1/2}\end{array}$$
(6.7)
holds with suitable coefficients depending on the spatial dimension $`s`$.
###### Proof.
We have to calculate the integral on the right-hand side of (6.3) and, to do so, it is split into two parts according to $`|𝒙|>𝑹`$ or $`|𝒙|𝑹`$ with an abitrary radius $`R`$ which is held fixed for the moment. For large $`|𝒙|`$ we use the estimate (2.2a) for the integrand and get in terms of the norm $`𝒬^{2s}`$:
$$\begin{array}{c}E\left(\mathrm{\Delta }^{}(\mathrm{\Delta },\mathrm{\Gamma })\right)[\alpha x(L_0),L_{0}^{}{}_{}{}^{}]E\left(\mathrm{\Delta }^{}(\mathrm{\Delta },\mathrm{\Gamma })\right)\hfill \\ \hfill [\alpha x(L_0),L_{0}^{}{}_{}{}^{}]4L_0L_0(L_0)_{2^1|𝒙|}+2L_0(L_0)_{2^1|𝒙|}^2\\ \hfill 4L_0\mathrm{\hspace{0.17em}2}^{2s}|𝒙|^{\mathit{2}𝒔}𝒬^{\mathit{2}𝒔}(𝑳_\mathit{0})+\mathit{2\hspace{0.17em}2}^{\mathit{4}𝒔}|𝒙|^{\mathit{4}𝒔}𝒬^{\mathit{2}𝒔}(𝑳_\mathit{0})^\mathit{2}\text{.}\end{array}$$
(6.8a)
Accordingly, the respective integral is subject to the inequality
$$\begin{array}{c}_{|𝒙|>𝑹}d^sxE\left(\mathrm{\Delta }^{}(\mathrm{\Delta },\mathrm{\Gamma })\right)[\alpha x(L_0),L_{0}^{}{}_{}{}^{}]E\left(\mathrm{\Delta }^{}(\mathrm{\Delta },\mathrm{\Gamma })\right)\hfill \\ \hfill 2^{2s+2}L_0𝒬^{2s}(L_0)_{|𝒙|>𝑹}d^sx|𝒙|^{\mathit{2}𝒔}+\mathit{2}^{\mathit{4}𝒔+\mathit{1}}𝒬^{\mathit{2}𝒔}(𝑳_\mathit{0})^\mathit{2}_{|𝒙|>𝑹}𝒅^𝒔𝒙|𝒙|^{\mathit{4}𝒔}\text{.}\end{array}$$
(6.8b)
The integrand for small $`|𝒙|`$ is evaluated observing the spectral projections arising on the right-hand side of (6.3), which are abbreviated as $`\mathrm{\Delta }^{}\mathrm{\Delta }^{}(\mathrm{\Delta },\mathrm{\Gamma })`$ with closure $`\overline{\mathrm{\Delta }^{}}`$. This leads to
$$\begin{array}{c}E\left(\mathrm{\Delta }^{}(\mathrm{\Delta },\mathrm{\Gamma })\right)[\alpha x(L_0),L_{0}^{}{}_{}{}^{}]E\left(\mathrm{\Delta }^{}(\mathrm{\Delta },\mathrm{\Gamma })\right)\hfill \\ \hfill E(\mathrm{\Delta }^{})\alpha x(L_0)E(\overline{\mathrm{\Delta }^{}}\mathrm{\Gamma })L_{0}^{}{}_{}{}^{}E(\mathrm{\Delta }^{})+E(\mathrm{\Delta }^{})L_{0}^{}{}_{}{}^{}E(\overline{\mathrm{\Delta }^{}}+\mathrm{\Gamma })\alpha x(L_0)E(\mathrm{\Delta }^{})\\ \hfill E\left(\mathrm{\Delta }^{}(\overline{\mathrm{\Delta }^{}}\mathrm{\Gamma })\right)L_0E\left(\mathrm{\Delta }^{}(\overline{\mathrm{\Delta }^{}}\mathrm{\Gamma })\right)^2+E\left(\mathrm{\Delta }^{}(\overline{\mathrm{\Delta }^{}}+\mathrm{\Gamma })\right)L_0E\left(\mathrm{\Delta }^{}(\overline{\mathrm{\Delta }^{}}+\mathrm{\Gamma })\right)^2\\ \hfill 2E\left(\mathrm{\Delta }^{\prime \prime }(\mathrm{\Delta },\mathrm{\Gamma })\right)L_0E\left(\mathrm{\Delta }^{\prime \prime }(\mathrm{\Delta },\mathrm{\Gamma })\right)^2\end{array}$$
(6.9a)
where $`\mathrm{\Delta }^{\prime \prime }(\mathrm{\Delta },\mathrm{\Gamma })\mathrm{\Delta }(\overline{\mathrm{\Delta }}+\mathrm{\Gamma })\left(\mathrm{\Delta }^{}(\overline{\mathrm{\Delta }^{}}\mathrm{\Gamma })\right)\left(\mathrm{\Delta }^{}(\overline{\mathrm{\Delta }^{}}+\mathrm{\Gamma })\right)`$—the inclusion of $`\mathrm{\Delta }(\overline{\mathrm{\Delta }}+\mathrm{\Gamma })`$ into this definition being required at the very end of the present argumentation. The corresponding integral satisfies the inequality
$$\begin{array}{c}_{|𝒙|𝑹}d^sxE\left(\mathrm{\Delta }^{}(\mathrm{\Delta },\mathrm{\Gamma })\right)[\alpha x(L_0),L_{0}^{}{}_{}{}^{}]E\left(\mathrm{\Delta }^{}(\mathrm{\Delta },\mathrm{\Gamma })\right)\hfill \\ \hfill 2E\left(\mathrm{\Delta }^{\prime \prime }(\mathrm{\Delta },\mathrm{\Gamma })\right)L_0E\left(\mathrm{\Delta }^{\prime \prime }(\mathrm{\Delta },\mathrm{\Gamma })\right)^2_{|𝒙|𝑹}d^sx\text{.}\end{array}$$
(6.9b)
The integrals remaining in (6.8b) and (6.9b) are known from calculus (cf. \[22, Section 4.11\]):
$`{\displaystyle _{|𝒙|>𝑹}}d^sx|𝒙|^{\mathit{2}𝒔}`$ $`=\omega _s{\displaystyle _R^{\mathrm{}}}𝑑rr^{s1}r^{2s}=s^1\omega _sR^s\text{,}`$ (6.10a)
$`{\displaystyle _{|𝒙|>𝑹}}d^sx|𝒙|^{\mathit{4}𝒔}`$ $`=\omega _s{\displaystyle _R^{\mathrm{}}}𝑑rr^{s1}r^{4s}=(3s)^1\omega _sR^{3s}\text{,}`$ (6.10b)
$`{\displaystyle _{|𝒙|𝑹}}d^sx`$ $`=\omega _s{\displaystyle _0^R}𝑑rr^{s1}=s^1\omega _sR^s\text{,}`$ (6.10c)
where the factor $`\omega _s`$ is defined via the $`\mathrm{\Gamma }`$-function as
$$\omega _s2\mathrm{\Gamma }(s/2)^1\sqrt{\pi ^s}\text{.}$$
(6.10d)
Collecting the results from (6.8b), (6.9b) and (6.10) one gets for the complete integral
$$\begin{array}{c}_^sd^sxE\left(\mathrm{\Delta }^{}(\mathrm{\Delta },\mathrm{\Gamma })\right)[\alpha x(L_0),L_{0}^{}{}_{}{}^{}]E\left(\mathrm{\Delta }^{}(\mathrm{\Delta },\mathrm{\Gamma })\right)\hfill \\ \hfill \omega _sR^s(a^{}(s)𝒬^{2s}(L_0)^2R^{4s}+b^{}(s)L_0𝒬^{2s}(L_0)R^{2s}\\ \hfill +c^{}(s)E\left(\mathrm{\Delta }^{\prime \prime }(\mathrm{\Delta },\mathrm{\Gamma })\right)L_0E\left(\mathrm{\Delta }^{\prime \prime }(\mathrm{\Delta },\mathrm{\Gamma })\right)^2)\end{array}$$
(6.11)
with suitable $`s`$-dependent factors. So far the value of $`R`$ has been left open. To get the concise formula (6.7) we deliberately choose
$$R^{2s}E\left(\mathrm{\Delta }^{\prime \prime }(\mathrm{\Delta },\mathrm{\Gamma })\right)L_0E\left(\mathrm{\Delta }^{\prime \prime }(\mathrm{\Delta },\mathrm{\Gamma })\right)^2𝒬^{2s}(L_0)\text{,}$$
(6.12)
so that (6.11) simplifies to
$$\begin{array}{c}_^sd^sxE\left(\mathrm{\Delta }^{}(\mathrm{\Delta },\mathrm{\Gamma })\right)[\alpha x(L_0),L_{0}^{}{}_{}{}^{}]E\left(\mathrm{\Delta }^{}(\mathrm{\Delta },\mathrm{\Gamma })\right)\hfill \\ \hfill \omega _sR^sE\left(\mathrm{\Delta }^{\prime \prime }(\mathrm{\Delta },\mathrm{\Gamma })\right)L_0E\left(\mathrm{\Delta }^{\prime \prime }(\mathrm{\Delta },\mathrm{\Gamma })\right)^2\left(a^{}(s)L_0^2+b^{}(s)L_0+c^{}(s)\right)\text{.}\end{array}$$
(6.13)
Inserting the square root of (6.12) into this estimate and carrying the result over to (6.3), we finally arrive at (6.7), where $`\omega _s`$ has been included in the definition of the coefficients.
Note, that the above argument is independent of the occurrence of $`𝒬^{2s}(L_0)=0`$ or $`E\left(\mathrm{\Delta }^{\prime \prime }(\mathrm{\Delta },\mathrm{\Gamma })\right)L_0E\left(\mathrm{\Delta }^{\prime \prime }(\mathrm{\Delta },\mathrm{\Gamma })\right)=0`$, for in this case $`q_\mathrm{\Delta }(L_0)=0`$, so that (6.7) is trivially fulfilled. This consequence is immediate from the norm property of $`𝒬^{2s}`$. As to the second of the above conditions, it turns out to be important that we have included $`\mathrm{\Delta }(\overline{\mathrm{\Delta }}+\mathrm{\Gamma })`$ into the definition of $`\mathrm{\Delta }^{\prime \prime }(\mathrm{\Delta },\mathrm{\Gamma })`$. Thence the named assumption implies $`E(\overline{\mathrm{\Delta }}+\mathrm{\Gamma })L_0E(\mathrm{\Delta })=L_0E(\mathrm{\Delta })=0`$, and $`q_\mathrm{\Delta }(L_0)=0`$ is a result of Lemma 2.11. ∎
Our next aim is to single out a convex subset in the class of all particle weights which turns out to be compact in a suitably chosen topology. To this end it will be assumed from now on that the underlying quantum field theory satisfies the Fredenhagen–Hertel Compactness Criterion, under the assumption of which the following result can be established.
###### Proposition 6.5.
In a quantum field theory which satisfies the Fredenhagen–Hertel Compactness Condition the subsequent mapping, defined for bounded Borel subsets $`\mathrm{\Delta }`$ and compact, convex subsets $`\mathrm{\Gamma }`$ of $`^{s+1}`$,
$$S_\mathrm{\Delta }^\mathrm{\Gamma }:𝔏_0(\mathrm{\Gamma })𝔅()L_0S_\mathrm{\Delta }^\mathrm{\Gamma }(L_0)L_0E(\mathrm{\Delta })\text{,}$$
sends balls in $`𝔏_0(\mathrm{\Gamma })`$ of finite radius with respect to the norm $`|||.|||_m`$, $`m`$, onto precompact subsets of $`𝔅()`$ in its uniform topology.
###### Proof.
Let $`𝔏_{0,R}^m(\mathrm{\Gamma })`$ denote the closed $`R`$-ball, $`R>0`$, in $`𝔏_0(\mathrm{\Gamma })`$ with respect to $`|||.|||_m`$. By Definition 6.2, the condition $`|||L_0|||_mR`$, $`L_0𝔏_0(\mathrm{\Gamma })`$, implies $`𝒬^m(L_0)<R`$, stating a property of uniform approximation. This means that, given $`\epsilon >0`$, there exists a radius $`r_0`$, take e. g. $`r_0(2R/\epsilon )^{m/2}`$, such that to any $`L_0𝔏_{0,R}^m(\mathrm{\Gamma })`$ we can find a local operator $`(L_0)_{r_0}𝔄(𝒪_{r_0})`$ with
$$L_0(L_0)_{r_0}<\epsilon \text{.}$$
(6.14a)
Again according to Definition 6.2, one also has $`L_0<R`$, so that the collection of local operators just introduced belongs to the closed ball of radius $`R+\epsilon `$ in $`𝔄(𝒪_{r_0})`$. Now, the Fredenhagen–Hertel Compactness Condition ensues that there exists a finite number of operators $`L_k`$, $`k=1\text{,}\mathrm{}\text{,}N(\epsilon )`$, in this ball such that any $`(L_0)_{r_0}`$ satisfies the condition
$$E(\overline{\mathrm{\Delta }}+\mathrm{\Gamma })\left((L_0)_{r_0}L_k\right)E(\mathrm{\Delta })<\epsilon $$
(6.14b)
for at least one $`k`$. Combining this with (6.14a), we see that for any $`L_0𝔏_{0,R}^m(\mathrm{\Gamma })`$ there exists a suitable operator $`L_k𝔄(𝒪_{r_0})`$ with
$$E(\overline{\mathrm{\Delta }}+\mathrm{\Gamma })(L_0L_k)E(\mathrm{\Delta })L_0(L_0)_{r_0}+E(\overline{\mathrm{\Delta }}+\mathrm{\Gamma })\left((L_0)_{r_0}L_k\right)E(\mathrm{\Delta })<2\epsilon \text{.}$$
(6.14c)
It is an immediate consequence that finitely many elements from $`𝔏_{0,R}^m(\mathrm{\Gamma })`$ can be selected, serving as centres of $`4\epsilon `$-balls which cover the set
$$E(\overline{\mathrm{\Delta }}+\mathrm{\Gamma })𝔏_{0,R}^m(\mathrm{\Gamma })E(\mathrm{\Delta })=𝔏_{0,R}^m(\mathrm{\Gamma })E(\mathrm{\Delta })\text{.}$$
By arbitrariness of $`\epsilon `$, we have thus established precompactness of the mapping $`S_\mathrm{\Delta }^\mathrm{\Gamma }`$ in the sense of the Proposition. ∎
The results presented thus far have only laid down the groundwork for the topological considerations concerning the set of particle weights proper. For the moment we return here to the special continuity properties of the asymptotic functionals resulting from the limiting procedure expounded in Chapter 3. According to Proposition 3.11 in connection with (2.20) of Lemma 2.12, one has for any $`L_1\text{,}L_2𝔏`$ and any $`A𝔄`$:
$$\left|\sigma (L_{1}^{}{}_{}{}^{}AL_2)\right|h_{\mathrm{}}E(\overline{\mathrm{\Delta }}+\mathrm{\Gamma }_1)AE(\overline{\mathrm{\Delta }}+\mathrm{\Gamma }_2)q_\mathrm{\Delta }(L_1)q_\mathrm{\Delta }(L_2)\text{.}$$
(6.15)
Specializing now to operators $`A`$ from the unit ball of a local $`C^{}`$-algebra and to vacuum annihilation operators $`L_1`$ and $`L_2`$ from the $`|||.|||_{2s}`$-unit balls of $`𝔏_0(\mathrm{\Gamma }_1)`$ and $`𝔏_0(\mathrm{\Gamma }_2)`$, respectively, with compact and convex $`\mathrm{\Gamma }_k`$, we infer from (6.15) by use of Lemma 6.4 in connection with Definition 6.2 that there exist constants $`C^{}(\mathrm{\Delta },\mathrm{\Gamma }_1)`$ and $`C^{}(\mathrm{\Delta },\mathrm{\Gamma }_2)`$ such that
$$\begin{array}{c}\left|\sigma (L_{1}^{}{}_{}{}^{}AL_2)\right|h_{\mathrm{}}C^{}(\mathrm{\Delta },\mathrm{\Gamma }_1)C^{}(\mathrm{\Delta },\mathrm{\Gamma }_2)\hfill \\ \hfill E(\overline{\mathrm{\Delta }}+\mathrm{\Gamma }_1)AE(\overline{\mathrm{\Delta }}+\mathrm{\Gamma }_2)L_1E\left(\mathrm{\Delta }^{\prime \prime }(\mathrm{\Delta },\mathrm{\Gamma }_1)\right)^{1/2}L_2E\left(\mathrm{\Delta }^{\prime \prime }(\mathrm{\Delta },\mathrm{\Gamma }_2)\right)^{1/2}\end{array}$$
(6.16)
with appropriate bounded Borel sets $`\mathrm{\Delta }^{\prime \prime }(\mathrm{\Delta },\mathrm{\Gamma }_1)`$ and $`\mathrm{\Delta }^{\prime \prime }(\mathrm{\Delta },\mathrm{\Gamma }_1)`$, depending on $`\mathrm{\Delta }`$ and both of the compact sets $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$. An inequality of type (6.16) can likewise be imposed on the corresponding sesquilinear form, which opens up the way to distinguish a certain subclass in the space $`\mathrm{S}`$ of *all* sesquilinear forms on $`𝔏\times 𝔏`$.
###### Definition 6.6.
The set $`\mathrm{S}^\mathrm{b}`$ of all sesquilinear forms $`W(.|.)`$ on $`𝔏\times 𝔏`$ which are characterized by the existence of constants $`C^{}(\mathrm{\Delta },\mathrm{\Gamma }_1)`$, $`C^{}(\mathrm{\Delta },\mathrm{\Gamma }_2)`$ and $`C^{\prime \prime }`$ such that the following condition (6.17) holds for any operator $`A𝔄_1(𝒪)`$ as well as for $`L_1𝔏_{0,1}^{2s}(\mathrm{\Gamma }_1)`$ and $`L_2𝔏_{0,1}^{2s}(\mathrm{\Gamma }_2)`$ with compact and convex $`\mathrm{\Gamma }_l`$ is a subspace of $`\mathrm{S}`$:
$$\begin{array}{c}\left|W\left(L_1\right|AL_2)\right|C^{\prime \prime }C^{}(\mathrm{\Delta },\mathrm{\Gamma }_1)C^{}(\mathrm{\Delta },\mathrm{\Gamma }_2)\hfill \\ \hfill E(\overline{\mathrm{\Delta }}+\mathrm{\Gamma }_1)AE(\overline{\mathrm{\Delta }}+\mathrm{\Gamma }_2)L_1E\left(\mathrm{\Delta }^{\prime \prime }(\mathrm{\Delta },\mathrm{\Gamma }_1)\right)^{1/2}L_2E\left(\mathrm{\Delta }^{\prime \prime }(\mathrm{\Delta },\mathrm{\Gamma }_2)\right)^{1/2}\text{.}\end{array}$$
(6.17)
Its intersection with the positive cone $`\mathrm{W}`$ of all particle weights according to Definition 3.14 is again a positive proper convex cone, denoted $`\mathrm{W}^\mathrm{b}`$, which obviously comprises the particle weights induced by asymptotic functionals.
Due to (6.17), the space $`\mathrm{S}^\mathrm{b}`$ can be furnished with various seminorm topologies.
###### Definition 6.7.
For any combination of bounded regions $`𝒪`$ with compact and convex $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$, a seminorm $`𝒫_{\mathrm{\Gamma }_1,\mathrm{\Gamma }_2}^𝒪`$ can be introduced on $`\mathrm{S}^\mathrm{b}`$ by
$$𝒫_{\mathrm{\Gamma }_1,\mathrm{\Gamma }_2}^𝒪(W)sup\{\left|W\left(L_1\right|AL_2)\right|:A𝔄_1(𝒪),L_1𝔏_{0,1}^{2s}(\mathrm{\Gamma }_1),L_2𝔏_{0,1}^{2s}(\mathrm{\Gamma }_2)\}\text{.}$$
(6.18)
The convex subset of $`\mathrm{W}^\mathrm{b}`$, which is to be used from now on and will turn out to be compact when furnished with a suitable topology, is introduced again in view of (6.17).
###### Definition 6.8.
$`\mathrm{W}_\mathrm{c}^\mathrm{b}`$ is the convex set of all particle weights in $`\mathrm{W}^\mathrm{b}`$ which satisfy the inequality
$$\begin{array}{c}\left|W\left(L_1\right|AL_2)\right|C^{}(\mathrm{\Delta },\mathrm{\Gamma }_1)C^{}(\mathrm{\Delta },\mathrm{\Gamma }_2)\hfill \\ \hfill E(\overline{\mathrm{\Delta }}+\mathrm{\Gamma }_1)AE(\overline{\mathrm{\Delta }}+\mathrm{\Gamma }_2)L_1E\left(\mathrm{\Delta }^{\prime \prime }(\mathrm{\Delta },\mathrm{\Gamma }_1)\right)^{1/2}L_2E\left(\mathrm{\Delta }^{\prime \prime }(\mathrm{\Delta },\mathrm{\Gamma }_2)\right)^{1/2}\end{array}$$
(6.19)
for all bounded regions $`𝒪`$ and all compact and convex $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$. The difference between this condition and (6.17) is that the only $`W`$-dependent constant $`C^{\prime \prime }`$ has been omitted.
###### Remark.
According to (6.16), all particle weights arising from asymptotic functionals with $`h_{\mathrm{}}1`$ satisfy (6.19) and are thus contained in $`\mathrm{W}_\mathrm{c}^\mathrm{b}`$.
Now, as a consequence of the Compactness Criterion of Fredenhagen and Hertel, we know that there exist in each case finitely many operators in $`𝔄_1(𝒪)`$ as well as in $`𝔏_{0,1}^{2s}(\mathrm{\Gamma }_l)`$, $`l=1\text{,}2`$, serving as centres of $`\delta `$-balls to cover the sets $`𝔏_{0,1}^{2s}(\mathrm{\Gamma }_l)E\left(\mathrm{\Delta }^{\prime \prime }(\mathrm{\Delta },\mathrm{\Gamma }_l)\right)`$ and $`E(\overline{\mathrm{\Delta }}+\mathrm{\Gamma }_1)𝔄_1(𝒪)E(\overline{\mathrm{\Delta }}+\mathrm{\Gamma }_2)`$. These operators can be used to span finite-dimensional subspaces in $`𝔏_0(\mathrm{\Gamma }_1)`$ and $`𝔄(𝒪)𝔏_0(\mathrm{\Gamma }_2)`$. The corresponding space of sesquilinear forms defined on these domains is again finite-dimensional, so that its unit ball with respect to the relative $`𝒫_{\mathrm{\Gamma }_1,\mathrm{\Gamma }_2}^𝒪`$-topology can be covered by a finite number of $`\epsilon `$-balls ( note that bounded sets in finite-dimensional vector spaces are relatively compact). The restriction of any element $`W`$ of $`\mathrm{W}_\mathrm{c}^\mathrm{b}`$ to the named subspaces of $`𝔏_0(\mathrm{\Gamma }_1)`$ and $`𝔄(𝒪)𝔏_0(\mathrm{\Gamma }_2)`$ is thus contained in one of these balls. This in turn means, that we can even select a finite number of elements $`W_k\mathrm{W}_\mathrm{c}^\mathrm{b}`$, $`k=1\text{,}\mathrm{}\text{,}N(2\epsilon )`$, such that any element of $`\mathrm{W}_\mathrm{c}^\mathrm{b}`$ is contained in a $`2\epsilon `$-ball around at least one of these chosen forms with respect to the aforementioned relative $`𝒫_{\mathrm{\Gamma }_1,\mathrm{\Gamma }_2}^𝒪`$-topology. But then it can be shown that the $`3\epsilon `$-balls with respect to the $`𝒫_{\mathrm{\Gamma }_1,\mathrm{\Gamma }_2}^𝒪`$-topology proper indeed cover all of $`\mathrm{W}_\mathrm{c}^\mathrm{b}`$. To see this, let $`W_K`$ be the element pertaining to $`W\mathrm{W}_\mathrm{c}^\mathrm{b}`$ and let $`L_l𝔏_{0,1}^{2s}(\mathrm{\Gamma }_l)`$, $`l=1\text{,}2`$, and $`A𝔄_1(𝒪)`$ be arbitrary. Then there exist operators $`L_l^\delta 𝔏_{0,1}^{2s}(\mathrm{\Gamma }_l)`$ and $`A^\delta 𝔄_1(𝒪)`$ which satisfy
$`(L_lL_l^\delta )E\left(\mathrm{\Delta }^{\prime \prime }(\mathrm{\Delta },\mathrm{\Gamma }_l)\right)`$ $`<\delta \text{,}`$ (6.20a)
$`E(\overline{\mathrm{\Delta }}+\mathrm{\Gamma }_1)(AA^\delta )E(\overline{\mathrm{\Delta }}+\mathrm{\Gamma }_2)`$ $`<\delta \text{.}`$ (6.20b)
Now, making use of condition (6.19) on elements of $`\mathrm{W}_\mathrm{c}^\mathrm{b}`$ in connection with (6.20) as well as of the defining property for $`W_K`$, we get
$$\begin{array}{c}\left|W\left(L_1|AL_2\right)W_K\left(L_1|AL_2\right)\right|\left|W\left(L_1|AL_2\right)W\left(L_1^\delta |A^\delta L_2^\delta \right)\right|\hfill \\ \hfill +\left|W\left(L_1^\delta |A^\delta L_2^\delta \right)W_K\left(L_1^\delta |A^\delta L_2^\delta \right)\right|+\left|W_K\left(L_1^\delta |A^\delta L_2^\delta \right)W_K\left(L_1|AL_2\right)\right|\\ \hfill <2(2\delta ^{1/2})+2\epsilon \text{.}\end{array}$$
(6.21)
Since we are free to choose $`\delta `$ appropriately small in dependence on a given $`\epsilon >0`$, this final relation shows, upon taking the supremum with respect to the operators appearing on the left-hand side, that for each $`W\mathrm{W}_\mathrm{c}^\mathrm{b}`$ there exists at least one $`W_K\mathrm{W}_\mathrm{c}^\mathrm{b}`$ such that $`𝒫_{\mathrm{\Gamma }_1,\mathrm{\Gamma }_2}^𝒪\left(WW_K\right)<3\epsilon `$ in accordance with our statement.
As in Chapter 4 we want to pass at this point to countable families $`\{𝒪_n\}_n`$ and $`\{\mathrm{\Gamma }_l\}_l`$ of bounded regions in space-time and of compact and convex subsets of $`\mathrm{}\overline{V}_+`$. By Definition 6.7 any triple taken from these sequences defines a seminorm on $`\mathrm{S}^\mathrm{b}`$. In this way $`\mathrm{S}^\mathrm{b}`$ can be topologized with a sequence $`\{𝒫_m\}_m`$ of seminorms and thereby becomes a locally convex (Hausdorff) space. This space is metrizable according to \[44, Chapter Four, § 18, 2.(2)\] and its topology can moreover be derived from the increasing sequence of seminorms
$$^k(W)\underset{1mk}{\mathrm{max}}𝒫_m(W)\text{,}W\mathrm{S}^\mathrm{b}\text{,}\mathrm{k}\text{.}$$
(6.22)
This countable system can then be used to define an (F)-norm \[44, p. 163\] on $`\mathrm{S}^\mathrm{b}`$ which furnishes this space with the same topology. It is given by
$$|||W|||_F\underset{k=1}{\overset{\mathrm{}}{}}\frac{1}{2^k}\frac{^k(W)}{1+^k(W)}\text{,}W\mathrm{S}^\mathrm{b}\text{,}$$
(6.23)
and generates a translation-invariant metric \[44, Chapter Four, § 18, 2.(3)\]. Given $`\epsilon >0`$ there exists, according to (6.23), an index $`M`$ such that for any $`W\text{,}W^{}\mathrm{S}^\mathrm{b}`$
$$\begin{array}{c}|||WW^{}|||_F=\underset{k=1}{\overset{\mathrm{}}{}}\frac{1}{2^k}\frac{^k(WW^{})}{1+^k(WW^{})}\hfill \\ \hfill <\frac{\epsilon }{2}+\underset{k=1}{\overset{M}{}}\frac{1}{2^k}\frac{^k(WW^{})}{1+^k(WW^{})}\frac{\epsilon }{2}+\underset{k=1}{\overset{M}{}}\frac{1}{2^k}^M(WW^{})\text{.}\end{array}$$
(6.24)
A consequence of the preceding paragraph in combination with the definition (6.22) is the fact that $`\mathrm{W}_\mathrm{c}^\mathrm{b}`$ can be covered by a finite number of balls with a given arbitrarily small radius with respect to the $`^k`$-topologies. It is then an immediate conclusion from the definitions involved that the sesquilinear functionals arising as limits with respect to $`|||.|||_F`$ of sequences in $`\mathrm{W}_\mathrm{c}^\mathrm{b}`$ are again elements of $`\mathrm{W}_\mathrm{c}^\mathrm{b}`$. This convex subset thus turns out to be closed. Since it has been seen above to be precompact, it is indeed compact in the $`|||.|||_F`$-topology.
###### Proposition 6.9.
The convex set $`\mathrm{W}_\mathrm{c}^\mathrm{b}`$ in the class of all particle weights is compact with respect to the metric derived from the (F)-norm $`|||.|||_F`$.
The above work has laid the foundation for an application of Choquet’s Theorem \[1, Corollary I.4.9\] which tells us that any particle weight $`.|.`$ in the metrizable compact convex set $`\mathrm{W}_\mathrm{c}^\mathrm{b}`$ can be represented by a positive and normalized measure vanishing off its extreme boundary $`_e\mathrm{W}_\mathrm{c}^\mathrm{b}`$:
$$.|.=_{_e\mathrm{W}_\mathrm{c}^\mathrm{b}}d\upsilon (\zeta ).|._\zeta \text{.}$$
(6.25)
The above result represents the present status of the Choquet approach to a disintegration theory for particle weights. The problem to be tackled at this point is the open question of how a base can be fitted into the cone $`\mathrm{W}^\mathrm{b}`$ which is completely contained in $`\mathrm{W}_\mathrm{c}^\mathrm{b}`$. This would allow for a disintegration of a particle weight on this base in terms of extremal points, defining extremal rays of the cone $`\mathrm{W}^\mathrm{b}`$ and representing pure particle weights. On this foundation a complete theory in parallel to that developed for states on a $`C^{}`$-algebra $`𝔄`$ in \[11, Sections 4.1 and 4.2\] still awaits its completion. The advantage of this approach in comparison to the spatial disintegration presented in Chapter 4 is that, apart from the somewhat intricate topological considerations, it is more direct and the resulting pure particle weights are no longer subject to the restrictive Definition 4.1. On the other hand, the mathematical problems concerning convex sets in infinite-dimensional spaces are far from being trivial. Therefore, a lot of work remains to be done until eventually the particle content of a quantum field theory is seen to be encoded in the geometrical structure (the set of extreme rays) of a positive cone of particle weights.
## Chapter 7 Summary and Outlook
The present work is based on the general point of view that the concept of ‘particles’ is asymptotic in nature and simultaneously has to be founded by making appropriate use of the notion of locality. This reflects our conviction that the long-standing problem of ‘asymptotic completeness’ of quantum field theory, i. e., the question if a quantum field theoretic model can be interpreted completely in terms of particles, has to be tackled by the aid of further restrictions on the general structure, which essentially are of a local character. The question is, what the local structure of a theory should be in order that it governs scattering processes in such a way that asymptotically the physical states appear to clot in terms of certain entities named particles. The compactness and nuclearity conditions discussed in and the references therein are examples of this kind of approach. We do not claim that they already give a complete answer, but believe that they indicate the right direction.
In this thesis we have constructed asymptotic functionals on a certain algebra of detectors giving rise to particle weights which can be interpreted as mixtures of particle states. A disintegration theory has been developed for restricted particle weights by means of a highly technical procedure in Chapter 4. This constitutes the basis for the definition of mass and spin even in the case of charged states . We are convinced that the technicalities involved can be dissolved by future research. In this connection the analysis of concrete models may be helpful. Such investigations are already under way. They concern the Schwinger model and an application of our formalism to quantum electrodynamics . It is expected that some insight may be gained with respect to the open questions mentioned in the various chapters. E. g., the convergence problem in connection with Theorem 3.10 can perhaps be solved with additional information at hand, and the direct integral decomposition of Chapter 4 might get more manageable, unfolding the connection between the intrinsic energy-momenta pertaining to the irreducible representations and the geometrical energy-momenta (velocities) that stem from the asymptotic functionals.
So far we have considered *single* particle weights. Another field of future research is the inspection of coincidence arrangements of detectors as in . In this respect, too, the analysis of concrete models is helpful.
As indicated by Chapter 5 and in view of the partial results presented in Chapter 6 phase space restrictions seem to be a key ingredient in the general analysis, in particular of the Choquet approach to disintegration. This theory is still in its initial stage. But, difficult as the mathematical problems concerning convex sets in infinite-dimensional spaces are, it deserves further efforts. Presumably, both the spatial disintegration and the Choquet decomposition will eventually turn out to be essentially equivalent, revealing relations similar to those encountered in the disintegration theory of states on $`C^{}`$-algebras \[11, Chapter 4\]. Further studies have to disclose the geometrical structure of the positive cone of particle weights, as the particle content of a theory seems to be encoded in this kind of information.
## Appendix A Concepts of Differentiability
Various notions of differentiability have to be used in this work and some of them take on a somewhat unusual shape. So it seems right to collect in this appendix a number of definitions and propositions, both to assign a precise meaning to the concepts proper and to their consequences as well as to fix the notation.
### A.1 Differentiation in Locally Convex Spaces
###### Definition A.1.
Let $`𝔛`$ be a (real or complex) normed space and let $`𝔙`$ be a locally convex space over the same field whose topology is defined by the family $`\{q_\lambda :\lambda L\}`$ of seminorms which separate the points of $`𝔙`$. Suppose further that we are given an open subset $`𝔊`$ of $`𝔛`$.
* A mapping $`F:𝔊𝔙`$ is called differentiable at the point $`𝔵𝔊`$ if there exists a continuous linear mapping $`T:𝔛𝔙`$ such that for any vector $`𝔥`$ in a certain $`\mathfrak{0}`$-neighbourhood $`𝔘𝔛`$ the increment $`F(𝔵+𝔥)F(𝔵)`$ allows for the linearized approximation
$$F(𝔵+𝔥)F(𝔵)=T𝔥+R[F,𝔵](𝔥)\text{,}$$
(A.1a)
where $`R[F,𝔵]`$ is a mapping on $`𝔘`$ to $`𝔙`$ subject to the condition
$$\underset{𝔥\mathfrak{0}}{lim}𝔥^1q_\lambda \left(R[F,𝔵](𝔥)\right)=0$$
(A.1b)
for any seminorm $`q_\lambda `$, $`\lambda L`$. The linear operator $`T`$ occurring in (A.1a) is signified as $`𝔇F(𝔵)`$ and called the derivative of $`F`$ at $`𝔵`$.
* The mapping $`F:𝔊𝔙`$ is called differentiable if it is differentiable at any $`𝔵𝔊`$.
* The differentiable mapping $`F:𝔊𝔙`$ is called continuously differentiable if the mapping $`𝔊𝔵𝔇F(𝔵)𝔥𝔙`$, which exists by assumption, is continuous with respect to the locally convex topology of $`𝔙`$ for any given $`𝔥𝔛`$.
###### Remark.
The definition of the continuous linear operator $`𝔇F(𝔵)`$ requires uniqueness of the corresponding $`T`$ in (A.1a), but this is easily established. Assume the existence of another $`\mathfrak{0}`$-neighbourhood $`𝔘^{}`$, a continuous linear operator $`T^{}:𝔛𝔙`$ and a mapping $`R^{}[F,𝔵]:𝔘^{}𝔙`$ which, upon insertion into (A.1a), represent the increment $`F(𝔵+𝔥)F(𝔵)`$ such that $`R^{}[F,𝔵]`$ fulfills a condition analogous to (A.1b). Then
$$T𝔥T^{}𝔥=R^{}[F,𝔵](𝔥)R[F,𝔵](𝔥)\text{,}𝔥𝔘𝔘^{}\text{.}$$
Let $`𝔶𝔛`$, $`𝔶\mathfrak{0}`$, be arbitrary but fixed, then for $`\alpha \{0\}`$ small enough we infer from the above equation due to the linearity of both $`T`$ and $`T^{}`$
$$\begin{array}{c}q_\lambda \left(T𝔶T^{}𝔶\right)=q_\lambda \left(\alpha ^1\left(R^{}[F,𝔵](\alpha 𝔶)R[F,𝔵](\alpha 𝔶)\right)\right)\hfill \\ \hfill =𝔶\alpha 𝔶^1q_\lambda \left(R^{}[F,𝔵](\alpha 𝔶)R[F,𝔵](\alpha 𝔶)\right)\text{,}\end{array}$$
where the right-hand side vanishes in the limit $`\alpha 0`$ for any seminorm $`q_\lambda `$, according to (A.1b). This yields $`q_\lambda \left(T𝔶\right)=q_\lambda \left(T^{}𝔶\right)`$, valid also for $`𝔶=\mathfrak{0}`$, and as a consequence $`T𝔶=T^{}𝔶`$ for any $`𝔶𝔛`$ since the seminorms $`q_\lambda `$ separate the points in $`𝔙`$.
An immediate consequence of the presumed continuity of the linear operators $`𝔇F(𝔵)`$, entering as derivatives the representation (A.1a) of the increment of $`F`$ at $`𝔵`$, is the fact that differentiability implies continuity.
###### Corollary A.2.
Let $`𝔛`$ be a normed space and let $`𝔙`$ be a locally convex space. If the mapping $`F:𝔊𝔙`$, $`𝔊𝔛`$ open, is differentiable at the point $`𝔵𝔊`$ then it is also continuous in $`𝔵`$.
The methods used in the standard theory of differentiable functions yield the following propositions when applied to the concept laid open in Definition A.1, the main modification being the occurrence of seminorms $`q_\lambda `$ on $`𝔙`$ in (A.1b).
###### Proposition A.3 (Product Rule for Derivatives).
Let $`𝔛`$ be a normed space and $`𝔊`$ an open subset of $`𝔛`$.
* Suppose that $`𝔙`$ is a locally convex space and that the mappings $`F:𝔊𝔙`$ and $`f:𝔊𝕂`$, $`𝕂`$ the scalar field of both $`𝔛`$ and $`𝔙`$, are differentiable at $`𝔵𝔊`$. Then their product $`fF`$ is differentiable at this point, too, and the derivative at $`𝔵`$ is given by
$$𝔇(fF)(𝔵)𝔥=𝔇f(𝔵)𝔥F(𝔵)+f(𝔵)𝔇F(𝔵)𝔥\text{,}𝔥𝔛\text{.}$$
* Let $`𝔜`$ be a normed algebra and assume that the mappings $`F:𝔊𝔜`$ and $`G:𝔊𝔜`$ are differentiable at $`𝔵𝔊`$. Then their product $`FG`$ is differentiable at $`𝔵`$, too, and the derivative is
$$𝔇(FG)(𝔵)𝔥=𝔇F(𝔵)𝔥G(𝔵)+F(𝔵)𝔇G(𝔵)𝔥\text{,}𝔥𝔛\text{.}$$
###### Proposition A.4 (Chain Rule for Derivatives).
Let $`𝔛`$ and $`𝔜`$ be normed spaces and let $`𝔙`$ be a locally convex space. Assume further that the mapping $`G:𝔊_1𝔜`$ is differentiable at $`𝔵𝔊_1`$ and that the mapping $`F:𝔊_2𝔙`$ is differentiable at $`G(𝔵)`$, where $`𝔊_1`$ and $`𝔊_2`$ are open subsets of $`𝔛`$ and $`𝔜`$, respectively, and $`G(𝔊_1)𝔊_2`$. Then the composition of $`F`$ and $`G`$: $`FG:𝔊_1𝔙`$, exists and is differentiable at $`𝔵`$ with a derivative connected to those of $`F`$ and $`G`$ through
$$𝔇(FG)(𝔵)=𝔇G\left(F(𝔵)\right)𝔇F(𝔵)\text{.}$$
The fundamental Mean Value Theorem which has to be formulated in the setting of Definition A.1 is based on the following two lemmas. Their proof as well as that of the theorem proper is an adaptation of the reasoning to be found in \[37, Kapitel XX, Abschnitt 175\].
###### Lemma A.5.
Let $`F:[a,b]𝔙`$ be a continuous mapping on the compact interval $`[a,b]`$ to the locally convex space $`𝔙`$ and suppose that it is differentiable on the interior of this set with $`𝔇F(x)=0`$ for any $`x]a,b[`$. Then $`F`$ is constant on $`[a,b]`$.
###### Proof.
Let $`s`$ and $`t`$ be arbitrary distinct points in $`]a,b[`$. We shall assume $`s<t`$ and want to show that $`F(s)=F(t)`$. Define $`u2^1(ts)`$ and consider one of the seminorms $`q_\lambda `$ topologizing $`𝔙`$. There are two possibilities:
$`q_\lambda \left(F(u)F(s)\right)`$ $`q_\lambda \left(F(t)F(u)\right)\text{,}`$ (A.2a)
$`q_\lambda \left(F(t)F(u)\right)`$ $`>q_\lambda \left(F(u)F(s)\right)\text{.}`$ (A.2b)
Depending on the actual situation we define an interval $`]s_1,t_1[[a,b]`$, choosing $`s_1s`$, $`t_1u`$ in case (A.2a) and $`s_1u`$, $`t_1t`$ in case (A.2b). Independent of this selection is the estimate
$$q_\lambda \left(F(t)F(s)\right)q_\lambda \left(F(t)F(u)\right)+q_\lambda \left(F(u)F(s)\right)2q_\lambda \left(F(t_1)F(s_1)\right)\text{.}$$
(A.3)
The same procedure can then be applied to the interval $`]s_1,t_1[`$, to the resulting interval $`]s_2,t_2[`$ and so on. In this way a sequence of intervals $`]s_n,t_n[`$ is constructed, which is decreasing with respect to the inclusion relation: $`]s_{n+1},t_{n+1}[]s_n,t_n[`$. Furthermore the lengths are explicitly known as $`t_ns_n=2^n(ts)`$ and the estimate (A.3) can be generalized to
$$q_\lambda \left(F(t)F(s)\right)2^nq_\lambda \left(F(t_n)F(s_n)\right)\text{.}$$
(A.4)
There exists exactly one point $`u_0]a,b[`$ belonging to all intervals of this sequence and by assumption $`𝔇F(u_0)=0`$, so that for $`h`$ in a small 0-neighbourhood $`𝒰`$ the increment of $`F`$ at $`u_0`$ is represented by
$$F(u_0+h)F(u_0)=hR(h)$$
(A.5a)
with a mapping $`R:𝒰𝔙`$ satisfying
$$\underset{h0}{lim}q_\lambda \left(R(h)\right)=0\text{.}$$
(A.5b)
Hence, given $`ϵ>0`$, there exists $`N`$ such that $`q_\lambda \left(R(u_0s_n)\right)`$ and $`q_\lambda \left(R(t_nu_0)\right)`$ are majorized by $`(ts)^1ϵ`$ for $`n>N`$. According to (A.5a) this implies
$$\begin{array}{c}q_\lambda \left(F(t_n)F(s_n)\right)q_\lambda \left(F(t_n)F(u_0)\right)+q_\lambda \left(F(s_n)F(u_0)\right)\hfill \\ \hfill |t_nu_0|q_\lambda \left(R(t_nu_0)\right)+|u_0s_n|q_\lambda \left(R(u_0s_n)\right)\\ \hfill (t_nu_0)\frac{ϵ}{ts}+(u_0s_n)\frac{ϵ}{ts}=(t_ns_n)\frac{ϵ}{ts}=\frac{ϵ}{2^n}\text{,}\end{array}$$
where we made use of the length formula for the interval $`]s_n,t_n[`$. From (A.4) one then infers
$$q_\lambda \left(F(t)F(s)\right)2^n\frac{ϵ}{2^n}=ϵ\text{,}$$
so that, by arbitrariness of $`ϵ`$ and $`q_\lambda `$ together with the separation property of the seminorms, we see that $`F(t)=F(s)=𝔳_0𝔙`$. This relation holds for any $`s`$, $`t]a,b[`$ and extends by the supposed continuity of $`F`$ to all of $`[a,b]`$, establishing $`F𝔳_0`$ as stated. ∎
###### Lemma A.6.
Let $`F:[a,b]𝔙`$ be a continuous mapping on the compact interval $`[a,b]`$ to the locally convex space $`𝔙`$ and define $`G:[a,b]\overline{𝔙}`$, $`\overline{𝔙}`$ the completion of $`𝔙`$, through the integral
$$G(x)_a^x𝑑\vartheta F(\vartheta )\text{,}x[a,b]\text{.}$$
Then the mapping $`G`$ is differentiable for any $`x_0]a,b[`$ and the action of the derivative $`𝔇G(x_0)`$ as a linear operator on $``$ is given by
$$𝔇G(x_0)h=hF(x_0)\text{,}h\text{.}$$
(A.6)
###### Proof.
By \[26, II.6.2\] $`G`$ is a well-defined $`\overline{𝔙}`$-valued mapping on the compact interval $`[a,b]`$. For $`x_0]a,b[`$ and $`h`$ satisfying $`x_0+h[a,b]`$ we have
$$G(x_0+h)G(x_0)=_{x_0}^{x_0+h}𝑑\vartheta F(\vartheta )\text{,}$$
hence
$$G(x_0+h)G(x_0)hF(x_0)=_{x_0}^{x_0+h}𝑑\vartheta \left(F(\vartheta )F(x_0)\right)\rho (h)\text{.}$$
(A.7a)
Now by assumption, $`\vartheta q_\lambda \left(F(\vartheta )F(x_0)\right)`$ is continuous on the compact interval $`I_h`$ of integration for any of the defining seminorms $`q_\lambda `$ of $`𝔙`$, and, according to \[26, II.6.2 in connection II.5.4\], one has for any $`h`$ the estimate
$$|h|^1q_\lambda \left(\rho (h)\right)|h|^1\left|_{x_0}^{x_0+h}𝑑\vartheta q_\lambda \left(F(\vartheta )F(x_0)\right)\right|\underset{\vartheta I_h}{\mathrm{max}}q_\lambda \left(F(\vartheta )F(x_0)\right)\text{,}$$
(A.7b)
where the right-hand side vanishes in the limit $`h0`$. Thus (A.7a) corresponds to the representation (A.1a) of Definition A.1 in terms of the increment $`G(x_0+h)G(x_0)`$ with a residual term $`\rho (h)`$ satisfying (A.1b). This proves differentiability of $`G`$ on $`]a,b[`$ along with relation (A.6). ∎
###### Theorem A.7 (Mean Value Theorem).
Let $`𝔛`$ be a normed space and $`𝔙`$ be a locally convex space. Let furthermore $`F:𝔊𝔙`$, $`𝔊𝔛`$ open, be a continuously differentiable mapping (in the sense of Definition A.1) and consider $`𝔵_0𝔊`$ and $`𝔥𝔛`$ small enough so that $`𝔵_0+\vartheta 𝔥𝔊`$ for $`0\vartheta 1`$. Then
$$F(𝔵_0+𝔥)F(𝔵_0)=_0^1𝑑\vartheta 𝔇F(𝔵_0+\vartheta 𝔥)𝔥\text{.}$$
(A.8)
###### Proof.
Given $`𝔵_0𝔊`$ and $`𝔥𝔛`$ as above we define two mappings $`F_1`$ and $`F_2`$ on the compact interval $`[0,1]`$ to $`𝔙`$ respectively $`\overline{𝔙}`$ through
$`s`$ $`F_1(s)F(𝔵_0+s𝔥)\text{,}`$ (A.9a)
$`s`$ $`F_2(s){\displaystyle _0^s}𝑑\vartheta 𝔇F(𝔵_0+\vartheta 𝔥)𝔥\text{.}`$ (A.9b)
From Lemma A.6 and Proposition A.4 we infer $`𝔇F_2(s)=𝔇F(𝔵_0+s𝔥)𝔥=𝔇F_1(s)`$ for any $`s]0,1[`$. This implies, according to Lemma A.5, that the mapping $`F_1F_2`$ is constant on the interval $`[0,1]`$ (Note, that $`F_1`$ as well as $`F_2`$ are continuous.). Hence
$$F(𝔵_0)=F_1(0)F_2(0)=F_1(1)F_2(1)=F(𝔵_0+𝔥)_0^1𝑑\vartheta 𝔇F(𝔵_0+\vartheta 𝔥)𝔥\text{,}$$
which is just equation (A.8) re-written. ∎
### A.2 Differentiation on Analytic Manifolds
Being of a local nature, the concept of differentiability set out in Definition A.1 can be generalized to $`𝔙`$-valued mappings on analytic manifolds in the following way.
###### Definition A.8.
Let $``$ be a (real or complex) analytic manifold of dimension $`d`$ and let $`𝔙`$ be a locally convex space over the same field. Let furthermore $`(𝒰,\varphi )`$ denote a local chart on $``$, which means that $`\varphi (𝒰)𝕂^d`$, $`𝕂=`$ or $`𝕂=`$.
* The mapping $`F:𝒰𝔙`$ is called differentiable (with respect to $`\varphi `$) at $`m_0𝒰`$ if $`F\varphi ^1:\varphi (𝒰)𝔙`$ is differentiable at $`\varphi (m_0)`$ in the sense of Definition A.1. The derivative is denoted $`𝔇_\varphi F(m_0)𝔇\left(F\varphi ^1\right)\left(\varphi (m_0)\right)`$.
* $`F:𝒰𝔙`$ is called (continuously) differentiable if $`F\varphi ^1`$ is (continuously) differentiable in the sense of Definition A.1.
* The mapping $`F:𝔙`$ is called (continuously) differentiable if to any $`m_0`$ there exists a local chart $`(𝒰,\varphi )`$ containing $`m_0`$ such that $`F𝒰`$ is (continuously) differentiable with respect to $`\varphi `$.
* Let $`\{𝒆_𝒊:𝒊=\mathit{1}\text{,}\mathrm{}\text{,}𝒅\}`$ be the canonical orthonormal basis of $`𝕂^d`$. Then $`F:𝒰𝔙`$ is said to have continuous partial derivatives if there exist $`d`$ continuous mappings $`F_\varphi ^i:𝒰𝔙`$, such that the increment of $`F`$ in direction $`𝒆_𝒊`$ at any $`m_0=\varphi ^1(𝒕_\mathit{0})𝒰`$ allows for the representation
$$F\varphi ^1(𝒕_\mathit{0}+𝒉𝒆_𝒊)𝑭\varphi ^\mathit{1}(𝒕_\mathit{0})=𝒉𝑭_\varphi ^𝒊(𝒎_\mathit{0})+𝑹[𝑭\varphi ^\mathit{1},𝒕_\mathit{0}](𝒉)$$
(A.10a)
if $`h𝕂`$ is small enough, where the residual term satisfies
$$\underset{h0}{lim}|h|^1q_\lambda \left(R[F\varphi ^1,𝒕_\mathit{0}](𝒉)\right)=\mathit{0}$$
(A.10b)
for any seminorm $`q_\lambda `$, $`\lambda L`$.
* Higher derivatives of the mapping $`F:𝒰𝔙`$ are defined recursively in terms of partial derivatives of the mappings $`F_\varphi ^i`$, $`i=1\text{,}\mathrm{}\text{,}d`$, and, if they happen to exist, are denoted $`F_\varphi ^\kappa `$ for multi-indices $`\kappa (k_1,\mathrm{},k_d)_0^d`$ in an obvious fashion (for given $`i`$ let $`F_\varphi ^{\kappa _i}F_\varphi ^i`$ where all entries in $`\kappa _i`$ apart from $`k_i=1`$ vanish). $`F:𝒰𝔙`$ is called $`N`$-fold (or infinitely often) continuously differentiable if the mappings $`F_\varphi ^\kappa `$ exist and are continuous for any $`|\kappa |_ik_iN`$ (or $`|\kappa |<\mathrm{}`$). These concepts apply equally to mappings $`F`$ defined on all of $``$.
###### Remark.
If $`F`$ is differentiable at $`m_0𝒰`$ with respect to the local chart $`(𝒰,\varphi )`$ it is also differentiable with respect to any other local chart $`(𝒱,\psi )`$ containing $`m_0`$, and according to Proposition A.4 one has
$$𝔇_\psi F(m_0)=𝔇_\varphi F(m_0)\left(\varphi \psi ^1\right)^{}\left(\psi (m_0)\right)\text{,}$$
(A.11)
where $`\left(\varphi \psi ^1\right)^{}`$ denotes the first derivative (Jacobi matrix) of the analytic function $`\varphi \psi ^1:\psi (𝒰𝒱)\varphi (𝒰𝒱)`$.
Strictly speaking, the definition of and notation for higher derivatives of a mapping $`F:𝒰𝔙`$ is justified only after the following two results are established.
###### Proposition A.9.
$`F:𝒰𝔙`$ is continuously differentiable if and only if it has continuous partial derivatives in all directions $`𝐞_𝐢`$, $`i=1\text{,}\mathrm{}\text{,}d`$.
###### Proof.
* If $`F`$ is continuously differentiable the mappings
$$𝒰m_0F_\varphi ^i(m_0)𝔇_\varphi F(m_0)𝒆_𝒊$$
(A.12)
are continuous for any $`i`$; furthermore (A.10a) and (A.10b) correspond for each $`i`$ exactly to (A.1a) and (A.1b) of Definition A.1 setting $`𝔥=h𝒆_𝒊`$, so that the first part of the statement is almost trivial.
* Let all the partial derivatives of $`F`$ exist as continuous mappings $`F_\varphi ^i:𝒰𝔙`$, then, for small $`𝒉=_𝒊𝒉_𝒊𝒆_𝒊𝕂^𝒅`$, we have through an application of the Mean Value Theorm A.7 for any $`m_0=\varphi ^1(𝒕_\mathit{0})𝒰`$
$$\begin{array}{c}F\varphi ^1(𝒕_\mathit{0}+𝒉)𝑭\varphi ^\mathit{1}(𝒕_\mathit{0})\hfill \\ \hfill =\underset{i=1}{\overset{d}{}}\left[F\varphi ^1\left(𝒕_\mathit{0}+\underset{𝒋=\mathit{1}}{\overset{𝒊}{}}𝒉_𝒋𝒆_𝒋\right)𝑭\varphi ^\mathit{1}\left(𝒕_\mathit{0}+\underset{𝒋=\mathit{1}}{\overset{𝒊\mathit{1}}{}}𝒉_𝒋𝒆_𝒋\right)\right]\\ \hfill =\underset{i=1}{\overset{d}{}}h_iF_\varphi ^i(m_0)+\underset{i=1}{\overset{d}{}}_0^1𝑑\vartheta h_i\left[F_\varphi ^i\varphi ^1\left(𝒕_\mathit{0}+\underset{𝒋=\mathit{1}}{\overset{𝒊\mathit{1}}{}}𝒉_𝒋𝒆_𝒋+\vartheta 𝒉_𝒊𝒆_𝒊\right)𝑭_\varphi ^𝒊(𝒎_\mathit{0})\right]\text{.}\end{array}$$
(A.13)
Due to continuity of the mappings $`F_\varphi ^i`$, the second term on the right-hand side multiplied with $`|𝒉|^\mathit{1}`$ can be estimated by
$$\begin{array}{c}|𝒉|^\mathit{1}𝒒_\lambda \left(\underset{𝒊=\mathit{1}}{\overset{𝒅}{}}_\mathit{0}^\mathit{1}𝒅\vartheta 𝒉_𝒊\left[𝑭_\varphi ^𝒊\varphi ^\mathit{1}\left(𝒕_\mathit{0}+\underset{𝒋=\mathit{1}}{\overset{𝒊\mathit{1}}{}}𝒉_𝒋𝒆_𝒋+\vartheta 𝒉_𝒊𝒆_𝒊\right)𝑭_\varphi ^𝒊(𝒎_\mathit{0})\right]\right)\hfill \\ \hfill |𝒉|^\mathit{1}\underset{𝒊=\mathit{1}}{\overset{𝒅}{}}|𝒉_𝒊|\underset{\mathit{0}\vartheta \mathit{1}}{\mathrm{max}}𝒒_\lambda \left(𝑭_\varphi ^𝒊\varphi ^\mathit{1}\left(𝒕_\mathit{0}+\underset{𝒋=\mathit{1}}{\overset{𝒊\mathit{1}}{}}𝒉_𝒋𝒆_𝒋+\vartheta 𝒉_𝒊𝒆_𝒊\right)𝑭_\varphi ^𝒊(𝒎_\mathit{0})\right)\\ \hfill \underset{i=1}{\overset{d}{}}\underset{0\vartheta 1}{\mathrm{max}}q_\lambda \left(F_\varphi ^i\varphi ^1\left(𝒕_\mathit{0}+\underset{𝒋=\mathit{1}}{\overset{𝒊\mathit{1}}{}}𝒉_𝒋𝒆_𝒋+\vartheta 𝒉_𝒊𝒆_𝒊\right)𝑭_\varphi ^𝒊(𝒎_\mathit{0})\right)\text{,}\end{array}$$
(A.14)
where the last expression of the above inequality is seen to vanish in the limit $`𝒉\mathit{0}`$ by assumption. Thus (A.13) in connection with (A.14) establishes continuous differentiability of the mapping $`F:𝒰𝔙`$ with
$$𝔇_\varphi F(m_0)𝒉=\underset{𝒊=\mathit{1}}{\overset{𝒅}{}}𝒉_𝒊𝑭_\varphi ^𝒊(𝒎_\mathit{0})\text{,}𝒉𝕂^𝒅\text{.}$$
###### Proposition A.10.
Assume that the mixed derivatives $`F_\varphi ^{ij}\left(F_\varphi ^i\right)_\varphi ^j`$ and $`F_\varphi ^{ji}\left(F_\varphi ^j\right)_\varphi ^i`$, $`i\text{,}j\{1,\mathrm{},d\}`$ of the mapping $`F:𝒰𝔙`$ exist and are continuous on $`𝒰`$. Then they coincide:
$$F_\varphi ^{ij}(m_0)=F_\varphi ^{ji}(m_0)\text{,}m_0𝒰\text{.}$$
###### Proof.
For $`m_0=\varphi ^1(𝒕_\mathit{0})𝒰`$ and sufficiently small $`h\text{,}k𝕂`$ consider the following expression which involves two increments of $`F\varphi ^1`$:
$$F\varphi ^1(𝒕_\mathit{0}+𝒉𝒆_𝒋+𝒌𝒆_𝒊)𝑭\varphi ^\mathit{1}(𝒕_\mathit{0}+𝒉𝒆_𝒋)𝑭\varphi ^\mathit{1}(𝒕_\mathit{0}+𝒌𝒆_𝒊)+𝑭\varphi ^\mathit{1}(𝒕_\mathit{0})\text{.}$$
By assumption on the existence and continuity of the mixed derivatives we can apply the Mean Value Theorem A.7 twice to the above expression: One can consider the increments with respect to $`𝒆_𝒊`$ and apply the Mean Value Theorem to them first and afterwards to the resulting integrand which takes on the form of an increment with respect to $`𝒆_𝒋`$, or one carries out the same procedure with the roles of $`𝒆_𝒊`$ and $`𝒆_𝒋`$ interchanged. Upon division by $`hk0`$ this yields the integrals
$$\begin{array}{c}_0^1𝑑\vartheta _0^1𝑑\vartheta ^{}F_\varphi ^{ij}\varphi ^1(𝒕_\mathit{0}+\vartheta 𝒉𝒆_𝒋+\vartheta ^{}𝒌𝒆_𝒊)\text{,}\hfill \\ \hfill =_0^1𝑑\vartheta _0^1𝑑\vartheta ^{}F_\varphi ^{ji}\varphi ^1(𝒕_\mathit{0}+\vartheta 𝒉𝒆_𝒋+\vartheta ^{}𝒌𝒆_𝒊)\text{,}\end{array}$$
for any $`h\text{,}k𝕂\{0\}`$. Specializing to sequences $`\{h_n\}_n`$ and $`\{k_n\}_n`$ in this set which converge to $`0`$, it is a consequence of Lebesgue’s Dominated Convergence Theorem (cf. \[26, II.5.6 and II.6.2\]) that for $`n\mathrm{}`$ the left-hand side converges to $`F_\varphi ^{ij}(m_0)`$ whereas the right-hand side approaches $`F_\varphi ^{ji}(m_0)`$ in the locally convex topology of $`𝔙`$. Since this topology separates the elements of $`𝔙`$, we conclude that these limits coincide and get the assertion by arbitrariness of $`m_0𝒰`$. ∎
### A.3 Differentiation on Automorphism Lie Groups
The concepts developed thus far can now be applied to the case where the underlying analytic manifold is a (real or complex) Lie group $`𝒢`$ acting via a strongly continuous group of automorphisms $`\{\alpha _g:g𝒢\}\mathrm{Aut}𝔅`$ on the $`C^{}`$-algebra $`𝔅`$. These automorphisms $`\alpha _g`$, when applied to a given element $`B𝔅`$, define a $`𝔅`$-valued mapping on $`𝒢`$, for which statements can be proved that go beyond the above results. In doing so we shall be concerned with the canonical coordinates $`(𝒰_0,\varphi _0)`$ of the first kind around the neutral element $`1`$ of $`𝒢`$ where $`1/\varphi _0^\mathrm{}1\mathrm{}\mathit{0}\mathrm{}`$ (cf. \[55, Section 2.10\]). Note also, that, for given $`g𝒢`$, the left and right translations $`l_g`$ and $`r_g`$ on $`𝒢`$ as well as their composition $`i_g=l_gr_{g^1}`$ are analytic diffeomorphisms, so that their application to $`(𝒰_0,\varphi _0)`$ yields local charts around $`g`$ and $`1`$, respectively (cf. \[55, Section 2.1\]).
###### Proposition A.11.
Let $`𝒢`$ be a $`d`$-dimensional real or complex Lie group and let $`𝔅`$ be a $`C^{}`$-algebra. For given $`B𝔅`$ define the mapping
$$\mathrm{\Xi }_B:𝒢𝔅g\mathrm{\Xi }_B(g)\alpha _g(B)\text{.}$$
* $`\mathrm{\Xi }_B`$ is continuously differentiable on $`𝒢`$ if and only if it is differentiable at $`1𝒢`$.
* If $`\mathrm{\Xi }_B`$ is differentiable at $`1𝒢`$, then $`\mathrm{\Xi }_{\alpha _g^{}(B)}`$ is differentiable for any $`g^{}𝒢`$ and the mapping
$$𝒢\times 𝒰(g^{},g)𝔇_\varphi \mathrm{\Xi }_{\alpha _g^{}(B)}(g)𝒉$$
is jointly continuous in $`g^{}`$ and $`g`$ for any local chart $`(𝒰,\varphi )`$ around $`g`$ and any $`𝒉𝕂^𝒅`$.
###### Proof.
* To prove the non-trivial part, suppose that $`g𝒢`$ is arbitrary but fixed. Then $`(g𝒰_0,\varphi _g)`$, $`\varphi _g\varphi _0l_{g^1}`$, is a local chart around $`g`$ with $`\varphi _g^1=l_g\varphi _0^1`$. According to the definition of $`\mathrm{\Xi }_B`$ we have
$$\mathrm{\Xi }_B\varphi _g^1=\mathrm{\Xi }_Bl_g\varphi _0^1=\alpha _g\mathrm{\Xi }_B\varphi _0^1$$
and, since the automorphisms are norm-preserving, the assumed differentiability of the mapping $`\mathrm{\Xi }_B\varphi _0^1`$ at $`\mathit{0}`$ carries over to $`\mathrm{\Xi }_B\varphi _g^1`$ which by Definition A.8 means that $`\mathrm{\Xi }_B`$ is differentiable at $`g=\varphi _g^1(\mathit{0})`$:
$$𝔇_{\varphi _g}\mathrm{\Xi }_B(g)=𝔇(\mathrm{\Xi }_B\varphi _g^1)(\mathit{0})=\alpha _𝒈𝔇(𝜩_𝑩\varphi _\mathit{0}^\mathit{1})(\mathit{0})=\alpha _𝒈𝔇_{\varphi _\mathit{0}}𝜩_𝑩(1\mathrm{}\text{.}$$
In view of (A.11) this relation can be re-written with respect to an arbitrary local chart $`(𝒰,\varphi )`$ on $`𝒢`$ containing $`g`$:
$$𝔇_\varphi \mathrm{\Xi }_B(g)=𝔇_{\varphi _g}\mathrm{\Xi }_B(g)(\varphi _g\varphi ^1)^{}\left(\varphi (g)\right)=\alpha _g𝔇_{\varphi _0}\mathrm{\Xi }_B(1\mathrm{}𝐌^\varphi \mathrm{}\mathrm{g}\mathrm{}\text{,}$$
(A.15)
where the matrix elements of $`𝐌^\varphi (g)\left(\varphi _g\varphi ^1\right)^{}\left(\varphi (g)\right)`$ are analytic in $`g𝒰`$. Since the automorphisms are norm-preserving and act stongly continuous on $`𝔅`$, it is evident that application of the above operator to any vector $`𝒉𝕂^𝒅`$ yields a continuous mapping on $`𝕂^d`$ to $`𝔅`$, thus establishing continuous differentiability of $`\mathrm{\Xi }_B`$ on $`𝒢`$ as stated.
* Let $`g^{}𝒢`$ be arbitrary and consider the local chart $`(𝒰_0g^{},\psi _g^{})`$, $`\psi _g^{}\varphi _0r_{g_{}^{}{}_{}{}^{1}}`$, around $`g^{}`$ with inverse $`\psi _g^{}^1=r_g^{}\varphi _0^1`$. Then
$$\mathrm{\Xi }_{\alpha _g^{}(B)}\varphi _0^1=\mathrm{\Xi }_Br_g^{}\varphi _0^1=\mathrm{\Xi }_B\psi _g^{}^1\text{,}$$
so that the assumed differentiablity of $`\mathrm{\Xi }_B`$ at $`1`$ and thus, according to the first part, at $`g^{}`$ with respect to the local chart $`(𝒰_0g^{},\psi _g^{})`$ implies differentiability of $`\mathrm{\Xi }_{\alpha _g^{}(B)}`$ at $`1`$. By an application of (A.15) we have
$$𝔇_{\varphi _0}\mathrm{\Xi }_{\alpha _g^{}(B)}(1\mathrm{}/𝔇_{\psi _\mathrm{g}^{}}\mathrm{\Xi }_\mathrm{B}\mathrm{}\mathrm{g}^{}\mathrm{}/\alpha _\mathrm{g}^{}𝔇_{\varphi _0}\mathrm{\Xi }_\mathrm{B}\mathrm{}1\mathrm{}𝐍\mathrm{}\mathrm{g}^{}\mathrm{}\text{,}$$
(A.16)
where the matrix elements of $`𝐍(g^{})\left(\varphi _g\psi _g^{}^1\right)^{}(\mathit{0})`$ are analytic in $`g^{}`$. This in turn can, again by use of (A.15), be generalized to any $`g𝒢`$ lying in the local chart $`(𝒰,\varphi )`$:
$$𝔇_\varphi \mathrm{\Xi }_{\alpha _g^{}(B)}(g)=\alpha _g𝔇_{\varphi _0}\mathrm{\Xi }_{\alpha _g^{}(B)}(1\mathrm{}𝐌^\varphi \mathrm{}\mathrm{g}\mathrm{}/\alpha _{\mathrm{g}\mathrm{g}^{}}𝔇_{\varphi _0}\mathrm{\Xi }_\mathrm{B}\mathrm{}1\mathrm{}𝐍\mathrm{}\mathrm{g}^{}\mathrm{}𝐌^\varphi \mathrm{}\mathrm{g}\mathrm{}\text{,}$$
(A.17)
an expression which is obviously continuous in both variables $`g^{}`$ and $`g`$ when applied to an arbitrary element $`𝒉`$ of $`𝕂^d`$. ∎
###### Remark.
Note, that in the case of differentiability of $`\mathrm{\Xi }_B`$ the mapping $`g𝔇_\varphi \mathrm{\Xi }_B(g)`$ need not be continuous in the operator-norm topology of the Banach space of linear operators on $`𝕂^d`$ to $`𝔅`$, since the automorphism group $`\{\alpha _g:g𝒢\}\mathrm{Aut}𝔅`$ is only supposed to be strongly continuous.
Consider those operators $`B𝔅`$ for which the mapping $`\mathrm{\Xi }_B`$ is continuously differentiable on $`𝒢`$. According to Proposition A.11 this is equivalent to differentiability at $`1`$ with respect to the canonical coordinates $`(𝒰_0,\varphi _0)`$. Therefore one can define mappings $`\delta ^i`$ corresponding to the partial derivatives of $`\mathrm{\Xi }_B`$ at $`1`$ (cf. (A.12)) by
$$\delta ^i(B)𝔇_{\varphi _0}\mathrm{\Xi }_B(1\mathrm{}𝒆_𝒊\text{,}𝒊/\mathit{1}\text{,}\mathrm{}\text{,}𝒅\text{.}$$
Since $`\mathrm{\Xi }_B`$ depends linearly on $`B`$, it is easily seen that
$`\delta ^i(B_1+B_2)`$ $`=\delta ^i(B_1)+\delta ^i(B_2)\text{,}`$ (A.18a)
$`\delta ^i(\lambda B_1)`$ $`=\lambda \delta ^i(B_1)\text{,}`$ (A.18b)
for any $`\lambda 𝕂`$ and $`B_1`$, $`B_2`$ in $`𝔅`$ subject to Proposition A.11. Moreover, $`\mathrm{\Xi }_{B_1B_2}=\mathrm{\Xi }_{B_1}\mathrm{\Xi }_{B_2}`$, so that Proposition A.3 yields
$$\delta ^i(B_1B_2)=\delta ^i(B_1)B_2+B_1\delta ^i(B_2)\text{.}$$
(A.18c)
Equations (A.18a) through (A.18c) show that the mappings $`\delta ^i`$ act as derivations of the $`C^{}`$-algebra $`𝔅`$ (cf. \[24, Chapter III.9\] and \[48, Section 8.6\]). Their domains are certain subalgebras which are invariant under transformations from the automorphism group $`\{\alpha _g:g𝒢\}`$, since by (A.16) for any $`g^{}𝒢`$ and any $`B𝔅`$ with differentiable $`\mathrm{\Xi }_B`$ one has
$$\delta ^i\left(\alpha _g^{}(B)\right)=𝔇_{\varphi _0}\mathrm{\Xi }_{\alpha _g^{}(B)}(1\mathrm{}𝒆_𝒊/\alpha _𝒈^{}\mathrm{}𝔇_{\varphi _\mathit{0}}𝜩_𝑩\mathrm{}1\mathrm{}𝐍\mathrm{}\mathrm{g}^{}\mathrm{}𝒆_𝒊\mathrm{}/\underset{𝒋/\mathit{1}}{\overset{𝒅}{}}𝐍_{\mathrm{𝒋𝒊}}\mathrm{}𝒈^{}\mathrm{}\alpha _𝒈^{}\mathrm{}\delta ^𝒋\mathrm{}𝑩\mathrm{}\mathrm{}\text{.}$$
(A.19)
Let $`\iota _M`$ denote the $`M`$-tuple $`(i_1,\mathrm{},i_M)`$ with integer entries $`1i_ld`$, then the corresponding products of derivations $`\delta ^{\iota _M}\delta ^{i_M}\mathrm{}\delta ^{i_1}`$ act as linear operators on certain subspaces of $`𝔅`$ which are again invariant with respect to $`\{\alpha _g:g𝒢\}`$, possibly the trivial space $`\{0\}`$ (note, that in general the derivations will not commute). Making use of the concepts of differentiability introduced above together with the fact that left and right translations act as analytic diffeomorphisms on the group $`𝒢`$, it is a matter of elementary considerations to establish the following connection between products $`\delta ^{\iota _M}`$ of the above kind and the partial derivatives of the mapping $`\mathrm{\Xi }_B`$ indexed by multi-indices $`\kappa `$:
$`\mathrm{\Xi }_{B,\varphi }^\kappa (g)`$ $`={\displaystyle \underset{\iota _M,M|\kappa |}{}}C_{\kappa ,\iota _M}^\varphi (g)\alpha _g\left(\delta ^{\iota _M}(B)\right)\text{,}`$ (A.20a)
$`\delta ^{\iota _M}(B)`$ $`={\displaystyle \underset{\kappa ,|\kappa |M}{}}D_{\iota _M,\kappa }^{\varphi _0}(1\mathrm{}\mathrm{\Xi }_{\mathrm{B}\mathrm{}\varphi _0}^\kappa \mathrm{}1\mathrm{}\text{.}`$ (A.20b)
Here the real or complex functions $`C_{\kappa ,\iota _M}^\varphi `$ and $`D_{\iota _M,\kappa }^{\varphi _0}`$ are analytic on the respective charts $`(𝒰,\varphi )`$ and $`(𝒰_0,\varphi _0)`$, containing $`g`$ and $`1`$ respectively. Implicit in (A.20a) and (A.20b) is the perception that the mapping $`\mathrm{\Xi }_B`$ is $`N`$-fold (or infinitely often) continuously differentiable if and only if the operator $`B`$ belongs to the domain of all $`\delta ^{\iota _M}`$ for $`MN`$ (or any $`M<\mathrm{}`$).
We formulate these results in the following definition and subsequent proposition.
###### Definition A.12.
Let $`\delta ^i`$, $`i=1\text{,}\mathrm{}\text{,}d`$, denote the partial derivations pertaining to the mappings $`𝒢g\mathrm{\Xi }_B(g)=\alpha _g(B)𝔅`$ for certain $`B𝔅`$ via
$$\delta ^i(B)𝔇_{\varphi _0}\mathrm{\Xi }_B(1\mathrm{}𝒆_𝒊\text{.}$$
(A.21)
For given $`N`$ the domain of arbitrary $`N`$-fold products $`\delta ^{\iota _N}`$ of these derivations is an invariant subspace of $`𝔅`$ with respect to the automorphism group $`\{\alpha _g:g𝒢\}`$ and denoted $`𝒟^{(N)}(𝔅)`$: the space of $`N`$-fold differentiable operators. The elements of the space $`𝒟^{(\mathrm{})}(𝔅)_N𝒟^{(N)}(𝔅)`$ in turn are called infinitely often differentiable with respect to $`\{\alpha _g:g𝒢\}`$. Accordingly, the resulting operators $`\delta ^{\iota _N}(B)`$ are designated as the derivatives of $`B𝔅`$, if this element happens to lie in their domain.
###### Proposition A.13.
For given $`B𝔅`$ the mapping
$$\mathrm{\Xi }_B:𝒢𝔅g\mathrm{\Xi }_B(g)\alpha _g(B)$$
is $`N`$-fold or infinitely often continuously differentiable if and only if the operator $`B`$ belongs to $`𝒟^{(N)}(𝔅)`$ respectively $`𝒟^{(\mathrm{})}(𝔅)`$.
### A.4 Differentiable Linear Mappings
In this section a special notion of differentiability for linear mappings on a locally convex space $`𝔙`$ is introduced, which is motivated by the following result that is valid under the assumption of continuity.
###### Proposition A.14.
Let $`𝔛`$ be a (real or complex) normed space, $`𝔊𝔛`$ open, and let $`𝔙`$ and $`𝔚`$ be locally convex spaces over the same field $`𝕂`$ or $`𝕂`$ with topologies defined by the families $`\{q_\lambda :\lambda L\}`$ and $`\{q_\mu ^{}:\mu M\}`$ of seminorms separating the points of $`𝔙`$ and $`𝔚`$, respectively. If $`F:𝔊𝔙`$ is differentiable at the point $`𝔵𝔊`$ and $`\mathrm{\Psi }:𝔙𝔚`$ is a continuous linear mapping then the composition
$$\mathrm{\Psi }F:𝔊𝔚$$
is differentiable at $`𝔵`$, too, and its derivative is given by
$$𝔇(\mathrm{\Psi }F)(𝔵)=\mathrm{\Psi }𝔇F(𝔵)\text{.}$$
(A.22)
If $`F`$ is differentiable on all of $`𝔊`$ the same holds true for $`\mathrm{\Psi }F`$ and (A.22) is valid for any $`𝔵𝔊`$.
###### Proof.
By assumption on $`F`$ (relations (A.1a) and (A.1b)) in connection with linearity of $`\mathrm{\Psi }`$, the increment of $`\mathrm{\Psi }F`$ at $`𝔵`$ allows for the representation
$$(\mathrm{\Psi }F)(𝔵+𝔥)(\mathrm{\Psi }F)(𝔵)=\mathrm{\Psi }𝔇F(𝔵)𝔥+\mathrm{\Psi }\left(R[F,𝔵](𝔥)\right)\text{,}$$
(A.23)
where
$$\underset{𝔥\mathfrak{0}}{lim}𝔥^1q_\lambda \left(R[F,𝔵](𝔥)\right)=0$$
for any seminorm $`q_\lambda `$, $`\lambda L`$. But, due to continuity of $`\mathrm{\Psi }`$, there exist to any seminorm $`q_\mu ^{}`$ on $`𝔚`$ a finite number of seminorms $`q_{\lambda _i}`$ on $`𝔙`$, $`i=1\text{,}\mathrm{}\text{,}N`$, and a positive constant $`C_\mu `$ such that for any $`𝔳𝔙`$
$$q_\mu ^{}\left(\mathrm{\Psi }(𝔳)\right)C_\mu \underset{1iN}{\mathrm{max}}q_{\lambda _i}(𝔳)\text{,}$$
and therefore
$$0𝔥^1q_\mu ^{}\left(\mathrm{\Psi }\left(R[F,𝔵](𝔥)\right)\right)C_\mu \underset{1iN}{\mathrm{max}}\left(𝔥^1q_{\lambda _i}\left(R[F,𝔵](𝔥)\right)\right)\underset{𝔥\mathfrak{0}}{\overset{}{}}0\text{.}$$
This is just the formulation of (A.1b) for $`\mathrm{\Psi }F`$ and thus proves, according to (A.23), differentiability of this mapping at $`𝔵`$ together with (A.22). The remainder of the assertion is a trivial consequence. ∎
The above results can easily be generalized to $`𝔙`$-valued mappings on an analytic manifold $``$.
###### Corollary A.15.
Let $``$ be a (real or complex) analytic manifold of dimension $`d`$ and let $`𝔙`$ and $`𝔚`$ be locally convex spaces over the same field. If $`F:𝒰𝔙`$ is differentiable at the point $`m_0𝒰`$, $`(𝒰,\varphi )`$ a local chart on $``$, and $`\mathrm{\Psi }:𝔙𝔚`$ is a continuous linear mapping then
$$\mathrm{\Psi }F:𝒰𝔚$$
is differentiable at $`m_0`$, and its derivative is given by
$$𝔇_\varphi (\mathrm{\Psi }F)(m_0)=\mathrm{\Psi }𝔇_\varphi F(m_0)\text{.}$$
(A.24)
Accordingly, $`\mathrm{\Psi }F`$ is differentiable on all of $``$ in case that $`F`$ is.
Proposition A.14 motivates the following definition which does no longer depend on the assumption of continuity.
###### Definition A.16.
* Let $`𝔛`$ be a normed space and $``$ a family of differentiable mappings on $`𝔛`$ with values in a locally convex space $`𝔙`$. A linear mapping $`\mathrm{\Psi }`$ on $`𝔙`$ to the locally convex space $`𝔚`$ is called $``$-differentiable if and only if $`\mathrm{\Psi }F:𝔛𝔚`$ is differentiable on $`𝔛`$ for any $`F`$ with
$$𝔇(\mathrm{\Psi }F)(𝔵)=\mathrm{\Psi }𝔇F(𝔵)\text{,}𝔵𝔛\text{.}$$
* Let $``$ be an analytic manifold and let $`𝔙`$, $`𝔚`$ and $`\mathrm{\Psi }`$ be as above. Assume furthermore that $``$ is a family of differentiable $`𝔙`$-valued mappings on $``$. Then $`\mathrm{\Psi }`$ is called $``$-differentiable if and only if $`\mathrm{\Psi }F:𝔚`$ is differentiable on $``$ for any $`F`$ and
$$𝔇_\varphi (\mathrm{\Psi }F)(m_0)=\mathrm{\Psi }𝔇_\varphi F(m_0)$$
for any chart $`(𝒰,\varphi )`$ around the arbitrary element $`m_0`$.
## Appendix B A Lemma on Norm-Separable $`𝑪^{}`$-Algebras
The following result is an adaptation of \[43, Lemma 14.1.17\] to our needs.
###### Lemma B.1.
Let $`𝔄`$ be a unital $`C^{}`$-subalgebra of $`𝔅()`$, where the Hilbert space $``$ is separable. There exists a norm-separable $`C^{}`$-algebra $`𝔄^0`$, containing the unit element $`\mathrm{𝟏}`$, that lies strongly dense in $`𝔄`$.
###### Proof.
Let $`\left\{\varphi _n\right\}_n`$ be a dense sequence of non-zero vectors in $``$ and let $`𝔐𝔄^{\prime \prime }`$ denote the von Neumann algebra generated by $`𝔄`$. According to von Neumann’s Density Theorem, $`𝔐`$ coincides with the strong closure $`𝔄^{}`$ of the algebra $`𝔄`$, which by assumption acts non-degenerately on $``$ (cf. \[24, Section I.3.4\], \[11, Corollary 2.4.15\]).
First we assume the existence of a separating vector for $`𝔐`$, which is thus cyclic for $`𝔐^{}`$ \[24, Section I.1.4\]. Then any normal functional on $`𝔐`$ is of the form $`\omega _{\psi ,\psi ^{}}𝔐`$ with $`\psi \text{,}\psi ^{}`$ \[54, Theorem V.3.15\]. Choose operators $`A_{j,k}𝔄_1`$ satisfying
$$\omega _{\varphi _j,\varphi _k}(A_{j,k})\omega _{\varphi _j,\varphi _k}𝔐2^1\text{,}$$
(B.1)
which is possible due to Kaplansky’s Density Theorem \[48, Theorem 2.3.3\]. Let $`𝔄^0`$ denote the norm-separable $`C^{}`$-algebra generated by the unit element $`\mathrm{𝟏}`$ together with all the operators $`A_{j,k}`$, $`j\text{,}k`$, and select a normal functional $`\omega _{\xi ,\theta }`$ on $`𝔐`$ with the properties $`\omega _{\xi ,\theta }𝔄^0=0`$ and $`\omega _{\xi ,\theta }𝔐>0`$. Without loss of generality we can assume $`\omega _{\xi ,\theta }𝔐=1`$. To any $`ϵ>0`$ there exist vectors $`\varphi _j\text{,}\varphi _k`$ from the dense sequence in $``$ rendering $`\varphi _j\xi `$ and $`\varphi _k\theta `$ small enough so that
$$(\omega _{\xi ,\theta }\omega _{\varphi _j,\varphi _k})𝔐<ϵ\text{.}$$
(B.2)
Making use of (B.1) we then get the estimate
$$\begin{array}{c}ϵ>(\omega _{\xi ,\theta }\omega _{\varphi _j,\varphi _k})𝔐(\omega _{\xi ,\theta }\omega _{\varphi _j,\varphi _k})(A_{j,k})\hfill \\ \hfill =\omega _{\varphi _j,\varphi _k}(A_{j,k})\omega _{\varphi _j,\varphi _k}𝔐2^1\text{,}\end{array}$$
which in connection with (B.2) implies
$$\omega _{\xi ,\theta }𝔐(\omega _{\xi ,\theta }\omega _{\varphi _j,\varphi _k})𝔐+\omega _{\varphi _j,\varphi _k}𝔐<2ϵ+2^1\text{.}$$
By arbitraryness of $`ϵ`$ we infer $`\omega _{\xi ,\theta }𝔐2^1`$ in contradiction to the assumption that $`\omega _{\xi ,\theta }𝔐`$ be normalized. Thus, $`\omega _{\xi ,\theta }𝔄^0=0`$ implies $`\omega _{\xi ,\theta }𝔐=0`$, i. e. any normal functional on $`𝔐`$ annihilating $`𝔄^0`$ annihilates $`𝔐`$ as well. Now, since the $`C^{}`$-algebra $`𝔄^0`$ acts non-degenerately on $``$, von Neumann’s Density Theorem tells us that its strong and $`\sigma `$-weak closures coincide with $`𝔄_{}^{0}{}_{}{}^{\prime \prime }=𝔄_{}^{0}{}_{}{}^{}`$, and this in turn is equal to the von Neumann algebra $`𝔐`$; for the existence of an element $`A𝔐`$ not contained in $`𝔄_{}^{0}{}_{}{}^{}`$ would, by the Hahn-Banach-Theorem, imply a $`\sigma `$-weakly continuous (normal) functional that vanishes on $`𝔄^0`$ but not on $`A𝔐𝔄_{}^{0}{}_{}{}^{}`$ in contradiction to the above result.
Now suppose that there does not exist a separating vector for the von Neumann algebra $`𝔐=𝔄^{}`$. Then the sequence
$$\left((n\varphi _n)^1\varphi _n\right)_n\underset{¯}{}\underset{n=1}{\overset{\mathrm{}}{}}$$
is such a vector for the von Neumann algebra $`\underset{¯}{𝔐}\left(_{n=1}^{\mathrm{}}\iota \right)(𝔐)`$, where $`\iota `$ denotes the identity representation of $`𝔐`$ in $``$. The result of the preceding paragraph thus applies to the $`C^{}`$-algebra $`\underset{¯}{𝔄}\left(_{n=1}^{\mathrm{}}\iota \right)(𝔄)`$ of operators on the separable Hilbert space $`\underset{¯}{}`$ which is weakly dense in $`\underset{¯}{𝔐}`$: $`\underset{¯}{𝔄}^{}=\underset{¯}{𝔐}`$. We infer that there exists a norm-separable $`C^{}`$-subalgebra $`\underset{¯}{𝔄}^0`$ of $`\underset{¯}{𝔄}`$ including its unit $`\underset{¯}{\mathrm{𝟏}}(\mathrm{𝟏})_n`$, which is strongly dense in $`\underset{¯}{𝔄}`$. Now, $`\underset{¯}{\iota }_{n=1}^{\mathrm{}}\iota `$ is a faithful representation of $`𝔄`$ on $`\underset{¯}{}`$ and its inverse $`\underset{¯}{\iota }^1:\underset{¯}{𝔄}𝔄`$ is a faithful representation of $`\underset{¯}{𝔄}`$ on $``$ which is continuous with respect to the strong topologies of $`\underset{¯}{𝔄}`$ and $`𝔄`$. Therefore $`𝔄^0\underset{¯}{\iota }^1\left(\underset{¯}{𝔄}^0\right)`$ is a norm-separable $`C^{}`$-subalgebra of $`𝔄`$, containing the unit element $`\mathrm{𝟏}`$ and lying strongly dense in $`𝔄`$. ∎
## Acknowledgements
My sincere gratitude is due to Prof. Dr. Detlev Buchholz for his support and forbearance. I not only learned a lot from his views on theoretical physics and its relationship to mathematics, but his ideas also constituted the basis on which I could erect my contributions to the topic presented in this thesis.
I should furthermore like to thank Prof. Dr. Klaus Fredenhagen for his immediate readiness to write the additional report on this work.
Financial support by the Deutsche Forschungsgemeinschaft is gratefully acknowledged which I obtained from the Graduiertenkolleg at the II. Institut für Theoretische Physik of the University of Hamburg.
I want to thank AnnA for all her support and encouragement, and for unrepiningly sharing the burden of my strain in the final stages of this project.
Eventually, I want to express my deepest gratitude to my parents for their patience and confidence. I am afraid that I shall be unable to pass back even only part of what they have done for me.
German translation of the quotation on page 4
Diodoros: Griechische Weltgeschichte IV, 59 (5)
(nach der Übersetzung von Otto Veh)
Theseus beseitigte auch bei Eleusis den Kerkyon, der die Passanten zum Ringkampf veranlaßte und den, der unterlag, umbrachte. Sodann tötete er auch den Prokrustes, wie er hieß, der am sogenannten Korydallos in Attika hauste. Der nötigte die vorüberziehenden Wanderer, sich auf ein Bett niederzulegen und war einer zu lang, dann schlug er ihm die herausragenden Körperteile ab; denen aber, die kleiner waren, zog er die Füße in die Länge, weshalb er den Namen Prokrustes erhielt.
|
warning/0005/cond-mat0005131.html
|
ar5iv
|
text
|
# Fluctuation and relaxation properties of pulled fronts: a possible scenario for non-Kardar-Parisi-Zhang behavior
## Abstract
We argue that while fluctuating fronts propagating into an unstable state should be in the standard KPZ universality class when they are pushed, they should not when they are pulled: The universal $`1/t`$ velocity relaxation of deterministic pulled fronts makes it unlikely that the KPZ equation is the appropriate effective long-wavelength low-frequency theory in this regime. Simulations in 2$`D`$ confirm the proposed scenario, and yield exponents $`\beta 0.29\pm 0.01`$, $`\zeta 0.40\pm 0.02`$ for fluctuating pulled fronts, instead of the KPZ values $`\beta =1/3`$, $`\zeta =1/2`$. Our value of $`\beta `$ is consistent with an earlier result of Riordan et al.
Over a decade ago, Kardar, Parisi and Zhang (KPZ) introduced their celebrated stochastic equation
$`{\displaystyle \frac{h}{t}}`$ $`=`$ $`\nu ^2h+{\displaystyle \frac{\lambda }{2}}(h)^2+\eta ,`$ (1)
$`\eta (𝐱,t)\eta (𝐱^{},t^{})`$ $`=`$ $`D\delta ^d(𝐱𝐱^{})\delta (tt^{}),`$ (2)
to describe the fluctuation properties of growing interfaces with height $`h`$ under the influence of the noise term $`\eta `$. A clear “derivation” of the KPZ equation is difficult to give, just as much as the Landau-Ginzburg-Wilson Hamiltonian can not straightforwardly be “derived” from the Ising model. However, one expects the KPZ equation to be the proper effective long wavelength low frequency theory for interfacial growth phenomena whose deterministic macroscopic evolution equation is of the form
$$\frac{h}{t}=v(h)+\text{curvature corrections}.$$
(3)
Here $`v(h)`$ is the deterministic growth velocity of a planar interface as a function of the orientation $`h`$. For, as long as the curvature corrections of the form $`^2h`$ are nonzero, the long wavelength expansion of (3) immediately yields the gradient term in (1). The philosophy is then that in the presence of noise, all the relevant terms in the KPZ equation (1) are generated, and that this is sufficient to yield the asymptotic KPZ scaling. In agreement with this picture, many interface growth models have been found to show the universal asymptotic scaling properties predicted by (1).
A dynamical interface equation of the form (3) is appropriate for interfaces whose long wavelength and slow time dynamics is essentially local in space and time, i.e., dependent on the local and instantaneous values of the slope and curvature. The applicability of the KPZ equation is therefore not limited to situations with a microscopically sharp interface: Many pattern forming systems of the reaction-diffusion type exhibit fronts whose intrinsic width $`l`$ is finite. For curvatures $`\kappa `$ small compared to $`l^1`$, $`\kappa l1`$, an effective interface approximation or moving boundary approximation of the form (3) can then be derived using standard techniques . These approximations apply whenever the internal stability modes of the fronts relax exponentially on a short time scale, so that an adiabatic decoupling becomes exact in the limit $`\kappa l0`$. The best known example of such a type of analysis is for the curvature driven growth in the Cahn-Hilliard equation, but moving boundary techniques have recently been applied successfully to many other such problems . In all these cases, the internal relaxation modes within the fronts or transition zones are indeed exponentially decaying on a short time scale.
From the above perspective, recent results for the relaxation properties of planar fronts propagating into an unstable state, suggest an interesting new scenario for non-KPZ behavior. Fronts propagating into unstable states generally come in two classes, so-called pushed fronts and pulled fronts . Pushed fronts propagating into an unstable state are the immediate analogue of fronts between two linearly stable states. In the thin interface limit, $`\kappa l1`$, the dynamics of such fronts becomes essentially local and instantaneous, and given by an equation of the form (3); according to the arguments given above, fluctuating pushed fronts should thus obey KPZ scaling: following standard practice by saying that the $`d`$+1D KPZ equation (where the +1 refers to the time dimension) describes the fluctuations of a $`d`$-dimensional interface, the conclusion is that fluctuations of $`d`$-dimensional pushed fronts in $`(d`$+$`1)`$ bulk dimensions are described by the $`d`$+1D KPZ equation.
Pulled fronts, however, behave very differently from pushed ones. A pulled front propagating into a linearly unstable state is basically “pulled along” by the linear growth dynamics of small perturbations spreading into the linearly unstable state. The crucial new insight for our discussion is the recent finding that pulled fronts can not be described by an effective interface equation like (3) that is local and instantaneous in space and time, even if they are weakly curved on the spatial scale. This just reflects the fact that the dynamically important region of pulled fronts is the semi-infinite leading edge region ahead of the front, not the nonlinear front region itself. Technically, the breakdown of an interfacial description is seen from the divergence of the solvability type integrals that arise in the derivation of a moving boundary approximation in dimensions $`d`$$``$$`2`$ . More intuitively, the result can be understood as follows: a deterministic pulled front in $`d`$=$`1`$ relaxes to its asymptotic speed $`v^{}`$ with a universal power law as $`v(t)=v^{}+c_1/t+c_{3/2}/t^{3/2}+\mathrm{}`$, where $`c_1(<0)`$ and $`c_{3/2}`$ are known coefficients . Clearly, this very slow power law relaxation implies that an adiabatic decoupling of the internal front dynamics and the large scale pattern dynamics can not be made and hence that there is no long-wavelength effective interface equation of the form (3) for pulled fronts. There is then a priori no reason to expect that fluctuating pulled fronts are in the KPZ universality class!
It is our aim to test this scenario by introducing a simple stochastic lattice model whose front dynamics can be changed from pushed to pulled by tuning a single parameter. Our results are consistent with our conjecture that pulled fluctuating fronts are not in the standard KPZ universality class, while pushed fronts are. In fact, our results put an earlier empirical finding of Riordan et al. into a new perspective: These authors obtained essentially the same growth exponent as we do for the non-KPZ case, but the connection with the transition from pushed to pulled front dynamics was not made.
Our stochastic model is motivated by and the results for deterministic planar fronts in the nonlinear diffusion equation
$$\rho /t=D^2\rho +k_1\rho +k_2\rho ^2k_3\rho ^3.$$
(4)
As discussed in , the planar fronts with $`\rho >0`$ propagating into the unstable state $`\rho =0`$ are pulled for all values $`k_2<\sqrt{k_1k_3/2}`$ and pushed for larger $`k_2`$. In the pulled regime, the asymptotic front velocity, is $`v^{}=2\sqrt{Dk_1}`$, while in the pushed regime the asymptotic front velocity equals $`v^{}=2\sqrt{Dk_1}[(K+\sqrt{K^2+4})/32]`$ where $`K=k_2/\sqrt{k_1k_3}`$. We confine ourselves here to studying two limits where the stochastic front dynamics can easily be understood intuitively.
We study the dynamics of particles on a square lattice, subject to the constraints that no more than one particle can occupy each lattice site. The stochastic moves are illustrated in Fig. 1. They consist of diffusive hops of particles to neighboring empty sites and of birth and death processes on sites neighboring an occupied site. In a mean field approximation, this stochastic model is equivalent to a discrete version of (4). We will study here the two cases indicated in Fig. 1. For $`k_1=0`$ (Fig. 1a), planar fronts are definitely pushed: Since the linear spreading speed $`v^{}=0`$ for $`k_1=0`$, the front must then be pushed, even if corrections to the mean field behavior are important in the front region or behind the front. Likewise, when $`k_3=0`$ and $`k_2<0`$ (Fig. 1b), the nonlinearities behind the front only limit the birth (growth) rate, so in this limit the stochastic planar front is definitely pulled.
Our simulations are done on 2D strips which are long in the $`y`$ direction and of width $`L`$ in the $`x`$ direction. In the $`x`$ direction, periodic boundary conditions are used. The Monte Carlo simulations are started with a configuration in which the first few rows ($`100`$) of the lattice are occupied with a probability equal to the equilibrium density. All other lattice sites are empty. After an initial transient, the scaling properties of the interface width are studied in the standard way using the following definition of the interface height $`h`$. We define a coarse-grained density variable at each lattice site as the average occupation of sites on a $`(2m+1)\times (2m+1)`$ grid centered at that site. We then define the position $`h(x_i)`$ of the interface as the first point where this coarse-grained density reaches half the equilibrium density value. Our results for ensemble averaged width of the interface (see below) are obtained by averaging over $`100`$ runs for the largest system $`L=2048`$ to about $`3200`$ runs for the smallest $`L=64`$. Although we have performed simulations with $`m=1,2`$ and $`m=3`$ almost all the data presented subsequently are those for a representative value of $`m=2`$. The coarse-grained density field and the corresponding interface position $`h`$ for a typical configuration is shown in Fig. 2.
The interface width $`w`$ of a given realization is defined in the usual way, $`w^2(t)=\overline{(h(x_i,t)\overline{h(x_i,t)})^2}`$, where the overbar denotes a spatial average, $`\overline{h}=L^1_{x_i}h(x_i,t)`$. The proper scaling to study is the ensemble averaged mean square interface width $`W^2w^2`$. As is well known, in the KPZ equation $`W`$ obeys a scaling form $`W(t)=t^\beta 𝒴\left(\frac{t}{L^z}\right)`$. Here the scaling function $`𝒴(u)`$ is about constant for $`u1`$ and $`𝒴u^\beta `$ for $`u1`$, with the KPZ exponents $`z=3/2`$ and $`\beta =1/3`$ in $`1+1`$D. For $`tL^z`$, the width saturates at $`W_{sat}L^\zeta `$ where $`\zeta =\beta z`$ is the roughness exponent. In Fig. 3, we show our data for stochastic pushed and pulled fronts by plotting $`W/L^\zeta `$ versus $`t/L^z`$ for a range of system sizes ($`L=128`$ to $`2048`$). Following standard practice, we always plot the subtracted width $`W^2(t)W^2(0)`$ to minimize the effect of the initial front width. The kinetic parameters are chosen to be $`k_2=0.5,k_3=1.0`$ for the pushed model and $`k_2=0,k_1=0.1`$ for the pulled model. The diffusion rate $`D=0.25`$ is the same in both cases. In Fig. 3a we use the 1+1D KPZ exponents to obtain a data collapse. Clearly, good scaling collapse of the pushed data confirms that the pushed fronts are in the universality class of the 1+1D KPZ equation. By contrast, use of 1+1D KPZ values does not lead to good scaling collapse of the pulled data. In Fig. 3b we show the same sets of data but now with exponents $`\zeta =0.4`$ and $`z=1.38`$ to obtain the best possible scaling collapse of the pulled data. It is clear that the two sets of exponents, though only moderately different from each other, are well beyond error-bars. More accurate estimates of the exponents for the pulled case were obtained as follows. In Fig. 4 we fit a power law to the non-saturated part of the width for the largest system $`L=2048`$ and obtain $`\beta 0.29\pm 0.01`$ for the growth exponent. Plotting the saturated width $`W_{sat}`$ as a function of system size $`L`$ (Fig. 4, inset) yields $`\zeta 0.4\pm 0.02`$ for the roughness exponent. Once $`\zeta `$ is known the dynamic exponent is obtained by requiring good scaling collapse of Fig. 3b., $`z=1.38\pm 0.06`$. The value of $`\beta `$ is consistent with that reported by Riordan et al. for this model, $`\beta =0.272\pm 0.007`$, but their apparent value of $`z1`$ is not the true dynamic exponent related to the interface roughness through $`\zeta =\beta z`$, since they studied the ensemble averaged width of the front .
Another way to investigate the possible difference with the 1+1D KPZ behavior is to study the distribution $`P(w^2/W^2)`$. For 1D interface models whose long time interface configurations are given by a Gaussian distribution, like the KPZ model, the distribution function $`P(w^2/W^2)`$ is uniquely determined, without adjustable parameters . As Fig. 5 shows, in the pushed regime our data are completely consistent with this distribution function, but in the pulled regime the measured distribution function deviates significantly from the universal prediction for Gaussian interface fluctuations.
The essential difference between pushed and pulled fronts is that for pushed fronts the dynamically important region is the finite transition zone between the two phases it separates, whereas for pulled fronts it is the semi-infinite leading edge ahead of the front itself . It is precisely for this reason that the wandering of stochastic pulled fronts in one bulk dimension with multiplicative noise was recently found to be subdiffusive and determined by the 1+1D KPZ equation, not by a “0+1D” stochastic Langevin equation . By extending this idea it has been recently conjectured that the scaling exponents of stochastic pulled fronts in $`d`$+$`1`$ bulk dimensions are generally given by the $`(d`$$`+`$$`1)`$$`+`$$`1`$D KPZ equation instead of the $`d`$+1D KPZ equation, essentially because the dimension perpendicular to the front can not be integrated out . The scaling exponents we find here in 2 bulk dimensions are indeed close to those reported for the $`2`$+$`1`$D KPZ equation , the supposedly most accurate values being $`\zeta =0.393(3)`$, $`\beta =0.245(3)`$ . Moreover, the probability distribution $`P(w^2/W^2)`$ of pulled fronts fits the $`P(w^2/W^2)`$ of the 2+1D KPZ equation quite well without adjustable parameters, see Fig. 5. For justification and further exploration of this conjecture, we refer to .
An interesting limit of our model is obtained when we further take $`k_2=0`$ in Fig. 1b. In this case only birth and diffusion occurs, leading to an equilibrium density $`\rho =1`$ behind the front. If we put $`D=0`$ as well, the result is an Eden-like model with the modification that the probability of adding a particle is proportional to the number of neighbors, not independent of it. Numerical simulations in this Eden-like limit indicate that the KPZ exponents are recovered, as it should, and hence that the model has a KPZ to non-KPZ transition at intermediate values of $`\rho _{eq}`$ and $`D`$.
In conclusion, even though one should always be aware of the possibility of a very slow crossover to asymptotic behavior in such studies — a problem that has plagued some earlier tests of KPZ scaling in, e.g., the Eden model — taken together our data as well as those of give, in our opinion, reasonably convincing evidence for our scenario that the absence of an effective interface description for deterministic pulled fronts also entails non-KPZ scaling of stochastic pulled fronts.
We thank J. Krug, T. Bohr and Z. Rácz for stimulating discussions. G. Tripathy is supported by the Dutch Foundation for Fundamental Research on Matter (FOM).
|
warning/0005/hep-ph0005081.html
|
ar5iv
|
text
|
# References
DFTT 18/2000
INFNCA-TH0008
VUTH 00-13
hep-ph/0005081
Phenomenology of transverse single spin
asymmetries in inclusive processes<sup>*</sup><sup>*</sup>*Talk delivered by M. Anselmino at the Fifth Workshop on Quantum Chromodynamics, January 3-7, 2000, Villefranche, France.
M. Anselmino<sup>1</sup>, M. Boglione<sup>2</sup> and F. Murgia<sup>3</sup>
<sup>1</sup>Dipartimento di Fisica Teorica, Università di Torino and
INFN, Sezione di Torino, Via P. Giuria 1, 10125 Torino, Italy
<sup>2</sup>Dept. of Physics and Astronomy, Vrije Universiteit Amsterdam,
De Boelelaan 1081, 1081 HV Amsterdam, The Netherlands
<sup>3</sup>Dipartimento di Fisica, Università di Cagliari and
INFN, Sezione di Cagliari, CP 170, I-09042 Monserrato (CA), Italy
## Abstract
A phenomenological description of single transverse spin asymmetries for the inclusive production of hadrons in $`pp`$ and $`\mathrm{}p`$ processes is discussed within pQCD and a straightforward generalization of the factorization theorem with the inclusion of parton intrinsic transverse motion. Fits to existing data, predictions for new processes and interpretation of recent results are presented.
QCD formalism for $`A^{\mathbf{}}B\mathbf{}CX`$ at leading twist, $`k_{\mathbf{}}\mathbf{=}\mathrm{𝟎}\mathbf{:}A_N\mathbf{=}\mathrm{𝟎}`$ It is well known that perturbative QCD and the factorization theorem at leading twist can be used to describe the large $`p_T`$ production of a hadron $`C`$ resulting from the interaction of two polarized hadrons $`A`$ and $`B`$:
$`{\displaystyle \frac{E_Cd^3\sigma ^{A,S_A+B,S_BC+X}}{d^3𝒑_C}}`$ $`=`$ $`{\displaystyle \underset{a,b,c,d;\{\lambda \}}{}}\rho _{\lambda _a^{},\lambda _a^{}}^{a/A,S_A}f_{a/A}(x_a)\rho _{\lambda _b^{},\lambda _b^{}}^{b/B,S_B}f_{b/B}(x_b)`$
$``$ $`\widehat{M}_{\lambda _c^{},\lambda _d^{};\lambda _a^{},\lambda _b^{}}\widehat{M}_{\lambda _c^{},\lambda _d^{};\lambda _a^{},\lambda _b^{}}^{}D_{\lambda _C^{},\lambda _C^{}}^{\lambda _c^{},\lambda _c^{}}(z),`$
where $``$ denotes the usual convolutions \[see, for example, Ref. for details\]. $`\rho _{\lambda _a^{},\lambda _a^{}}^{a/A,S_A}(x_a)`$ is the helicity density matrix of parton $`a`$ inside the polarized hadron $`A`$; similarly for $`\rho _{\lambda _b^{},\lambda _b^{}}^{b/B,S_B}(x_b)`$. The $`\widehat{M}_{\lambda _c^{},\lambda _d^{};\lambda _a^{},\lambda _b^{}}`$’s are the helicity amplitudes for the elementary process $`abcd`$; if one wishes to consider higher order (in $`\alpha _s`$) contributions also elementary processes involving more partons should be included. $`D_{\lambda _C^{},\lambda _C^{}}^{\lambda _c^{},\lambda _c^{}}(z)`$ is the product of fragmentation amplitudes for the $`cC+X`$ process
$$D_{\lambda _C^{},\lambda _C^{}}^{\lambda _c^{},\lambda _c^{}}={\displaystyle }_{X,\lambda _X}𝒟_{\lambda _X^{},\lambda _C^{};\lambda _c^{}}𝒟_{\lambda _X^{},\lambda _C^{};\lambda _c^{}}^{},$$
(2)
where the $`{\displaystyle }_{X,\lambda _X}`$ stands for a spin sum and phase space integration of the undetected particles, considered as a system $`X`$. The usual unpolarized fragmentation function $`D_{C/c}(z)`$, i.e. the density number of hadrons $`C`$ resulting from the fragmentation of an unpolarized parton $`c`$ and carrying a fraction $`z`$ of its momentum, is given by
$$D_{C/c}(z)=\frac{1}{2}\underset{\lambda _c^{},\lambda _C^{}}{}D_{\lambda _C^{},\lambda _C^{}}^{\lambda _c^{},\lambda _c^{}}(z).$$
(3)
For simplicity of notations we have not shown in Eq. (S0.Ex1) the $`Q^2`$ scale dependences in the unpolarized distribution and fragmentation functions, $`f`$ and $`D`$.
In the case in which only one hadron, say $`A`$, is polarized (orthogonally to the scattering plane) Eq. (S0.Ex1) reads
$`{\displaystyle \frac{E_Cd^3\sigma ^{A^{}BCX}}{d^3𝒑_C}}`$ $`=`$ $`{\displaystyle \underset{a,b,c,d;\{\lambda \}}{}}\rho _{\lambda _a^{},\lambda _a^{}}^{a/A^{}}f_{a/A}(x_a){\displaystyle \frac{1}{2}}f_{b/B}(x_b)`$
$``$ $`\widehat{M}_{\lambda _c^{},\lambda _d^{};\lambda _a^{},\lambda _b^{}}\widehat{M}_{\lambda _c^{},\lambda _d^{};\lambda _a^{},\lambda _b^{}}^{}D_{\lambda _C^{},\lambda _C^{}}^{\lambda _c^{},\lambda _c^{}}(z).`$
Eqs. (S0.Ex1) and (S0.Ex2) hold at leading twist and large $`p_T`$ values of the produced hadron; the intrinsic $`𝒌_{}`$ of the partons have been integrated over and collinear configurations dominate both the distribution functions and the fragmentation processes; one can then see that, in this case, there cannot be any single spin asymmetry. In fact, total angular momentum conservation in the (forward) fragmentation process \[see Eq. (2)\] implies $`\lambda _c^{}=\lambda _c^{}`$; this, in turns, together with helicity conservation in the elementary process, implies $`\lambda _a^{}=\lambda _a^{}`$. If we further notice that, by parity invariance, $`_{\lambda _C^{}}D_{\lambda _C^{},\lambda _C^{}}^{\lambda _c^{},\lambda _c^{}}`$ does not depend on $`\lambda _c^{}`$ and that $`_{\lambda _b^{},\lambda _c^{},\lambda _d^{}}|\widehat{M}_{\lambda _c^{},\lambda _d^{};\lambda _a^{},\lambda _b^{}}|^2`$ is independent of $`\lambda _a^{}`$, we remain with $`_{\lambda _a^{}}\rho _{\lambda _a^{},\lambda _a^{}}^{a/A^{}}=1`$. Moreover, in the absence of intrinsic $`𝒌_{}`$ and initial state interactions, the parton density numbers $`f_{a/A}(x_a)`$ cannot depend on the spin of $`A`$ and any spin dependence disappears from Eq. (S0.Ex2), so that:
$$d\sigma ^{A^{}BCX}d\sigma ^{A^{}BCX}=0.$$
(5)
Therefore one does not expect any sizeable single spin (or left-right) asymmetry:
$$A_N\frac{d\sigma ^{}(𝒑_T)d\sigma ^{}(𝒑_T)}{d\sigma ^{}(𝒑_T)+d\sigma ^{}(𝒑_T)}=\frac{d\sigma ^{}(𝒑_T)d\sigma ^{}(𝒑_T)}{2d\sigma ^{unp}}$$
(6)
This is contradicted by data showing large values of $`A_N`$ in $`p^{}p\pi X`$ and $`\overline{p}^{}p\pi X`$ processes.
Intrinsic $`k_{}`$ effects in fragmentation and/or distribution functions: $`A_N\mathbf{}\mathrm{𝟎}`$
The above conclusion may be avoided by considering the transverse motion of the quarks relatively to the parent hadron or of the observed hadron relatively to the fragmenting quark. The original suggestion that the intrinsic $`𝒌_{}`$ of the quarks in the distribution functions might give origin to single spin asymmetries was first made by Sivers ; such an effect is not forbidden by QCD time reversal invariance provided one takes into account soft initial state interactions among the colliding hadrons . A similar suggestion for a possible origin of single spin asymmetries was later made by Collins , concerning transverse momentum effects in the fragmentation of a polarized quark. A consistent phenomenological application of these ideas can be found in a series of papers . More recently, a further possible source of single spin effects, related to distribution functions, was discussed by Boer .
When allowing for parton intrinsic motion spin effects can remain in new – spin and $`𝒌_{}`$ dependent – distribution or fragmentation functions. We list these functions in the sequel.
$`\mathrm{\Delta }^Nf_{q/p^{}}(x,𝒌_{})`$ is the difference between the density numbers $`\widehat{f}_{q/p^{}}(x,𝒌_{})`$ and $`\widehat{f}_{q/p^{}}(x,𝒌_{})`$ of quarks $`q`$, with all possible polarization, longitudinal momentum fraction $`x`$ and intrinsic transverse momentum $`𝒌_{}`$, inside a transversely polarized proton with spin $``$ or $``$:
$`\mathrm{\Delta }^Nf_{q/p^{}}(x,𝒌_{})`$ $``$ $`\widehat{f}_{q/p^{}}(x,𝒌_{})\widehat{f}_{q/p^{}}(x,𝒌_{})`$
$`=`$ $`\widehat{f}_{q/p^{}}(x,𝒌_{})\widehat{f}_{q/p^{}}(x,𝒌_{})`$
where the second line follows from the first one by rotational invariance. Notice that $`\mathrm{\Delta }^Nf_{q/p^{}}(x,𝒌_{})`$ vanishes when $`𝒌_{}0`$; parity invariance also requires $`\mathrm{\Delta }^Nf`$ to vanish when the proton transverse spin has no component perpendicular to $`𝒌_{}`$, so that
$$\mathrm{\Delta }^Nf_{q/p^{}}(x,𝒌_{})k_{}\mathrm{sin}\alpha $$
(8)
where $`\alpha `$ is the angle between $`𝒌_{}`$ and the $``$ direction.
$`\mathrm{\Delta }^Nf`$ by itself is a leading twist distribution function, but its $`𝒌_{}`$ dependence, when convoluted with the elementary partonic cross-section, results in twist-3 contributions to single spin asymmetries. This same function (up to some factors) has also been introduced in Ref. – where it is denoted by $`f_{1T}^{}`$ – as a leading twist $`T`$-odd distribution function. The exact relation between $`\mathrm{\Delta }^Nf`$ and $`f_{1T}^{}`$ is discussed in Ref. .
A function analogous to $`\mathrm{\Delta }^Nf_{q/p^{}}(x,𝒌_{})`$ can be defined for the fragmentation process of a transversely polarized parton , giving the difference between the density numbers $`\widehat{D}_{h/q^{}}(z,𝒌_{})`$ and $`\widehat{D}_{h/q^{}}(z,𝒌_{})`$ of hadrons $`h`$, with longitudinal momentum fraction $`z`$ and transverse momentum $`𝒌_{}`$ inside a jet originated by the fragmentation of a transversely polarized quark with spin $``$ or $``$:
$`\mathrm{\Delta }^ND_{h/q^{}}(z,𝒌_{})`$ $``$ $`\widehat{D}_{h/q^{}}(z,𝒌_{})\widehat{D}_{h/q^{}}(z,𝒌_{})`$
$`=`$ $`\widehat{D}_{h/q^{}}(z,𝒌_{})\widehat{D}_{h/q^{}}(z,𝒌_{}).`$
A closely related function is denoted by $`H_1^{}`$ in Refs. and its correspondence with $`\mathrm{\Delta }^ND`$ is discussed in Ref. . Again we expect
$$\mathrm{\Delta }^ND_{h/q^{}}(z,𝒌_{})k_{}\mathrm{sin}\beta $$
(10)
where $`\beta `$ is the angle between $`𝒌_{}`$ and the $``$ direction.
Similarly one can introduce the difference $`\mathrm{\Delta }^Nf_{q^{}/p}(x,𝒌_{})`$ between the density numbers $`\widehat{f}_{q^{}/p}(x,𝒌_{})`$ and $`\widehat{f}_{q^{}/p}(x,𝒌_{})`$ of quarks $`q`$, with spin $``$ or $``$, longitudinal momentum fraction $`x`$ and intrinsic transverse momentum $`𝒌_{}`$, inside an unpolarized proton:
$`\mathrm{\Delta }^Nf_{q^{}/p}(x,𝒌_{})`$ $``$ $`\widehat{f}_{q^{}/p}(x,𝒌_{})\widehat{f}_{q^{}/p}(x,𝒌_{})`$
$`=`$ $`\widehat{f}_{q^{}/p}(x,𝒌_{})\widehat{f}_{q^{}/p}(x,𝒌_{}).`$
A closely related function, denoted by $`h_1^{}`$, was discussed in Ref. .
Finally, although not relevant for the single spin asymmetries we consider here, one could introduce the difference between the density numbers $`\widehat{D}_{h^{}/q}(z,𝒌_{})`$ and $`\widehat{D}_{h^{}/q}(z,𝒌_{})`$ of hadrons $`h`$, with longitudinal momentum fraction $`z`$ and transverse momentum $`𝒌_{}`$ inside a jet originated by the fragmentation of an unpolarized quark $`q`$:
$`\mathrm{\Delta }^ND_{h^{}/q}(z,𝒌_{})`$ $``$ $`\widehat{D}_{h^{}/q}(z,𝒌_{})\widehat{D}_{h^{}/q}(z,𝒌_{})`$
$`=`$ $`\widehat{D}_{h^{}/q}(z,𝒌_{})\widehat{D}_{h^{}/q}(z,𝒌_{}).`$
Such a function might prove useful in tackling the longstanding problem of hyperon polarization in inclusive $`pN`$ processes . A closely related function is denoted by $`D_{1T}^{}`$ in Refs. .
Assuming that the factorization theorem, Eq. (S0.Ex2), holds when parton intrinsic motion is taken into account, and using the new functions (S0.Ex3), (S0.Ex4) and (S0.Ex5), one has, at leading order in $`k_{}`$, for the $`p^{}p\pi X`$ process:
$`{\displaystyle \frac{E_\pi d^3\sigma ^{}}{d^3𝒑_\pi }}{\displaystyle \frac{E_\pi d^3\sigma ^{}}{d^3𝒑_\pi }}`$ $`=`$ $`{\displaystyle \underset{a,b,c,d}{}}{\displaystyle }{\displaystyle \frac{dx_adx_b}{\pi z}}\times \{`$
$`{\displaystyle d^2𝒌_{}\mathrm{\Delta }^Nf_{a/p^{}}(x_a,𝒌_{})f_{b/p}(x_b)\frac{d\widehat{\sigma }}{d\widehat{t}}(x_a,x_b,𝒌_{})D_{\pi /c}(z)}`$
$`+`$ $`{\displaystyle d^2𝒌_{}^{}h_1^{a/p}(x_a)f_{b/p}(x_b)\mathrm{\Delta }_{NN}\widehat{\sigma }(x_a,x_b,𝒌_{}^{})\mathrm{\Delta }^ND_{\pi /c}(z,𝒌_{}^{})}`$
$`+`$ $`{\displaystyle }d^2𝒌_{}^{\prime \prime }h_1^{a/p}(x_a)\mathrm{\Delta }^Nf_{b^{}/p}(x_b,𝒌_{}^{\prime \prime })\mathrm{\Delta }_{NN}^{}\widehat{\sigma }(x_a,x_b,𝒌_{}^{\prime \prime })D_{\pi /c}(z)\},`$
where the second line corresponds to the so-called Sivers effect , the third to Collins effect and the fourth one to the mechanism recently proposed by Boer . Apart from the new functions $`\mathrm{\Delta }^Nf`$ and $`\mathrm{\Delta }^ND`$, the other quantities appearing in Eq. (S0.Ex9) are the unpolarized quark distribution and fragmentation functions, $`f`$ and $`D`$; the unpolarized cross-section for the elementary process $`abcd`$, $`d\widehat{\sigma }/d\widehat{t}`$; the transverse spin content of the proton:
$$h_1^{q/p}=f_{q^{}/p^{}}(x)f_{q^{}/p^{}}(x)$$
(14)
and the elementary double spin asymmetries, computable in pQCD:
$$\mathrm{\Delta }_{NN}\widehat{\sigma }=\frac{d\widehat{\sigma }^{a^{}bc^{}d}}{d\widehat{t}}\frac{d\widehat{\sigma }^{a^{}bc^{}d}}{d\widehat{t}},$$
(15)
$$\mathrm{\Delta }_{NN}^{}\widehat{\sigma }=\frac{d\widehat{\sigma }^{a^{}b^{}cd}}{d\widehat{t}}\frac{d\widehat{\sigma }^{a^{}b^{}cd}}{d\widehat{t}}$$
(16)
An equation with the same structure and a similar physical meaning as Eq. (S0.Ex9) can be found in Ref. , where the new functions $`\mathrm{\Delta }^Nf`$ and $`\mathrm{\Delta }^ND`$ are replaced by higher twist parton correlation functions.
Eq. (S0.Ex9), or some of its terms, can be used for a phenomenological description of single spin asymmetries. We next summarize what has been done.
Fits to existing data, predictions Sivers effect only In Refs. and the scheme of Eq. (S0.Ex9) was adopted, taking into account only Sivers effect; to simplify things it was assumed – in this very first application of the idea of intrinsic quark motion and spin dependence – that the integral over $`𝒌_{}`$ is dominated by configurations in which $`𝒌_{}`$ lies in the scattering plane \[that is, $`\mathrm{sin}\alpha =1`$ in Eq. (8)\] and its magnitude equals some average value :
$$\frac{1}{M}k_{}^0(x_a)=0.47x_a^{0.68}(1x_a)^{0.48},$$
(17)
where $`M`$ = 1 GeV/$`c^2`$.
The residual $`x_a`$ dependence in $`\mathrm{\Delta }^Nf_{a/p^{}}`$ not coming from $`k_{}^0`$ was taken to be of the simple form
$$N_ax_a^{\alpha _a}(1x_a)^{\beta _a},$$
(18)
where $`N_a`$, $`\alpha _a`$ and $`\beta _a`$ are free parameters. Only $`u`$ and $`d`$ quark contributions to $`\mathrm{\Delta }^Nf_{a/p^{}}`$ were considered.
One then ends up with the simple expression:
$`{\displaystyle d^2𝒌_{}\mathrm{\Delta }^Nf_{a/p^{}}(x_a,𝒌_{})\left[\frac{d\widehat{\sigma }}{d\widehat{t}}(𝒌_{})\frac{d\widehat{\sigma }}{d\widehat{t}}(𝒌_{})\right]}`$ (19)
$``$ $`{\displaystyle \frac{k_{}^0(x_a)}{M}}N_ax_a^{\alpha _a}(1x_a)^{\beta _a}\left[{\displaystyle \frac{d\widehat{\sigma }}{d\widehat{t}}}(𝒌_{}^0){\displaystyle \frac{d\widehat{\sigma }}{d\widehat{t}}}(𝒌_{}^0)\right],`$
where $`k_{}^0(x_a)`$ is given by Eq. (17) and, choosing $`xz`$ as the scattering plane and $`z`$ as the direction of the incoming polarized proton, $`𝒌_{}^0=(k_{}^0,0,0)`$.
Inserting this result into the first two lines of Eq. (S0.Ex9), one can compute the single spin asymmetry (6) in terms of the parameters $`\alpha _a,\beta _a,N_a`$; only the leading valence quark contributions to $`\mathrm{\Delta }^Nf`$ were considered in the numerator of $`A_N`$, while all leading order pQCD elementary processes were included in the denominator. The parameters were fixed by performing a best fit to the experimental data, with the results shown in Fig. 1, corresponding to the $`\mathrm{\Delta }^Nf`$ functions:
$$\mathrm{\Delta }^Nf_{u/p^{}}(x,k_{}^0)=6.92x^{2.02}(1x)^{4.06}\mathrm{\Delta }^Nf_{d/p^{}}(x,k_{}^0)=2.33x^{1.44}(1x)^{4.62}$$
(20)
It is clear from Fig. 1 how Sivers effect alone allows a good fit of the experimental data. Let us add a few more comments:
* The resulting expressions of $`\mathrm{\Delta }^Nf_{q/p^{}}`$, Eq. (20), are reasonable and quite acceptable; in particular they satisfy the positivity condition $`|\mathrm{\Delta }^Nf_{q/p^{}}|2f_{q/p}`$.
* The opposite sign of $`\mathrm{\Delta }^Nf_{u/p^{}}`$ and $`\mathrm{\Delta }^Nf_{d/p^{}}`$ is expected from transverse momentum conservation.
* It is easy to explain, within this scheme, why almost opposite values of $`A_N^{\pi ^+}`$ and $`A_N^\pi ^{}`$ do not imply, as one might naively expect, $`A_N^{\pi ^0}0`$ .
* This same scheme, with the same $`\mathrm{\Delta }^Nf`$ functions (20), was used – without any free parameter – to compute $`A_N`$ in $`\overline{p}^{}p\pi X`$; the agreement with data is good .
Collins effect only A similar analysis was performed in Ref. taking into account only Collins effect. Again, it was assumed that the main contribution comes from an average $`k_{}^0`$ value, with a simple parametrization of $`\mathrm{\Delta }^ND_{\pi /c}(z,k_{}^0)`$ for all leading valence quarks:
$$\mathrm{\Delta }^ND_{\pi /c}(z,k_{}^0)=\frac{k_{}^0(z)}{M}N_cz^{\alpha _c}(1z)^{\beta _c},$$
(21)
where $`N_c`$, $`\alpha _c`$ and $`\beta _c`$ are free parameters and
$$\frac{k_{}^0(z)}{M}=0.61z^{0.27}(1z)^{0.20},$$
(22)
with $`M`$ = 1 GeV/$`c^2`$.
Another free parameter is contained in the expression of $`h_1^{q/p}(x)=P^{q/p^{}}f_{q/p}(x)`$ where $`P^{q/p^{}}`$ is the transverse polarization of quark $`q`$ inside the transversely polarized proton.
The experimental data were best fitted as in Fig. 2, with the resulting $`\mathrm{\Delta }^ND`$ function and $`P^{q/p^{}}`$ values :
$`\mathrm{\Delta }^ND_{\pi /q}(z,k_{}^0)`$ $`=`$ $`0.13z^{2.60}(1z)^{0.44}z0.977`$
$`\mathrm{\Delta }^ND_{\pi /q}(z,k_{}^0)=2D_{\pi /q}(z)`$ $`=`$ $`2.20z^1(1z)^{1.2}z>0.977`$ (23)
$`P^{u/p^{}}={\displaystyle \frac{2}{3}}`$ $`P^{d/p^{}}=0.88.`$
The quality of the fit is comparable to that obtained using only Sivers effect, shown in Fig. 1. However, some comments are now necessary:
* In order to fit the data, $`\mathrm{\Delta }^ND_{\pi /q}`$ has to saturate at large $`z`$ the positivity constraint $`|\mathrm{\Delta }^ND_{\pi /q}|2D_{\pi /q}`$. Otherwise, the values of $`A_N`$ at large $`x_F`$ would be much too small.
* The resulting value of $`h_1^{d/q}`$ is
$$h_1^{d/p}(x)=0.88f_{d/p}(x),$$
(24)
which violates the Soffer’s bound
$$|h_1^{d/p}|\frac{1}{2}(f_{d/p}+\mathrm{\Delta }d).$$
(25)
$`\mathrm{\Delta }d`$ is the $`d`$ quark helicity distribution, which, being in most parametrization of data negative, makes bound (25) very strict.
* A parametrization of $`\mathrm{\Delta }^ND`$ which satisfies Soffer’s bound was used in Ref. , with a good resulting fit, provided one allows $`\mathrm{\Delta }d`$ to become positive at large $`x`$ values; however, also in this case the positivity constraint has to be saturated at large $`z`$.
We can at this point conclude that the phenomenological approach to single spin asymmetries based on the generalization of the factorization theorem with the inclusion of quark intrinsic motion and on the introduction of the new spin and $`𝒌_{}`$ dependent functions (S0.Ex3), (S0.Ex4) and (S0.Ex5) is indeed promising and worth being pursued. The relative contributions of the several terms in Eq. (S0.Ex9) is still unknown; a first guess would indicate that Sivers effect is indeed necessary, while Collins effect alone leads to unreasonably large values of $`|\mathrm{\Delta }^ND|`$, the Collins function. The third effect, originating from the quarks in the unpolarized hadron has not been studied yet, although it should not be relevant for large and positive $`x_F`$ values. We turn now to a study of Collins effect in DIS.
Fragmentation of a polarized quark in semi-inclusive DIS The inclusive production of hadrons in DIS with transversely polarized nucleons, $`\mathrm{}N^{}\mathrm{}hX`$, is the ideal process to study Collins effect; in such a case, in fact, Sivers effect, which requires initial state interactions , is negligible and any single spin asymmetry must originate from spin dependences in the fragmentation of a polarized quark.
If one looks at the $`\gamma ^{}N^{}hX`$ process in the $`\gamma ^{}N`$ c.m. system, the elementary interaction is simply a $`\gamma ^{}`$ hitting a transversely polarized quark, which bounces back and fragments into a jet containing the detected hadron. The hadron $`p_T`$ in this case coincides with its $`k_{}`$ inside the jet; the fragmenting quark polarization can be computed from the initial quark one.
The spin and $`𝒌_{}`$ dependent fragmentation function for a quark with momentum $`𝒑_q`$ and a transverse polarization vector $`𝑷_q`$ ($`𝒑_q𝑷_q=0`$) which fragments into a hadron with momentum $`𝒑_h=z𝒑_q+𝒑_T`$ ($`𝒑_q𝒑_T=0`$) can be written as:
$$D_{h/q}(𝒑_q,𝑷_q;z,𝒑_T)=\widehat{D}_{h/q}(z,p_T)+\frac{1}{2}\mathrm{\Delta }^ND_{h/q}(z,p_T)\frac{𝑷_q(𝒑_q\times 𝒑_T)}{|𝒑_q\times 𝒑_T|}$$
(26)
where $`\widehat{D}_{h/q}(z,p_T)`$ is the unpolarized fragmentation function and $`\mathrm{\Delta }^ND`$ is the same function as introduced in Eq. (S0.Ex4). Notice that – as required by parity invariance – the only component of the polarization vector which contributes to the spin dependent part of $`D`$ is that perpendicular to the $`qh`$ plane; in general one has:
$$𝑷_q\frac{𝒑_q\times 𝒑_T}{|𝒑_q\times 𝒑_T|}=P_q\mathrm{sin}\mathrm{\Phi }_C,$$
(27)
where $`P_q=|𝑷_q|`$ and we have defined the Collins angle $`\mathrm{\Phi }_C`$.
Eq. (26) leads to a possible single spin asymmetry in semi-inclusive DIS off nucleons with polarization $`\pm 𝑷`$,
$$A_N^h\frac{d\sigma ^{\mathrm{}+p,𝑷\mathrm{}^{}+h+X}d\sigma ^{\mathrm{}+p,𝑷\mathrm{}^{}+h+X}}{d\sigma ^{\mathrm{}+p,𝑷\mathrm{}^{}+h+X}+d\sigma ^{\mathrm{}+p,𝑷\mathrm{}^{}+h+X}},$$
(28)
which is given by
$$A_N^h(x,y,z,\mathrm{\Phi }_C,p_T)=\frac{_qe_q^2h_1^{q/p}(x)\mathrm{\Delta }^ND_{h/q}(z,p_T)}{2_qe_q^2f_{q/p}(x)D_{h/q}(z,p_T)}\frac{2(1y)}{1+(1y)^2}P_T\mathrm{sin}\mathrm{\Phi }_C,$$
(29)
where $`P_T`$ is the nucleon polarization vector component transverse with respect to the $`\gamma ^{}`$ direction; $`x`$ and $`y`$ are the usual DIS variables. If one collects data at different kinematical values one should integrate over the relevant $`x`$ and $`y`$ regions so that Eq. (29) reads:
$$A_N^h(z,\mathrm{\Phi }_C,p_T)=\frac{_q𝑑x𝑑ye_q^2h_{1q}(x)\mathrm{\hspace{0.25em}2}(1y)/(xy^2)\mathrm{\Delta }^ND_{h/q}(z,p_T)P_T\mathrm{sin}\mathrm{\Phi }_C}{2_q𝑑x𝑑ye_q^2f_{q/p}(x)(1+(1y)^2)/(xy^2)D_{h/q}(z,p_T)}$$
(30)
Some preliminary data on $`A_N^\pi `$ have recently appeared and have been analysed in Ref. using Eq. (29). By saturating the unknown values of $`h_1^{u,d/p}`$ with the Soffer’s bound \[see Eq. (25)\] lower bounds for the Collins function have been obtained for the fragmentation of a $`u`$ quark into a $`\pi ^+`$. From SMC data one has
$$\frac{|\mathrm{\Delta }^ND_{\pi /q}(z,p_T)|}{2D_{\pi /q}(z,p_T)}\stackrel{>}{}(0.24\pm 0.15)z0.45,p_T0.65\text{GeV}/c,$$
(31)
and from HERMES data
$$\frac{|\mathrm{\Delta }^ND_{\pi /q}(z,p_T)|}{2D_{\pi /q}(z,p_T)}\stackrel{>}{}\mathrm{\hspace{0.25em}0.20}\pm 0.04(stat.)\pm 0.04(syst.)z0.2.$$
(32)
If confirmed, HERMES and SMC data indicate a large value of the Collins function, which might then play a significant role in other processes. In particular, it would be of great interest to compare the single spin asymmetry measured in the inclusive process $`\mathrm{}p^{}\pi X`$ with that measured by E704 Collaboration in $`p^{}p\pi X`$ processes: if the origin of the asymmetry is mainly in Collins mechanism similar results should be found in both cases. More data on single transverse spin asymmetries will be available in the future from operating or progressing facilities like JLAB, RHIC and COMPASS.
Acknowledgements One of us (M.A.) would like to thank the organizers of the Workshop for the kind invitation.
|
warning/0005/gr-qc0005063.html
|
ar5iv
|
text
|
# A Consistent Histories approach to the Unruh Effect
## 1 Introduction
### 1.1 Consistent Histories
The consistent histories approach to quantum theory originated in the pioneering work of Griffiths and Omnes . Initially the formalism was developed in an attempt to escape the familiar difficulties of the Copenhagen interpretation. More recently, Gell-Mann and Hartle suggested that generalised history theories may be useful in tackling the problems of quantum cosmology and quantum gravity, in particular the problem of time.
The basic ingredient of ‘conventional’ consistent histories is a time-ordered sequence of propositions about the system represented by a class operator:
$$C_\alpha :=\alpha _{t_1}(t_1)\alpha _{t_2}(t_2)\mathrm{}\alpha _{t_n}(t_n)$$
(1)
where $`\alpha _{t_i}(t_i)`$ is a Heisenberg picture projection operator representing a proposition made about the system at time $`t_i`$. To make physical predictions we must use the decoherence functional to identify (strongly) consistent sets of histories, *i.e.* sets $`\{\alpha _i\}`$ such that,
$`d(\alpha _i,\alpha _j)`$ $`:=`$ $`Tr_{}[C_{\alpha _i}^{}\rho C_{\alpha _j}]`$ (2)
$`=`$ $`0\text{if}ij`$ (3)
Within such consistent sets, the probability of a particular history $`\alpha _i`$ ‘occuring’ is $`d(\alpha _i,\alpha _i)`$. The consistency condition guarantees that the Kolmolgorov sum rules are satisfied.
If generalised history theories are to be useful in formulating quantum gravity, then it is important to understand how more conventional theories such as non-relativistic quantum mechanics and quantum field theory (QFT) can be formulated in history language. While non-relativistic quantum mechanics has been extensively studied within the formalism, there are very few results concerning QFT. This is the motivation for this paper in which we re-derive a well-known result in the theory of QFT on curved spaces, from a histories perspective. The Unruh effect is an analogue of Hawking radiation, but the gravitational field that induces the radiation is ‘apparent’ rather than ‘real’, *i.e.* it is measured by an observer accelerating through empty space rather than by an observer in the gravitational field of a black hole.
### 1.2 The HPO Approach
Isham proposed an algebraic scheme for generalised history theories of the type suggested by Gell-Mann and Hartle. The algebraic axioms are set up in analogy with the logical approach to single time quantum theory which is concerned with the pair $`(,𝒮)`$ where $``$ is the lattice of projection operators on a Hilbert space and $`𝒮`$ is the set of density matrices. Isham proposed that a generalised history theory should be composed of the pair $`(𝒰𝒫,𝒟)`$ where $`𝒰𝒫`$ is an *orthoalgebra* of propositions about possible histories and $`𝒟`$ is the space of decoherence functionals.
To fit conventional consistent histories into these axioms, we would like to interpret the class operators as logical propositions; however, the product of non-commuting projection operators is not a projection operator. This means it is difficult to define conjuctions, disjunctions and negations consistently. However, the *tensor product* of two projectors *is* a projector on the tensor product space. This is the central idea of the history projection operator (HPO) approach to consistent histories. The tensor product of Schrödinger picture projection operators, $`\alpha _{t_1}\alpha _{t_2}\mathrm{}\alpha _{t_n}`$, which is a projector on the n-time history space, $`𝒱^n:=_{t_1}_{t_2}\mathrm{}_{t_n}`$, represents the proposition “$`\alpha _{t_1}`$ is true at time $`t_1`$ *and then* $`\alpha _{t_2}`$ is true at time $`t_2`$ $`\mathrm{}`$ *and then* $`\alpha _{t_n}`$ is true at time $`t_n`$.” Now we can define the logical operations as we would for projection operators in any Hilbert space. So in this case, the orthoalgebra $`𝒰𝒫`$ is in fact the lattice $`𝒫(𝒱^n)`$ of projection operators on the history space.
The decoherence functional (2) can be written as
$$d(\alpha _i,\alpha _j)=Tr_{𝒱^n𝒱^n}(\alpha _i\alpha _jX)$$
(4)
for some $`X(𝒱^n𝒱^n)`$ where $`()`$ is defined as the set of bounded operators on $``$. Conversely Gleason’s theorem can be used to show that any decoherence functional that satisfies certain natural conditions can be written in this form . Therefore $`𝒟`$, the space of decoherence functionals, is the set of all functionals of this form. This result also holds in the continuous time case .
## 2 The Simple Harmonic Oscillator
### 2.1 Continuous Times
In extending HPO theory to the case of continuous time, which we anticipate to be important for QFT, we encounter the continuous tensor product of the single-time Hilbert space: $`𝒱^{cts}:=_t_t`$. To deal with this object it is useful to confine ourselves for the moment to the simple harmonic oscillator (SHO), where $`_t=L^2()`$ and to consider the *history group* . We can view $`𝒱^n`$ arising as the representation space for the n-fold direct product of the Weyl group of single time quantum theory on the line:
$`[x_{t_i},x_{t_j}]`$ $`=`$ $`0`$ (5)
$`[p_{t_i},p_{t_j}]`$ $`=`$ $`0`$ (6)
$`[x_{t_i},p_{t_j}]`$ $`=`$ $`i\mathrm{}\delta _{ij}`$ (7)
The advantage of this perspective is that it can be readily generalised to the case of continuous time. For this we consider the algebra:
$`[x_f,x_g]`$ $`=`$ $`0`$ (8)
$`[p_f,p_g]`$ $`=`$ $`0`$ (9)
$`[x_f,p_g]`$ $`=`$ $`i\mathrm{}(f,g)`$ (10)
where $`f,gL^2()`$ ; $`x_f:=𝑑tf(t)x_t`$ and $`(f,g):=𝑑tf(t)g(t)`$. This algebra is clearly isomorphic to the algebra of a one-dimensional QFT and suggests that field theory techniques will be useful in studying the theory. It is well-known that this algebra has a representation on the Fock space over $`L^2()`$, denoted $`[L^2()]`$. Indeed it can be shown that ,
$$𝒱^{cts}:=_t_t[L^2()]$$
(11)
and again $`𝒰𝒫`$ is a lattice, now it is the set of projection operators on the continuous history space, $`𝒫(𝒱^{cts})`$. The condition that the time-averaged Hamiltonian is self-adjoint is sufficient to select a unique representation of the history algebra . This representation is defined by the Fock basis associated with the creation operator,
$$a_f^{}:=\sqrt{\frac{m\omega }{2\mathrm{}}}x_fi\sqrt{\frac{1}{2m\omega \mathrm{}}}p_f$$
(12)
### 2.2 Time averaged propositions
The physical interpretation of a continuous time HPO theory is based on the assumption that projectors onto the spectrum of self-adjoint operators on $`𝒱`$represent propositions about the time-averages of physical quantities. So projections onto the eigenvectors of the $`x_f`$ operators introduced in (8) represent propositions about the average position of the particle over time. As $`[x_f,x_g]=0`$, these operators have common eigenvectors for any smearing function. We denote these eigenvectors $`|x()`$ and they can be interpreted as fine-grained histories or trajectories of the particle. In the single-time theory, $`x_t|x=x|x`$ so, formally, $`x_t|x()=x(t)|x()`$, which suggests;
$$x_f|x():=𝑑tf(t)x_t|x()=(f,x)|x()$$
(13)
If this is to make sense then $`x()`$ must be a member of $`L^2()`$. However, it is likely that the eigenvectors $`x()`$ will be distributions rather than functions. The natural procedure now would be to interpret the symbol $`(f,x)`$ to be the real number obtained from the pairing of the distribution $`x`$ with the function $`f`$. This implies that the allowed functions $`f`$ should really be members of Schwartz space rather than $`L^2()`$. We will not confront this issue here, and just consider functions which are members of some unspecified space, $`\tau `$.
For each $`f\tau `$ we have an equivalence relation, $`_f`$, on trajectories if we define $`x()_fy()`$ if $`(f,x)=(f,y)`$. We denote these equivalence classes by $`[(f,x)]`$. Now we consider projections onto the spectrum of $`x_f`$. We denote the operator which projects onto the eigenvector of $`x_f`$ with eigenvalue $`(f,x)`$ as $`P_{(f,x)}`$; it projects onto the equivalence class of trajectories $`[(f,x)]`$, *i.e.* onto a coarse-grained history. Similar remarks obviously apply to operators $`P_{(f,p)}`$ which project onto coarse grained momentum trajectories $`[(f,p)]`$.
Another operator of physical significance is the smeared Hamiltonian:
$`H_f`$ $`:=`$ $`{\displaystyle 𝑑tf(t)(\frac{1}{2m}p_tp_t+\frac{m\omega ^2}{2}x_tx_t)}`$ (14)
$`=`$ $`\mathrm{}\omega {\displaystyle 𝑑tf(t)(a_t^{}a_t+\frac{1}{2})}`$ (15)
Projections onto its spectrum represent propositions about the time-averaged energy of the system.
For our purposes, the average number operator $`N`$ will be of prime importance. We can formally define it as follows:
$$N:=𝑑ta_t^{}a_t$$
(16)
The eigenvectors of this operator are vectors of the form
$$|n_f:=(n!)^{1/2}𝑑tf(t_1\mathrm{}t_n)a_{t_1}^{}\mathrm{}a_{t_n}^{}|0$$
(17)
These are also eigenvectors of the Hamiltonian. The average number operator has a highly degenerate spectrum as vectors of the above form have eigenvalue $`n`$ for all smearing functions $`f`$, as can be easily checked. We will denote the projection operator onto $`|n_f`$ as $`P_{n_f}`$; it represents a proposition about the average number of quanta present in a particular time interval.
### 2.3 Propositions within a finite time interval
We can write the average number operator defined above in the form, $`N=N_{f=1}`$ where $`N_f:=𝑑tf(t)a_t^{}a_t`$. This shows that there is a problem with the definition because the constant function $`f=1`$ is not a member of $`L^2()`$. However it is a member of $`L^2[a,b]`$ where $`[a,b]`$ is a finite interval of the real line. This suggests that we should really be dealing with propositions in a finite interval of time.
Consider again the proposition $`P_{n_f}`$. Intuitively the support of $`f`$ affects the time period in which the proposition is made. In other words if $`supp(f)[a,b]`$ then the proposition $`P_{n_f}`$ refers to the average number of particles during the time period $`[a,b]`$. We can formulate this rigorously by splitting up $`𝒱^{cts}`$ as follows:
$`𝒱^{cts}`$ $`:=`$ $`_t_t`$ (18)
$`=`$ $`𝒱^{[\mathrm{},a]}𝒱^{[a,b]}𝒱^{[b,\mathrm{}]}`$ (19)
where $`𝒱^{[a,b]}:=_{t[a,b]}_t`$. Now we can use the isomorphisms:
$$_{t[a,b]}e^{L_t^2[a,b]}e^{{}_{}{}^{}_{a}^{b}L_t^2[a,b]}[L^2[a,b]]$$
(20)
. Here, $`{}_{}{}^{}_{a}^{b}L_t^2[a,b]`$ is the direct integral Hilbert space over the interval $`[a,b]`$. An element of this Hilbert space, $`F`$, can be considered as a one-parameter family of elements of $`L^2[a,b]`$ which we denote by $`f_t`$, where $`t[a,b]`$. The inner product is defined as
$$(F,G)_{{}_{}{}^{}_{a}^{b}L_t^2[a,b]}:=_a^b𝑑t(f_t,g_t)_{L_t^2[a,b]}$$
(21)
From the right hand side of (20) we can see that $`𝒱^{[a,b]}`$ naturally carries a representation of the Lie algebra
$`[x_f,x_g]`$ $`=`$ $`0`$ (22)
$`[p_f,p_g]`$ $`=`$ $`0`$ (23)
$`[x_f,p_g]`$ $`=`$ $`i(f,g)`$ (24)
where $`f,gL^2[a,b]`$. The natural interpretation of these operators is that they are associated with time averaged propositions about position and momentum in the finite time interval $`[a,b]`$. We can form complex combinations of these operators in the usual way to define creation and annihilation operators. Projections onto the eigenvectors of the average number operator associated with these correspond to propositions about the average number of particles in the time interval $`[a,b]`$.
We can now see that propositions on $`𝒱^{cts}`$ smeared by functions in a finite time interval are isomorphic with propositions on $`𝒱^{[a,b]}`$ by
$$P_{n_f}𝕀_{𝒱^{[\mathrm{},a]}}P_{n_f}^{[a,b]}𝕀_{𝒱^{[b,\mathrm{}]}}$$
(25)
where $`fL^2[a,b]`$ and $`P_{n_f}^{[a,b]}𝒱^{[a,b]}`$. So from now on when we use the average number operator $`N`$ it should be understood that in fact we are averaging over a finite time interval *i.e.* we are smearing with functions $`fL^2[a,b]`$.
This is consistent with the definition of finite time interval projectors for coherent states given by Isham *et al* .
### 2.4 The Decoherence Functional
Isham *et al* and Anastopolous have defined decoherence functionals for continuous time projectors in the HPO scheme by considering projections onto coherent states. However, we are interested in propositions concerning the average number of quanta. These cannot be simply related to coherent states, so we will take a different approach and require our decoherence functional to respect the dynamical time translation symmetry of quantum theory. As the projectors onto eigenstates of $`N`$ commute with the Hamiltonian, we would expect the probability of any such proposition to be decided by its probability in the initial state. We will see that this is indeed the case and that these propositions also form a ‘canonical’ consistent set. We shall then require these conditions to hold in the HPO formalism to obtain a condition on the decoherence functional. Analogous remarks apply to any symmetry of the system, *i.e.* propositions regarding the spectral projectors of any operator which commutes with the Hamiltonian will form a consistent set and their probabilities will be decided by their value in the initial state.
Let us first examine the discrete time case for the SHO with single-time number operator defined by $`N^{st}:=a^{}a`$. Here we have time translation symmetry $`[H,N^{st}]=0`$, which corresponds to the conservation of the number of quanta. We begin by considering a 2-time history in the conventional set-up. It has eigenvectors $`|n:=(n!)^{1/2}(a^{})^n|0`$ and we denote Schrödinger picture projectors onto these vectors by $`P_n`$. The class operator takes a particularly simple form, $`C_{n_1n_2}:=P_{n_1}(t_1)P_{n_2}(t_2)=P_{n_1}P_{n_2}=\delta _{n_1n_2}P_{n_1}`$. The decoherence functional is then,
$`d_{SHO}(m_1m_2,n_1n_2)`$ $`:=`$ $`Tr_{}[C_{m_1m_2}\rho C_{n_1n_2}^{}]`$ (26)
$`=`$ $`\delta _{m_1m_2}\delta _{n_1n_2}Tr_{}[P_{m_1}\rho P_{n_1}]`$ (27)
$`=`$ $`\delta _{m_1m_2}\delta _{n_1n_2}\delta _{m_1n_1}\rho _{m_1n_1}`$ (28)
We can see that the fact that the projectors commute with the Hamiltonian means that they must all project onto the same state for the answer to be non-zero. This shows that propositions about the average number of particles,or more generally propositions about any symmetry of a system, always make up a consistent set. It is also clear that the probabilities assigned to these propositions depend on the initial state alone.
Now we examine this in the HPO scheme. The history space is $`𝒱^2=_{t_1}_{t_2}`$. We can write the above decoherence functional as a trace over $`^5:=_{t_0}_{t_1}_{t_2}_{t_1}_{t_2}`$ using the trick in :
$$d_{SHO}(m_1m_2,n_1n_2)=Tr_^5[\rho P_{m_1}P_{m_2}P_{n_1}P_{n_2}S_5]$$
(29)
Tracing over the initial Hilbert space we obtain,
$$d_{SHO}(m_1m_2,n_1n_2)=Tr_{𝒱^2𝒱^2}[P_{m_1}P_{m_2}P_{n_1}P_{n_2}Z]$$
(30)
where $`Z(𝒱^2𝒱^2)`$ and is defined in terms of its matrix elements in the energy basis as,
$$i_1\mathrm{}i_4|Z|j_1\mathrm{}j_4=\delta _{i_1j_2}\delta _{i_2j_3}\delta _{i_3j_4}\rho _{i_4j_1}$$
(31)
Now it is the operator $`Z`$ that contains the initial conditions and forces all the projectors to project onto the same state. In fact, by using these energy eigenstates we have removed the dynamics from the decoherence functional and are left only with the initial conditions and temporal structure encoded in the operator $`Z`$. Note that this does not uniquely define $`Z`$ as any $`Z^{}`$ defined by $`Z^{}=U^{}ZU`$ where $`U`$ is of the form $`e^{if(H)}e^{ig(H)}`$ has the same matrix elements if $`H`$ is the time-averaged Hamiltonian; $`H:=𝑑tH_t`$ .
Consider a continuous time energy proposition in standard history theory, represented by the class operator $`C_{\{n_t\}}:=\mathrm{\Pi }_tP_{n_t}(t)`$. Heuristically, this is going to be zero unless all of the $`n_t`$ are equal. If they are all equal, to $`n`$ say, then the infinite product will equal $`P_n`$. In this case the decoherence functional will give the same result as before:
$`d_{SHO}(\{m_s\},\{n_t\})`$ $`=`$ $`\delta _{mn}\rho _{mn}ifm_s=ms,n_t=nt`$ (32)
$`=`$ $`0\text{otherwise}`$ (33)
We can now understand the degeneracy in the spectrum of the average number operator in the HPO approach. It corresponds to the fact that the number of quanta is conserved and must be an integer. Therefore the time-averaged number of quanta must be an integer over any time period.
From the above discussion we require that the continuous time HPO decoherence functional satisfies
$$d_{SHO}^{cts}(n_f,m_g)=\delta _{mn}\rho _{mn}$$
(34)
for all functions $`f,g`$. This guarantees that:
1. The functional $`d_{SHO}^{cts}`$ assigns the correct probabilities to average number propositions. $`P_{n_f}`$ corresponds to the proposition “There are an average of $`n`$ quanta over the time interval $`tsupp(f)`$”. However, we know that the number of quanta is constant in time so the smearing function is irrelevant and that the probability of finding $`n`$ particles at any time is simply $`\rho _{nn}`$.
2. Number propositions still form a consistent set.
There is a class of operators $`Z^{cts}(𝒱^{cts}𝒱^{cts})`$ such that the decoherence functional can be written in the form:
$$d_{SHO}^{cts}(n_f,m_g):=Tr_{𝒱^{cts}𝒱^{cts}}[P_{n_f}P_{m_g}Z^{cts}]$$
(35)
such $`Z^{cts}`$ must satisfy,
$$m_fn_g|Z^{cts}|m_f^{}^{}n_g^{}^{}:=\delta _{mn}\delta _{nm^{}}\delta _{m^{}n^{}}\rho _{mn}$$
(36)
for all functions $`f,f^{},g,g^{}`$ as can be easily shown by taking the trace over energy eigenstates:
$$Tr_{𝒱^{cts}𝒱^{cts}}[X]=𝒟\mu [m_f]𝒟\mu [n_g]m_fn_g|X|m_fn_g$$
(37)
The measure $`𝒟\mu [m_f]`$ can be assumed to exist because there is a well-defined measure on $`𝒱^{cts}`$ defined in terms of coherent states . The condition (36) only defines $`Z^{cts}`$ up to a unitary transformation.
## 3 Quantum Field Theory
### 3.1 The HPO approach to QFT
We use throughout the signature $`(+,)`$. To construct an HPO version of canonical QFT on Minkowski space-time, $``$, we must first foliate $``$with a one parameter family of space-like surfaces using some timelike vector $`n^\mu `$, normalised by $`\eta _{\mu \nu }n^\mu n^\nu =1`$. Note that this corresponds to a choice of time direction as seen by some inertial observer. This choice obviously breaks Lorentz covariance and an important unsolved problem in the HPO programme is to show the equivalence of theories based on all such slicings. See for a relevant discussion. In this paper however, we will not consider this problem and just consider slices orthogonal to the vector $`n:=_{x^0}`$ where $`x^\mu `$ is the coordinate system on $``$in which our inertial observer is at rest. Now we consider a canonical 3-dimensional QFT to be defined on each Cauchy surface $`𝒞_t`$, where $`𝒞_t`$is defined by
$$𝒞_t:=\{m|x^0(m)=t\}$$
(38)
$``$is a globally hyperbolic space-time so these Cauchy surfaces are all isomorphic. In fact they are all homeomorphic to $`^3`$ so $`_t=[L^2(^3,d^3x)]`$ for all times $`t`$. We define the history algebra to be (in non-rigorous unsmeared form),
$`[\varphi _{t_1}(𝐱_\mathrm{𝟏}),\varphi _{t_2}(𝐱_\mathrm{𝟐})]`$ $`=`$ $`0`$ (39)
$`[\pi _{t_1}(𝐱_\mathrm{𝟏}),\pi _{t_2}(𝐱_\mathrm{𝟐})]`$ $`=`$ $`0`$ (40)
$`[\varphi _{t_1}(𝐱_\mathrm{𝟏}),\pi _{t_2}(𝐱_\mathrm{𝟐})]`$ $`=`$ $`i\mathrm{}\delta (t_1t_2)\delta ^3(𝐱_\mathrm{𝟏}𝐱_\mathrm{𝟐})`$ (41)
with $`𝐱_\mathrm{𝟏}𝒞_{t_1}`$. As shown in , the requirement that the Hamiltonian is self-adjoint is sufficient to select a representation of this algebra on the history space,
$$𝒱^{}:=_t_t[L^2(),d^4x)]$$
(42)
This representation is defined by the annihilation operator
$$a_t(𝐱):=\frac{1}{\sqrt{2}}\left(K_M^{1/4}\varphi _t(𝐱)+iK_M^{1/4}\pi _t(𝐱)\right)$$
(43)
where $`K_M`$ is defined by $`(K_Mf)(t,𝐱):=(_x^2+m^2)f(t,𝐱)`$. Equation (43) is a familiar equation in an unusual form. If we write $`\varphi _t(𝐱)`$ in terms of $`a_t(𝐤)`$ and $`a_{t}^{}{}_{}{}^{}(𝐤)`$ (defined as the 3 dimensional Fourier transforms of $`a_t(𝐱)`$ and $`a_{t}^{}{}_{}{}^{}(𝐱)`$ respectively) then we have
$$\varphi _t(𝐱)=\frac{d^3k}{(2\omega _k)^{1/2}}(e^{i𝐤.𝐱}a_{t}^{}{}_{}{}^{}(𝐤)+e^{i𝐤.𝐱}a_t(𝐤))$$
(44)
However, we must not let the familiar form of these equations make us forget that we are dealing with a history theory. In particular we must remember that the $`\varphi _t(𝐱)`$ operator is in the *Schrödinger* picture and the $`t`$ label that it carries represents the time that a particular proposition is made, *i.e.* it is a *logical* time quite separate from *dynamical* time. We can introduce dynamical time by using a one parameter unitary group as usual, but this must involve the introduction of a second time label:
$`\varphi _t(s,𝐱)`$ $`:=`$ $`e^{isH}\varphi _t(𝐱)e^{isH}`$ (45)
$`=`$ $`{\displaystyle }{\displaystyle \frac{d^3k}{(2\omega _k)^{1/2}}}(e^{i(𝐤.𝐱\omega _ks)}a_{t}^{}{}_{}{}^{}(𝐤)+e^{i(𝐤.𝐱\omega _ks)}a_t(𝐤)`$ (46)
where $`H:=𝑑tH_t(𝒱^{})`$ is the time-averaged Hamiltonian. Another difference with the canonical theory is that only projection operators have any meaning. Here we will be interested in propositions about the number of particles in a particular mode so we now define these:
$$N_𝐤:=𝑑ta_t^{}(𝐤)a_t(𝐤)$$
(47)
This operator has a highly degenerate spectrum as vectors of the form
$$|n_f^𝐤:=(n!)^{1/2}𝑑tf(t_1\mathrm{}t_n)a_{t_1}^{}{}_{}{}^{}(𝐤)\mathrm{}a_{t_n}^{}{}_{}{}^{}(𝐤)|0^M$$
(48)
are eigenvectors, with eigenvalue $`n`$ for all functions $`f`$. This degeneracy is the result of the fact that we are considering a free theory, so each $`N_𝐤`$ is separately conserved ($`[N_𝐤,H]=0`$) and must be an integer. Projectors $`P_{n_f^𝐤}`$ which project onto these vectors represent propositions about the average number of particles in mode $`𝐤`$ in the interval $`tsupp(f)`$. Symmetry implies that the propositions $`P_{n_f^𝐤}`$ form a canonical consistent set and that the probability of these propositions is decided by the probability in the initial state:
$$d^{}(m_f^𝐤,n_g^𝐤^{})=\delta _{mn}\delta ^3(𝐤𝐤^{})\rho _{m^𝐤n^𝐤^{}}^M$$
(49)
for all $`f,g`$, where $`\rho ^M(_{t_0})`$ is defined by its matrix elements:
$$\rho _{m^𝐤n^𝐤}^M:=m^𝐤|\rho ^M|n^𝐤$$
(50)
and $`|n^𝐤:=(a_{t_0}^{}(𝐤))^n|0^M`$.
We can write the decoherence functional in the form
$$d^{}(m_f^𝐤,n_g^𝐤^{})=Tr_{𝒱^{}𝒱^{}}[P_{m_f^𝐤}P_{n_g^𝐤^{}}𝒵^{}]$$
(51)
if $`𝒵^{}(𝒱^{}𝒱^{})`$ satisfies
$$m_f^{𝐤_\mathrm{𝟏}}n_g^{𝐤_\mathrm{𝟐}}|𝒵^{}|m_f^{}^{𝐤_\mathrm{𝟏}^{}}n_g^{}^{𝐤_\mathrm{𝟐}^{}}=\delta _{mn}\delta _{nm^{}}\delta _{m^{}n^{}}\delta (𝐤_\mathrm{𝟏}𝐤_\mathrm{𝟐})\delta (𝐤_\mathrm{𝟐}𝐤_\mathrm{𝟏}^{})\delta (𝐤_\mathrm{𝟏}^{}𝐤_\mathrm{𝟐}^{})\rho _{m^{𝐤_\mathrm{𝟏}}n^{𝐤_\mathrm{𝟏}}}$$
(52)
for all $`f,f^{},g,g^{}`$, which only defines $`𝒵^{}`$ up to a unitary transformation as before.
### 3.2 Canonical QFT on Rindler space-time
Consider an observer accelerating with constant acceleration, $`\alpha `$, through $``$. Let $`\xi ^\mu `$ denote the coordinates in which this observer is at rest. Then $`\xi ^\mu `$ are related to the coordinates $`x^\mu `$ by
$$(x^1)^2(x^0)^2=(\xi ^1)^2,x^0/x^1=tanh(\alpha \xi ^0),x^2=\xi ^2,x^3=\xi ^3$$
(53)
So constantly accelerating observers follow hyperbolae in $``$. These hyperbolae split into 2 sets depending on the sign of $`\xi ^1`$. Rindler space, $``$ , is defined to be the space covered by the coordinates $`\xi ^\mu `$ with $`\xi ^1>0`$. It corresponds to the wedge $`x>|t|`$ in ordinary Minkowski coordinates. Similarly, $``$ is defined to be the space covered by $`\xi ^\mu `$ with $`\xi ^1<0`$. It corresponds to the the wedge $`x<|t|`$. The metric in these coordinates takes the form,
$$ds^2:=g_{\mu \nu }d\xi ^\mu d\xi ^\nu =(\alpha \xi ^1)^2(d\xi ^0)^2(d\xi ^1)^2(d\xi ^2)^2(d\xi ^3)^2$$
(54)
The vector $`_{\xi ^0}`$ is a globally time-like Killing vector field in $``$. Therefore $``$is globally hyperbolic and we can formulate QFT canonically by using $`_{\xi ^0}`$ to select a particular representation of the canonical commutation relations. On non-globally hyperbolic space-times there is no globally time-like vector field and therefore no way to select one of the infinite number of unitarily inequivalent representations. This is the major difficulty in the theory of QFT in curved spaces. However, this does not concern us here and we proceed by solving the classical Klein-Gordon equation in curved space-time:
$$(g^{\mu \nu }_\mu _\nu +m^2)\varphi ^R(\xi )=0$$
(55)
Here, $`_\mu `$ is the covariant derivative associated with the metric (54). As shown in , equation (55) can be reduced to a Bessel equation with solutions $`u_\kappa ^R(\xi )`$. Following the canonical procedure we now second quantise and expand the quantum field in terms of creation and annihilation operators,
$$\varphi ^R(\xi ):=\frac{d^3\kappa }{(2\omega _\kappa )^{1/2}}\left(u_\kappa ^R(\xi )b^R(\stackrel{}{\kappa })^{}+u_\kappa ^R(\xi )b^R(\stackrel{}{\kappa })\right)$$
(56)
We can write down a similar equation for the field in $``$and because $`𝒞_\tau ^{}𝒞_\tau ^{}`$ is a Cauchy surface for $``$we can expand the field on $``$as :
$`\varphi (x)={\displaystyle }{\displaystyle \frac{d^3\kappa }{(2\omega _\kappa )^{1/2}}}(b^R(\stackrel{}{\kappa })\overline{u}_\kappa ^R(x)+b^R(\stackrel{}{\kappa })^{}\overline{u}_\kappa ^R(x)`$
$`+b^L(\stackrel{}{\kappa })\overline{u}_\kappa ^L(x)+b^L(\stackrel{}{\kappa })^{}\overline{u}_\kappa ^L(x))`$
where
$`\overline{u}_\kappa ^R(x)`$ $`:=`$ $`u_\kappa ^R(x)ifx`$ (57)
$`:=`$ $`0\text{otherwise}`$ (58)
and similarly for $`\overline{u}_\kappa ^L(x)`$.
Unruh used the analytic properties of the eigenfunctions $`u_\kappa ^R(x)`$ to find the Bogoliubov transformation between the above expansion and the usual one:
$$\varphi (x)=\frac{d^3k}{(2\omega _k)^{1/2}}\left(a(𝐤)e^{ik.x}+a^{}(𝐤)e^{ik.x}\right)$$
(59)
Unruh showed that the inertial vacuum can be written as a thermal density matrix in the Fock basis associated with the accelerating observer. It is this result that leads to the claim that an accelerating observer appears to be immersed in a thermal bath.
### 3.3 The Histories Approach
We now formulate QFT on Rindler space-time using the HPO approach and show how the the Unruh effect appears within the formalism.
Firstly we use the time coordinate of our accelerating observer to foliate $``$with a one parameter family of spacelike Cauchy surfaces $`𝒞_\tau ^{}`$where
$$𝒞_\tau ^{}:=\{r|\xi ^0(r)=\tau \}$$
(60)
The single time Hilbert space for the theory is then $`_\tau :=[L^2(𝒞_\tau ^{},d\mu )]`$ where $`d\mu (\xi )=(\alpha \xi ^1)^1d^3\xi `$ . The History space is
$$𝒱^{}:=_\tau _\tau [L^2(,d\mu d\tau )]$$
(61)
By analogy with equations (39) we define the history algebra to be
$`[\varphi _{\tau _1}(\stackrel{}{\xi _1}),\varphi _{\tau _2}(\stackrel{}{\xi _2})]`$ $`=`$ $`0`$ (62)
$`[\pi _{\tau _1}(\stackrel{}{\xi _1}),\pi _{\tau _2}(\stackrel{}{\xi _2})]`$ $`=`$ $`0`$ (63)
$`[\varphi _{\tau _1}(\stackrel{}{\xi _1}),\pi _{\tau _2}(\stackrel{}{\xi _2})]`$ $`=`$ $`i\mathrm{}\delta (\tau _1\tau _2)\delta ^3(\stackrel{}{\xi _1}\stackrel{}{\xi _2})`$ (64)
with $`\stackrel{}{\xi _1}𝒞_{\tau _1}`$.
The Hamiltonian of the real scalar field in $``$is
$$H_\tau ^R=\frac{1}{2}d^3\xi \alpha \xi ^1(\pi _\tau ^R(\stackrel{}{\xi })^2+_\xi \varphi _\tau ^R(\stackrel{}{\xi })._\xi \varphi _\tau ^R(\stackrel{}{\xi })+m^2\varphi _\tau ^R(\stackrel{}{\xi })^2)$$
(65)
where the vector field $`_\xi `$ is defined by $`_\xi :=_{\xi ^1}+_{\xi ^2}+_{\xi ^3}`$, and the dot product is taken using the 3-metric on $`𝒞_\tau ^{}`$; $`g^3=diag(1,1,1)`$ . Equation (65) has the same form for all $`\tau `$ so the representation of the history algebra in which $`H_\tau ^R`$ is self-adjoint is isomorphic on each $`_\tau `$ . The commutation relations of the smeared Hamiltonian with $`\varphi _\tau ^R(\stackrel{}{\xi })`$ and $`\pi _\tau ^R(\stackrel{}{\xi })`$ are
$`[H_f^R,\varphi _\tau ^R(\stackrel{}{\xi })]`$ $`=`$ $`i\mathrm{}\alpha \xi ^1f(\tau )\pi _\tau ^R(\stackrel{}{\xi })`$ (66)
$`[H_f^R,\pi _\tau ^R(\stackrel{}{\xi })]`$ $`=`$ $`i\mathrm{}f(\tau )K_R\varphi _\tau ^R(\stackrel{}{\xi })`$ (67)
where $`K_R`$ is defined by $`(K_Rf)(\tau ,\stackrel{}{\xi }):=(_\xi (\alpha \xi ^1_\xi )+\alpha \xi ^1m^2)f(\tau ,\stackrel{}{\xi })`$. Now we can follow the analysis of Isham *et al* to show that there is a unitary representation of the exponentiated commutation relations and that therefore the Hamiltonian exists as a self-adjoint operator in this representation. We can deduce the associated annihilation operators to be
$$b_\tau ^R(\stackrel{}{\xi })=\frac{1}{\sqrt{2}}\left(K_R^{1/4}\varphi _\tau ^R(\stackrel{}{\xi })+i\frac{\alpha \xi ^1}{K_R^{1/4}}\pi _\tau ^R(\stackrel{}{\xi })\right)$$
(68)
This defines a particular complexification of the test function space which is equivalent to a choice of positive and negative frequencies consistent with the Killing field $`_{\xi ^0}`$. Using these creation and annihilation operators we can build the Fock basis for the history theory. These equations can be written in a more familiar form by taking the spectral transform of the $`b_\tau ^R(\stackrel{}{\xi })`$ and $`b_{\tau }^{R}{}_{}{}^{}(\stackrel{}{\xi })`$, that is by expanding them in terms of the eigenfunctions of $`K_R`$, $`u_\kappa ^R(\stackrel{}{\xi })`$ <sup>*</sup><sup>*</sup>*these are just the functions $`u_\kappa ^R(\xi )`$, but with the time dependent part set to 1:
$$\varphi _\tau ^R(\stackrel{}{\xi }):=\frac{d^3\kappa }{(2\omega _\kappa )^{1/2}}\left(u_\kappa ^R(\stackrel{}{\xi })b_{\tau }^{R}{}_{}{}^{}(\stackrel{}{\kappa })+u_\kappa ^R(\stackrel{}{\xi })b_\tau ^R(\stackrel{}{\kappa })\right)$$
(69)
as before. There is obviously a strong similarity between the histories version of this problem and the canonical version. But from the histories perspective the result of Unruh shows nothing because a thermal density matrix is not a projection operator and so has no meaning when defined on $`𝒱^{}`$. Only elements of $`𝒫(𝒱^{})`$ and $`𝒫(𝒱^{})`$ are meaningful in a history theory as these can be considered as propositions about histories, *i.e.* as elements of $`𝒰𝒫^{}`$ and $`𝒰𝒫^{}`$. We have to change our approach so that we are talking about projectors onto eigenvectors of the average Rindler particle number operator:
$$N_\kappa :=𝑑\tau b_{\tau }^{R}{}_{}{}^{}(\stackrel{}{\kappa })b_\tau ^R(\stackrel{}{\kappa })$$
(70)
These vectors are of the form
$$|n_f^\kappa :=(n!)^{1/2}𝑑\tau f(\tau _1\mathrm{}\tau _2)b_{\tau _1}^{R}{}_{}{}^{}(\stackrel{}{\kappa })\mathrm{}b_{\tau _n}^{R}{}_{}{}^{}(\stackrel{}{\kappa })|\mathrm{\hspace{0.17em}\hspace{0.17em}0}^R$$
(71)
and have a degenerate spectrum in the same way as those for the inertial observer because we are still considering a free theory. Projectors onto these vectors represent propositions about the time-averaged number of particles in each mode, as seen by the accelerating observer.
The space of propositions about possible histories is not the same for the accelerating observer as for the inertial observer, but this is not the only difference. The decoherence functional associated with a quantum system depends on both the initial conditions and the Hamiltonian. The accelerating observer has a different Hamiltonian to the inertial observer and so has a different decoherence functional.
As before, the fact that $`[N_\kappa ,H^R]=0`$ implies that,
$$d^{}(m_f^\kappa ,n_g^\kappa ^{})=\delta _{mn}\delta ^3(\kappa \kappa ^{})\rho _{m^\kappa n^\kappa }^R$$
(72)
for all $`f,g`$, in notation which parallels that of (49) but now, $`\rho ^R(_{\tau _0}^R)`$ and $`|n^\kappa _{\tau _0}^R`$ is defined by:
$$|n^\kappa :=(b_{\tau _0}^{R}{}_{}{}^{}(\kappa ))^n|0^R$$
(73)
We can write this in the form,
$$d^{}(m_f^\kappa ,n_g^\kappa ):=Tr_{𝒱^{}𝒱^{}}[P_{m_f^\kappa }P_{n_g^\kappa }𝒵^{}]$$
(74)
for $`𝒵^{}(𝒱^{}𝒱^{})`$ defined similarly to the Minkowski case, (52).
### 3.4 The Unruh Effect
Finally we can see how the Unruh effect arises in the HPO formalism. Let us consider the situation in the inertial vacuum, *i.e.* the initial density matrix is
$$\rho _{n^𝐤n^𝐤}^M=\delta _{0n}$$
(75)
for all $`𝐤^\mathrm{𝟑}`$, where the matrix elements are taken in the Fock representation associated with the inertial observer. Note that this density matrix means that the probability of the inertial observer detecting $`n`$ particles in any mode is zero unless $`n=0`$:
$$d^{}(m_f^𝐤,n_g^𝐤^{})=\delta _{mn}\delta (𝐤𝐤^{})\delta _{0n}$$
(76)
The density matrix $`\rho ^M`$ is defined on some initial Hilbert space $`_{t_0}`$, but we can choose our Cauchy surfaces so that:
$$_{t_0}=_{\tau _0}^{}_{\tau _0}^{}$$
(77)
Using Unruh’s result on this initial Hilbert space we can write the inertial vacuum as a thermal density matrix in the representation associated with the accelerating observer. Tracing over $`_{\tau _0}^{}`$ we obtain the initial condition for the accelerating observer
$$\rho _{n^\kappa n^\kappa }^R=N_{\frac{2\pi }{\alpha }}(n\omega _\kappa )$$
(78)
where $`N_\beta (E)`$ is the thermal distribution giving the probability of a scalar particle having energy $`E`$ in a heat bath of inverse temperature $`\beta `$. Finally,
$$d^{}(n_f^\kappa ,m_g^\kappa ^{})=\delta _{mn}\delta (\kappa \kappa ^{})N_{\frac{2\pi }{\alpha }}(n\omega _\kappa ),$$
(79)
which shows that the accelerating observer detects a thermal spectrum at inverse temperature $`\beta =\frac{2\pi }{\alpha }`$, in agreement with the result of Unruh.
## 4 Conclusion
We have shown that it is possible to consider average number propositions within the continuous time HPO formalism. We have postulated a condition on the decoherence functional which ensures that energy propositions form a consistent set, as they do in the conventional theory, and which gives the correct probabilities for such propositions. This condition is defined for the SHO and for QFT but can easily be generalised to any system with symmetries as its construction involves only the matrix elements of the initial density matrix in the basis associated with the symmetry.
We have shown that the HPO scheme allows the construction of QFT in curved space-time and have re-derived the well-known result of Unruh within this scheme. In fact, the general nature of the HPO formalism - in particular its ability to cope with very general temporal support strucures and the associated non-unitary evolution - means that it can potentially be used to formulate QFT on much more general space-times such as non-globally hyperbolic space-times or those with topology change. This remains a task for future research.
Another potentially interesting avenue of research is to attempt to apply the formalism to other problems in conventional QFT such as scattering. Scattering type questions typically involve propositions such as ”there are $`n_1`$ particles of type $`1`$ at time $`t_1`$ and then $`n_2`$ particles of type $`2`$ at time $`t_2`$”. We cannot pose such questions in the formalism as presented here because we cannot embed discrete time propositions into the continuous time history space. The best we can do is to use propositions with support in a neighbourhood of $`t_1`$ and $`t_2`$ which we can choose arbitrarily small. Non-trivial scattering questions necessarily involve interactions and we haven’t considered these here, but in principle there is no reason why perturbation theory could not be developed.
## 5 Acknowledgements
I would like to thank Chris Isham for many useful discussions.
|
warning/0005/astro-ph0005151.html
|
ar5iv
|
text
|
# A Look at What Is (and Isn’t) Known About Quasar Broad Line Regions and How Narrow-Line Seyfert 1 Galaxies Fit In
## 1 Introduction
I was asked by the organizers if I would give “some kind of review of the properties of the BLR.” Why, at a workshop on NLS1s, should we consider BLRs in non-NLS1s? For me, and perhaps most people who attended this workshop, the interest of NLS1s has been that these extrema in the distributions of AGN properties might tell us more about how all AGNs work. This is, however, a two-way street: if NLS1 are “just” extrema of a continuum of BLR properties, then NLS1 BLRs are not fundamentally different from the BLRs of other AGNs. Therefore, a model of NLS1s must have the BLRs be consistent with those of non-NLS1s. I present here an overview of some general BLR results with which NLS1 models need to be consistent.
In BLR research the big question is “what role does the BLR play in the quasar phenomenon?” Some very basic specific questions include:
* Is the BLR one thing or more than one thing?
* Where is the BLR located? Or where are the components of the BLR located? (e.g., what is the distribution in $`R,\mathrm{\Theta },\mathrm{\Phi }`$?)
* Where is the BLR coming from or going to?
A survey of participants in the 1998 Nebraska BLR conference revealed a complete lack of consensus on the answers to these basic questions (Gaskell 1999).
## 2 Basic BLR Parameters
Observationally an AGN presents the astronomer with the four Stokes parameters as a function of projected velocity ($`v_{\mathrm{proj}}`$) and time ($`t`$). If we ignore polarization and subtract out various problems such as the continuum, absorption lines, and blended emission lines, then, for each line, $`i`$, we get $`L_i(v_{\mathrm{proj}},t)`$. If we ignore $`t`$ we get $`L_i(v_{\mathrm{proj}})`$, the line profiles. If we ignore $`v`$ we get $`L_i`$, the integrated intensities of the lines.
The traditional approach (going back to the late 1960s) has been to try to explain $`L_i`$ with a “typical” photoionized cloud. Obviously this assumes that the profiles are the same. We will see below that they are not, but this traditional approach has none the less been quite successful in explaining the stronger high-ionization lines and has produced a number of important results.
* The broad lines are produced mainly by photoionization. We can say this with confidence because the lines vary with the continuum, and the lines can vary in intensity by an order of magnitude. Not only do the high-ionization lines vary with the continuum, but the low ionization lines (such as Mg II and Fe II) do too. Other energy input sources might be present, but the dominant energy input is from photoionization.
* The mean physical conditions in the BLR do not vary a lot from object to object.
* The covering factor, $`\mathrm{\Omega }10\%`$.
* The density is $`10^{10}`$$`10^{12}`$ (or $`10^{13}`$) cm<sup>-3</sup>.
* The mass of gas seen in the BLR is of the order of a solar mass.
* If the BLR emission arises in free clouds with dimensions comparable the the Strömgren length, the number of clouds is enormous ($`10^{16}`$)
* Abundances are $``$ solar to a few times solar (e.g., Gaskell, Wampler, & Shields 1981; Hamann & Ferland 1993) which is consistent with what one expects in the centers of galaxies.
## 3 Optical Fe II Strengths in NLS1s
One of the major challenges for photo-ionization models has always been to explain the great strengths of low-ionization lines. In particular, explaining the origin of Fe II emission is a long-standing problem. This is reviewed by Suzy Collin-Souffrin elsewhere in this volume so I will not discuss it here, but I will point out one result (Gaskell 1985) which is not widely appreciated in discussion of the BLR of NLS1s.
One supposed characteristic of NLS1s is the great strength of the optical Fe II emission (e.g., Boller, Brandt, & Fink 1996). Wills (1982) pointed out for a heterogenous collection of AGNs that Fe II/H$`\beta `$ seemed to be roughly inversely proportional to the FWHM. As is well known, optical Fe II is hard to measure. I therefore looked at H$`\beta `$ and optical Fe II emission in the Lick Observatory sample of Seyfert 1 spectra of Osterbrock and his graduate students – a sample with very high signal-to-noise ratio spectra where the lines had been carefully de-blended. From this homogeneous sample I was able to confirm the FWHM vs. Fe II/H$`\beta `$ correlation of Wills (1982). However, the interesting thing is that while there is considerable object-to-object variation in the equivalent width of the optical Fe II (something that needs to be explained), there is no correlation of the equivalent width of the Fe II emission with FWHM (see figure 2 of Gaskell 1985). Instead, the FWHM vs. Fe II/H$`\beta `$ correlation arises because H$`\beta `$ gets weaker as the lines get narrower. So rather than trying to explain a (non-existent) anomalous strength of Fe II in NLS1s, we should be trying to explain the anomalous weakness of H$`\beta `$ in NLS1s. H$`\beta `$ is probably being thermalized because of the high density and this led to the prediction (Gaskell 1985) that UV spectra of narrow-line objects would show an increase in the strength of Si III\] relative to C III\]. This prediction has now been confirmed (Kuraszkiewicz et al. 1998; Wills et al. 1999).
## 4 The Need For (At Least) Two BLR Components
To be valid the traditional “typical photoionized cloud” analysis requires the $`L_i(v)`$ to be the same. If we look at the first moments of the line profiles (the line centroids) we find that the high-ionization lines are blueshifted relative to both the low-ionization lines and the rest frame of the host galaxy (Gaskell 1982). If we look at the second moments of the line profiles (the line widths) we find that FWHM $``$ ionization potential (Shuder 1982). These differences and other considerations have led some workers to argue that there are at least two fundamentally different BLR components (Gaskell 1987; Collin-Souffrin & Lasota 1988). Using the terminology of Gaskell (1987) these are:
* BLR I—a fairly traditional H II region producing the strong UV lines (e.g., as in the models of Davidson 1973). Collin-Souffrin & Lasota (1988) refer to BLR I as the “HIL” (high-ionization) BLR.
* BLR II—a large partially ionized zone (PIZ) which produces strong Balmer emission, Mg II, Fe II, Ca II, O I, etc.. Collin-Souffrin & Lasota (1988) refer to this as the “LIL” (low-ionization) BLR.
There are several reasons why I believe we need two (or more) components:
1. To Explain the Integrated Intensities
Collin-Souffrin et al. (1979, 1980, 1981) have argued that single photo-ionized clouds are unable to simultaneously explain the integrated line intensities of both the high-ionization lines and the hydrogen lines (especially Ly$`\alpha `$/H$`\beta `$). This is true a fortiori when one considers $`L_i(v)`$ (Snedden & Gaskell 1999ab, 2000). By constructing a grid of models using Gary Ferland’s photoionization code CLOUDY (Ferland 2000) and looking at the profiles of the strong BLR I lines, we get a hydrogen density $`n_H`$ const. $`(10^{11}`$ cm$`{}_{}{}^{3})`$, independent of the projected velocity, $`v`$, and an ionization parameter, $`U_1`$ const. $`(10^{1.5})`$, again independent of $`v`$. We found that the BLR I lines alone can be explained by either optically thick or optically thin models.
When we take the physical conditions we deduce from our analysis of the BLR I lines and try to predict the strengths of the BLR II lines we find that optically thick photoionization models can only explain the hydrogen-line ratios at low $`v`$; the wings of the Balmer lines are seriously overpredicted (see Figure 2 in Gaskell & Snedden 1999$`a`$). This means that BLR I is mostly optically thin. This has been already suggested by Ferland, Korista, & Peterson (1990) from emission-line variability considerations. When we try to use our grids of CLOUDY models to deduce conditions from the BLR II lines alone we find we need very optically thick clouds with $`n_H10^{13}`$ cm<sup>-3</sup> and a very low ionization parameter, $`U_110^310^4`$ to satisfy the constraints (Snedden & Gaskell 2000).
2. To Explain the Different Line Profiles
(a) As already noted above, BLR I lines are broader than BLR II lines.
(b) Not only are the line widths different, but Mathews & Wampler (1985) found the C IV (BLR I) and Mg II (BLR II) FWHMs in general to be uncorrelated for both radio-loud and radio-quiet AGNs. Further analysis (Gaskell & Mariupolskaya 2000), while confirming the independence of the C IV and Mg II FWHMs, shows that the relationship between BLR I and BLR II FWHMs is complicated (for example, Corbin 1993 found the FWHMs of C IV and H$`\beta `$ to be quite well correlated).
3. Because BLR I is Blueshifted
The standard explanation of the blueshifting (Gaskell 1982) is that BLR I is (at least approximately) spherical and radially outflowing but something (the accretion disc or inner torus) is blocking our view of the far side. There are, however, problems with this and other explanations have been offered (see below).
4. Because BLR I and BLR II Emission Come From Different Radii
One of the first results of what is often called “reverberation mapping” of BLRs (using light echoes to probe the structure of BLRs) was that the higher-ionization lines come from closer in to the central source (Gaskell & Sparke 1986). This has now been widely confirmed by major observational campaigns studying a number of objects. Not only are the responsivity weighted radii different, but the “transfer functions” (the light echoes seen in response to a $`\delta `$-function in the photoionizing continuum) of BLR I and BLR II have difference shapes (compare Krolik et al. 1991 with Horne, Welsh, & Peterson 1991). While we now recognize that BLR I is “stratified” and seems to have a different structure from BLR II, it is important to note that the distances of the different emitting regions from the black hole are not very different.
## 5 Spectropolarimetric Results
Spectropolarimetry is proving to be a powerful tool for probing the structure of AGNs. Three results stand out in particular:
1. The percentage polarization can be different for the BLR and the continuum.
2. The position angle (PA) of the polarization can also be different for the BLR and the continuum.
3. Both the percentage polarization and the PA can vary across the BLR line profiles (see, for example, Martel 1998).
The first two results imply that the size of the scattering region is comparable to the size of the BLR (and perhaps the scatterer is mixed in with the BLR?). The third result is very important for BLR kinematics because it implies that there is organized bulk motion in the BLR.
## 6 How is the BLR Gas Moving?
There is no consensus as to what the BLR is doing. Almost every kind of model of BLR kinematics is still under active consideration: random motions, Keplerian discs, infall, and various sorts of outflows or winds. I believe that understanding BLR kinematics is crucial for understanding the role of the BLR in AGNs and why NLS1s are different. My confident pre-1987 expectation was that BLR I at least was moving radially outwards. This was because of
* The blueshift of BLR I.
* The existence of broad absorption line quasars. Gas causing blueshifted absorption lines must be moving away from the quasar.
If the BLR is moving radially outwards then as the emission lines vary in response to changes in the ionizing continuum we expect the blue wings of the line to vary first, since this gas would be closest to our line of sight. I was therefore quite surprised to find (Gaskell 1988$`b`$) that the wings of C IV vary almost together and the red wing leads the blue wing slightly. This has now been confirmed for many objects (see, for example, Koratkar & Gaskell 1991). There are several possible solutions to the dilemma these results present:
1. The Shifts Are Not Caused by Bulk Radial Motion of BLR Clouds
Electron scattering can cause a blueshift of line profiles (Edmonds 1950). This has been modeled more recently by Kallman & Krolik (1986) and Ferrara & Pietrini (1993). At least a modest optical depth to electron scattering is needed ($`\tau _{es}0.5)`$. This mechanism has the advantage that we know the electrons are there because we see them in polarized light as the so-called “mirrors” in some Seyfert 2s. Continuity requires that the density of electrons increase inwards from the observed “mirror” to the BLR as $`r^2`$. The stellar wind model of Taylor (1998; see also these proceedings) offers another possible explanation of the blueshifting.
2. There is Bulk Radial Outflow, but it is Not Causing the Blue Wing to Vary First
There have been a couple of proposals in this category. These include the disc/wind model of Chiang & Murray (1996) and the hydromagnetic outflow model of Bottorff et al. (1999).
## 7 Line Emission from Discs
If the dominant motion of the BLR is in a plane, the line width depends on the orientation of the plane of motion (e.g., the disc) to the observer’s line of sight. A face-on viewing position has been a widely discussed explanation of the narrowness of the BLR lines in NLS1s. Wills & Browne (1986) showed that, for radio-loud AGNs at least, there is indeed a strong correlation between the FWHM of H$`\beta `$ and orientation.
If we are not viewing the disc face on, the Keplerian rotation will make us see double-humped profiles. Such displaced humps have long been seen, especially in the so-called ‘3C 390.3 objects,’ but the disc explanation of this line profile structure has been controversial (see Gaskell & Snedden 1999 and Sulentic et al. 1999). Problems arise because the disc model makes some definite predictions that were not verified:
* The blueshifted hump should always be stronger than the redshifted hump. In fact it has long been known that there are cases where the redshifted hump is much stronger than the blueshifted hump (e.g., Osterbrock & Cohen 1979) and our statistical study of line profiles (Gaskell & Snedden 1999) shows that red peaks are about equally likely to be stronger than blue peaks.
* The BLR responds to changes in the photoionizing continuum. If the source of ionizing radiation is located on the axis of symmetry, the observer sees the red and the blue sides of the line going up and down in intensity together. Gaskell (1988$`a`$) pointed out that this was not the case.
Despite my earlier objections I have now become convinced that we are seeing the signature of disc emission in some BLR line profiles. The International AGN Watch (IAW) monitoring of 3C 390.3 (Dietrich et al. 1998) showed that on a light-crossing timescale both the red and blue sides of the Balmer lines vary up and down together as predicted by the disc model. The profile changes reported earlier were on a much longer timescale and are not related to ionizing continuum changes. Interestingly, profile changes in typical (non-NLS1) radio-quiet AGNs are also independent of what the ionizing continuum is doing (see Peterson, Pogge, & Wanders 1999).
The line-profile statistics and long-term profile variability observations can all be explained if there is emission from a disc that is not azimuthally symmetric. The asymmetries might be hot spots or spiral wave patterns on the disc (see Gilbert et al. 1999 for some possible models). In at least one object (3C 390.3), a drift in wavelength consistent with orbital motion has been seen (Gaskell 1996; Eracleous et al. 1997).
It might be said that NLS1s are very different from broad-line radio galaxies, such as 3C 390.3, but it is important to recognize that some objects displaying disc-like emission line profiles are radio-quiet Seyfert galaxies. Also, difference spectra of “typical” Seyferts often show the same disc-like signatures seen in 3C 390.3 objects. We have carried out a study of the statistics of apparent structure in Balmer line profiles for samples of radio-loud and radio-quiet objects (Gaskell & Snedden 1999), and the results are consistent with the hypothesis that a disc-like component of BLR emission is present in all AGNs, regardless of type. Disc-like profiles are harder to detect in narrower-line objects whether they be radio-loud or radio-quiet. We believe the disc component is harder to detect in radio-quiet objects simply because they tend to have narrower line profiles.
## 8 Putting it All Together
There are many results and complications that space does not permit me to cover, but, while the jury is still out on many key issues (notably the kinematics), I think a picture is emerging: the BLR consists of two main components. One component (BLR II) arises predominantly in a very optically thick disc and the other (BLR I) arises predominantly from a more spherical distribution of more optically thin gas. If this is correct the leading model is therefore one of the disc-plus-atmosphere models proposed by Shields (1977), Collin-Souffrin et al. (1980), and others.
At a conference like this we are focusing in on differences between objects but it is important to realize the great similarities between between members of the quasar family, especially in continuum and BLR properties. The differences are mostly subtle, and it took a lot of work by many people over many years to discover them. There is also a continuum in the distribution of differences.
Why do BLR properties change as the BLR II lines get narrower? (i.e., as we go to NLS1s)? The thermalization of H$`\beta `$ and the increase in Si III\]/C III\] tells us that the density must be higher in NLS1s. Since the FWHM difference is most pronounced for the BLR II lines, and variations in the FWHM of BLR II are demonstrably influenced by orientation for non-NLS1s (Wills & Browne 1986), NLS1s are presumably seen face-on. I offer the following as a tentative attempt at an integrated picture of the NLS1 phenomenon:
1. The driving factor is a denser circumnuclear environment.
2. This denser environment provides an enhanced black hole fueling rate, as has been argued elsewhere in this conference.
3. The denser environment produces more obscuration of the central engine and produces a narrower opening angle through which we see BLR II more face-on than in non-NLS1s.
I would like to thank all the conference organizers and especially Thomas Boller for all their efforts in putting on this particularly well focused workshop.
|
warning/0005/hep-ph0005239.html
|
ar5iv
|
text
|
# A SEE-SAW MECHANISM FOR LARGE NEUTRINO MIXING FROM SMALL QUARK AND LEPTON MIXINGS
## 1 Introduction
Recent impressive results from Superkamiokande $`^\mathrm{?}`$ have cofirmed that the atmospheric neutrino anomaly can be successfully interpreted in terms of neutrino oscillations. Also the solar neutrino deficit, observed by several experiments $`^\mathrm{?}`$, is probably an indication of a different sort of neutrino oscillations. Since neutrino oscillations imply neutrino masses, one is forced to look for viable extensions of the Standard Model. In doing this, it is worth to stress that the extreme smallness of neutrino masses in comparison with quark and charged lepton masses seems to point in favour of a different nature of the former, maybe linked to lepton-number violation. Experimental facts on neutrino masses could then provide an indication on the very large energy scale where lepton-number is violated. Grand Unified Theories (GUTs) are certainly a very attractive framework where neutrino masses can be analyzed, because they predict - besides, of course, unification of the three gauge coupling constants - lepton and baryon-number violation. Since in GUTs all fermion masses are related, neutrino masses and mixings could also provide an insight on the mechanism for the generation of charged fermion masses. In particular the observation of a nearly maximal mixing angle for $`\nu _\mu \nu _\tau `$ is particularly impressive. At present solar neutrino mixings can be either large or very small, depending on which particular solution will eventually be established by the data. Large mixings in the neutrino sector are very interesting because a first guess was in favour of small mixings, in analogy to what is observed for left mixings in the quark sector. If confirmed, single or double maximal mixings can provide an important hint on the mechanisms that generate neutrino masses.
In the context of GUTs, many theoretical descriptions of large neutrino mixing(s) have been discussed (see Ref. 4 for a review). In most models large mixings are already present at the level of Dirac and/or Majorana matrices for neutrinos. Instead, here I discuss the interesting class of models where large, possibly maximal, neutrino mixings are generated by the see-saw mechanism starting from nearly diagonal Majorana and Dirac matrices for neutrinos, without fine-tuning or stretching small parameters into becoming large. A more complete discussion of the subject can be found in Ref. 1.
With neutrino masses settled, observation of proton decay will be the next decisive challenge remained to support or eventually put in crisis GUTs. In fact, SuperKamiokande $`^\mathrm{?}`$ is giving lower bounds on the proton life-time which, for certain decay modes, yet exclude part of the range generally predicted by GUTs. Then it seems interesting to see if it is possible to construct a simple but realistic GUT model $`^\mathrm{?}`$ which not only correctly reproduces the informations on neutrino masses and the actual bounds on proton decay, but also overcome the typical problem of these theories, that is the doublet-triplet splitting problem, without destroying gauge coupling unification.
## 2 Starting Assumptions
Since the experimental status of neutrino oscillations is still very preliminary, one has to make a number of assumptions on how the data will finally look like. Here I assume that only two distinct oscillation frequencies exist, the largest being associated with atmospheric neutrinos and the smallest with solar neutrinos. I assume that the hint of an additional frequency from the LSND experiment $`^\mathrm{?}`$ will disappear, thus avoiding the introduction of new sterile neutrino species. Dealing with only the three known species of light neutrinos, the atmospheric neutrino oscillations are interpreted as nearly maximal $`\nu _\mu \nu _\tau `$ oscillations while the solar neutrino oscillations correspond to the disappearance of $`\nu _e`$ into nearly equal fractions of $`\nu _\mu `$ and $`\nu _\tau `$. One has to be open minded to all the three most likely solutions for solar neutrino oscillations $`^\mathrm{?}`$: the two MSW solutions with small (SA) or large (LA) mixing angle, or the vacuum oscillation solution (VO). Assuming only two frequencies, given by
$$\mathrm{\Delta }_{sun}m_2^2m_1^2,\mathrm{\Delta }_{atm}m_3^2m_{1,2}^2,$$
(1)
there are two extreme possibilities for the mass eigenvalues:
$$\mathrm{A}:m_3>>m_{2,1}\mathrm{B}:m_1m_2m_3.$$
(2)
Configuration B imply a very precise near degeneracy of squared masses: it would need a relative splitting $`|\mathrm{\Delta }m/m|\mathrm{\Delta }m_{atm}^2/2m^210^3`$$`10^4`$ and a much smaller one for solar neutrinos, especially if explained by vacuum oscillations: $`|\mathrm{\Delta }m/m|10^{10}`$$`10^{11}`$. Foreseeing a GUT framework, it is reasonable to assume that the Dirac neutrino matrix has a strongly hierarchical structure, as is the case for charged fermions. So, it seems quite implausible that, starting from hierarchical Dirac matrices, one end up via the see-saw mechanism into a nearly perfect degeneracy of squared masses. As a consequence, here I will focus on models of type A with large effective light neutrino mass splittings and large mixings.
## 3 A $`2\times 2`$ Example
Reconciling large splittings with large mixing would seem difficult. Indeed, one could guess that, in analogy to what is observed for quarks, large splittings correspond to small mixings because only close-by states are strongly mixed. At the contrary, via the see-saw mechanism $`^\mathrm{?}`$, there are two particularly simple ways in which this can be realized.
Without loss of generality, leaving apart for the moment the eventual presence of flavor symmetries, one can go to the basis where both the charged lepton Dirac mass matrix $`m_D^l`$ and the Majorana matrix $`M`$ for the right-handed neutrinos are diagonal. For simplicity, let’s start assuming that the role of the first generation is not crucial in the mechanism for the generation of neutrino masses, so than one can, with good approximation, work in the 2 by 2 case (in the next section I will however relax this condition). If one writes $`m_D`$ (defined by $`\overline{R}m_DL`$) and $`M`$ in the most general way:
$$m_D=v\left[\begin{array}{cc}a& b\\ c& 1\end{array}\right],M=\left[\begin{array}{cc}M_2& 0\\ 0& M_3\end{array}\right],$$
(3)
where $`v`$ is a vacuum expectation value, $`a,b`$ and $`c`$ are Yukawa couplings, then, via the see-saw, one obtains:
$$m_\nu =\frac{v^2}{M_3}\left[\begin{array}{cc}\frac{a^2}{M_2}+\frac{c^2}{M_3}& \frac{ab}{M_2}+\frac{c}{M_3}\\ \frac{ab}{M_2}+\frac{c}{M_3}& \frac{b^2}{M_2}+\frac{1}{M_3}\end{array}\right].$$
(4)
The request of large splittings among the light neutrino’s eigenvalues is equivalent to demanding that the determinant of the previous matrix is much smaller than its trace. It is then possibile to see at first sight that two very natural cases arise, respectively when the terms with $`M_3`$ or $`M_2`$ at the denominator are dominant.
One simple example of the first case is realized if $`M_2M_3`$ and $`a,b1`$. In order to have a large splitting, one must have $`c1`$, that is the right-handed neutrino of the third generation couples with the same strenght $`^\mathrm{?}`$ to left-handed $`\nu _\mu `$ and $`\nu _\tau `$. The heaviest mass for light neutrinos results $`m_3v^2/M_3`$. Since in the hierarchical case, the data from SuperKamiokande suggest $`m_30.05eV`$, if one assumes that $`v`$ is a typical weak scale, namely 250 GeV, then $`M_310^{15}GeV`$, just the order of a GUT scale. It is worth to stress that this first mechanism is based on asymmetric Dirac matrices with, in the case of the example, a large left-handed mixing already present in the Dirac matrix. It has been observed $`^{\mathrm{?},\mathrm{?}}`$ that in SU(5) GUT left-handed mixings for leptons tend to correspond to right-handed mixings for $`d`$ quarks (in a basis where $`u`$ quarks are diagonal). Since large right-handed mixings for quarks are not in contrast with experiment, viable GUT models that correctly reproduce the data on fermion masses and mixings can be constructed following this mechanism.
An alternative possibility $`^\mathrm{?}`$ is to have the dominance of the terms with $`M_2`$ at the denominator. This is achieved for any $`c<1`$ if $`a^2,b^2>M_2/M_3`$. The request for large splitting is then equivalent to require also $`ab`$. Now it is the second generation right-handed neutrino which is particularly light and which couples with the same strenght $`^\mathrm{?}`$ to left-handed $`\nu _\mu `$ and $`\nu _\tau `$. In order to be more specific, consider one particular example with symmetric matrices. These matrices are interesting because, for instance, one could want to preserve left-right symmetry at the GUT scale. Then, the observed smallness of left-handed mixings for quarks would also demand small right-handed mixings. Starting from
$$m_D=v\left[\begin{array}{cc}ϵ& xϵ\\ xϵ& 1\end{array}\right],M^1=\frac{1}{M_3}\left[\begin{array}{cc}r_2& 0\\ 0& 1\end{array}\right],$$
(5)
where $`ϵ`$ is a small number, $`x`$ is of O(1) and $`r_2M_3/M_2`$, then, via the see-saw, it is sufficient that $`ϵ^2r_21`$ in order to have approximately:
$$m_\nu =\frac{v^2}{M_3}ϵ^2r_2\left[\begin{array}{cc}1& x\\ x& x^2\end{array}\right].$$
(6)
The determinant is naturally vanishing so that the mass eigenvalues are widely split and for $`x1`$ the mixing is nearly maximal. It is exactly maximal if $`x=1`$. The see-saw mechanism has created large mixing from almost nothing: all relevant matrices entering the see-saw mechanism are nearly diagonal $`^{\mathrm{?},\mathrm{?}}`$, that is they are diagonalized by transformations that go into the identity in the limit of vanishing $`ϵ`$. Clearly, the crucial factorization of the small parameter $`ϵ^2`$ only arises if the light Majorana eigenvalue is coupled to $`\nu _\mu `$ and $`\nu _\tau `$ with comparable strength, that is $`x1`$. An interesting feature of this second case, in connection with a possible realization within a GUT scheme, is that it requires $`M_3>v^2/m_3`$, so that one can push $`M_3`$, the scale of lepton-number violation, beyond the GUT scale. This is desirable because, for instance, this is expected in SU(5) if right-handed neutrinos are present and also in the breaking of SO(10) to SU(5).
Summarizing, the second case require a peculiar hierarchy in the Majorana eigenvalues in order to work, but it has however the good characteristic that it can be realized even with nearly diagonal matrices.
## 4 Generatization to the $`3\times 3`$ Case
It is straightforward to extend the previous model to the 3 by 3 case. One simple class of examples with symmetric mass matrices is the following one. Starting from
$$m_D=v\left[\begin{array}{ccc}ϵ^{\prime \prime }& ϵ^{}& yϵ^{}\\ ϵ^{}& ϵ& xϵ\\ yϵ^{}& xϵ& 1\end{array}\right],M^1=\frac{1}{\mathrm{\Lambda }}\left[\begin{array}{ccc}r_1& 0& 0\\ 0& r_2& 0\\ 0& 0& r_3\end{array}\right],$$
(7)
where, unless otherwise stated, $`x`$ and $`y`$ are O(1); $`ϵ`$, $`ϵ^{}`$ and $`ϵ^{\prime \prime }`$ are independent small numbers and $`r_iM_3/M_i`$. One expects $`ϵ^{\prime \prime }ϵ^{}ϵ1`$ and, perhaps, also $`r_1r_2r_3=1`$, if the hierarchy for right-handed neutrinos follows the same pattern as for known fermions. Depending on the relative size of the ratios $`r_i/r_j`$, $`ϵ/ϵ^{}`$ and $`ϵ^{}/ϵ^{\prime \prime }`$, it is possible to have models with dominance of any of the $`r_{1,2}`$. For example, setting $`x=1`$ (keeping $`y`$ of O(1)) and assuming $`r_2ϵ^2r_1ϵ^2,r_3`$, together with $`r_2ϵ^2r_1ϵ^{\prime \prime 2}`$ and $`r_2ϵr_1ϵ^{\prime \prime }`$, with good accuracy we obtain:
$$m_\nu =\frac{v^2}{\mathrm{\Lambda }}r_2ϵ^2\left[\begin{array}{ccc}\frac{ϵ^2}{ϵ^2}& \frac{ϵ^{}}{ϵ}& \frac{ϵ^{}}{ϵ}\\ \frac{ϵ^{}}{ϵ}& 1+\frac{r_1ϵ^2}{r_2ϵ^2}& 1\\ \frac{ϵ^{}}{ϵ}& 1& 1+\frac{r_3}{r_2ϵ^2}\end{array}\right].$$
(8)
Since the subdeterminant of the 23 block is vanishing, the eigenvalues are widely split. Having set $`x=1`$ the atmospheric neutrino mixing is nearly maximal. The solar neutrino mixing is instead generically small in these models, being proportional to $`ϵ^{}/ϵ`$. Thus the SA-MSW solution is obtained. It is easy to find set of parameter values that lead to an acceptable phenomenology within these solutions. As an illustrative example take:
$$ϵ\lambda ^4,ϵ^{}\lambda ^6,ϵ^{\prime \prime }\lambda ^{12},r_1\lambda ^{12},r_2\lambda ^9,$$
(9)
where $`\lambda \mathrm{sin}\theta _C`$, $`\theta _C`$ being the Cabibbo angle. The neutrino mass matrices than become
$$m_D=v\left[\begin{array}{ccc}\lambda ^{12}& \lambda ^6& \lambda ^6\\ \lambda ^6& \lambda ^4& \lambda ^4\\ \lambda ^6& \lambda ^4& 1\end{array}\right],M=\mathrm{\Lambda }\left[\begin{array}{ccc}\lambda ^{12}& 0& 0\\ 0& \lambda ^9& 0\\ 0& 0& 1\end{array}\right]$$
(10)
and, in units of $`v^2/\mathrm{\Lambda }`$, we obtain: $`m_31/\lambda ,m_21,m_1\lambda ^4`$. The solar mixing angle $`\theta _{12}`$ is of order $`\lambda ^2`$, suitable to the SA-MSW solution. Also $`\theta _{13}\lambda ^2`$.
Models based on symmetric matrices are directly compatible with left-right symmetry and therefore are naturally linked with SO(10). This is to be confronted with models that have large right-handed mixings for quarks, which, in SU(5), can be naturally translated into large left-handed mixings for leptons. In this connection it is interesting to observe that the proposed textures for the neutrino Dirac matrix can also work for up and down quarks. For example, the matrices
$$m_D^u\left[\begin{array}{ccc}0& \lambda ^6& \lambda ^6\\ \lambda ^6& \lambda ^4& \lambda ^4\\ \lambda ^6& \lambda ^4& 1\end{array}\right],m_D^d\left[\begin{array}{ccc}0& \lambda ^3& \lambda ^3\\ \lambda ^3& \lambda ^2& \lambda ^2\\ \lambda ^3& \lambda ^2& 1\end{array}\right],$$
(11)
where for each entry the order of magnitude is specified in terms of $`\lambda \mathrm{sin}\theta _C`$, lead to acceptable mass matrices and mixings. In fact $`m_u:m_c:m_t=\lambda ^8:\lambda ^4:1`$ and $`m_d:m_s:m_b=\lambda ^4:\lambda ^2:1`$. The $`V_{CKM}`$ matrix receives a dominant contribution from the down sector in that the up sector angles are much smaller than the down sector ones. The same kind of texture can also be adopted in the charged lepton sector.
## 5 Bimixing
The solar mixing angle is generically small in the class of models explicitly discussed above. However, small mixing angles in the Dirac and Majorana neutrino mass matrices do not exclude a large solar mixing angle. For instance, this is generated from the asymmetric, but nearly diagonal mass matrices:
$$m_D=v\left[\begin{array}{ccc}\lambda ^6& \lambda ^6& 0\\ 0& \lambda ^4& \lambda ^4\\ 0& 0& 1\end{array}\right],M=\mathrm{\Lambda }\left[\begin{array}{ccc}\lambda ^{12}& 0& 0\\ 0& \lambda ^{10}& 0\\ 0& 0& 1\end{array}\right].$$
(12)
They give rise to a light neutrino mass matrix of the kind:
$$m_\nu =\left[\begin{array}{ccc}\lambda ^2& \lambda ^2& 0\\ \lambda ^2& 1& 1\\ 0& 1& 1\end{array}\right]\frac{v^2}{\lambda ^2\mathrm{\Lambda }},$$
(13)
which is diagonalized by large $`\theta _{12}`$ and $`\theta _{23}`$ and small $`\theta _{13}`$. The mass hierarchy is suitable to the large angle MSW solution.
## 6 Outlook and Conclusions
Other comments about this mechanism are contained in Ref. 1. For instance here we show that:
i) the results obtained are stable under renormalization from the high energy scale where the mass matrices are produced down to the electroweak scale; in fact, if light neutrino masses are hierarchical, one always expects renormalization effects to be negligible $`^\mathrm{?}`$;
ii) it is possible to construct specific realizations of the mechanism sketched here, e.g. in the context of SU(5) $`\times `$ broken horizontal flavour symmetries.
Summarizing, in most models $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$ that describe neutrino oscillations with nearly maximal mixings, there appear large mixings in at least one of the matrices $`m_D`$, $`m_D^l`$, $`M`$ (i.e. the neutrino and charged-lepton Dirac matrices and the right-handed Majorana matrix). In this contribution I have discussed the peculiar possibility that large neutrino mixing is only produced by the see-saw mechanism starting from all nearly diagonal matrices. Although this possibility is certainly rather special, models of this sort can be constructed without an unrealistic amount of fine-tuning and are well compatible with grand unification ideas and the related phenomenology for quark and lepton masses.
## Acknowledgments
Many thanks go to the organizers for the stimulating and relaxed atmosphere of this conference, held in such beautiful and inspiring surroundings. It is pleasure to thank Guido Altarelli and Ferruccio Feruglio for the enjoyable collaboration on which this talk is based.
## References
|
warning/0005/hep-ph0005243.html
|
ar5iv
|
text
|
# Two-loop matching conditions for (MS)̄ parton densities
## Abstract
We discuss how the the operator product expansion (OPE) can be used to derive asymptotic expressions for certain integrals. This yields operator matrix elements (OME’s) which determine the matching conditions for $`\overline{\mathrm{MS}}`$ parton densities across heavy flavour thresholds. Then we construct four and five-flavour densities from a three-flavour set via the evolution of the AP equation using LO and NLO splitting functions.
It is well known that $`\alpha _s(\mu ^2,n_f,\mathrm{\Lambda }(n_f))`$ in pQCD requires matching conditions as the scale $`\mu `$ crosses flavour matching points. At these points the number of light-flavours $`n_f`$ changes by unity so the QCD scale parameter $`\mathrm{\Lambda }(n_f)`$ in the solution of the differential equation for the $`\beta `$ function is redefined to make the running coupling continuous. When heavy quarks are included another scale $`m`$, the mass of the heavy quark, enters and the matching conditions are more complicated. The precise relations which need to be satisfied are given in , . In order lowest order pQCD one can choose to make the $`\alpha _s`$ continuous across heavy flavour thresholds at $`\mu =m`$. However this does not hold in higher order pQCD as the matching conditions in the $`\overline{\mathrm{MS}}`$ scheme then contain non-logarithmic terms. Hence there is a discontinuity in $`\alpha _s`$ at $`\mu =m`$.
Recently the analogous problem of deriving the two-loop matching conditions on parton densities as the mass factorization scale crosses the heavy flavour thresholds has been solved in . The way this was done is as follows. We examined the large $`Q^2`$ limit of the heavy quark coefficient functions which appear in NLO perturbation expressions for heavy quark extrinsic pair production in deep inelastic scattering. These quantities are functions of the virtuality of the photon probe $`\sqrt{Q^2}`$, the mass of the heavy quark $`m`$, the renormalization scale $`\mu `$, which is chosen equal to the mass factorization scale, and the partonic Bjorken scaling variable $`z`$. The number of heavy $`D^{}`$ mesons produced in deep inelastic scattering can be derived by convoluting these heavy ($`c\overline{c}`$) quark coefficient functions with appropriate combinations of three-flavour light parton densities (u,d,s and g) and with heavy ($`c\overline{c}`$) quark fragmentation functions . Note that the heavy $`c\overline{c}`$ pair only appears in the final state. Unfortunately we do not have analytic expressions for all these heavy quark coefficient functions. Some only exist as two-dimensional integrals over very complicated expressions. However there are convenient tables for all of them in .
In the limit $`Q^2m^2`$ the complicated integrals in the heavy quark coefficient functions reduce to terms with powers of $`\mathrm{ln}(Q^2/\mu ^2)`$ and $`\mathrm{ln}(\mu ^2/m^2)`$ multiplied by functions of the variable $`z`$. These results can be reexpressed as convolutions of light-mass coefficient functions $`𝒞(z,Q^2/\mu ^2)`$ which contain the terms with powers in $`\mathrm{ln}(Q^2/\mu ^2)`$ and OME’s $`A(z,\mu ^2/m^2)`$ which contain the powers in $`\mathrm{ln}(\mu ^2/m^2)`$ . The way we evaluated these OME’s is described in so we only give an outline here. We wrote the heavy quark coefficient functions in terms of dispersion integrals for off-shell forward Compton scattering as is normally done for the OPE in deep inelastic scattering. We then changed variables to write the dispersion integral in terms of a variable $`z^{}`$ which is between zero and unity. Next we expanded the denominator in a Taylor series in $`z^{}`$. To take the limit $`Q^2m^2`$ of the dispersion integral we add and subtract the same dispersion integral where we take the limit $`Q^2m^2`$ in the integrand. This integrand contains the OPE of the standard heavy quark (Q) nonsinglet and singlet operators in pQCD taken between states with momentum $`k`$, namely
$`<Q(k)|O_{Q,\mu _1,\mu _2,\mathrm{}.\mu _n}(0)|Q(k)>.`$ (1)
The heavy quark operator
$`O_{Q,\mu _1,\mu _2,\mathrm{}.\mu _n}(x)=\overline{\psi }(x)\gamma _{\mu _1}D_{\mu _2}\mathrm{}\mathrm{}D_{\mu _n}\psi (x),`$ (2)
is a gauge invariant operator containing the heavy quark field $`\psi (x)`$ and the covariant derivative $`D_\mu =_\mu +igA_\mu `$. It can be shown that the original integral minus the integral involving the OPE does not contain any mass singularities as $`m0`$ so it cannot depend on the heavy quark mass $`m`$ and therefore only contains terms with powers of $`\mathrm{ln}(Q^2/\mu ^2)`$. Hence the integrals which contain the evaluation of the OME’s in the OPE yield all the terms containing powers of $`\mathrm{ln}(\mu ^2/m^2)`$. This means that analytic expressions for the two-loop OME’s with one heavy quark loop and light-quark or gluon incoming and outgoing states contain the information we require to extract the $`A(z,\mu ^2/m^2)`$. Of course the actual evaluation of the five OME’s which exist in order $`\alpha _s^2`$ requires the introduction of infrared and ultraviolet regulators, the use of gauge invariant operators, contractions with light-like four vectors to make the projections and $`\overline{\mathrm{MS}}`$ renormalization. The results of this analysis are encapsulated in expressions like
$`\stackrel{~}{A}_{Qq}^{\mathrm{PS},(2)}(z,\mu ^2/m^2)=A_1(z)\mathrm{ln}^2(\mu ^2/m^2)`$
$`+A_2(z)\mathrm{ln}(\mu ^2/m^2)+A_3(z),`$ (3)
where
$`A_1(z)=C_FT_f[8(1+z)\mathrm{ln}z`$
$`16/(3z)4+4z+16z^2/3],`$ (4)
$`A_2(z)=C_FT_f[8(1+z)\mathrm{ln}^2z`$
$`+(8+40z+64z^2/3)\mathrm{ln}z`$
$`+160/(9z)16+48z448z^2/9],`$ (5)
$`A_3(z)=C_FT_f\{(1+z)[32\mathrm{S}_{1,2}(1z)`$
$`+16\mathrm{ln}z\mathrm{Li}_2(1z)`$
$`16\zeta (2)\mathrm{ln}z4/3\mathrm{ln}^3z]`$
$`+(32/(3z)+88z32z^2/3)\mathrm{Li}_2(1z)`$
$`+(32/(3z)8+8z+32z^2/3)\zeta (2)`$
$`+(2+10z+16z^2/3)\mathrm{ln}^2z`$
$`(56/3+88z/3+448z^2/9)\mathrm{ln}z448/(27z)`$
$`4/3124z/3+1600z^2/27\}.`$ (6)
The tilde indicates that an overall factor of $`n_f`$ has been extracted from the function and the (2) in the superscript means this is the second order term in an expansion in $`a_s=\alpha _s/(4\pi )`$. The five functions $`\stackrel{~}{A}_{Qq}^{\mathrm{PS},(2)}(z,\mu ^2/m^2)`$, where PS denotes pure singlet under the flavour group (i.e., no non-singlet projection exists), $`\stackrel{~}{A}_{Qg}^{\mathrm{S},(2)}(z,\mu ^2/m^2)`$, where S denotes singlet under the flavour group, $`A_{gg,Q}^{\mathrm{S},(2)}(z,\mu ^2/m^2)`$, $`A_{gq,Q}^{\mathrm{S},(2)}(z,\mu ^2/m^2)`$, and $`A_{qq,Q}^{\mathrm{NS},(2)}(z,\mu ^2/m^2)`$ where NS denotes non-singlet under the flavour group, which exist in order $`\alpha _s^2`$ pQCD are given in . Alternative discussions of their derivation and use are given in . Note that they contain nonlogarithmic terms such as $`A_3(z)`$ in Eq.(6) in order $`\alpha _s^2`$ so there is no scale $`\mu `$ where we can make them all vanish. Since we know the four-flavour light mass coefficient functions $`𝒞(z,Q^2/\mu ^2)`$ in order $`\alpha _s^2`$ we can analytically evaluate the convolutions with the appropriate $`A`$’s to obtain asymptotic expressions for the heavy quark coefficient functions. They were given in . As far as this workshop is concerned we would like to point out that this ”inverse mass factorization method” is an elegant use of the OPE to obtain asymptotic expansions of integrals.
Normally parton densities are fitted to specific functions of $`x`$ at a scale $`\mu `$ and the AP equations then govern the evolution of these densities to other scales. Suppose one begins with a three-flavour set containing densities for u,d,s quarks and the gluon g. Then the above results allow one to define four-flavour parton densities at scales $`\mu m_c`$ from the input set of three-flavour densities in fixed-order perturbation theory (FOPT). Let the $``$ symbol denote the convolution integral $`fg=f(x/y)g(y)𝑑y/y`$, where $`xy1`$, then we define the charm density
$`f_{c+\overline{c}}(n_f+1,\mu ^2)=`$
$`a_s(n_f,\mu ^2)\stackrel{~}{A}_{Qg}^\mathrm{S}\left({\displaystyle \frac{\mu ^2}{m_c^2}}\right)f_g^\mathrm{S}(n_f,\mu ^2)`$
$`+a_s^2(n_f,\mu ^2)[\stackrel{~}{A}_{Qq}^{\mathrm{PS}}\left({\displaystyle \frac{\mu ^2}{m_c^2}}\right)f_q^\mathrm{S}(n_f,\mu ^2)`$
$`+\stackrel{~}{A}_{Qg}^\mathrm{S}\left({\displaystyle \frac{\mu ^2}{m_c^2}}\right)f_g^\mathrm{S}(n_f,\mu ^2)],`$ (7)
the singlet gluon density
$`f_g^\mathrm{S}(n_f+1,\mu ^2)=f_g^\mathrm{S}(n_f,\mu ^2)`$
$`+a_s(n_f,\mu ^2)A_{gg,Q}^\mathrm{S}({\displaystyle \frac{\mu ^2}{m_c^2}})f_g^\mathrm{S}(n_f,\mu ^2)`$
$`+a_s^2(n_f,\mu ^2)[A_{gq,Q}^\mathrm{S}({\displaystyle \frac{\mu ^2}{m_c^2}})f_q^\mathrm{S}(n_f,\mu ^2),`$
$`+A_{gg,Q}^\mathrm{S}({\displaystyle \frac{\mu ^2}{m_c^2}})f_g^\mathrm{S}(n_f,\mu ^2)],`$ (8)
and the light mass quark densities
$`f_{k+\overline{k}}(n_f+1,\mu ^2)=f_{k+\overline{k}}(n_f,\mu ^2)`$
$`+a_s^2(n_f,\mu ^2)A_{qq,Q}^{\mathrm{NS}}\left({\displaystyle \frac{\mu ^2}{m_c^2}}\right)f_{k+\overline{k}}(n_f,\mu ^2),`$ (9)
for $`n_f=3`$ and $`m_c^2\mu ^2<m_b^2`$. Note that we have suppressed the $`x`$ dependence to make the notation more compact. These expressions were used in to construct a variable flavour number scheme (VFNS) for the heavy quark contributions to the deep inelastic structure functions.
Note however that the above procedure does not resum the potentially large terms in $`\mathrm{ln}(\mu ^2/m_c^2)`$ which are explicitly left in the parton densities. To do this we need to evolve the above densities via the AP equation rather than using FOPT. This is new work in using three-flavour densities at small scales from . The latter LO and NLO densities are started at very small scales $`\mu _0`$ below the mass of the charmed quark. Hence three flavor evolution proceeds from the initial $`\mu _0^2`$ to the scale $`\mu ^2=m_c^2=1.96`$ $`(\mathrm{GeV}/\mathrm{c}^2)^2`$. In this region $`\alpha _s`$ is large so we had to be very careful to get numerically accurate solutions of the evolution equation. Fortunately there are standard inputs and tables in with which we could compare the parton densities from our evolution code. We chose the matching scale $`\mu `$ at the mass of the charm quark $`m_c`$ so that all the $`\mathrm{ln}(\mu ^2/m_c^2)`$ terms in the OME’s vanish at this point leaving only the nonlogarithmic pieces in the order $`\alpha _s^2`$ OME’s to contribute to the right-hand-sides of Eqs. (7), (8) and (9). Note that the LO and NLO charm densities vanish at the scale $`\mu =m_c`$ since
$`\stackrel{~}{A}_{Qg}^{\mathrm{S},(1)}(z,\mu ^2/m^2)=4T_f(z^2+(1z)^2)\mathrm{ln}(\mu ^2/m^2),`$
does not have a non-logarithmic term. The NNLO charm density starts off with a finite $`x`$-dependent shape in order $`a_s^2`$ determined by
$`f_{c+\overline{c}}(n_f+1,m_c^2)=`$
$`a_s^2(n_f,m_c^2)[\stackrel{~}{A}_{Qq}^{\mathrm{PS}}(1)f_q^\mathrm{S}(n_f,m_c^2)`$
$`+\stackrel{~}{A}_{Qg}^\mathrm{S}(1)f_g^\mathrm{S}(n_f,m_c^2)],`$ (11)
with $`n_f=3`$. Hence the OME’s provide the boundary condition for the evolution of the (massless) charm density. Also note that we ordered the terms on the right-hand-side of Eq.(11) in powers of $`\alpha _s`$ so that the result contains a product of NLO OME’s and LO parton densities, although this is not evident here. The result is then strictly order $`a_s^2`$ and should be multiplied by order $`a_s^0`$ coefficient functions when forming the zero-mass variable flavour number scheme (ZM-VFNS) charm density contribution to the deep inelastic structure functions.
The four-flavour gluon density is also generated at the matching point in the same way. At $`\mu =m_c`$ we define
$`f_g^\mathrm{S}(n_f+1,m_c^2)=f_g^\mathrm{S}(n_f,m_c^2)`$
$`+a_s^2(n_f,m_c^2)[A_{gq,Q}^\mathrm{S}(1)f_q^\mathrm{S}(n_f,m_c^2),`$
$`+A_{gg,Q}^\mathrm{S}(1)f_g^\mathrm{S}(n_f,m_c^2)].`$ (12)
The four-flavor light quark (u,d,s) densities are generated using
$`f_{k+\overline{k}}(n_f+1,m_c^2)=f_{k+\overline{k}}(n_f,m_c^2)`$
$`+a_s^2(n_f,m_c^2)A_{qq,Q}^{\mathrm{NS}}(1)f_{k+\overline{k}}(n_f,m_c^2).`$ (13)
The total four-flavor singlet quark density follows from the sum of Eqs.(11) and (13).
Next the resulting four-flavor densities are evolved from their boundary values using the four-flavor evolution kernels in the AP equations in either LO or NLO up to the scale $`\mu ^2=20.25`$ $`(\mathrm{GeV}/\mathrm{c}^2)^2`$. The bottom quark density is then generated at this point using
$`f_{b+\overline{b}}(n_f+1,m_b^2)=`$
$`a_s^2(n_f,m_b^2)[\stackrel{~}{A}_{Qq}^{\mathrm{PS}}(1)f_q^\mathrm{S}(n_f,m_b^2)`$
$`+\stackrel{~}{A}_{Qg}^{(\mathrm{S})}(1)f_g^\mathrm{S}(n_f,m_b^2)],`$ (14)
and the five-flavour gluon and light quark densities (which now include charm) are generated using Eqs. (12) and (13) with $`n_f=4`$ and replacing $`m_c^2`$ by $`m_b^2`$. Therefore only the nonlogarithmic terms in the order $`a_s^2`$ OME’s contribute to the matching conditions on the bottom quark density. Then all the densities are evolved up to higher $`\mu ^2`$ as a five-flavor set with either LO or NLO splitting functions.
The above formulae and their evolution with LO and NLO splitting functions have been implemented in a C++ computer code to yield the CS parton density set. They were used in the construction of two VFNS for the charm quark contribution to the deep inelastic structure functions in . Note that approximate expressions for the three loop splitting functions are now available in . When NNLO parton densities are available from fits to experimental data we can incorporate them into our computer program.
As an illustration we would like to compare the charm and bottom quark densities in the CS , MRST98 and CTEQ5 sets. The latter two sets work with order $`\alpha _s`$ matching conditions so the parton densities are continuous across heavy flavour thresholds. The MRST98 sets use a procedure proposed in , while the CTEQ5 sets use the different ACOT procedure in . Here we show the five-flavor densities. In the CS set they start at $`\mu ^2=m_b^2=20.25`$ $`\mathrm{GeV}^2`$. At this scale the charm densities in the CS, MRST98 (set 1) and CTEQ5HQ sets are shown in Figs.1,2,3 respectively. Since the CS charm density starts off negative for small $`x`$ at $`\mu ^2=m_c^2=1.96`$ $`\mathrm{GeV}^2`$ (see the plots in ) it is smaller than the corresponding CTEQ5HQ density. At larger $`\mu ^2`$ all the CS curves in Fig.1 are below those for CTEQ5HQ in Fig.3 although the differences are small. In general the CS c-quark densities are more equal to those in the MRST98 (set 1) in Fig.2.
At the matching point $`\mu ^2=20.25`$ $`\mathrm{GeV}^2`$ the b-quark density also starts off negative at small $`x`$ as can be seen in Fig.4, which is a consequence of the explicit form of the OME’s in . At $`O(\alpha _s^2)`$ the nonlogarithmic terms do not vanish at the matching point and yield a finite function in $`x`$, which is the boundary value for the evolution of the b-quark density. This negative start slows down the evolution of the b-quark density at small $`x`$ as the scale $`\mu ^2`$ increases. Hence the CS densities at small $`x`$ in Fig.4 are smaller than the MRST98 (set 1) densities in Fig.5 and the CTEQ5HQ densities in Fig.6 at the same values of $`\mu ^2`$. The differences between the sets are still small, of the order of five percent at small $`x`$ and large $`\mu ^2`$. This will lead to differences in cross sections for processes involving incoming b-quarks at the Tevatron.
We suspect that the differences between these results for the c and b-quark densities are primarily due to the different gluon densities in the three sets rather to than the effects of the different boundary conditions. This could be checked theoretically if both LO and NLO three-flavor sets were provided by MRST and CTEQ at small scales. We note that CS uses the GRV98 LO and NLO gluon densities, which are rather steep in $`x`$ and generally larger than the latter sets at the same values of $`\mu ^2`$. Since the discontinuous boundary conditions suppress the charm and bottom densities at small $`x`$, they enhance the gluon densities in this same region (in order that the momentum sum rules are satisfied). Hence the GRV98 three flavour gluon densities and the CS four and five flavor gluon densities are generally larger than those in MRST98 (set 1) and CTEQ5HQ. Unfortunately experimental data are not yet precise enough to decide which set is the best one.
Acknowledgements.
I am grateful to B. Harris, E. Laenen and W.L van Neerven for helpful comments on this report.
|
warning/0005/cond-mat0005355.html
|
ar5iv
|
text
|
# Ground state energy of a non-integer number of particles with 𝛿 attractive interactions
## 1 Introduction
We consider a system of $`n`$ identical quantum particles on a ring of size $`L`$ with $`\delta `$ attractive interactions. If we call $`x_\alpha `$ (for $`1\alpha n`$) the positions of the particles, the Hamiltonian of this system is
$$=\frac{1}{2}\underset{\alpha }{}\frac{^2}{x_\alpha ^2}\gamma \underset{\alpha <\beta }{}\delta (x_\alpha x_\beta ),$$
(1)
where $`\gamma `$ is the strength of the attractive ($`\gamma 0`$) interactions. The main goal of the present work is to define and to calculate the ground state energy $`E_0(n,L,\gamma )`$ of (1) when $`n`$ is not an integer (especially when $`n`$ is small).
This system of particles in one dimension with $`\delta `$ interactions has a long history in the theory of exactly soluble models. It was first introduced to describe a Bose gas by Lieb and Liniger who calculated by Bethe ansatz the ground state energy and the excitations for repulsive interactions (that is for negative $`\gamma `$) in the thermodynamic limit ($`n`$ and $`L`$ go to infinity keeping $`n/L`$ constant).
The problem arose also in the theory of disordered systems: the calculation of the free energy of a directed polymer in a random medium in $`1+1`$ dimensions by the replica method reduces to finding the ground state energy of (1): if $`Z(x,t)`$ is the partition function of a directed polymer joining the points $`(0,0)`$ and $`(x,t)`$ on a cylinder with periodic boundary conditions ($`x+Lx`$)
$$Z(x,t)=_{(0,0)}^{(x,t)}𝒟y(s)\mathrm{exp}\left(_0^t𝑑s\left[\frac{1}{2}\left(\frac{dy(s)}{ds}\right)^2+\eta (y(s),s)\right]\right),$$
(2)
where the random medium is characterised by a Gaussian white noise $`\eta (y,t)`$
$$\eta (y,t)\eta (y^{},t^{})=\gamma \delta (yy^{})\delta (tt^{}),$$
(3)
then the integer moments of $`Z(x,t)`$ are given for large $`t`$ by
$$\underset{t\mathrm{}}{lim}\frac{1}{t}\mathrm{ln}\left[\frac{Z^n(x,t)}{Z(x,t)^n}\right]=E_0(n,L,\gamma ),$$
(4)
where $``$ denotes the average over the random medium and $`E_0(n,L,\gamma )`$ is the ground state energy of (1). The knowledge of $`E_0(n,L,\gamma )`$ determines for large $`t`$ the whole distribution of $`\mathrm{ln}Z(x,t)`$. For example, the variance of $`\mathrm{ln}Z(x,t)`$ is
$$\underset{t\mathrm{}}{lim}\frac{\mathrm{ln}^2Z(x,t)\mathrm{ln}Z(x,t)^2}{t}=\frac{^2E_0(n,L,\gamma )}{n^2}|_{n=0}.$$
(5)
Of course, to obtain this variance or other characteristics of the distribution of $`\mathrm{ln}Z(x,t)`$, one should be able to define and to calculate the ground state energy of (1) not only for integer $`n`$, but for *any value of $`n`$*. Moreover, because of (3), the interactions in (1) must be attractive ($`\gamma 0`$); So in contrast to the Bose gas initially studied, the interactions are attractive and the interesting limit is no longer the thermodynamic limit $`n\mathrm{}`$ but rather the limit $`n0`$.
For integer $`n`$ and $`L=\mathrm{}`$, the $`n`$ particles form a bound state at energy
$$E_0(n,\mathrm{},\gamma )=\gamma ^2\frac{n(n^21)}{24}.$$
(6)
Using this formula for non-integer $`n`$ helped to understand several properties of the distribution of $`\mathrm{ln}Z(t)`$ when $`L`$ is infinite. There are however a number of difficulties with (6) for non-integer $`n`$, in particular a problem of convexity: $`d^2\mathrm{ln}Z^n/dn^2`$ should be positive for all $`n`$, and so (4, 6) cannot be valid at least for negative $`n`$. We believe that these difficulties are due to the exchange of limit $`t\mathrm{}`$ and $`L\mathrm{}`$ and this is why we try in the present work to determine $`E_0(n,L,\gamma )`$ for finite $`L`$.
The paper is organised as follows: in section 2, we recall the Bethe ansatz equations which give the ground state energy of (1) for an (integer) number $`n`$ of particles and we write the integral equation (13) which is a way of solving the coupled non-linear equations of the Bethe ansatz. The main advantage of this integral equation is that both the strength $`\gamma `$ of the interactions and the number $`n`$ of particles appear as continuous parameters. In section 3 we solve (13) perturbatively in $`c`$ (where $`c=\gamma L/2`$) for arbitrary $`n`$. We notice that the coefficients in the small $`c`$ expansion of $`E_0(n,L,\gamma )`$ are all polynomials in $`n`$, thus allowing to define the expansion even for non-integer $`n`$. In section 4, we show how to generate a small $`n`$ expansion of the solution of (13). We give explicit expressions up to order $`n^3`$ of $`E_0(n,L,\gamma )`$ and we notice that the coefficients of the small $`n`$ expansion have in general a zero radius of convergence in $`c`$.
## 2 The Bethe ansatz equations
The Bethe ansatz consists in looking in the region $`0x_1\mathrm{}x_nL`$ for a ground state wave function $`\mathrm{\Psi }(x_1,\mathrm{},x_n)`$ of the form
$$\mathrm{\Psi }(x_1,\mathrm{},x_n)=\underset{P}{}A_Pe^{\frac{2}{L}(q_1x_{P(1)}+\mathrm{}+q_nx_{P(n)})},$$
(7)
where the sum in (7) runs over all the permutations $`P`$ of $`\{1,\mathrm{},n\}`$. The value of $`\mathrm{\Psi }`$ in other regions can be deduced from (7) by symmetries. One can show that (7) is an eigenstate of (1) at energy
$$E(n,L,\gamma )=\frac{2}{L^2}\underset{1\alpha n}{}q_\alpha ^2,$$
(8)
if the $`\{q_\alpha \}`$ are solutions of the $`n`$ coupled equations
$$e^{2q_\alpha }=\underset{\beta \alpha }{}\frac{q_\alpha q_\beta +c}{q_\alpha q_\beta c},$$
(9)
where
$$c=\frac{\gamma L}{2}.$$
(10)
(A derivation of (9) can be found in . Note that $`ik_j`$ and $`c`$ in become here respectively $`\frac{2}{L}q_j`$ and $`\gamma `$; the $`c`$ in and our $`c`$ defined in (10) are thus different.) Moreover, for $`\gamma 0`$, all the $`q_\alpha `$ are distinct.
There are *a priori* many solutions of (9). We look for the ground state, that is the solution $`\{q_\alpha \}`$ for which (8) is minimal. When $`c=0`$, the problem reduces to $`n`$ non-interacting particles $`\{q_\alpha \}=\{0\}`$ (we have periodic boundary conditions). Because the ground state solution is not degenerate, the solution $`\{q_\alpha \}`$ of (9) must have the symmetry $`\{q_\alpha \}=\{q_\alpha \}`$, depend continuously on $`c`$, and vanish as $`c0`$.
Let us introduce the following function of $`\{q_\alpha \}`$:
$$B(u)=\frac{1}{n}e^{\frac{c}{4}(u^21)}\underset{q_\alpha }{}\rho (q_\alpha )e^{q_\alpha (u1)},$$
(11)
where the parameters $`\rho (q_\alpha )`$ are defined by
$$\rho (q_\alpha )=\underset{q_\beta q_\alpha }{}\frac{q_\alpha q_\beta +c}{q_\alpha q_\beta }.$$
(12)
The function $`B(u)`$ is a rather complicated (but easier to manipulate than the $`q_\alpha `$) symmetric function of the ground state solution $`\{q_\alpha \}`$ of (9). As shown in the appendix, it satisfies the integral equation
$`B(1+u)B(1u)=nc{\displaystyle _0^u}e^{\frac{c}{2}(v^2uv)}B(1v)B(1+uv)𝑑v,`$ (13)
and the following two conditions
$`B(1)`$ $`=1,`$ (14)
$`B(u)`$ $`=B(u).`$ (15)
Moreover, the energy (8) can be deduced from of $`B(u)`$ by
$$E_0(n,L,\gamma )=\frac{2}{L^2}\left[\frac{n^3c^2}{6}+\frac{nc^2}{12}+\frac{nc}{2}nB^{\prime \prime }(1)\right].$$
(16)
The derivation of (13, 14, 15, 16) is given in the appendix. How these relations lead to small $`c`$ or small $`n`$ expansions is explained in sections 3 and 4.
## 3 Expansion in powers of $`c`$
One could try to solve the equations (9) perturbatively in $`c`$ but the approach turns out to be quickly complicated. Instead we are going to see that the integral equation (13) is very convenient to obtain $`E_0(n,L,\gamma )`$ for small $`c`$.
It is known (and easy to check from (9)) that the $`q_\alpha `$ scale like $`\sqrt{c}`$ for small $`c`$. Therefore each coefficient $`B_i(u)`$ of the small $`c`$ expansion of $`B(u)`$ defined by (11) is a polynomial in $`u`$.
$$B(u)=B_0(u)+cB_1(u)+c^2B_2(u)+\mathrm{}$$
(17)
Moreover, conditions (14) and (15) impose that all the $`B_i(u)`$ are even, that $`B_0(1)=1`$ and $`B_i(1)=0`$ for any $`i1`$.
At zero-th order in $`c`$, we find, using (13):
$$B_0(1+u)B_0(1u)=0.$$
(18)
Thus $`B_0(u)`$ and $`B_0(1+u)`$ are both even functions of $`u`$. As $`B_0(u)`$ is a polynomial and $`B_0(1)=1`$, the only solution is
$$B_0(u)=1.$$
(19)
We put this back into (13) and get at first order in $`c`$
$$B_1(1+u)B_1(1u)=nu.$$
(20)
Again, using the fact that $`B_1(u)`$ is an even polynomial such that $`B_1(1)=0`$, the only possible solution is:
$$B_1(u)=\frac{n}{4}(u^21).$$
(21)
It is easy to see from (13) that at any order in $`c`$, we have to solve
$`B_i(1+u)B_i(1u)=\text{“some polynomial odd in }u\text{},`$ (22)
and that there is a unique even polynomial solution satisfying $`B_i(1)=0`$. One can generate as many $`B_i(u)`$ as needed to obtain $`B(u)`$ up to any desired order in $`c`$.
$`B(u)=1+{\displaystyle \frac{cn(1+u^2)}{4}}+{\displaystyle \frac{c^2n(1+2n)(1+u^2)^2}{96}}+O(c^3).`$ (23)
Relation (16) then gives the energy. Up to the fourth order in $`c`$, we find
$`{\displaystyle \frac{L^2}{2}}E_0(n,L,\gamma )=n(n1)(`$ $`{\displaystyle \frac{c}{2}}+{\displaystyle \frac{c^2}{12}}+{\displaystyle \frac{n}{180}}c^3+({\displaystyle \frac{n^2}{1512}}{\displaystyle \frac{n}{1260}})c^4+\mathrm{}).`$ (24)
## 4 Solution for small $`n`$
It is clear from section 3 that if we stop the small $`c`$ expansion of $`B(u)`$ at a given order, $`B(u)`$ and $`E_0(n,L,\gamma )`$ are polynomials in $`n`$. This allows to define the small $`c`$ expansion of $`B(u)`$ or of $`E_0(n,L,\gamma )`$ for an arbitrary value of $`n`$.
Moreover, we can write a small $`n`$ expansion of $`B(u)`$ by collecting all the terms proportional to $`n^k`$ in the small $`c`$ expansion of $`B(u)`$ and calling the series obtained $`b_k(u)`$. Then,
$$B(u)=1+nb_1(u)+n^2b_2(u)+\mathrm{}$$
(25)
Conditions (14, 15) impose that $`b_k(u)=b_k(u)`$ and $`b_k(1)=0`$ for all $`k1`$.
We are now going to describe a procedure which leads to a recursion on the $`b_k(u)`$ and allows to write them not only as power series in $`c`$ but as explicit functions of $`c`$ and $`u`$. If we insert (25) into (13), we get at first order in $`n`$
$$b_1(1+u)b_1(1u)=c_0^ue^{\frac{c}{2}(v^2uv)}𝑑v.$$
(26)
It can be checked that a solution of (26) compatible with the conditions $`b_1(u)=b_1(u)`$ and $`b_1(1)=0`$ is
$$b_1(u)=\sqrt{c}_0^+\mathrm{}𝑑\lambda \frac{\mathrm{cosh}\frac{\lambda u\sqrt{c}}{2}\mathrm{cosh}\frac{\lambda \sqrt{c}}{2}}{\mathrm{sinh}\frac{\lambda \sqrt{c}}{2}}e^{\frac{\lambda ^2}{2}}.$$
(27)
There are however other solutions to the difference equation (26): one could add to (27) an arbitrary function $`F(u,c)`$ even and periodic in $`u`$ of period 2 and vanishing at $`u=1`$. If we require that each term in the small $`c`$ expansion of $`b_1(u)`$ is polynomial in $`u`$ (as justified in section 3), we see that *all the terms of the small $`c`$ expansion of $`F(u,c)`$ must be identically zero*. For example $`F(u,c)=\mathrm{exp}(c^{1/4})(\mathrm{cos}(\pi u)+1)`$ is an acceptable function. This already shows that $`b_1(u)`$ given by (27) has indeed for small $`c`$ expansion the series obtained by collecting all the terms proportional to $`n`$ in the small $`c`$ expansion of section 3.
If the solution $`b_1(u)`$ (27) had a non-zero radius of convergence in $`c`$, it would be natural to choose the only $`b_1(u)`$ which is analytic in $`c`$ at $`c=0`$ by taking $`F(u,c)=0`$. Unfortunately, this is not the case: by making the change of variable $`\lambda ^2=\mu `$, expression (27) appears as the Borel sum of a divergent series.
We found no conclusive reasons why (27) is the solution of (26) we should select. However, one can notice that, when $`n`$ is an integer, $`B(u)`$ is analytic in $`u`$ and goes to zero when $`u\pm i\mathrm{}`$ (see (11)). Here, the $`b_1(u)`$ given by (27) grows like $`\mathrm{ln}|u|`$ when $`u\pm i\mathrm{}`$. Adding a non-zero periodic $`F(u,c)`$ would either lead to an exponential growth in the imaginary direction or introduce singularities in the complex $`u`$ plane. So (27) is the solution of (26), analytic in the whole $`u`$ plane, which has the slowest growth in the imaginary direction.
The same difficulty of selecting the right solution appears at every order in the expansion in powers of $`n`$. We are now going to explain the procedure we have used to select one solution. If we insert (25) into (13), we have to solve at any order $`k`$ in the small $`n`$ expansion
$$b_k(1+u)b_k(1u)=\varphi _k(u),$$
(28)
where $`\varphi _k(u)`$ is a function odd in $`u`$ which can be calculated from the previous orders
$$\varphi _k(u)=c\underset{i=0}{\overset{k1}{}}_0^ue^{\frac{c}{2}(v^2uv)}b_i(1v)b_{ki1}(1+uv)𝑑v.$$
(29)
(For consistency, we use $`b_0(u)=1`$.) It can be checked that a solution to (28) is
$$b_k(u)=_0^+\mathrm{}𝑑\lambda \frac{\mathrm{cosh}\frac{\lambda u\sqrt{c}}{2}\mathrm{cosh}\frac{\lambda \sqrt{c}}{2}}{\mathrm{sinh}\frac{\lambda \sqrt{c}}{2}}a_k(\lambda ),$$
(30)
where $`a_k(\lambda )`$ is given by
$$a_k(\lambda )=\frac{1}{2i\pi }_0^+\mathrm{}𝑑u\mathrm{sin}\frac{\lambda u}{2}\varphi _k\left(\frac{iu}{\sqrt{c}}\right).$$
(31)
Indeed, the verification is a simple matter of algebra. (We convinced ourselves that $`b_k(u)\mathrm{ln}^k|u|`$ and $`\varphi _k(u)\mathrm{ln}^{k1}|u|/u`$ as $`u\pm i\mathrm{}`$, and that $`a_k(\lambda )\mathrm{ln}^{k1}|\lambda |`$ for $`\lambda 0`$ and $`a_k(\lambda )\mathrm{exp}(\lambda ^2(k+1)/4k)`$ for $`\lambda \mathrm{}`$, so that all the integrals in (30, 31) converge.)
As for $`b_1(u)`$, one could add to the $`b_k(u)`$ given by (30, 31) an arbitrary function $`F_k(u,c)`$, even and periodic in $`u`$ of period 2 to obtain the general solution of (28). However, to be consistent with the small $`c`$ expansion of section 3, the small $`c`$ expansion of $`F_k(u,c)`$ should be identically zero. Moreover if we want the solution of (28) to be analytic in the whole $`u`$ plane and not to grow too fast when $`u\pm i\mathrm{}`$, we must take $`F_k(u,c)=0`$.
Expression (27) for $`b_1(u)`$ is in fact a particular case of the procedure (30, 31); when applied to (26), it gives indeed $`a_1(\lambda )=\sqrt{c}\mathrm{exp}(\lambda ^2/2)`$.
At second order in the small $`n`$ expansion, we find for $`\lambda >0`$
$$a_2(\lambda )=ce^{\frac{\lambda ^2}{2}}\left[_0^\lambda 𝑑\mu e^{\frac{\mu ^2}{2}}\frac{2\mathrm{cosh}\frac{\lambda \mu }{2}2}{\mathrm{tanh}\frac{\mu \sqrt{c}}{2}}+_\lambda ^+\mathrm{}𝑑\mu e^{\frac{\mu ^2}{2}}\frac{e^{\frac{\lambda \mu }{2}}2}{\mathrm{tanh}\frac{\mu \sqrt{c}}{2}}\right].$$
(32)
The expressions of $`b_3(u)`$ or $`a_3(\lambda )`$ would be much longer to write and higher orders even more complicated. Recursion (29, 30, 31) allows nevertheless to calculate in principle the whole expansion in powers of $`n`$.
Using relation (16) and the expressions (27) and (32) of $`b_1(u)`$ and $`a_2(\lambda )`$, we find that the energy $`E_0(n,L,\gamma )`$ is given up to order $`n^3`$:
$`{\displaystyle \frac{L^2}{2}}E_0(n,L,\gamma )=n\left({\displaystyle \frac{c}{2}}+{\displaystyle \frac{c^2}{12}}\right)n^2{\displaystyle \frac{c^{3/2}}{4}}{\displaystyle _0^+\mathrm{}}𝑑\lambda {\displaystyle \frac{\lambda ^2}{\mathrm{tanh}\frac{\lambda \sqrt{c}}{2}}}e^{\frac{\lambda ^2}{2}}+n^3{\displaystyle \frac{c^2}{6}}`$ (33)
$`n^3{\displaystyle \frac{c^2}{4}}{\displaystyle _0^+\mathrm{}}𝑑\lambda {\displaystyle \frac{\lambda ^2e^{\frac{\lambda ^2}{2}}}{\mathrm{tanh}\frac{\lambda \sqrt{c}}{2}}}\left[{\displaystyle _0^\lambda }𝑑\mu e^{\frac{\mu ^2}{2}}{\displaystyle \frac{2\mathrm{cosh}\frac{\lambda \mu }{2}2}{\mathrm{tanh}\frac{\mu \sqrt{c}}{2}}}+{\displaystyle _\lambda ^+\mathrm{}}𝑑\mu e^{\frac{\mu ^2}{2}}{\displaystyle \frac{e^{\frac{\lambda \mu }{2}}2}{\mathrm{tanh}\frac{\mu \sqrt{c}}{2}}}\right].`$
As explained in the introduction, this small $`n`$ expansion of $`E_0(n,L,\gamma )`$ gives the cumulants of the free energy in the directed polymer problem. Of course, if we expand (33) in powers of $`c`$, we recover (24).
## 5 Conclusion
In this work, we have developed a method allowing to calculate perturbatively the ground state energy of (1) for a non-integer number $`n`$ of particles. We first generated for integer $`n`$ a perturbation series in powers of the interaction $`c`$. Each term of this series is polynomial in $`n`$, allowing to define a small $`c`$ expansion of the energy for non-integer $`n`$. This series, at least for small $`n`$, has in general a zero radius of convergence, in contrast to integer $`n`$ for which the radius of convergence of the perturbation theory is non-zero. (For $`n=2`$, the closest singularities of $`E_0(2,L,\gamma )`$ in the $`c`$ plane lie at $`c3.30\pm i\mathrm{\hspace{0.17em}4.12}`$.)
We believe that the fact that each term in the perturbation theory is polynomial in $`n`$ is generic and would be true for an arbitrary pair interaction and in any dimension. As the link (4) to directed polymers is valid in any dimension, it would be useful and interesting to try to recover our results by doing a direct perturbation theory of the Hamiltonian instead of our Bethe ansatz approach (which is limited to $`1+1`$ dimensions and to a $`\delta `$ potential) in order to see whether the calculations could be extended to higher dimensions.
Our calculation of the ground state energy for non-integer $`n`$ is based on the integral equation (13) and the conditions (14, 15). When we tried to solve the problem for small $`n`$, at each order we had to select a particular solution of a difference equation. We did not find a conclusive reason to justify the solution we selected, apart from some analyticity properties and growth criterion in the complex plane of the variable $`u`$. It would certainly be interesting to justify our choice (33) by calculating the second and the third cumulants of $`\mathrm{ln}Z`$ directly (and not only perturbatively to all orders in $`c`$).
In our small $`n`$ expansion of section 4, the terms become quickly very complicated. There is however a regime, which corresponds to the large $`c`$ limit of recursion (29, 30, 31) where one can handle all orders in the small $`n`$ expansion. This allows one to calculate the whole distribution of $`\mathrm{ln}(Z(x,t))/t`$ when $`t`$ is very large and $`(1/t)\mathrm{ln}(Z(x,t)/Z(x,t))`$ of order $`1/L`$. One can then recover the same large deviation function as found for the asymmetric exclusion process, as expected since the directed polymer problem in $`1+1`$ dimensions and the asymmetric exclusion process are both representatives of the KPZ equation in dimension $`1`$. This strengthens the conjecture that the solutions to the difference equations we selected in section 4 give indeed the right non-integer moments of the partition function.
From the point of view of the theory of disordered systems, our results give one of the very few examples for which the distribution of $`Z`$ can be calculated exactly. In particular they could provide a good test of the replica approach and of other variational methods.
A simple and interesting phenomenon visible in the present work (which is generic of all kinds of disordered systems with Gaussian disorder) is that the weak disorder expansion (here small $`c`$ expansion) of non-integer moments of the partition function has a zero radius of convergence whereas integer moments have a non-zero radius of convergence. This is already visible in the trivial example of a single Ising spin $`\sigma =\pm 1`$ in a random Gaussian field $`h`$; the partition function at temperature $`T`$ is $`Z=2\mathrm{cosh}(h/T)`$, and it is easy to check that all non-integer moments of the partition function have a zero radius of convergence in $`1/T`$.
Acknowledgements We thank François David, Michel Gaudin, Vincent Pasquier, Leonid Pastur, Herbert Spohn and André Voros for useful discussions.
## A Derivation of (13, 14, 15, 16)
Let us first establish some useful properties of the numbers $`\rho (q_\alpha )`$ defined by (12). If the $`q_\alpha `$ are the $`n`$ roots of the polynomial $`P(X)`$ defined as
$$P(X)=\underset{q_\alpha }{}(Xq_\alpha ),$$
(A.1)
it is easy to see that the $`\rho (q_\alpha )`$ defined in (12) satisfy
$$\frac{P(X+c)}{P(X)}=1+c\underset{q_\alpha }{}\frac{\rho (q_\alpha )}{Xq_\alpha }.$$
(A.2)
(The two sides have the same poles with the same residues and coincide at $`X\mathrm{}`$.) Expanding the right hand side of (A.2) for large $`X`$, we get
$`{\displaystyle \frac{P(X+c)}{P(X)}}=1+c{\displaystyle \underset{q_\alpha }{}}{\displaystyle \frac{\rho (q_\alpha )}{X}}\left(1+{\displaystyle \frac{q_\alpha }{X}}+{\displaystyle \frac{q_\alpha ^2}{X^2}}\right)+O\left({\displaystyle \frac{1}{X^4}}\right).`$ (A.3)
On the other hand, using (8, A.1) and the symmetry $`\{q_\alpha \}=\{q_\alpha \}`$ we have
$$P(X)=X^n+\frac{L^2}{4}E_0(n,L,\gamma )X^{n2}+O(X^{n4}),$$
(A.4)
so that
$$\frac{P(X+c)}{P(X)}=1+\frac{nc}{X}+\frac{c^2\left(\genfrac{}{}{0pt}{}{n}{2}\right)}{X^2}+\frac{c^3\left(\genfrac{}{}{0pt}{}{n}{3}\right)cE_0(n,L,\gamma )L^2/2}{X^3}+O\left(\frac{1}{X^4}\right).$$
(A.5)
Comparing (A.3) and (A.5), we get the relations
$`{\displaystyle \underset{q_\alpha }{}}\rho (q_\alpha )`$ $`=n,`$ (A.6)
$`{\displaystyle \underset{q_\alpha }{}}q_\alpha \rho (q_\alpha )`$ $`=c\left({\displaystyle \genfrac{}{}{0pt}{}{n}{2}}\right),`$ (A.7)
$`{\displaystyle \underset{q_\alpha }{}}q_\alpha ^2\rho (q_\alpha )`$ $`=c^2\left({\displaystyle \genfrac{}{}{0pt}{}{n}{3}}\right){\displaystyle \frac{E_0(n,L,\gamma )L^2}{2}}`$ (A.8)
Moreover, by letting $`X=\pm q_\beta c`$ in (A.2) we get for any $`q_\beta `$ root of $`P(X)`$
$$\frac{1}{c}=\underset{q_\alpha }{}\frac{\rho (q_\alpha )}{q_\alpha q_\beta +c}=\underset{q_\alpha }{}\frac{\rho (q_\alpha )}{q_\alpha +q_\beta +c}.$$
(A.9)
Lastly using the symmetry $`\{q_\alpha \}=\{q_\alpha \}`$ and the definition (12), the Bethe ansatz equations (9) reduce to
$$e^{q_\alpha }\rho (q_\alpha )e^{q_\alpha }\rho (q_\alpha )=0.$$
(A.10)
From the definition (11) of $`B(u)`$ and the properties (A.6A.10), it is straightforward to establish (1316): the integral equation (13) is a direct consequence of (11) and (A.9). Properties (14, 15) follow from (11, A.6) and (11, A.10) respectively. Lastly (16) is a consequence of (11, A.6A.8).
|
warning/0005/astro-ph0005287.html
|
ar5iv
|
text
|
# Large Scale Pressure Fluctuations and Sunyaev-Zel’dovich Effect
## I Introduction
In recent years, increasing attention has been given to the physical properties of the intergalactic warm and hot plasma gas distribution associated with large scale structure and the possibility of its detection (e.g., ). It is now widely believed that at least $``$ 50% of the present day baryons, when compared to the total baryon density through big bang nucleosynthesis, are present in this warm gas distribution and have remained undetected given its nature (e.g., ). Currently proposed methods for the detection of this gas with include observations of the thermal diffuse X-ray emission (e.g., ), associated X-ray and UV absorption and emission lines (e.g., ) and resulting Sunyaev-Zel’dovich (SZ; ) effect (e.g., ).
The SZ effect arises from the inverse-Compton scattering of CMB photons by hot electrons along the line of sight. This effect has now been directly imaged towards massive galaxy clusters (e.g., ), where temperature of the scattering medium can reach as high as 10 keV producing temperature changes in the CMB of order 1 mK at Rayleigh-Jeans wavelengths. Previous analytical predictions of the resulting SZ effect due to large scale structure have been based on either through a Press-Schechter (PS; ) description of the contributing galaxy clusters (e.g., ) or using a biased description of the pressure power spectrum with respect to the dark matter density field (e.g., ). Numerical simulations (e.g., ) are beginning to improve some of these analytical predictions, but are still limited to handful of simulations with limited dynamical range and resolution. Therefore, it is important that one consider improving analytical models of the large scale structure SZ effect, and provide predictions which can be easily tested through simulations.
Our present study on the large scale baryon pressure and the resulting SZ effect is timely for several reasons, including the fact that improving numerical simulations have recently begun to make detailed predictions for the pressure power spectrum and SZ effect such that those predictions can be extended and improved with analytical models . Also, several studied have considered the possibility that large scale baryon distribution can be probed with upcoming CMB missions using SZ effect (e.g., ). Our calculations can be used to further refine these predictions and to investigate the possibility how such analytical model as the one presented here can be tested with observations.
As part of this study, we extend previous studies by considering the full power spectrum and bispectrum, the Fourier space analog of the three-point function, of pressure fluctuations. The pressure power spectrum and bispectrum contains all necessary information on the large scale distribution of temperature weighted baryons, whereas, the SZ power spectrum is only a projected measurement of the pressure power spectrum. This can be compared to weak gravitational lensing, where lensing is a direct probe of the projected dark matter density distribution. The bispectrum of pressure fluctuations, and SZ bispectrum, contains all the information present at the three-point level, whereas conventional statistics, such as skewness, do not. An useful advantage of using the 3d statistics, such as the pressure power spectrum, is that they can directly compared to numerical simulations, while only 2d statistics, such as the projected pressure power spectrum along the line of sight, basically the SZ power spectrum, can be observed. Our approach here is to consider both such that our calculations can eventually be compared to both simulations and observations.
The calculation of pressure power spectrum and bispectrum requires detailed knowledge on the baryon distribution, which can eventually be obtained numerically through hydrodynamical simulations. Here, we provide an analytical technique to obtain the pressure power spectrum and bispectrum by describing the baryon distribution in the universe as (1) present in virialized halos with overdensities $`200`$ with respect to background densities (2) unshocked diffuse baryons in overdensities $`10`$ that trace a Jeans-smoothed dark matter density field (3) the intermediate overdensity region, which is likely to be currently undergoing in shock heating and falling on to structures such as filaments. In the present paper we discuss the first two regimes, while a useful approach to include the latter, through simulations, is discussed.
Our description of baryons present in virialized halos follow recent studies on the dark matter density field through halo contributions following and applied to lensing statistics in and . For the description of baryons, the critical ingredients are: the PS formalism for the mass function; the NFW profile of , and the halo bias model of . The baryons are assumed to be in hydrostatic equilibrium with respect to dark matter distribution, which is a valid assumption, at least for the high mass halos that have been observed with X-ray instruments, given the existence of regularity relations between cluster baryon and dark matter physical properties (e.g., ). We take two descriptions of the temperature structure: (1) virial temperature and (2) virial temperature plus an additional source of minimum energy. The latter consideration allows the possibility for a secondary source of energy for baryons, such as due to preheating through stellar formation and feedback processes. Numerical simulations (e.g., ), as well observations (e.g., ), suggest the existence of such an energy source. The low photoionized overdensity baryons are described following the analytical description of . The fraction of baryons present in such low overdensities are assumed to follow what has been measured in numerical simulations of . We suggest that such baryons provide a lower limit to the SZ effect in the absence of any contribution from baryons present in virialized halos.
Throughout this paper, we will take $`\mathrm{\Lambda }`$CDM as our fiducial cosmology with parameters $`\mathrm{\Omega }_c=0.30`$ for the CDM density, $`\mathrm{\Omega }_b=0.05`$ for the baryon density, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.65`$ for the cosmological constant, $`h=0.65`$ for the dimensionless Hubble constant and a scale invariant spectrum of primordial fluctuations, normalized to galaxy cluster abundances ($`\sigma _8=0.9`$ see ) and consistent with COBE . For the linear power spectrum, we take the fitting formula for the transfer function given in .
The paper is organized as following: In §II, we review the dark matter halo approach to modeling the density field and extend it to model properties associated with large scale baryon distribution, mainly the pressure fluctuations that contributes to the observable SZ effect. We suggest recent papers by Seljak , Ma & Fry , Cooray & Hu , and Scoccimarro et al for details on the dark matter halo approach and applications to other observable statistics such as galaxy properties and weak gravitational lensing. As necessary, we use techniques developped in these papers for our current calculation. In §III we apply the formalism to the convergence power spectrum, skewness, and bispectrum. We conclude in §IV with a summary of our main results.
## II Density and Pressure Power Spectra
### A General Definitions
In order to calculate the contribution to temperature anisotropies through SZ effect associated with large scale structure, we divide the LSS with overdensities $`200`$ as collapsed and virialized halos with a gas distribution following hydrostatic equilibrium and with virial temperatures. The pressure power spectrum can be calculated using an extension to the dark matter halo approach by assuming a physical relation between baryons and dark matter.
The baryons with overdensities $`10`$ track the dark matter distribution and their power spectrum has been studied by . These baryons have temperatures similar to photoionization energies of Hydrogen and Helium. Current numerical simulations, e.g., , suggest that most of the baryons are in such low overdensities at $`z>1`$, while at present day, are in virialized halos. The calculation of SZ effect due to such baryons follow , except that we include the redshift dependence of mass fraction within such low density halos following the numerical results of , and modify the mean temperature of such baryons to be consistent with photoionization energies.
First we discuss the pressure and related power spectra due to collapsed halos.
The dark matter profile of collapsed halos are taken to be the NFW with a density distribution
$$\rho _\delta (r)=\frac{\rho _s}{(r/r_s)(1+r/r_s)^2}.$$
(1)
The density profile can be integrated and related to the total dark matter mass of the halo within $`r_v`$
$$M_\delta =4\pi \rho _sr_s^3\left[\mathrm{log}(1+c)\frac{c}{1+c}\right]$$
(2)
where the concentration, $`c`$, is $`r_v/r_s`$. Choosing $`r_v`$ as the virial radius of the halo, spherical collapse tells us that $`M=4\pi r_v^3\mathrm{\Delta }(z)\rho _b/3`$, where $`\mathrm{\Delta }(z)`$ is the overdensity of collapse (see e.g. ) and $`\rho _b`$ is the background matter density today. We use comoving coordinates throughout. By equating these two expressions, one can eliminate $`\rho _s`$ and describe the halo by its mass $`M`$ and concentration $`c`$.
Following , we take the concentration of dark matter halos to be
$$c(M,z)=a(z)\left[\frac{M}{M_{}(z)}\right]^{b(z)},$$
(3)
where $`a(z)=10.3(1+z)^{0.3}`$ and $`b(z)=0.24(1+z)^{0.3}`$. Here $`M_{}(z)`$ is the non-linear mass scale at which the peak-height threshold, $`\nu (M,z)=1`$. The above concentration is chosen so that dark matter halos provide a reasonable match to the the non-linear density power spectrum as predicted by the ; it extends the treatment of to the redshifts of interest for SZ effect. We caution the reader that eqn. (3) is only a good fit for the $`\mathrm{\Lambda }`$CDM model assumed.
The gas density profile, $`\rho _g(r)`$, is calculated assuming the hydrostatic equilibrium between the gas distribution and the dark matter density field with in a halo. This is a valid assumption given that current observations of halos, mainly galaxy clusters, suggest the existence of regularity relations, such as size-temperature (e.g., ), between physical properties of dark matter and baryon distributions.
The hydrostatic equilibrium implies,
$$\frac{kT_e}{\mu m_p}\frac{d\mathrm{log}\rho _g}{dr}=\frac{GM_\delta (r)}{r^2},$$
(4)
where now the $`M_\delta (r)`$ is the mass only out to a radius of $`r`$. Note that we have assumed here an isothermal temperature for the gas distribution. Solving for the the equations above, we can analytically calculate the baryon density profile $`\rho _g(r)`$
$$\rho _g(r)=\rho _{g0}e^b\left(1+\frac{r}{r_s}\right)^{br_s/r},$$
(5)
where $`b`$ is a constant, for a given mass,
$$b=\frac{4\pi G\mu m_p\rho _sr_s^3}{k_BT_e},$$
(6)
with the Boltzmann constant, $`k_B`$.
In general, the halos are described with virial temperatures
$$k_BT_e=\frac{\gamma G\mu m_pM_\delta (r_v)}{3r_v},$$
(7)
with $`\gamma =3/2`$ and $`\mu =0.59`$. In addition, we also consider the possibility for the existence of a constant non-gravitational energy in small mass halos, consistent with observations of galaxy groups, due to what is commonly known as “preheating”. The possibility for such a minimum energy for baryons today comes from heating before virialization due to energy injection and feedback processes such as processes associated with stellar formation. The total gas mass present in a dark matter halo within $`r_v`$ is
$$M_g(r_v)=4\pi \rho _{g0}e^br_s^3_0^c𝑑xx^2(1+x)^{b/x}.$$
(8)
The physical properties of the profile defined in Eq. (5) for baryons within dark matter halos, and a comparison to commonly used profiles such as isothermal and so-called beta-profiles, can be found in and . Compared to conventional profiles, this profile has the advantage that it is directly related to the dark matter profile parameters, such as central density $`\rho _s`$ and concentration via scale radius $`r_s`$, thus, any changes to the dark matter distribution produces resulting changes in the baryon distribution. Also, one can study the effect of temperature variations on the gas distribution as the parameter $`b`$ defined in Eq. (6) depends on it. A proper normalization for the dark matter halo distribution containing baryons comes through the $`c(M,z)`$ relation in Eq. (3) such that the non-linear dark matter power spectrum is produced in numerical simulations by same halos. Note that our hydrostatic equilibrium ignores the self-gravity contribution from baryons to the total potential as we only include the dark matter contribution to total mass. Since the baryon mass is expected to be $`10`$% of the total mass, we can safely ignore, as a first approximation, the contribution to total mass from baryons themselves.
Roughly speaking, the perturbative aspect of the clustering of the dark matter and baryons is described by the correlations between halos, whereas the nonlinear aspect is described by the correlations within halos, i.e. the halo profiles. We will consider the Fourier analogies of the two and three point correlations of the dark matter density, $`\delta `$, baryon pressure, $`\mathrm{\Pi }`$, and galaxy distribution, $`g`$, defined in the usual way
$$\delta _i^{}(𝐤)\delta _i(𝐤^{})=(2\pi )^3\delta (𝐤𝐤^{})P_i^\mathrm{t}(k),$$
(9)
$$\delta _i(𝐤_1)\delta _i(𝐤_2)\delta _i(𝐤_3)=(2\pi )^3\delta (𝐤_1+𝐤_2+𝐤_3)B_i^\mathrm{t}(k_1,k_2,k_3),$$
(10)
with $`i`$ representing $`\delta `$, $`\mathrm{\Pi }`$ or $`g`$. We will also consider cross-correlations between the two, such as the dark matter density-pressure power spectrum $`P_{\delta \mathrm{\Pi }}^\mathrm{t}(k)`$, which is what one probes by correlating, say, the SZ effect and weak gravitational lensing observations. Here and throughout, we occasionally suppress the redshift dependence where no confusion will arise.
As presented in , these spectra are related to the linear density power spectrum $`P(k)`$ through the bias parameters and the normalized 3d Fourier transform of the density profile $`\rho _i(r,M)`$
$$y_i(k,M)=\frac{1}{M_i}_0^{r_v}𝑑r\mathrm{\hspace{0.17em}4}\pi r^2\rho _i(r,M)\frac{\mathrm{sin}(kr)}{kr},$$
(11)
where $`i`$ represents either the density, $`\delta `$, or the gas, $`g`$, profile and associated masses respectively given in equations 2 and 8. With an increase in temperature relative to the virial temperature of the halo, especially for halos with masses $`10^{13}`$ M$`_{}`$, the gas profile is such that it does not fall rapidly at the virial radius, leading to an arbitrary cut off when doing the Fourier transformation. We included an additional filter to the gas density profile such that the gas density profile decreases smoothly but promptly to zero at the virial radius: $`\rho ^{}(r)=\rho (r)[\mathrm{erfc}(rr_v/\sqrt{2}\mathrm{\Delta }r)1]`$ with $`\mathrm{\Delta }rr_s`$. Detailed aspects of the pressure and SZ statistics due to medium to small mass halos ($`10^{13}`$ M$`_{}`$) are sensitive to sharpness of this transition, but these issues do not change our primary results. Here, we concentrate mostly on the statistics due to massive and rare halos. Another possibility not considered here is to include the role of baryons at the outskirts of halos. The physical properties of such baryons can be semi-analytically calculated following the assumption that virial radius provides a shock boundary for the equilibrium of baryons within and outside virialized regimes. The baryons outside halos are likely to preheated and trace the Jeans-smoothed version of the dark matter density field. The proper inclusion of such baryons require the aid of numerical simulations or semi-analytical models. These baryons are likely to include the ones present in overdensities between 200 and 10, which we have neglected in the present calculation.
Following , it is convenient to define a general integral over the halo mass function $`dn/dM`$,
$`I_{\mu ,i_1\mathrm{}i_\mu }^{\beta ,\eta ,\gamma }(k_1,\mathrm{},k_\mu ;z){\displaystyle 𝑑M\left(\frac{M}{\rho _b}\right)^\mu \frac{dn}{dM}(M,z)b_\beta (M)}`$ (12)
$`\left({\displaystyle \frac{\rho _b}{M}}{\displaystyle \frac{N_g}{\overline{n}_g}}\right)^\gamma \times T_e(M,z)^\eta y_{i_1}(k_1,M)\mathrm{}y_{i_\mu }(k_\mu ,M),`$ (13)
where $`b_01`$. gives the following analytic predictions for the bias parameters which agree well with simulations:
$$b_1(M;z)=1+\frac{\nu ^2(M;z)1}{\delta _c},$$
(14)
and
$`b_2(M;z)`$ $`=`$ $`{\displaystyle \frac{8}{21}}[b_1(M;z)1]+{\displaystyle \frac{\nu ^2(M;z)3}{\sigma ^2(M;z)}}.`$ (15)
Here, $`T_e(M,z)`$ is the electron temperature of the baryon distribution of given halo when pressure power spectrum is considered, $`\nu (M,z)=\delta _c/\sigma (M,z)`$, where $`\sigma (M,z)`$ is the rms fluctuation within a top-hat filter at the virial radius corresponding to mass $`M`$, and $`\delta _c`$ is the threshold overdensity of spherical collapse (see for useful fitting functions).
The terms related to the galaxy power spectrum includes the average number of galaxies per dark matter halo, $`N_g`$, the mean number density of galaxies in the universe $`\overline{n}_g`$. These parameters are discussed in § II E involving the galaxy-pressure power spectrum.
We use the Press-Schechter (PS; ) mass function to describe $`dn/dM`$. We take the minimum mass to be $`10^3`$ M$`_{}`$ while the maximum mass is varied to study the effect of massive halos on related statistics. In general, masses above $`10^{16}`$ M$`_{}`$ do not contribute to low order statistics due to the exponential decrease in the number density of such massive halos.
To summarize, in comparison to previous work on the SZ effect from virialized halos, our model has following advantages:
(1) A physically motivated profile for the distribution of baryons in virialized dark matter halos, instead of an assumed profile such as the isothermal model or so-called $`\beta `$-profile.
(2) A mass function for the virialized structures with the dark matter distribution of such halos reproducing the numerically simulated dark matter power spectrum and higher order correlations.
(3) Easily extendable variations to the baryon physics so as to account for issues such as pre-heating.
(4) Direct calculation of 3d properties of the large scale baryon distribution, such as the pressure power spectrum and bispectrum, which can be compared easily in numerical simulations.
We now discuss the calculation of properties related to the dark matter and baryons. In the case of baryons, we discuss pressure as this is the property that leads to the SZ effect allowing a useful probe of them.
### B Density Power Spectrum
Following , we can decompose the density power spectrum, as a function of redshift, into contributions from single halos (shot noise or “Poisson” contributions),
$$P_{\delta \delta }^{\mathrm{PP}}(k)=I_{2,\delta \delta }^{0,0,0}(k,k),$$
(16)
and correlations between two halos,
$$P_{\delta \delta }^{\mathrm{hh}}(k)=\left[I_{1\delta }^{1,0,0}(k)\right]^2P(k),$$
(17)
such that
$$P_{\delta \delta }^\mathrm{t}=P_{\delta \delta }^{\mathrm{PP}}+P_{\delta \delta }^{\mathrm{hh}}.$$
(18)
As $`k0`$, $`P_{\delta \delta }^{\mathrm{hh}}P(k)`$.
### C Pressure Power Spectrum
As above, we can decompose the pressure power spectrum, as relevant for the SZ effect, into contributions from single halos
$$P_{\mathrm{\Pi }\mathrm{\Pi }}^{PP}(k)=I_{2,gg}^{0,2,0}(k,k),$$
(19)
and correlations between halos
$$P_{\mathrm{\Pi }\mathrm{\Pi }}^{hh}(k)=\left[I_{1,g}^{1,1,0}(k)\right]^2P(k),$$
(20)
such that
$$P_{\mathrm{\Pi }\mathrm{\Pi }}(k)^\mathrm{t}=P_{\mathrm{\Pi }\mathrm{\Pi }}^{PP}(k)+P_{\mathrm{\Pi }\mathrm{\Pi }}^{hh}(k)$$
(21)
For the pressure power spectrum, since $`\eta `$ in Eq. (13) is non-zero, there is additional mass weighing arising from the fact that $`T_eM^{2/3}`$ resulting in an additional mass dependence. The dependence is such that most of the contributions to the pressure, and thus, to the SZ, power spectrum comes from most massive and rarest halos. This dependence has already been observed in numerical simulations by .
As we discuss later, this dependance on high mass halos to produce most of the pressure fluctuations also leads to several interesting results with regards to the detection and observability of SZ effect, among which are
(1) Most of the contribution to large scale SZ effect results from massive clusters of galaxies, while smaller mass halos and structures at low electron temperatures such as filaments do not contribute significantly
(2) Since massive halos dominate the SZ effect, and the distribution of these halos are Poisson and highly non-Gaussian, most of the contribution to two point and higher-order statistics of SZ effect will be dominated by Poisson term and there will be a significant non-Gaussianity associated with large scale SZ effect
In fact, as we find later, the large scale correlations only contribute at a level of 10% to the SZ power spectrum suggesting that such correlations can be mostly disregarded. The non-Gaussianity associated with SZ effect may become a useful tool to separate out its contribution from other sources of foregrounds in CMB anisotropy data, though this task can be efficiently carried out using frequency information (see, ). We will return to all these issues in later sections.
### D Density-Pressure Power Spectrum
The cross-correlation between the density and gas field, as appropriate for lensing-SZ cross-correlation can be decomposed as a single halo
$$P_{\mathrm{\Pi }\delta }^{PP}(k)=I_{2,g\delta }^{0,1,0}(k,k),$$
(22)
and
$`P_{\mathrm{\Pi }\delta }^{hh}(k)=\left[I_{1,g}^{1,1,0}(k)\right]\left[I_{1,\delta }^{1,0,0}(k)\right]P(k),`$ (23)
such that
$$P_{\mathrm{\Pi }\delta }(k)^\mathrm{t}=P_{\mathrm{\Pi }\delta }^{PP}(k)+P_{\mathrm{\Pi }\delta }^{hh}(k)$$
(24)
With the pressure and density field power spectra, one can define a bias associated with the large scale pressure, relative to density field,
$$b_\mathrm{\Pi }(k)\overline{T_e}=\sqrt{\frac{P_\mathrm{\Pi }(k)}{P_\delta (k)}},$$
(25)
and the dimensionless correlation coefficient between the dark matter and baryon distributions
$$r_\mathrm{\Pi }(k)=\frac{P_{\mathrm{\Pi }\delta }(k)}{\sqrt{P_\delta (k)P_\mathrm{\Pi }(k)}}.$$
(26)
In Eq. (25), the average density weighted temperature is
$$\overline{T_e}=𝑑M\frac{M}{\rho _p}\frac{dn}{dM}(M,z)T_e(M,z).$$
(27)
Following , one can define a covariance matrix in Fourier space containing the full information on scale dependence of bias and correlations:
$$\widehat{𝐂}(k)\left(\begin{array}{cc}P_{\delta \delta }(k)& P_{\mathrm{\Pi }\delta }(k)\\ P_{\mathrm{\Pi }\delta }(k)& P_{\mathrm{\Pi }\mathrm{\Pi }}(k)\end{array}\right)=P_{\delta \delta }(k)\left(\begin{array}{cc}1& b_\mathrm{\Pi }r_\mathrm{\Pi }\\ b_\mathrm{\Pi }r_\mathrm{\Pi }& b_\mathrm{\Pi }^2\end{array}\right).$$
(28)
The observation measurement of $`b_\mathrm{\Pi }`$ and $`r_\mathrm{\Pi }`$ can be considered by an inversion of the SZ-SZ, lensing-lensing and SZ-lensing power spectra as a function of redshift bins in which lensing-lensing or SZ-lensing power spectra are constructed. We discuss these possibilities later.
### E Galaxy-Pressure Power Spectrum
The cross-correlation between the galaxy distribution and gas field, as appropriate for galaxy-SZ cross-correlation can be decomposed as a single halo
$$P_{\mathrm{\Pi }\delta }^{PP}(k)=I_{2,g\delta }^{0,1,1}(k,k),$$
(29)
and
$`P_{\mathrm{\Pi }\delta }^{hh}(k)=\left[I_{1,g}^{1,1,0}(k)\right]\left[I_{1,\delta }^{1,0,1}(k)\right]P(k),`$ (30)
such that
$$P_{\mathrm{\Pi }\delta }(k)^\mathrm{t}=P_{\mathrm{\Pi }\delta }^{PP}(k)+P_{\mathrm{\Pi }\delta }^{hh}(k)$$
(31)
The calculation of galaxy-pressure power spectrum requires knowledge on the galaxy distribution within dark matter halos. Following, , we describe the average number of galaxies per halo, $`N_g`$ in Eq. (13), such that
$$N_g=\{\begin{array}{cc}\left(\frac{M}{M_{\mathrm{min}}}\right)^{0.6}& MM_{\mathrm{min}}\\ 0& M<M_{\mathrm{min}}\end{array}$$
(32)
where $`M_{\mathrm{min}}`$, the minimum dark matter halo mass in which a galaxy is found, is taken to be $`5.3\times 10^{11}h^1`$ M$`_{}`$ for our fiducial $`\mathrm{\Lambda }`$CDM cosmological model following . The above relation is consistent with semi-analytical models. however, we ignore scatter in the observed distribution on the mean number of galaxies per halo.
With the the average number of galaxies per halo, as a function of mass, the mean number density of galaxies can be written as an integral over the PS mass function
$$\overline{n}_g=𝑑MN_g\frac{dn}{dM}(M,z).$$
(33)
In practice, the cross-correlation between galaxy distribution and any other field requires the knowledge on the observable galaxy properties such as the magnitude limit, relation between luminosity and mass etc. The same restriction arising from observing conditions can be introduced as part of the weight function that takes into account the redshift projection of galaxies.
### F Pressure Bispectrum
Following , we can write the pressure bispectrum as
$`B_\mathrm{\Pi }^\mathrm{t}`$ $`=`$ $`B_\mathrm{\Pi }^{\mathrm{PPP}}+B_\mathrm{\Pi }^{\mathrm{Phh}}+B_\mathrm{\Pi }^{\mathrm{hhh}},`$ (34)
where
$`B_\mathrm{\Pi }^{\mathrm{PPP}}(k_1,k_2,k_3)=I_{3,ggg}^{0,3,0}(k_1,k_2,k_3),`$ (35)
for single halo contributions,
$`B_\mathrm{\Pi }^{\mathrm{Phh}}(k_1,k_2,k_3)=I_{2,gg}^{1,2,0}(k_1,k_2)I_{1,g}^{0,1,0}(k_3)P(k_3)+\mathrm{Perm}.`$ (36)
(37)
for double halo contributions, and
$`B_\mathrm{\Pi }^{\mathrm{hhh}}(k_1,k_2,k_3)`$ $`=`$ $`\left[2J(k_1,k_2,k_3)I_{1,g}^{1,1,0}(k_3)+I_{1,g}^{2,1,0}(k_3)\right]`$ (39)
$`\times I_{1,g}^{1,1,0}(k_1)I_{1,g}^{1,1,0}1(k_2)P(k_1)P(k_2)+\mathrm{Perm}.`$
for triple halo contributions. Here the 2 permutations are $`k_3k_1`$, $`k_2`$. Second order perturbation theory tells us that
$`J(k_1,k_2,k_3)`$ $`=`$ $`1{\displaystyle \frac{2}{7}}\mathrm{\Omega }_m^{2/63}+\left({\displaystyle \frac{k_3^2k_1^2k_2^2}{2k_1k_2}}\right)^2`$ (41)
$`\times \left[{\displaystyle \frac{k_1^2+k_2^2}{k_3^2k_1^2k_2^2}}+{\displaystyle \frac{2}{7}}\mathrm{\Omega }_m^{2/63}\right].`$
In addition to the pressure bispectrum, one can also define cross-correlation bispectra such as the pressure-pressure-density or pressure-density-density. These bispectra are relevant to the calculation of non-Gaussianities present in CMB through secondary anisotropies (e.g., ), and is necessary to determine the higher order moments associated with cross-correlations between individual effects such as SZ and weak lensing.
### G Baryons in small overdensities
The power spectrum of baryons that trace the Jeans-scale smoothed dark matter density field can be calculated following (GH), where they studied simple schemes to approximate the effect of gas pressure. One such scheme that has fractional errors on the 10% level for overdensities $`10`$ is to filter the density perturbations in the linear regime as $`P_b^2=f_b^2(k/k_\mathrm{F})P_{\delta \delta }^2`$ and treat the system as collisionless baryonic particles. Their best fit is obtained with the filter
$$f_b=\frac{1}{2}[e^{x^2}+\frac{1}{(1+4x^2)^{1/4}}]$$
(42)
and suggests $`k_\mathrm{F}=34\mathrm{\Omega }_m^{1/2}h`$Mpc<sup>-1</sup> as a reasonable choice for the thermal history dependent filtering scale.
For such baryons, we assume that their temperature is related to photoionization energy ($`25`$ eV). The mass fraction of baryons present in such small overdensities as a function of redshift is obtained through the numerical simulations of .
### H Results & Discussion
In Fig. 1(a-b), we show the pressure power power spectrum today ($`z=0`$), written such that $`\mathrm{\Delta }^2(k)=k^3P(k)/2\pi ^2`$ is the power per logarithmic interval in wavenumber. In Fig 1(a), we show individual contributions from the single and double halo terms and a comparison to the Jeans-scale smoothed dark matter density field power spectra, both linear and non-linear, following and using fitting function. Here, we have taken the electron temperature to be the virial temperature given in Eq. (7). Shown here is also the gas bias $`b_\mathrm{\Pi }(k)`$; at large scales, $`b_\mathrm{\Pi }(k)3`$ as $`k0`$, consistent with numerically measured bias for gas (e.g., ) and analytical estimates in . The pressure power spectrum is such that at scales below $``$ few $`h`$ Mpc<sup>-1</sup>, the pressure fluctuations are suppressed relative to the dark matter power spectrum; The resulting power spectrum can also be described as a smoothed, but biased, version of the dark matter power spectrum. The scale at which smoothing enters in to the power spectrum is determined by the scale radius of the dark matter and gas profiles. Thus, the direct measurement of the pressure power spectrum, to some extent, can be used as a probe of halo profiles.
In the same figure, we also show the measured pressure power spectrum in hydrodynamical simulations by for their $`\mathrm{\Lambda }`$CDM model. For comparative purposes, we have appropriately corrected their pressure power spectra based on the mean temperature of baryons as tabulated since our definition of the pressure power spectrum includes the temperature. The resolution of simulations limit the accuracy of power spectrum to the range in wavenumbers of $`0.2k2.0`$ h Mpc<sup>-1</sup>, and is only based on a single realization. In this range, we find that our analytical models predict more power than what is measured, while agreement is observed at scales of a $``$ few h Mpc<sup>-1</sup>. The extent to which our analytical calculations agree with simulations is encouraging; this is the first time that a detailed analytical model for the pressure power spectrum has been compared with a numerically measured one. Numerical simulations with improving dynamical range and resolution will eventually test the reliability of models such as the one presented here as a useful description of the pressure fluctuations of the universe. Till then, we consider the present model as an appropriate descrption of the large scale pressure fluctuations.
In Fig. 1(b), we show the dependence of pressure power as a function of maximum mass used in the calculation with maximum mass ranging from $`10^{16}`$, $`10^{15}`$, $`10^{14}`$ and $`10^{13}`$ M$`_{}`$ from top to bottom. Here, we have shown the single halo contribution. Also shown are the total contribution to pressure power spectrum when there is an additional source of energy. Here, we have taken the minimum temperature to be $``$ 0.75 keV; power spectra, in general, scale as the square of this energy if the real preheating energy is higher or lower than the one considered here. There are clear differences between the pressure power spectra with and without an additional source of energy. With increasing such additional non-gravitational energy, note that $`b`$ in Eq. (6) $`0`$ such that $`\rho _g(r)\rho _{g0}`$. Thus, there is no longer a clear turn over in the pressure power spectrum since the effect of smoothing resulting from scale radius $`r_s`$ is not present. The changes suggest the possibility that physical properties associated with large scale structure baryons can be probed with pressure power spectrum. In fact, the combined study of dark matter and pressure power spectra may allow a consistent determination of halo properties, and to break certain degeneracies associated with dark matter halo profile and concentration as noted in , while at the same time investigating presence of additional sources of energy.
In Fig. 2 and 3, we study the cross-correlation power spectra between pressure and density field and pressure and galaxy distribution, respectively. These power spectra are relevant to the study of correlations present between, say, SZ and weak gravitational lensing and SZ and galaxies, or a similar tracer of large scale structure, such as radio sources. The presence of additional source of energy clearly affects the cross-correlation power spectra, suggesting the possibility that such effects may be investigated using cross-correlations between a tracer of pressure fluctuations and a tracer of matter density fluctuations.
Since the bispectrum generally scales as the square of the power spectrum, it is useful to define
$$\mathrm{\Delta }_{\mathrm{eq}}^2(k)\frac{k^3}{2\pi ^2}\sqrt{B(k,k,k)},$$
(43)
which represents equilateral triangle configurations. In Fig. 4, we show the pressure bispectrum as produced by baryons present in virialized halos. Here, most of the contributions at relevant scales come from the single halo term. Given the additional dependence on temperature, and thus mass, the bispectrum is more strongly sensitive to the presence of rare and most massive halos. Thus, an increase in energy of such rare halos does not significantly change the pressure bispectrum, but such energy changes contribute when halos of mass $`10^{14}`$ M$`_{}`$.
## III SZ effect
The temperature decrement along the line of sight due to SZ effect can be written as the integral of pressure along the same line of sight
$$y\frac{\mathrm{\Delta }T}{T_{\mathrm{CMB}}}=g(x)𝑑ra(r)\frac{k_B\sigma _T}{m_ec^2}n_e(r)T_e(r)$$
(44)
where $`\sigma _T`$ is the Thomson cross-section, $`n_e`$ is the electron number density, $`r`$ is the comoving distance, and $`g(x)=x\mathrm{coth}(x/2)4`$ with $`x=h\nu /k_BT_{\mathrm{CMB}}`$ is the spectral shape of SZ effect. At Rayleigh-Jeans (RJ) part of the CMB, $`g(x)=2`$.
The spectral dependance of SZ effect is unique in that it can be separated from most other contributors to CMB temperature fluctuations, including the primary anisotropy itself. As discussed in Cooray et al. , the upcoming multifrequency CMB satellite and ballone-borne data, among which Planck providing the greatest information on SZ, allow the possibility for detailed studies on the SZ effect including its higher order correlations such as bispectrum and skewness. Since these observations are projected measurements of the pressure power spectrum and bispectrum along the line of sight, we now provide analytical predictions based for SZ effect based on our model for the pressure fluctuations.
### A SZ Power Spectrum
The angular power spectrum of the SZ effect is defined in terms of the multipole moments $`y_{lm}`$ of temperature fluctuations as
$$y_{lm}^{}y_{l^{}m^{}}=C_l^{\mathrm{SZ}}\delta _{ll^{}}\delta _{mm^{}}.$$
(45)
$`C_l^{\mathrm{SZ}}`$ is numerically equal to the flat-sky power spectrum in the flat sky limit. The SZ power spectrum can be written as a redshift projection of the pressure power spectrum
$$C_l^{\mathrm{SZ}}=𝑑r\frac{W^{\mathrm{SZ}}(r)^2}{d_A^2}P_{\mathrm{\Pi }\mathrm{\Pi }}^\mathrm{t}(\frac{l}{d_A},r),$$
(46)
where $`d_A`$ is the angular diameter distance. At RJ part of the frequency spectrum, the SZ weight function is
$$W^{\mathrm{SZ}}(r)=2\frac{k_B\sigma _T\overline{n}_e}{a(r)^2m_ec^2}$$
(47)
where $`\overline{n}_e`$ is the mean electron density today. In deriving Eq. (46), we have used the Limber approximation by setting $`k=l/d_A`$ and flat-sky approximation. In previous studies (e.g., and references therein), the SZ power spectrum due to massive halos have been calculated following projected $`y`$ parameter of individual halos. The two approaches are essentially same since the order in which the projection is taken does not matter, except that our approach allows us to calculate intermediate 3d properties of baryons, mainly pressure.
In Fig. 5(a), we show the SZ power spectrum due to baryons present in virialized halos compared with our previous prediction for SZ effect using a biased power spectrum for pressure fluctuations following non-linear dark matter power spectrum. As shown, most of the contributions to SZ power spectrum comes from individual massive halos, while the halo-halo correlations only contribute at a level of 10% at large angular scales. This is contrary to, say, the lensing convergence power spectrum discussed in , where most of the power at large angular scales is due to the halo-halo correlations. The difference can be understood by noting that the SZ effect is strongly sensitive to the most massive halos due to $`TM^{2/3}`$ dependence in temperature and to a lesser, but somewhat related, extent that its weight function increases towards low redshifts. The lensing weight function selectively probes the large scale dark matter density power spectrum at comoving distances half to that of background sources ($`z0.2`$ to 0.5 when sources are at a redshift of 1), but has no extra dependence on mass. We have also shown current upper limits on the temperature fluctuations at arcminute scale angular scales where potentially the physical properties of baryons can be studies. These limits come from (BIMA) and (ATCA).
Also shown is the contribution to SZ effect from baryons present in overdensities $`10`$ (curve labeled GH). The SZ power spectrum here was calculated by replacing the pressure power spectrum in Eq. (46) with the unbiased Jeans-scale smoothed dark matter power spectrum following and assuming a mean temperature of 25 eV for these baryons. The mass fraction of baryons present in such small overdensities were taken from numerical simulations of and roughly follows $`0.25(1+z)`$, such that at a redshift of 3 and above all of the baryons are present in such small overdensities. The power spectrum due to such baryons are roughly three orders of magnitude lower than the power spectrum predicted for SZ effect from baryons in virialized halos, but as shown in Fig. 5(b), this level is consistent with what is predicted for SZ effect when halos with mass greater than $`10^{13}`$ M$`_{}`$ is not present in observed fields.
As shown in Fig. 5(b), the lack of massive halos leads to a strong suppression of power, and halos with masses greater than $`10^{15}`$ M$`_{}`$ are needed to obtain the full power spectrum. The lack of massive halos not only lead to a change in the power spectrum at large angular scales, the lack of mases also affect the contribution at small angular scales. Increasing the minimum temperature of electrons from the values determined by virial theorem to a minimum energy value of 0.75 keV significantly affects the change resulting from the lack of massive halos. In fact, with a minimum energy for baryons, the change is smaller when halos with masses less than $`10^{14}`$ M$`_{}`$ is considered. At the higher end of masses, the minimum energy do not significantly affect the power spectrum; the resulting change is less than 30% compared to the power spectrum with electron temperature based on the virial theorem. The variations suggest several observational possibilities, including the determination of minimum electron temperature, ie. the energy related to preheating if it exists, by calculating the power spectrum with massive halos substratced in a wide-field SZ map such as the one that will be eventually made with Planck.
For less area surveys, such as planned interferometric observations of the wide-field SZ effect (e.g., the few square degree survey of Carlstrom et al. ), the sample variance due to lack of massive halos in observed fields can be problematic in the interpretation of the observed signal. The problem arises from the fact that massive halos which contribute to the SZ power spectrum are rare and that observations in small fields will not contain such adeqaute masses to provide the fully expected SZ signal. The dependance of SZ effect on massive halos is even problematic for numerical simulations with limited box sizes. As pointed out by , the measured power spectrum in their simulation varies significantly based on the considred line of sight.
The dependance of signal on massive halos is also present in other observables of large scale structure, such as weak gravitational lensing. Compared to weak lensing surveys, studied in and , the SZ effect depends more strongly on rare halos. Most of these halos are at low redshifts, thus, surveys which avoid regions with known clusters will inherently also include an additional bias. As an example, the contribution to 1-$`\sigma `$ detection of temperature anisotropies at arcminute scales by due to SZ effect requires detailed knowledge on the distribution of halo masses in the observed fields. For a measurement of the SZ power spectrum, with a sample variance less than 20%, requires observations of a field $``$ 1000 deg<sup>2</sup>, while the same can be achieved in an area of $``$ 100 deg<sup>2</sup> for lensing. As discussed in , however, the sample variance due to lack of massive and rare halos, which dominate the SZ power, does not directly imply a systematic bias as long as one uses an approach similar to the one suggested by carefully accounting for the sample variance that may be present from lack of massive halos. Such an approach requires a reliable model for the SZ effect and detailed numerical simulations will be required for such a study.
These issues can be ignored for the upcoming wide field CMB experiments, such as Planck, where the frequency coverage will allow the recovery of SZ effect over 65% of the sky not confused by galactic emissions, thereby, providing an accurate measurement of its power spectrum and higher order statistics (see, ). Such a wide-field SZ map is also highly desirable for several reasons including the presence of adequate mass distribution of the universe such that a fair sample is considered and the possibility to use such a wide-field map for various cross-correlations purposes, such as against Sloan galaxy distribution or a wide-field weak lensing survey (see, § III C and III D).
### B SZ Bispectrum & Skewness
The angular bispectrum of the SZ effect is defined as
$$y_{l_1m_1}y_{l_2m_2}y_{l_3m_3}=\left(\begin{array}{ccc}l_1& l_2& l_3\\ m_1& m_2& m_3\end{array}\right)B_{l_1l_2l_3}^{\mathrm{SZ}}.$$
(48)
and can be written following Limber approximation as
$`B_{l_1l_2l_3}^{\mathrm{SZ}}`$ $`=`$ $`\sqrt{{\displaystyle \frac{(2l_1+1)(2l_2+1)(2l_3+1)}{4\pi }}}\left(\begin{array}{ccc}l_1& l_2& l_3\\ 0& 0& 0\end{array}\right)`$ (52)
$`\times \left[{\displaystyle 𝑑r\frac{[W^{\mathrm{SZ}}(r)]^3}{d_A^4}B_\mathrm{\Pi }^\mathrm{t}(\frac{l_1}{d_A},\frac{l_2}{d_A},\frac{l_3}{d_A};r)}\right].`$
The more familiar flat-sky bispectrum is simply the expression in brackets . The basic properties of Wigner-3$`j`$ symbol introduced above can be found in .
Similar to the density field bispectrum, we define
$$\mathrm{\Delta }_{\mathrm{eq}l}^2=\frac{l^2}{2\pi }\sqrt{|B_{lll}^{\mathrm{SZ}}|},$$
(53)
involving equilateral triangles in $`l`$-space. The absolute value of $`B^{\mathrm{SZ}}`$ is considered in above since $`B^{\mathrm{SZ}}g(x)^3`$, which is a negative quantity at RJ part of the frequency spectrum with $`g(x)=2`$.
In Fig. 6(a), we show individual contributions $`\mathrm{\Delta }_{\mathrm{eq}l}^2`$ with a maximum mass of $`10^{16}`$ M$`_{}`$. As shown, the main contribution to bispectrum comes from individual halo term, while other terms involving correlations between halos contribute $`10`$%. In Fig. 6(b), we show bispectrum as a function of maximum mass used in the calculation. Here, we have shown the total contribution to the bispectrum in solid lines while the total bispectrum in the presence of a minimum temperature of 0.75 keV is shown with a dot-dashed line. The bispectrum, as discussed in § II H, is strongly sensitive to the single halo term due to additional mass weighing. Almost all of the contributions to the SZ bispectrum comes from the single halo term. The same dependence on mass massive rare halos decreases the effect of a temperature increase when maximum mass in $`10^{15}`$ M$`_{}`$. This is in contrast to the power spectrum, where differences are still present with an increase in temperature from virial to a minimum of 0.75 keV.
The measurement of full bispectrum in million to billion pixel data of a wide-field SZ map as the one that’ll be produced with Planck can in general be difficult. In fact, there is no algorithm yet to measure the full bispectrum in such wide-field data in a reasonable time and computational requirements. Given such a possibility, it is interesting to consider a collapsed measurement of the bispectrum; real space statistics such as the third moment and skewness allow this possibility. In fact, the skewness has now been measured for the COBE data by , while the bispectrum measurements have only been limited to specific configurations of the bispectrum such as equilateral triangles in $`l`$-space . The skewness allows an easily measurable aspect of the bispectrum and will be porbably be one of the first measurements of non-Gaussianity in a wide-field SZ map. The skewness can be calculated using the second, $`y^2(\sigma )`$, and third, $`y^3(\sigma )`$, moments of the SZ effect:
$$S_3(\sigma )=\frac{y^3(\sigma )}{y^2(\sigma )^2},$$
(54)
where the two moments are
$`y^3(\sigma )`$ $`=`$ $`{\displaystyle \frac{1}{4\pi }}{\displaystyle \underset{l_1l_2l_3}{}}\sqrt{{\displaystyle \frac{(2l_1+1)(2l_2+1)(2l_3+1)}{4\pi }}}`$ (58)
$`\times \left(\begin{array}{ccc}l_1& l_2& l_3\\ 0& 0& 0\end{array}\right)B_{l_1l_2l_3}^{\mathrm{SZ}}W_{l_1}(\sigma )W_{l_2}(\sigma )W_{l_3}(\sigma ).`$
and
$$y^2(\sigma )=\frac{1}{4\pi }\underset{l}{}(2l+1)C_l^{\mathrm{SZ}}W_l^2(\sigma ).$$
(60)
In Fig. 7(a), we show the absolute value of skewness, $`|S_3(\sigma )|`$, as a function of smoothing scale $`\sigma `$ when the maximum mass included in the calculation ranges from $`10^{16}`$ to $`10^{13}`$ M$`_{}`$ (from top to bottom). The absolute value of skewness is considered since $`S_3g(x)^1`$, which is a negative quantity at RJ part of the frequency spectrum with $`g(x)=2`$. As shown, the SZ skewness is heavily dependent on the presence of most massive and rare halos. The introduction of a minimum energy of 0.75 keV leads to decrease in skewness, which results from the fact that the power spectrum is more affected than the bispectrum by such an increase.
Given the dependence on rare and most massive halos, the SZ effect is considerably non-Gaussian. As discussed in , the non-Gaussianity can be used as a useful tool for the identification and separation of most rare and massive halos from the wide-field CMB data. An optimised algorithm that utilised both the non-Gaussianity of SZ effect and its frequency dependance will be useful for constructing a catalog of SZ clusters in upcoming wide-field data. With massive clusters separated out, the remaining contribution to SZ effect will be from halos of mass $`10^{14}`$ M$`_{}`$, such as galaxy groups. In , we considered contribution from such small halos as the one due to large scale structure. The extent to which such small halos contribute to the SZ power spectrum and higher order statistics clearly depend on the role of additional energy in baryons. Therefore, any measurement of the power spectrum with known massive clusters removed, can be in return used as a probe of physical properties related to baryons, mainly the extent to which prehating affects the electron temperature of low mass halos.
Instead of individual non-Gaussian statistics as skewness, one can construct the probability distribution function (pdf) using a SZ map smoothed on some scale $`\sigma `$. The use of pdf as a probe of cosmology was first suggested by for weak gravitational lensing convergence. The same technique can be easily extended to SZ. Using the Edgeworth expansion to capture small deviations from Gaussianity, one can write the pdf of SZ effect to second order as
$`p(y)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\pi y^2(\sigma )}}}e^{y(\sigma )^2/2y^2(\sigma )}`$ (61)
$`\times `$ $`\left[1+{\displaystyle \frac{1}{6}}S_3(\sigma )\sqrt{y^(\sigma )}H_3\left({\displaystyle \frac{y(\sigma )}{\sqrt{y^2(\sigma )}}}\right)\right],`$ (62)
where $`H_3(x)=x^33x`$ is the third order Hermite polynomial (see, for details).
In Fig. 7(b), we show the pdf of SZ effect at 12 as a function of maximum mass used in the calculation. As shown, the greatest departure from non-Gaussianity occur when the maximum mass of halos are greater than $`10^{14}`$ M$`_{}`$. Given that we have only constructed the pdf using terms only out to skewness, the presented pdfs should only be consider as approximate; With increasing non-Gaussianity behavior, the approximated pdfs are likely to depart from true distributions especially in the tails. Observationally, the pdf can be constructed easily by considering a histogram of the pixel temperature values of the SZ map. Though such a construction sounds straightforward, there is likely to be complications on the interpretation of such a pdf in the presence of instrumental noise and other foregrounds. Techniques which do no directly make a wide-field map, especially interferometric observations, will again require special techniques to construct the pdf. Therefore, the extent to which the full pdf, or such an histogram, can be used a probe of cosmology and the accuracy to which pdfs can be constructed from upcoming wide-field CMB anisotropy data, such as Planck and planned interferometric surveys (Carlstrom, private communication), need to be investigated in detail. We leave these issues for further study.
### C SZ-Weak Lensing Cross-correlation
Similar to the SZ power spectrum, the angular power spectrum of weak lensing convergence can be defined in terms of the multipole moments $`\kappa _{lm}`$ as
$$\kappa _{lm}^{}\kappa _{l^{}m^{}}=C_l^\kappa \delta _{ll^{}}\delta _{mm^{}},$$
(63)
and can be written in terms of the dark matter power spectrum by
$$C_l^\kappa =𝑑r\frac{W^\kappa (r)^2}{d_A^2}P_\delta ^\mathrm{t}(\frac{l}{d_A};r),$$
(64)
When all background sources are at a distance of $`r_s`$, the lensing weight function becomes
$$W^\kappa (r)=\frac{3}{2}\mathrm{\Omega }_m\frac{H_0^2}{c^2a}\frac{d_A(r)d_A(r_sr)}{d_A(r_s)}.$$
(65)
The detail properties of lensing statistics, under the dark matter halo approach, is discussed in and .
The cross-correlation between the SZ effect and weak gravitational lensing can be similarly defined in terms of the individual multipole moments as
$$\kappa _{lm}^{}y_{l^{}m^{}}=C_l^{\mathrm{SZ}\kappa }\delta _{ll^{}}\delta _{mm^{}}.$$
(66)
This is now related to the dark matter-pressure power spectrum by
$$C_l^{\mathrm{SZ}\kappa }=𝑑r\frac{W^{\mathrm{SZ}}(r)W^\kappa (r)}{d_A^2}P_{\mathrm{\Pi }\delta }^\mathrm{t}(\frac{l}{d_A};r).$$
(67)
Finally the cross-correlation coefficient between SZ and weak lensing is
$$C(\mathrm{SZ},\kappa )_l=\frac{C_l^{\mathrm{SZ}\kappa }}{\sqrt{C_l^{\mathrm{SZ}}C_l^\kappa }}$$
(68)
### D SZ-Galaxy Cross-correlation
Similar to the SZ-weak lensing cross-correlation, one can study the cross-correlation between the galaxy distribution, which traces the large scale structure, and the SZ effect. The power spectrum of galaxy distribution can be defined terms of the multipole moments $`g_{lm}`$ as
$$g_{lm}^{}g_{l^{}m^{}}=C_l^g\delta _{ll^{}}\delta _{mm^{}},$$
(69)
and can be written as a projection of the dark matter power spectrum
$$C_l^g=𝑑r\frac{W^g(r)^2}{d_A^2}P_\delta ^\mathrm{t}(\frac{l}{d_A};r).$$
(70)
Here, $`C_l^g`$ should be understood as the 2d Fourier transform of the galaxy correlation function, generally referred to as $`w(\theta )`$ in the literature. The weight function for galaxy projection involves the redshift distribution
$$W^\kappa (r)=\frac{dN}{dr}$$
(71)
normalized such that $`𝑑r(dN/dr)=1`$.
As before. the cross-correlation between the SZ effect and galaxy distribution can be similarly defined in terms of the individual multipole moments as
$$g_{lm}^{}y_{l^{}m^{}}=C_l^{\mathrm{SZ}g}\delta _{ll^{}}\delta _{mm^{}}.$$
(72)
This is now related to the dark matter-pressure power spectrum by
$$C_l^{\mathrm{SZ}g}=𝑑r\frac{W^{\mathrm{SZ}}(r)W^g(r)}{d_A^2}P_{\mathrm{\Pi }\delta }^\mathrm{t}(\frac{l}{d_A};r).$$
(73)
Finally the cross-correlation coefficient between SZ and weak lensing is
$$C(\mathrm{SZ},g)_l=\frac{C_l^{\mathrm{SZ}\kappa }}{\sqrt{C_l^{\mathrm{SZ}}C_l^g}}$$
(74)
In Fig. 8, we show the SZ-weak lensing and SZ-galaxy cross-correlation power spectra as a function of maximum mass used in the calculation, while the correlation coefficients are shown in Fig. 9. In order to describe the galaxy distribution, we have considered a survey at low redshifts, similar to the Sloan Digital Sky Survey (SDSS)<sup>*</sup><sup>*</sup>*http://www.sdss.org. Such a low redshift tracer is desirable since contributions to SZ effect primarily comes from large scale structures at redshifts $`<1`$. Here, we have assumed the redshift distribution of Sloan galaxies follow $`dN/dzz^2\mathrm{exp}[(z/z_c)^{3/2}]`$ with a mean redshift $`z_m`$, of 0.2 ($`z_c1.412z_m`$). For SZ-weak lensing, the cross-correlation is such that SZ and lensing traces each other out to angular scales of $``$ 1000 when most massive and rarest halos are involved with a decrease in cross-correlation between the two at small angular scales. For SZ-galaxy cross-correlation, there is additional correlation at large angular scales, when compared to lensing, while the correlation is suppressed at small angular scales. The decrease at small angular scales is due to the fact that small halos that contribute to SZ and lensing do not necessarily contribute to the galaxy power spectrum. With an increase in additional energy for halos, the cross-correlation increases by few percent, however, this small increase unlikely to be determined accurately through observations. The correlation is sensitive to the redshift distribution of galaxies, which depends mostly on the selection criteria imposed by observations. The selection function of weak lensing can be considered well understood, however, a straight forward interpretation of any observed SZ-galaxy cross-correlation will require a clear understanding of observable related to galaxy distribution.
The SZ-SZ, lensing-lensing and SZ-lensing power spectra dependent different on the bias and correlation parameters. Since the bias and correlation are scale and redshift dependent, the measurement of these power spectra, which are projected along the redshift distributions, do not allow a direct probe of these quantities. A useful approach would be to consider the inversion of these power spectra in redshift bins by considering the measured lensing-lensing and SZ-lensing power spectra as a function of redshift. Note that the SZ-SZ power spectrum cannot be easily separated in redshift space as we do not have the ability to separate individual redshift contributions, unlike say in lensing, where one can use the redshifts of background sources to construct convergence as a function of redshift. With adequate signal-to-noise from wide-field surveys, it is likely that such an approach will allow studies to be carried out on the extent to which temperature weight baryons trace the dark matter and their correlation properties.
Similarly, as studied in , the cross-correlation of SZ against galaxy data, as a function of redshift, is expected to provide information on the properties of clustering of galaxies with respect to the temperature weighted baryon field represented by SZ effect. Such a cross-correlation will help understand the extend to which hot/warm gas is present in the outskirts of individual galaxies. With individual sets of SZ, weak lensing and galaxy maps, it is likely that a tremendous amount of information on physical properties associated with dark matter, baryon and galaxy distribution will be obtained through both a comparison of individual power spectra and higher order moments and cross-correlations and higher order moments associated with such cross-correlations. We hope to study some of these possibilities in detail in future studies.
## IV Summary & Conclusions
Using an extension of the dark matter halo approach, we have presented an efficient method to calculate the large scale structure pressure power spectrum and its high order moments, such as bispectrum. We have divided the contribution to large scale pressure power spectrum based on the overdensities in which contributing baryons are present with (1) baryons present in virialized halos with overdensities greater than $``$ 200 and in hydrostatic equilibrium with the density field of such halos, (2) photoionized baryons in overdensities less than $``$ 10 and which trace the Jeans-scale smoothed dark matter density field, and (3) baryons present in the mid overdensity regime which are likely to be undergoing collapse and shock heating.
Our approach allows us to calculate not only 2d statistics such as the projected pressure power spectrum, or the SZ effect, which will be observed, but also the 3d statistics that can be directly compared to predictions based on numerical simulations. We have performed such a comparison to recently published numerical simulations by and found good agreement between our analytical calculations and their simulations. The current simulations are limited to a handful of realizations and limited dynamical range and resoltion. With improving resolution and accuracy, analytical models such as the one presented here will be tested in detail against numerical calculations. Analytical calculations, aided by numerical simulations, will eventually allow detailed studies of large scale baryon distribution using observations such as the wide-field SZ effect.
The projected pressure power spectrum along the line of sight, provides a direct calculation of the large scale structure SZ effect and its higher order correlations. In the absence of massive and rare halos, we have suggested that baryons present in small overdensities provide a lower limit to any contribution to SZ effect. The extent to which baryons present in overdensities between 10 and 200 contribute to the correlations in large scale pressure and, from it, the SZ effect, requires additional studies, preferably with numerical simulations. Presently, our Understanding the role of preheating and its effect of baryons will also be another challenge as the SZ observations will clearly depend on such additional energy contributions to large scale baryon distribution. We have suggested the possibility of using SZ power spectrum and higher order correlations, such as the SZ skewness, as a probe of preaheating. Such a study will require a wide-field SZ map and this task will be completed with Planck observations. The unique frequency dependance of the SZ effect, together with its non-Gaussian behavior, will allow the construction of a reliable SZ cluster catalog with will aid in cosmologicalo studies of structure formation.
Our approach allows one to study possible systematic effects that may be present in upcoming SZ observations of small area fields due to the presence or absence of rare massive halos in such fields that will be observed. We have shown that the SZ effect as well as its non-Gaussian properties are mainly due to the most massive and rarest virialized halos in the universe. The lack of massive halos in observed SZ fields can introduce a systematic bias in the power spectrum, but the sample variance introduced by the lack of such masses, can be easily corrected based on the prior knowledge on mass distribution of observed fields. Due to additional mass dependence through temperature, the effect of mass is such that the SZ effect is more dependent on the rare halos than weak gravitational lensing convergence. The same SZ halos also contribute to lensing convergence and the cross-correlation between SZ and lensing can be used a probe of clustering properties between density and temperature weighted baryon fields. Given the great potential to study baryon distribution using SZ, various issues suggested here involving such correlations merit further study.
###### Acknowledgements.
The author greatly thanks Wayne Hu for useful discussions and helpful suggestions that led to the calculations presented in this paper. We also acknowledge useful discussions with Gil Holder, Lloyd Knox and Joe Mohr and thank Alexandre Refregier, Ue-Li Pen and their collaborators for providing results from numerical simulation presented in .
|
warning/0005/quant-ph0005128.html
|
ar5iv
|
text
|
# Concurrent Quantum Computation
## Abstract
A quantum computer is a multi-particle interferometer that comprises beam splitters at both ends and arms, where the $`n`$ two-level particles undergo the interactions among them. The arms are designed so that relevant functions required to produce a computational result is stored in the phase shifts of the $`2^n`$ arms. They can be detected by interferometry that allows us to utilize quantum parallelism. Quantum algorithms are accountable for what interferometers to be constructed to compute particular problems. A standard formalism for constructing the arms has been developed by the extension of classical reversible gate arrays. By its nature of sequential applications of logic operations, the required number of gates increases exponentially as the problem size grows. This may cause a crucial obstacle to perform a quantum computation within a limited decoherence time. We propose a direct and concurrent construction of the interferometer arms by one-time evolution of a physical system with arbitrary multi-particle interactions. It is inherently quantum mechanical and has no classical analogue. Encoding the functions used in Shor’s algorithm for prime factoring, Grover’s algorithm and Deutsch-Jozsa algorithm requires only one-time evolution of such a system regardless of the problem size $`n`$ as opposed to its standard sequential counterpart that takes $`O(n^3)`$, $`O(n)`$ and $`O(n2^n)`$.
A computation entails encoding of a function whether classically or quantum mechanically. The encoding of a function has been carried out in the form of a bit-flip. Such an “oracle" $`U_\mathrm{c}`$ in reversible classical computers and also in proposed quantum computers is a transformation of an $`n`$-bit input $`x`$ and a work bit $`w`$,
$$U_\mathrm{c}:|x|w|x|wf(x),$$
(1)
where $``$ stands for exclusive-OR. When the work bit is initially set to be 0, the oracle returns the function value $`f(x)`$ in the work bit. The customary construction of such an oracle pertains to sequential application of reversible gates. In quantum computation, however, the oracle is often converted into a transformation of the form,
$$U_\mathrm{q}:|x(1)^{f(x)}|x,$$
(2)
where the information about $`f(x)`$ is encoded in the phase of a linear superposition state so that quantum parallelism could be utilized. The conversion of $`U_\mathrm{c}`$ into $`U_\mathrm{q}`$ is performed by supplementary transformation or automatically by appropriately initialized work bit. The construction of the oracle (1) by sequential application of one-bit and two-bit gates demand an exponential or polynomial number of operations as the problem size grows as shown in Table 1. In order to perform a quantum computation, we need to complete the construction of the oracle while the coherence of the system is maintained, although classical computation does not bring the issue of a limited decoherence time into a question. We propose a concurrent construction of the transformation $`U_\mathrm{q}`$ by only one-time evolution of a physical system that has arbitrary multi-particle interactions. The exponentially hard work that is expressed by the complexity of gate arrays to be applied to a physical system in the case of the sequential construction will be replaced by adjustment of an exponentially large number of coupling strengths in the system before the coherence of the system is required.
The system required to concurrently implement an arbitrary $`n`$-bit Boolean function $`f(x)`$ for an $`n`$-bit string $`x`$ consists of $`n`$ two-level particles with arbitrary multi-particle interactions among them. The Hamiltonian of the system is,
$``$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}\mathrm{}\omega _i\sigma _{iz}+{\displaystyle \underset{i<j}{}}\mathrm{}\omega _{ij}\sigma _{iz}\sigma _{jz}+{\displaystyle \underset{i<j<k}{}}\mathrm{}\omega _{ijk}\sigma _{iz}\sigma _{jz}\sigma _{kz}`$ (3)
$`+\mathrm{}+\mathrm{}\omega _{12\mathrm{}n}\sigma _{1z}\sigma _{2z}\mathrm{}\sigma _{nz},`$
where $``$ stands for tensor product. The eigenstates of the Pauli matrix $`\sigma _{iz}=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)`$ for the $`i`$-th particle is used as the computational basis $`\{|0,|1\}`$ for the $`i`$-th bit $`x_i`$, where $`\sigma _{iz}|x_i=(1)^{x_i}|x_i`$ ($`x_i=\{0,1\}`$). The computational basis of $`x`$ is defined in the $`2^n`$-dimensional state space spanned by the two-basis states of the $`n`$ particles, $`|x=|x_n\mathrm{}x_1=|x_n\mathrm{}|x_1`$. The state $`|x`$ represents a number $`x=_{i=1}^n2^{i1}x_i`$. The Hamiltonian is diagonal in the computational basis $`\{|00\mathrm{}0,|00\mathrm{}1,\mathrm{},|11\mathrm{}1\}`$. Diagonal elements of the $`(2^n1)`$ terms in the Hamiltonian and a $`2^n`$-dimensional vector $`(1,1,\mathrm{},1)`$ (diagonal elements of the $`2^n\times 2^n`$ identity matrix) constitute an orthonormal basis to expand $`2^n`$-dimensional vectors. It suggests that the transformation $`U_\mathrm{q}`$ in (2), which is a ($`2^n\times 2^n`$)-diagonal matrix, can be constructed by a global phase shift and one-time evolution of the system, $`U_{}(\tau )=e^{\frac{i}{\mathrm{}}\tau }`$ for time $`\tau `$ as
$$U_\mathrm{q}=e^{i\varphi }U_{}(\tau ).$$
(4)
The concurrent construction of the transformation $`U_\mathrm{q}`$ for a given function comprises a preparation of the Hamiltonian (by adjusting the coefficients, $`\mathrm{}\omega _i`$, $`\mathrm{}\omega _{ij}`$, $`\mathrm{}`$, $`\mathrm{}\omega _{12\mathrm{}n}`$) so that (4) is satisfied and a time-evolution by that Hamiltonian. Only the latter process requires the conserved quantum coherence of the system.
The diagonal elements of (4) are decomposed as
$`e^{i\pi f(x)}`$ $`=`$ $`e^{i\varphi }\times {\displaystyle \underset{i}{}}e^{i\omega _i\tau (1)^{x_i}}\times {\displaystyle \underset{i<j}{}}e^{i\omega _{ij}\tau (1)^{x_i+x_j}}`$ (5)
$`\times {\displaystyle \underset{i<j<k}{}}e^{i\omega _{ijk}\tau (1)^{x_i+x_j+x_k}}`$
$`\times \mathrm{}\times e^{i\omega _{12\mathrm{}n}\tau (1)^{x_1+x_2+\mathrm{}x_n}},`$
which gives the solution,
$`\varphi ={\displaystyle \frac{\pi }{2^n}}{\displaystyle \underset{x}{}}f(x),`$
$`\omega _i\tau ={\displaystyle \frac{\pi }{2^n}}{\displaystyle \underset{x}{}}(1)^{x_i}f(x),`$
$`\omega _{ij}\tau ={\displaystyle \frac{\pi }{2^n}}{\displaystyle \underset{x}{}}(1)^{x_i+x_j}f(x),`$
$`\mathrm{}`$
$`\omega _{12\mathrm{}n}\tau ={\displaystyle \frac{\pi }{2^n}}{\displaystyle \underset{x}{}}(1)^{x_1+x_2+\mathrm{}+x_n}f(x),`$ (6)
and determines the coefficients in the Hamiltonian to be prepared so that only one-time evolution for time $`\tau `$ by the Hamiltonian will construct $`U_\mathrm{q}`$ by itself for a given function $`f(x)`$. Note that when $`f(x)`$ is a constant function, all coefficients in (6) except $`\varphi `$ are zero.
The formalism for concurrently constructing the transformation $`U_\mathrm{q}`$ can be applied to encode functions and calculated values that are used in existing quantum algorithms. Examples are shown below.
Deutsch-Jozsa algorithm. The algorithm solves the following problem by quantum parallelism. Given the oracle for an $`n`$-bit Boolean function $`f:\{0,1\}^n\{0,1\}`$, determine either (A) $`f`$ is a constant function (at 0 or 1) or (B) $`f`$ is a balanced function (the sequence $`f(0)`$, $`\mathrm{}`$, $`f(2^n1)`$ contains exactly $`2^{n1}`$ zeros and $`2^{n1}`$ ones). In the original algorithm, the oracle that has a form of $`U_\mathrm{c}`$ is used twice together with another unitary operation on the work bit $`S:|w(1)^w|w`$ between the two applications. The three unitary operations applied to the system in series are designed so that the information about $`f(x)`$ is transferred from the work bit to the phase of the control register $`|x`$ as in (2). Now we have a method to construct the transformation $`U_\mathrm{q}`$ concurrently, and the work bit can be removed.
When $`f(x)`$ takes only 0 or 1, $`f(x)`$ in (6) can be replaced by $`f(x)`$ since $`\pi f(x)\pi f(x)`$ (mod $`2\pi `$). Therefore a set of parameters obtained by replacing $`f(x)`$ by $`(1)^{x_1+x_2+\mathrm{}+x_n}(f(x)2N_x)`$ ($`N_x`$ is an integer) in (6) is also solution to (5). For a balanced function, we choose $`N_x=x_1x_2`$, so that $`\omega _{12\mathrm{}n}\tau =\frac{\pi }{2^n}\left(2^{n1}_x2N_x\right)=0`$. Thus, to encode a balanced $`n`$-bit Boolean function, we need the multi-particle interactions up to $`(n1)`$-th order (the $`n`$-particle interaction is not required). Other parameters are determined as,
$`\varphi ={\displaystyle \frac{\pi }{2^n}}{\displaystyle \underset{x}{}}(1)^{x_1+x_2+\mathrm{}+x_n}(f(x)2N_x),`$
$`\omega _i\tau ={\displaystyle \frac{\pi }{2^n}}{\displaystyle \underset{x}{}}(1)^{x_1+\mathrm{}+x_{i1}+x_{i+1}\mathrm{}+x_n}(f(x)2N_x),`$
$`\mathrm{}.`$ (7)
Grover’s algorithm. The algorithm explains how a data can be found out of $`2^n`$ random data entries. The data search problem is described by a function $`f(x)`$ that returns 1 for a single unknown value of $`x`$, say $`x=\tau `$, and 0 for the rest of $`x`$. The algorithm uses the information about $`f(x)`$ encoded in phase of the control register as in (2), which can be implemented concurrently, when the coefficients in the Hamiltonian are chosen as $`\varphi =\frac{\pi }{2^n}`$, $`\omega _i\tau =\frac{\pi }{2^n}(1)^{\tau _i}`$, $`\omega _{ij}\tau =\frac{\pi }{2^n}(1)^{\tau _i+\tau _j}`$, $`\mathrm{}`$, and $`\omega _{12\mathrm{}n}\tau =\frac{\pi }{2^n}(1)^{\tau _1+\tau _2+\mathrm{}\tau _n}`$.
(iii) Simon’s algorithm determines whether a given function $`f:\{0,1\}^n\{0,1\}^m`$ with $`mn`$ is periodic $`f(x)=f(x^{})x^{}=xs`$ (with a nontrivial string $`s`$) or one-to-one. Encoding the function $`f(x)`$ in the oracle of the form $`U_\mathrm{c}`$ in the original algorithm can be replaced by
$$|xe^{i\frac{\pi }{2^{m1}}f(x)}|x,$$
(8)
which can be constructed concurrently. The transformation (8) has the same form as $`U_\mathrm{q}`$, where $`(1)=e^{i\pi }`$ is replaced by $`e^{i\pi /2^{m1}}`$. The formalism for constructing the transformation $`U_\mathrm{q}`$ works exactly the same way. In order to implement an arbitrary function in this problem, all multi-particle interactions up to $`n`$-particle interaction in the Hamiltonian are necessary.
Shor’s algorithm for prime factorization. In order to factorize an odd number $`N`$, we randomly choose $`a`$ ($`N`$ and $`a`$ need to be relatively prime) and find the order $`r`$ of $`a`$, the least $`r`$ that satisfies $`a^r1\mathrm{mod}N`$. Finding the order is the prime part of Shor’s algorithm and starts with the transformation on two $`n`$-bit registers,
$$\frac{1}{\sqrt{q}}\underset{x=0}{\overset{q1}{}}|x|1\frac{1}{\sqrt{q}}\underset{x=0}{\overset{q1}{}}|x|a^x\mathrm{mod}N,$$
(9)
where $`q=2^n`$ satisfies $`N^2q<2N^2`$. Replacing the transformation (9) by,
$$\frac{1}{\sqrt{q}}\underset{x=0}{\overset{q1}{}}|x\frac{1}{\sqrt{q}}\underset{x=0}{\overset{q1}{}}e^{i\frac{\pi }{2N}(a^x\mathrm{mod}N)}|x,$$
(10)
which encodes $`a^x\mathrm{mod}N`$ in the phase of the control register instead of in the work bit, leaves the algorithm unchanged. The phase factors $`a^x\mathrm{mod}N`$ in (10) can be calculated classically as,
$$e^{i\frac{\pi }{2N}(a^x\mathrm{mod}N)}=e^{i\frac{\pi }{2N}_{i=1}^n(a^{2^{i1}}\mathrm{mod}N)^{x_i}},$$
(11)
where products refer to multiplication mod ($`N`$). Equation (11) has the same form as (5) if $`\lambda _i^{x_i}=\frac{1+\lambda _i}{2}+\frac{1\lambda _i}{2}(1)^{x_i}`$ ($`\lambda _i=a^{2^{i1}}\mathrm{mod}N`$) is used. Therefore, the time-evolution of the system for time $`\tau `$ constructs the transformation (11) by itself when the coefficients are chosen as, $`\varphi =_{i=1}^n\frac{1+\lambda _i}{2}`$, $`\omega _i\tau =\varphi \times \frac{1\lambda _i}{1+\lambda _i}`$, $`\omega _{ij}\tau =\varphi \times \frac{1\lambda _i}{1+\lambda _i}\times \frac{1\lambda _j}{1+\lambda _j}`$, $`\mathrm{}`$, and $`\omega _{12\mathrm{}n}\tau =_{i=1}^n\frac{1\lambda _i}{2}`$.
Quantum Fourier transform. To perform Quantum Fourier transform $`A_q`$ for an $`n`$-bit register ($`q=2^n`$),
$$A_q:|x\underset{y=0}{\overset{q1}{}}e^{2\pi ixy/q}|y,$$
(12)
as used in Shor’s algorithm, we need the Walsh-Hadamard transformation on each bit, $`H=\frac{1}{\sqrt{2}}\left(\begin{array}{cc}1& 1\\ 1& 1\end{array}\right)`$ in the computational basis, and controlled-phase-shift operators on pairs of bits, defined as $`S_{j,k}=e^{i\theta _{kj}x_jx_k}`$ on the $`j`$-th bit and $`k`$-th bit with $`j<k`$ where $`\theta _{kj}=\pi /2^{kj}`$. The right-hand side of (5) is equal to $`S_{j,k}`$ if we choose $`2\omega _{jk}\tau =\theta _{kj}`$, $`2\omega _j\tau =\theta _{kj}`$, $`2\omega _k\tau =\theta _{kj}`$ and the rest of the parameters are zero, aside from the global phase factor. A product of multiple controlled-phase-shift operators $`S_{l,l^{}}S_{l,l^{}1}\mathrm{}S_{l,l+1}`$ is also diagonal in the computational basis and can be implemented concurrently by setting $`2\omega _{l,m}\tau =\theta _{ml}`$ ($`m=l+1,\mathrm{},l^{}`$), $`2\omega _m\tau =\theta _{ml}`$, $`2\omega _l\tau =_m\theta _{ml}`$. One-bit rotations and two-particle interactions in the Hamiltonian are sufficient to implement controlled-phase-shift operators and products of those concurrently.
The time evolution of a system that concurrently implements functions by our proposed scheme and the number of necessary elementary gates for the function implementation by means of the standard sequential scheme are compared in Table 1 for existing quantum algorithm. It is challenging to find a system that has multi-particle interactions with reasonably large and controllable strengths in order to implement arbitrary functions, but if found, many functions and quantum algorithms will be implemented by only one-time evolution of the system, which may cross out the current biggest obstacle to quantum computation, the short decoherence time of a quantum system.
Scaling laws of evolution time and number of elementary gates necessary to implement functions for existing algorithms by means of the proposed concurrent implementation and the standard sequential implementation, respectively. In the most general case, an $`n`$-bit Boolean function $`f(x):\{0,1\}^n\{0,1\}`$ is encoded in phases of the control register $`|x`$, as $`U_\mathrm{q}:|x(1)^{f(x)}|x`$, in order to utilize quantum interference effect. To implement such an $`n`$-bit Boolean function, we prepare the $`n`$-particle system that has multi-particle interactions among them, from two-particle interactions $`I^{(2)}`$ up to $`n`$-particle interactions $`I^{(n)}`$ in addition to one-bit rotations, $`R`$. A one-time evolution by the Hamiltonian with properly chosen coefficients of the terms in it builds the transformation $`U_\mathrm{q}`$ by itself. A Boolean function $`f(x)`$ used in Deutsch-Jozsa problem has a constraint that $`f(x)`$ is either constant (at 0 or 1) or balanced. Because of the constraint, the concurrent implementation does not call for the $`n`$-particle interaction $`I^{(n)}`$. In Grover’s data search algorithm, the function to be implemented $`f(x)`$ is zero for all $`x`$ but $`\tau `$ ($`f(\tau )=1`$) which we search for. As a general Boolean function, this function requires the system that has all multi-particle interactions in order to be implemented in a concurrent fashion. The transformation $`U_\mathrm{q}`$ is constructed by successive applications of elementary gates (one-bit gates and two-bit gates) by the standard sequential implementation with the help of a work bit prepared in the state $`\frac{1}{\sqrt{2}}(|0|1)`$ and another qubit. The number of required gate operations to implement an $`n`$-bit Boolean function in the sequential manner scales as $`O(n2^n)`$ as opposed to a one-time evolution of a system in the case of the concurrent implementation. The function implemented in Grover’s algorithm is special in that it only needs one $`(n+1)`$-bit gate that is constructed by $`O(n)`$ elementary gates. In Shor’s algorithm for prime factoring an odd number $`N`$, $`a^x\mathrm{mod}N`$ ($`a`$ is a randomly chosen integer relatively prime to $`N`$) is encoded in a work bit, as $`|x|0|x|a^x\mathrm{mod}N`$. The construction of a gate array for this transformation requires $`O(n^3)`$ elementary gates by means of the sequential implementation. Instead of encoding $`|a^x\mathrm{mod}N`$ is the work bit, encoding it in the phase of the control register $`|x`$, as $`|x\mathrm{exp}\left[i\frac{\pi }{2N}(a^x\mathrm{mod}N)\right]|x`$, still works. The factor $`\frac{\pi }{2N}`$ is determined to differentiate all possible values $`|a^xN`$ takes (between 0 and $`N1`$) and also to avoid a destructive interference that ruins the algorithm. In Simon’s algorithm, a function to be implemented is $`f:\{0,1\}^n\{0,1\}^m(mn)`$. The sequential implementation of such a function in a work bit is in need of $`O(mn2^n)`$ elementary gates. As in Shor’s algorithm, the function $`f`$ can be implemented in the phase of the control register $`|x`$, as $`|x\mathrm{exp}\left[\frac{\pi }{2^{m1}}f(x)\right]`$. To distinguish all $`2^m`$ possible values $`f`$ takes, the phase space of $`2\pi `$ is divided by $`2^m`$ (a Boolean function, $`f:\{0,1\}^n\{0,1\}`$, is a special case where $`m=1`$. In both Shor’s and Simon’s algorithms, the concurrent implementation of necessary functions involves all multi-particle interactions $`I^{(2)}`$, $`\mathrm{}`$, $`I^{(n)}`$ in the system. A controlled-phase-shift operator $`S_{j,k}`$ acts on a pair of bits, in this case $`j`$-th and $`k`$-th qubits of the control register $`|x`$. It adds a phase factor $`\mathrm{exp}\left(i\frac{\pi }{2^{kj}}\right)`$ only when both $`x_k`$ and $`x_j`$ are ones. Only two-particle interactions $`I^{(2)}`$ between $`k`$-th and $`j`$-th particles to implement this operator. Such controlled-phase-shift operators on all pairs of qubits ($`n(n1)/2`$ paris in total) compose Quantum Fourier transform, which is used in Shor’s algorithm, along with the Walsh-Hadamard transformation on each qubit ($`H_l`$ on $`l`$-th qubit). The Walsh-Hadamard transformations are applied to all qubits, from $`x_1`$ to $`x_n`$, and controlled-phase-shift operators $`S_{l,n}S_{l,n1}\mathrm{}S_{l,l+1}`$ are interleaved between $`H_l`$ and $`H_{l+1}`$. A controlled-phase-shift operator and a product of multiple of them can be concurrently implemented by a one-time evolution of a system that has only two-particle interactions and one-bit rotations. In total, Quantum Fourier transform is constructed by $`n`$ Walsh-Hadamard transformations and $`n`$ time-evolutions of a system in a concurrent manner, whereas in the case of the sequential implementation, $`n(n1)/2`$ controlled-phase-shift operators are implemented in series.
|
warning/0005/math0005150.html
|
ar5iv
|
text
|
# Banach embedding properties of non-commutative 𝐿^𝑝-spaces
## 1. Introduction
Let $`𝒩`$ be a finite von Neumann algebra and $`1p<2`$. Our main theorem yields that $`C_p`$ is not linearly isomorphic to a subspace of $`L^p(𝒩)`$ (where $`C_p`$ denotes the Schatten $`p`$-class). It follows immediately that for any infinite von Neumann algebra $``$, $`L^p()`$ is not isomorphic to a subspace of $`L^p(𝒩)`$, since $`C_p`$ is then isomorphic to a subspace of $`L^p()`$. (It is proved in \[S1\] that also $`C_p`$ does not embed in $`L^p(𝒩)`$ for any $`2<p<\mathrm{}`$.)
For $`𝒩`$ a semi-finite von-Neumann algebra and $`\tau `$ a faithful normal semi-finite trace on $`𝒩`$, $`L^p(\tau )`$ denotes the non-commutative $`L^p`$ space associated with $`(𝒩,\tau )`$ (see e.g., \[FK\]). The particular choice of trace $`\tau `$ is unimportant, for if $`\beta `$ is another such trace, $`L^p(\beta )`$ is isometric to $`L^p(\tau )`$. We also denote this (isometrically unique) Banach space by $`L^p(𝒩)`$.
Given $`C1`$ and non-negative reals $`a`$ and $`b`$, let $`a\stackrel{C}{}b`$ denote the equivalence relation $`\frac{1}{C}abCa`$. Sequences $`(x_j)`$ and $`(y_j)`$ in Banach spaces $`X`$ and $`Y`$ respectively all called $`C`$-equivalent if
(1.1)
$$\underset{i=1}{\overset{n}{}}\alpha _ix_i\stackrel{C}{}\underset{i=1}{\overset{n}{}}\alpha _iy_i\text{ for all }n\text{ and scalars }\alpha _1,\mathrm{},\alpha _n.$$
(Equivalently, there exists an invertible linear map $`T:[x_i][y_i]`$ with $`T,T^1C`$, where $`[x_i]`$ denotes the closed linear span of $`(x_i)`$.) $`(x_j)`$ is called unconditional if there is a constant $`u`$ so that for any $`n`$ and scalars $`c_1,\mathrm{},c_n`$ and $`\epsilon _1,\mathrm{},\epsilon _n`$ with $`|\epsilon _i|=1`$ for all $`i`$, $`_{i=1}^n\epsilon _ic_ix_iuc_ix_i`$ (then one says $`(x_j)`$ is $`u`$-unconditional). The usual $`\mathrm{}^p`$-basis refers to the unit vector basis $`(e_j)`$ of $`\mathrm{}^p`$, where $`e_j(i)=\delta _{ji}`$ for all $`i`$ and $`j`$.
Our main result goes as follows.
###### Theorem 1.1.
Let $`𝒩`$ be a finite von Neumann algebra, $`1p<2`$, and let $`(x_{ij})`$ be an infinite matrix in $`L^p(\tau )`$ where $`\tau `$ is a fixed faithful normal tracial state on $`𝒩`$. Assume for some $`C1`$ that every row and column of $`(x_{ij})`$ is $`C`$-equivalent to the usual $`\mathrm{}^2`$-basis and that $`(x_{i_k,j_k})_{k=1}^{\mathrm{}}`$ is unconditional, whenever $`i_1<i_2<\mathrm{}`$ and $`j_1<j_2<\mathrm{}`$. Then there exist $`i_1<i_2<\mathrm{}`$ and $`j_1<j_2<\mathrm{}`$ so that setting $`y_k=x_{i_k,j_k}`$ for all $`k`$, then
(1.2)
$$\underset{n\mathrm{}}{lim}n^{1/p}\underset{i=1}{\overset{n}{}}y_i^{}_{L^p(\tau )}=0$$
for all subsequences $`(y_k^{})`$ of $`(y_k)`$.
###### Corollary 1.2.
Let $`p`$ and $`𝒩`$ be as in 1.1. Let $`X`$ be a Banach space spanned by an infinite matrix of elements $`(x_{ij})`$ so that for some $`\lambda 1`$,
* every row and column of $`(x_{ij})`$ is $`\lambda `$-equivalent to the usual $`\mathrm{}^2`$ basis
* $`(x_{i_n,j_n})_{n=1}^{\mathrm{}}`$ is $`\lambda `$-equivalent to the usual $`\mathrm{}^p`$-basis, for all $`i_1<i_2<\mathrm{}`$ and $`j_1<j_2<\mathrm{}`$.
Then $`X`$ is not Banach isomorphic to a subspace of $`L^p(\tau )`$. In particular, $`C_p`$ does not embed in $`L^p(\tau )`$.
The Corollary yields its final statement since the standard matrix units $`(x_{ij})`$ for $`C_p`$ satisfy (i) and (ii) with $`\lambda =1`$.
To see why 1.1 $``$ 1.2, suppose to the contrary that $`T:XX^{}L^p(\tau )`$ were an isomorphic embedding, where $`X`$ is as in 1.2. Then $`(Tx_{ij})`$ satisfies the hypotheses of 1.1 with $`C=\lambda TT^1`$. However if $`(i_k),(j_k)`$ satisfies the conclusion of Theorem 1.1, $`(Tx_{i_k,j_k})`$ and hence $`(x_{i_k,j_k})`$ cannot be equivalent to the usual $`\mathrm{}^p`$-basis, a contradiction.
Let $`\mathrm{Rad}C_p`$ denote the “Rademacher unconditionalized version” of $`C_p`$ $`(1p<\mathrm{})`$. That is, letting $`(r_{ij})`$ be an independent matrix of $`\{1,1\}`$-valued random variables with $`P(r_{ij}=1)=P(r_{ij}=1)=\frac{1}{2}`$ for all $`i,j`$, and letting $`(c_{ij})`$ be a matrix of scalars with only finitely many non-zero terms, then
(1.3)
$$(c_{ij})_{\mathrm{Rad}_{C_p}}=𝔼_\omega (r_{ij}(\omega )c_{ij})_{C_p}.$$
###### Corollary 1.3.
Let $`p`$ and $`𝒩`$ be as in 1.1. Then $`\mathrm{Rad}C_p`$ is not isomorphic to a subspace of $`L^p(\tau )`$.
###### Proof.
The standard matrix units basis $`(x_{ij})`$ of $`\mathrm{Rad}C_p`$ also satisfies the hypotheses of Corollary 1.2 with $`\lambda =1`$. ∎
Corollary 1.3 yields new information in the classical, commutative case of $`L^p`$. (Throughout, $`L^p`$ refers to $`L^p`$ on the unit interval, endowed with Lebesgue measure; i.e., $`L^p=L^p(𝒩)`$ where $`𝒩=L^{\mathrm{}}`$ acting on $`L^2`$ via multiplication.) This also reveals a remarkable difference in the structure of $`L^p`$-spaces, $`p<2`$ or $`p>2`$, for $`\mathrm{Rad}C_p`$ is isometric to a subspace of $`L^p`$ for $`2<p<\mathrm{}`$ (cf. Theorem 5 of \[L-P\]). Also, let us note that $`\mathrm{Rad}C_p`$ is isometric to a subspace of $`L^p`$ $`(C_p)`$ for $`1p<2`$, so we obtain an unconditionalized version of $`C_p`$ in $`L^p()`$ which also does not embed in $`L^p(𝒩)`$, for $`𝒩`$ finite, where $`=L^{\mathrm{}}B(H)`$. (Throughout, $`L^p(X)`$ refers to the Bochner-Lebesgue space $`L^p(X,m)`$, where $`m`$ is Lebesgue measure.)
It is a classical result of C.A. McCarthy that $`C_p`$ does not “locally” embed in $`L^p`$, for $`1p<\mathrm{}`$ \[McC\]. Corollary 1.2 yields an “infinite” dimensional proof of this result for $`1p<2`$, as well as the apparently new discovery that also $`\mathrm{Rad}C_p`$ does not locally embed in $`L_p`$ for these $`p`$. To see this, we give the following.
###### Definition.
Let $`1p<\mathrm{}`$, $`n`$, and $`\lambda 1`$. A finite-dimensional Banach space $`X`$ is called a $`\lambda `$-$`GC_p^n`$-space provided there is an $`(n\times n)`$-matrix $`(x_{ij})`$ spanning $`X`$ so that
* any row and column of $`(x_{ij})`$ is $`\lambda `$-equivalent to the usual $`\mathrm{}_n^2`$-basis
* $`(x_{i_k,j_k})_{k=1}^m`$ is $`\lambda `$-equivalent to the usual $`\mathrm{}_m^p`$ basis for any $`m`$,
$$1i_1<\mathrm{}<i_mn\text{ and }1j_1<j_2<\mathrm{}<j_mn.$$
An infinite-dimensional space $`X`$ is called a $`\lambda `$-$`GC_p`$-space provided it admits a spanning matrix $`(x_{ij})`$ satisfying (i) and (ii) of Corollary 1.2; finally $`X`$ is called a $`GC_p`$-space if it is a $`\lambda `$-$`GC_p`$-space for some $`\lambda 1`$.
$`C_p^n`$ refers to the $`n^2`$-dimensional Schatten $`p`$-class consisting of $`n\times n`$ matrices in the $`C_p`$ norm; “$`G`$” stands for “Generalized”. For example, $`\mathrm{Rad}C_p^n`$ is a 1-$`GC_p^n`$ space. The next result yields that $`\lambda `$-$`GC_p^n`$-spaces cannot be uniformly embedded in $`L^p`$, hence in particular, we recapture the classical fact mentioned above that $`L^p`$ does not contain $`C_p^n`$’s uniformly. (For isomorphic Banach spaces $`X`$ and $`Y`$, $`d(X,Y)=inf\{TT^1:T`$ is a surjective isomorphism from $`X`$ to $`Y\}`$).
###### Corollary 1.4.
Let $`1p<2`$ and $`\lambda 1`$. Define:
$$\beta _{n,\lambda }=inf\{d(X,Y):X\text{ is a }\lambda \text{-}GC_p^n\text{-space and }YL^p\}.$$
Then $`lim_n\mathrm{}\beta _{n,\lambda }=\mathrm{}`$.
###### Proof.
Suppose this were false. Then we could choose $`\lambda 1`$ and $`X_1,X_2,\mathrm{}`$ subspaces of $`L^p`$ so that $`X_n`$ is a $`\lambda `$-$`GC_p^n`$-space for all $`n`$. Choose then $`(x_{ij}^n)`$ an $`n\times n`$ matrix of elements of $`X_n`$, satisfying (i) and (ii) of the definition, for all $`n`$. Let $`M_{00}`$ denote the linear space of all infinite matrices of scalars with only finitely many non-zero entries. Let $`U`$ be a free ultrafilter on $``$. Define a semi-norm $``$ on $`M_{00}`$ by
(1.4)
$$(c_{ij})=\underset{nU}{lim}c_{ij}x_{ij}^n.$$
It is easily checked that $``$ is indeed a semi-norm; let $`W`$ be its null space; $`W=\{(c_{ij})M_{00}:(c_{ij})=0\}`$, and let $`X`$ denote the completion of $`(M_{00},)/W`$. It follows easily that $`X`$ is a $`\lambda `$-$`GC_p`$-space. By standard ultraproduct techniques, it follows that $`X`$ is finitely representable in $`L^p`$. But then (since ultraproducts of (commutative) $`L^p(\mu )`$ spaces are (commutative) $`L^p(\nu )`$ spaces and any separable subspace of an $`L^p(\nu )`$ space is isometric to a subspace of $`L^p`$), $`X`$ isometrically embeds in $`L^p`$. This contradicts Corollary 1.2. ∎
###### Remark.
Theorem 1.1 may easily be extended to the case of general finite von Neumann algebras $`𝒩`$, and not just the finite, $`\sigma `$-finite ones covered by its statement. Corollaries 1.2 and 1.3 also hold in this setting, as well as the general formulations of Theorems 4.1 and 4.2. Indeed, in general, one has that $`L^p(𝒩)`$ is isometrically isomorphic to $`L^p(\tau )`$ for some semi-finite faithful normal trace $`\tau `$ on $`𝒩`$. Let $`(x_{ij})`$ be a matrix of elements of $`L^p(\tau )`$ satisfying the assumptions of Theorem 1.1, and let $`P`$ be the supremum of all the support projections of $`x_{ij}`$ and $`x_{ij}^{}`$, $`i,j=1,2,\mathrm{}`$. Then $`P`$ is a $`\sigma `$-finite projection in $`𝒩`$, and thus $`P𝒩P`$ is both finite and $`\sigma `$-finite. Moreover all the $`x_{ij}`$’s belong to $`L^p(P𝒩P,\tau ^{})=PL^p(𝒩,\tau )P`$, where $`\tau ^{}=\tau |P𝒩P`$. But in turn, $`L^p(P𝒩P,\tau ^{})`$ is isometrically isomorphic to $`L^p(P𝒩P,\tau ^{\prime \prime })`$ for some faithful finite normal trace $`\tau ^{\prime \prime }`$ on $`P𝒩P`$. This reduces the proof of Theorem 1.1 in the case of general finite von Neumann algebras, to those with a finite trace.
We now give a description of the results and proof-order of the paper.
If a matrix satisfies the hypotheses of Theorem 1.1, then every row and column has the property that the $`p^{th}`$ powers of absolute values of the terms form a uniformly integrable sequence. We develop the basic technical tools to explain and exploit this, in Section 2, through the device of the $`p`$-modulus of an element of $`L^p(𝒩)`$ with respect to a normal tracial state $`\tau `$ on $`𝒩`$. We give several useful inequalities for this modulus in Lemma 2.3. Although many of these can be obtained from the literature (e.g., \[FK\]), we give full proofs for the sake of completeness. We also obtain equivalences for relative weak compactness in $`L^1(𝒩)`$ in terms of uniform integrability in Proposition 2.5, and a useful non-commutative truncation equivalence for general $`p`$, in Corollary 2.7.
We give technical information concerning general unconditional sequences in $`L^p(𝒩)`$ for $`p<2`$ in Lemmas 3.13.3, yielding in particular the following definitive equivalences obtained in Corollaries 3.4 and 3.5. Let $`(f_n)`$ be a bounded unconditional sequence in $`L^p(𝒩)`$. Then the following are equivalent.
* $`(f_n)`$ has no subsequence equivalent to the usual $`\mathrm{}^p`$ basis.
* $`(|f_n|^p)`$ is uniformly integrable.
* $`lim_n\mathrm{}n^{1/p}_{i=1}^nf_i^{}_{L^p(\tau )}=0`$ for all subsequences $`(f_n^{})`$ of $`(f_n)`$.
The proof of Theorem 1.1 is then completed, using the standard ultraproduct construction of the finite ultrapower of a finite von Neumann algebra $`𝒩`$, and a result giving the connection between its associated $`L^p`$ space and the Banach ultrapower of $`L^p(𝒩)`$ (Lemma 3.6).
Section 4 yields results considerably stronger than Theorem 1.1. The arguments here do not use the ultraproduct construction in Section 3, and are thus more elementary (but also more delicate). Theorem 4.1 gives the following result (which immediately implies Theorem 1.1).
If a semi-normalized matrix in $`L^p(𝒩)`$ is such that all columns and “generalized” diagonals are unconditional while all rows are $`u`$-unconditional for some fixed $`u`$, then three alternatives occur: Either some column has an $`\mathrm{}^p`$-subsequence, or $`\mathrm{}_n^p`$’s are finitely represented in the terms of the rows, or the matrix has a “generalized diagonal” $`(y_k)`$ satisfying (1.2) of Theorem 1.1.
Using results from Banach space theory, we obtain in Theorem 4.2 that if $`p=1`$ or if $`p>1`$ and $`𝒩`$ is hyperfinite, the unconditionality assumption in 4.1 may be dropped. The case $`p>1`$ also uses recent non-commutative martingale inequalities (see \[SF\], \[PX1\]). The case $`p=1`$ uses techniques from \[R1\], which yield results for sequences in the preduals of arbitrary von Neumann algebras which may be independent interest (see Lemmas 4.8 and 4.9). The proof in this case also requires an apparently new elementary finite disjointness result (Lemma 4.10B).
Section 5 contains rather quick applications of our main results and the techniques of their proofs. For example, Proposition 5.1 asserts that neither the Row nor Column operator spaces completely embed in the predual of a finite von Neumann algebra; this is a quick consequence of our main result. Theorem 5.4 shows that for $`1p<2`$ and $`𝒩`$ finite, a subspace of $`L^p(𝒩)`$ contains $`\mathrm{}_n^p`$’s uniformly iff it contains an almost disjointly supported sequence (which of course is then almost isometric to $`\mathrm{}^p`$), extending the previously known commutative case \[R2\]). We give the concepts of the $`p`$-Banach-Saks and strong $`p`$-Banach-Saks properties in Definition 5.5, and extend the classical results of Banach-Saks \[BS\] and Szlenk \[Sz\] in Proposition 5.6. This result also yields that for $`p`$ and $`𝒩`$ as above, a weakly null sequence in $`L^p(𝒩)`$ has the property that every subsequence has a strong $`p`$-Banach-Saks subsequence if and only if the $`p^{th}`$ powers of absolute values of its terms are uniformly integrable.
The main result of Section 6 shows that there are precisely thirteen Banach isomorphism types among the spaces $`L^p(𝒩)`$ for $`𝒩`$ hyperfinite semi-finite, $`1p<\mathrm{}`$, $`p2`$. The embedding properties of the various types for $`p<2`$ are given in an eight-level Hasse diagram, in Theorem 6.2. This work completes the classification and embedding properties of the type I case given in \[S2\]. The main work in establishing this Theorem is found in the non-embedding results given in Theorems 6.3 and 6.9; we also give a new proof of a non-embedding result in the type I case, established in \[S2\], in our Proposition 6.5. The most delicate of these is Theorem 6.9, which yields that if $``$ is a type II von-Neumann algebra, and $`L^p()`$ embeds in $`L^p(𝒩)`$, then also $`𝒩`$ must have a type II or type III summand ($`1p<2`$). Of course this reduces directly to the case where $``$ is the hyperfinite type II factor; the proof requires our Theorem 4.1, and also rests upon recent discoveries of M. Junge \[J\] and Pisier-Xu \[PX2\].
Our methods do not cover the following case, which remains a fascinating open problem: Is it so that the predual of a type III von-Neumann algebra does not Banach embed in the predual of one of type II? In fact, we do not know if the predual of the injective type II factor can be Banach isomorphic to the predual of an injective type III-factor. We show in Theorem 7.2 that such factors cannot in general be distinguished by the Banach space isomorphism class (or even operator space isomorphism class) of their preduals. Letting $`R_\lambda `$ denote the Powers injective factor of type III<sub>λ</sub> and $`R_{\mathrm{}}`$ denote the Araki-Woods injective factor of type III<sub>1</sub>, we show that $`(R_\lambda )_{}`$ is completely isomorphic to $`(R_{\mathrm{}})_{}`$ for all $`0<\lambda <1`$. (For a von Neumann algebra $`𝒩`$, $`𝒩_{}`$ denotes its predual, also denoted here by $`L^1(𝒩)`$.) Thus there are uncountably many isomorphically distinct injective factors, all of whose preduals are completely isomorphic. We also show in Theorem 7.2 that there are uncountably many isomorphically distinct injective type III<sub>0</sub>-factors, all of whose preduals are completely isomorphic to $`(R_{\mathrm{}})_{}`$.
We show in Theorem 7.3 that the famous open isomorphism problem for free group von Neumann algebras cannot be resolved by the Banach (or even operator) space structure of the predual. Namely, we prove that the preduals of the $`L(F_n)`$’s are all completely isomorphic, for $`2n\mathrm{}`$, where $`F_n`$ is the free group on $`n`$ generators and $`L(F_n)`$ its associated von Neumann algebra. This extends the result of A. Arias \[Ar\], showing that the $`L(F_n)`$’s themselves are completely isomorphic as operator spaces. The proof of Theorem 7.3 relies basically on the deep result of D. Voiculescu that $`L(F_{\mathrm{}})M_k(L(F_{\mathrm{}}))`$ as von Neumann algebras, for $`k=2,3,\mathrm{}`$ (cf. \[Vo\] or \[VDN\]).
The results in Section 7 also extend to the case of the non-commutative spaces $`L^p(𝒩)`$, for $`1<p<\mathrm{}`$ (see Theorem 7.5). These isomorphism results (as well as the “positive” isomorphism results in Section 6) rely on the operator space version of the so-called Pełczyński decomposition method (see Lemma 6.13). Thus, one actually shows for von Neumann algebras $`𝒩`$ and $``$, that each of the spaces $`L^p(𝒩)`$ and $`L^p()`$ is completely isometric to a completely contractively complemented subspace of the other, and also (e.g., in the free group case $`=L(F_{\mathrm{}}))`$, that say $`L^p()`$ also has the property that $`(L^p()\mathrm{}L^p()\mathrm{})_\mathrm{}^p`$ completely contractively factors through $`L^p()`$, which then implies the operator space isomorphism of these two spaces. Thus the proofs of these operator space isomorphism results are actually based on natural isometric embedding properties of the $`L^p(𝒩)`$ spaces themselves.
## 2. The modulus of uniform integrability and weak compactness in $`L^1(𝒩)`$
Let $`𝒩`$ be a finite von Neumann algebra, acting on a Hilbert space $`H`$. Let $`𝒫=𝒫(𝒩)`$ denote the set of all (self-adjoint) projections in $`𝒩`$. We shall assume that $`𝒩`$ is endowed with a faithful normal tracial state $`\tau `$, which is atomless. That is, for all $`P𝒫`$ with $`P0`$, there is a $`QP`$, $`Q𝒫`$, with $`0<\tau (Q)<\tau (P)`$. (Equivalently, $`0QP`$, since $`\tau `$ is faithful.)
These assumptions cause no loss in generality. Indeed, if $`𝒩`$ has a faithful normal trace $`\gamma `$, then simply replace $`𝒩`$ by $`\stackrel{~}{𝒩}=𝒩\overline{}L^{\mathrm{}}`$, where $`\stackrel{~}{𝒩}`$ is equipped with the atomless trace $`\gamma =\tau m`$, with $`m`$ the trace on $`L^{\mathrm{}}`$ given by integration with respect to Lebesgue measure on $`[0,1]`$. $`𝒩`$ is ($``$-isomorphic to) a subalgebra of $`\stackrel{~}{𝒩}`$, and hence $`L^p(𝒩)`$ is isometric to a subspace of $`L^p(\stackrel{~}{𝒩})`$, so we may as well assume our space $`X`$ in Theorem 1.1 is already contained in $`L^p(\stackrel{~}{𝒩})`$.
Now if $`𝒩`$ is a MASA, it follows easily that also $`\tau |`$ is atomless. Indeed, were this false, we could choose $`P0`$, $`P`$ so that $`0QP`$, $`Q`$ implies $`Q=0`$ or $`Q=P`$. But then choosing $`Q𝒫(𝒩)`$, $`0QP`$ with $`0<\tau (Q)<\tau (P)`$, we obtain that if $`\stackrel{~}{}`$ is the von Neumann algebra generated by $``$ and $`Q`$, $`\stackrel{~}{}`$ is also commutative and $`\stackrel{~}{}`$, a contradiction.
###### Definition 2.1.
Given $`f𝒩_{}=L^1(\tau )`$, we define the modulus of uniform integrability of $`f`$ as the function on $`^+`$, $`\epsilon \omega (f,\epsilon )`$ given by
(2.1)
$$w(f,\epsilon )=sup\{\tau (|fP|),P𝒫,\tau (P)\epsilon \}.$$
We also define the lower modulus of $`f`$, $`\epsilon \underset{¯}{\omega }(f,\epsilon )`$, as
(2.2)
$$\underset{¯}{\omega }(f,\epsilon )=sup\{|\tau (fP)|:P𝒫,\tau (P)\epsilon \}.$$
To handle the case $`p1`$ in our Main Theorem, we also use the following $`p`$-moduli. (When $`\tau `$ is fixed, we set $`f_p=f_{L^p(\tau )}=(\tau (|f|^p))^{1/p}`$. Also, for $`f𝒩`$, we set $`f_{\mathrm{}}=f_𝒩`$.)
###### Definition 2.2.
Let $`0<p<\mathrm{}`$ and $`fL^p(\tau )`$. The $`p`$-modulus of $`f`$, $`\omega _p(f,)`$, the symmetric $`p`$-modulus of $`f`$, $`\omega _p^s(f,)`$, and the spectral $`p`$-modulus of $`f`$, $`\stackrel{~}{\omega }_p(f,)`$ are given, for $`0\epsilon 1`$, by
(2.3) $`\omega _p(f,\epsilon )`$ $`=sup\{fP_p:P𝒫,\tau (P)\epsilon \},`$
(2.4) $`\omega _p^s(f,\epsilon )`$ $`=sup\{PfP_p:P𝒫,\tau (P)\epsilon \},`$
(2.5) $`\stackrel{~}{\omega }_p(f,\epsilon )`$ $`=sup\{\left({\displaystyle _{(r,\mathrm{})}}t^pd(\tau E_{|f|}(t))\right)^{1/p}:\tau E_{|f|}((r,\mathrm{}))\epsilon \}`$
where for $`g`$ self-adjoint, $`E_g`$ denotes the spectral measure for $`g`$.
It is trivial that all these moduli are increasing (i.e., non-decreasing) functions on $`^+`$, which are continuous at $`0`$, thanks to the assumption that $`fL^p(\tau )`$. Actually, the assumption that $`\tau `$ is atomless yields that $`\omega _p(f,)`$, $`\underset{¯}{\omega }(f,)`$ and $`\omega _p^s(f,)`$ are absolutely continuous on $`[0,1]`$.
We now give some basic properties of these moduli. The most important of these is that several of them reduce to the uniform integrability modulus given in Definition 2.1. In particular, we obtain for $`fL^p(\tau )`$ and $`\epsilon >0`$ that
$$\omega _p^s(f,\epsilon )\omega _p(f^{},\epsilon )=\omega _p(f,\epsilon )=(\omega (|f|^p,\epsilon ))^{1/p}2\omega _p^s(|f|,\epsilon ).$$
For any $`f`$ affiliated with $`𝒩`$, we let $`t\mu (f,t)`$ denote the decreasing rearrangement of $`|f|`$ on $`[0,1]`$; $`\mu (f,t)=inf\{r0:\tau E_{|f|}((r,\mathrm{}))t\}`$.
###### Lemma 2.3.
Let $`1p<\mathrm{}`$, $`f,gL^p(\tau )`$, and $`\epsilon >0`$.
(2.6) $`\omega _p(f+g,\epsilon )`$ $`\omega _p(f,\epsilon )+\omega _p(g,\epsilon )`$
and
$`\omega _p^s(f+g,\epsilon )`$ $`\omega _p^s(f,\epsilon )+\omega _p^s(g,\epsilon ).`$
If $`f`$ is self-adjoint, then
(2.7)
$$\begin{array}{cc}\hfill \omega _p(f,\epsilon )=\omega _p^s(f,\epsilon )& =(\underset{¯}{\omega }(|f|^p,\epsilon ))^{1/p}\hfill \\ & =\mathrm{max}\{fP_p:Pf=fP,P𝒫,\text{ and }\tau (P)=\epsilon \}\hfill \\ & =\left(_0^\epsilon \mu ^p(f,t)𝑑t\right)^{1/p}\hfill \end{array}$$
and
(2.8)
$$\omega (f,\epsilon )2\underset{¯}{\omega }(f,\epsilon )\text{ when }p=1.$$
In general,
(2.9)
$$\begin{array}{cc}\hfill \omega _p^s(f,\epsilon )\omega _p(f,\epsilon )& =\omega _p(f^{},\epsilon )\hfill \\ & =\omega _p(|f|,\epsilon )=(\underset{¯}{\omega }(|f|^p,\epsilon ))^{1/p}2\omega _p^s(f,\epsilon )\hfill \end{array}$$
and in case $`p=1`$,
(2.10)
$$\underset{¯}{\omega }(f,\epsilon )\omega (f,\epsilon )4\underset{¯}{\omega }(f,\epsilon ).$$
Finally, let $`r=\epsilon ^{1/p}f_p`$. There exists a spectral projection $`P`$ for $`|f|`$ so that $`fP𝒩`$ with
(2.11)
$$fP_{\mathrm{}}r\text{ and }f(IP)_p\stackrel{~}{\omega }_p(f,\epsilon )\omega _p(f,\epsilon ).$$
The case $`p>1`$ uses the following classical submajorization inequality, due to H. Weyl \[W\].
###### Sublemma.
Let $`f`$ and $`g`$ be decreasing non-negative functions on $`(0,1]`$ so that
$$_0^xf(t)𝑑t_0^xg(t)𝑑t\text{ for all }0<x1.$$
Then also
$$_0^xf^p(t)𝑑t_0^xg^p(t)𝑑t\text{ for all }1<p<\mathrm{},$$
all $`0<x1`$.
###### Remarks.
1. This follows easily from the corresponding “discrete” formulation, cf. \[GK\]. Also, the result holds in greater generality; one does not need the functions to be non-negative, and moreover the conclusion generalizes to assert that
$$_0^x\mathrm{\Phi }f(t)d_0^x\mathrm{\Phi }g(t)𝑑t\text{ for all }0<x1$$
all continuous convex functions $`\mathrm{\Phi }`$.
2. All the assertions of Lemma 2.3 hold for semi-finite von Neumann algebras $`𝒩`$ that are atomless (i.e., have no minimal projections), endowed with a faithful normal trace $`\tau `$. Several of its assertions can also be deduced from results in \[FK\] and \[CS\]. For example, once one proves the equality of the first and last terms in (2.7), one may apply Lemma 4.1 of \[FK\] to obtain several of the other equalities in (2.7), for $`p=1`$; one then has that $`\omega (T,\epsilon )=\mathrm{\Phi }_\epsilon (T)`$ in the notation of \[FK\], and some other results in Lemma 2.3 follow from Theorem 4.4 of \[FK\]. However we prefer to give a “self-contained” treatment, in part because we take the modulus $`\omega (f,\epsilon )`$ as the primary concept in our development.
###### Proof of Lemma 2.3.
Let $`p,f,g`$ and $`\epsilon `$ be as in the statement. (2.6) is a trivial consequence of the fact that $`_p`$ is a norm (i.e., the triangle inequality). Also, we easily obtain that
(2.12) $`\omega _p^s(f,\epsilon )\omega _p(f,\epsilon )=\omega _p(|f|,\epsilon )`$
(2.13) $`\stackrel{~}{\omega }_p(f,\epsilon )\omega _p(f,\epsilon )`$
and in case $`p=1`$,
(2.14)
$$\underset{¯}{\omega }(f,\epsilon )\omega (f,\epsilon ).$$
Indeed, if $`P𝒫`$, then
(2.15)
$$|fP|=(Pf^{}fP)^{1/2}=(P|f|^2P)^{1/2}=\left||f|P\right|$$
which immediately yields the equality in (2.12). Since compression reduces the $`L^p(\tau )`$ norm, we have
(2.16)
$$PfP_p=P(fP)P_pfP_p$$
which gives the inequality in (2.12). If $`0r`$ and $`\tau E_{|f|}((r,\mathrm{}))\epsilon `$, then setting $`P=E_{|f|}((r,\mathrm{}))`$,
(2.17)
$$\left(_{(r,\mathrm{})}t^p𝑑\tau E_{|f|}(t)\right)^{1/p}=|f|P_p\omega _p(f,t),$$
yielding the inequality in (2.13). (2.14) is trivial, since for any $`P𝒫`$,
(2.18)
$$|\tau (fP)|\tau (|fP|)=fP_1.$$
For the non-trivial assertions of the Lemma, we need the following basic identities (cf. \[FK\], \[CS\]).
(2.19)
$$f_p^p=_0^{\mathrm{}}t^p𝑑\tau E_{|f|}(t)_0^1\mu ^p(f,t)𝑑t.$$
(The final inequality is also an equality, but this follows from the conclusion of our Lemma.)
Now let $`f`$ be self-adjoint. Let $`𝒩(f)`$ denote the von Neumann algebra generated by $`f`$, and let $``$ be a MASA contained in $`𝒩`$ with $`𝒩(f)`$. Then by our initial remarks, $`\tau |`$ is atomless. Let us identify (as we may), $``$ and $`\tau |`$ with an atomless probability space $`(\mathrm{\Omega },𝒮,\nu )`$. It follows that we may choose a countably generated $`\sigma `$-subalgebra $`𝒮_0`$ of $`𝒮`$ so that $`f`$ is $`𝒮_0`$-measurable and also $`\nu |𝒮_0`$ is atomless. Denote the corresponding von-Neumann algebra by: $`L^{\mathrm{}}(\nu |𝒮_0)=_0`$.
It then follows that $`(\mathrm{\Omega },𝒮_0,\nu )`$ is measure-isomorphic to $`([0,1],,m)`$ (where $``$ denotes the Borel subsets of $`[0,1]`$ and $`m`$ denotes Lebesgue measure on $``$), and moreover the measure-isomorphism may be so chosen that the “random-variable” $`f`$ is carried over to the decreasing function $`t\mu (f,t)`$ (cf. Lemma 4.1 of \[CS\]). It now follows that
(2.20)
$$_0^x\mu ^p(f,t)𝑑t\omega _p^p(f,x).$$
Indeed, it follows that there exists a set $`S𝒮_0`$ with $`\nu (S)=x`$ and $`_S|f|^p𝑑\nu =\tau (|\chi _Sf|^p)=_0^x\mu ^p(f,t)𝑑t`$ (where $`\chi _S`$ may be interpreted as the projection in $`_0`$ obtained via multiplication). Now we define a quantity $`\beta `$ (depending on $`x`$) by
(2.21)
$$\beta =sup\{f\psi _1:\psi 𝒩,\psi _{\mathrm{}}1,|\tau (\psi )|x\}.$$
We are going to prove that there exists a $`G𝒫(_0)`$ with $`\tau (G)=x`$ and
(2.22)
$$\tau (|fG|)=\tau (|f|G)=\beta .$$
Note that the first equality in (2.22) is trivial, since $`Gf`$. But then all the equalities in (2.7) for the case $`p=1`$, follow immediately, for we have also that then $`|f|G=G|f|G=|GfG|`$ and so trivially $`\tau (|f|G)\underset{¯}{\omega }(|f|,x)\beta `$ and $`\tau (|f|G)\omega _1^s(f,x)\beta `$; of course also $`\omega (f,x)\beta `$, hence by (2.22), $`\beta =\omega (f,x)`$. Moreover by the argument for (2.20) and (2.22) we have that $`\beta =\tau (|f|G)=_0^x\mu (f,t)𝑑t`$.
Before proving this basic claim, let us see why it also yields (2.7) for $`p>1`$ (via the Sublemma). Still keeping $`x`$ fixed, assume $`0<x\epsilon 1`$, and suppose $`P𝒫`$ with $`\tau (P)\epsilon `$. Now setting $`g=|fP|`$, $`g`$ is self-adjoint and “supported” on $`P`$, whence it easily follows that $`\mu (g,t)=0`$ for $`t>\epsilon `$.
But now we obtain that
(2.23)
$$_0^x\mu (g,t)𝑑t_0^x\mu (f,t)𝑑t.$$
Indeed,
(2.24)
$$\begin{array}{cc}\hfill _0^x\mu (g,t)𝑑t& \omega (g,x)=\omega (fP,x)\hfill \\ & =sup\{fPQ_1:\tau (Q)x\}\hfill \\ & =sup\{|\tau (fPQ\phi )|:\phi 𝒩,\phi _{\mathrm{}}1\}\text{ (by duality)}\hfill \\ & \beta \hfill \end{array}$$
(since $`PQ𝒩`$, $`PQ_{\mathrm{}}1`$, and $`|\tau (PQ)|\tau (Q)x`$).
Now (temporarily) unfixing $`x`$, we also have that (2.23) holds for $`x>\epsilon `$, since $`\mu (g,t)=0`$ for all $`t>\epsilon `$. Thus the Sublemma yields that
(2.25)
$$_0^\epsilon \mu ^p(g,t)𝑑t_0^\epsilon \mu ^p(f,t)𝑑t.$$
Hence in view of (2.19),
(2.26)
$$fP_p^p_0^\epsilon \mu ^p(f,t)𝑑t,$$
and so at last
(2.27)
$$\omega _p(f,\epsilon )\left(_0^\epsilon \mu ^p(f,t)𝑑t\right)^{1/p}.$$
Of course (2.20) combined with (2.27) now yields that
(2.28)
$$\omega _p(f,\epsilon )=\left(_0^\epsilon \mu ^p(f,t)𝑑t\right)^{1/p},$$
and now all the equalities in (2.7) follow for $`p>1`$ as well.
We now establish (2.22). Using the polar decomposition of $`f`$ and duality, we have that
(2.29)
$$\begin{array}{cc}\hfill \beta & =sup\{|\tau (f\psi \phi )|:\psi ,\phi 𝒩,\psi _{\mathrm{}},\phi _{\mathrm{}}1\text{ and }|\tau (\psi )|x\}\hfill \\ & =sup\{\tau (|f|\psi ):\psi 𝒩,\mathrm{\hspace{0.17em}0}\psi 1,\tau (\psi )x\}\hfill \\ & =sup\{\tau (|f|\psi ):\psi ,\mathrm{\hspace{0.17em}0}\psi 1,\tau (\psi )x\}.\hfill \end{array}$$
The last equality follows by a conditional expectation argument from classical probability theory.
Indeed, given $`0\psi 1`$ in $`𝒩`$ with $`\tau (\psi )x`$, there exists a unique $`\stackrel{~}{\psi }_0`$ such that
(2.30)
$$\tau (g\psi )=\tau (g\stackrel{~}{\psi })\text{ for all }gL^1(_0).$$
It follows that then $`0\stackrel{~}{\psi }1`$ and $`\tau (\stackrel{~}{\psi })x`$; this yields the desired equality.
Now let $`K`$ be defined:
(2.31)
$$K=\{\psi _0:0\psi 1\text{ and }\tau (\psi )x\}.$$
Then $`K`$ is a weak\* compact convex set, thus
(2.32)
$$K=\omega ^{}\overline{\text{co}}\{\phi :\phi \mathrm{Ext}K\}$$
and moreover
(2.33)
$$\beta =sup\{\tau (|f|\phi ):\phi \mathrm{Ext}K\}.$$
Now we claim that if $`\phi \mathrm{Ext}K`$, $`\phi `$ is a projection. To see this, again identifying $`_0`$ with $`L^{\mathrm{}}(\mathrm{\Omega },𝒮_0,\nu |𝒮_0)`$, we regard $`\phi `$ as an $`𝒮_0`$-measurable function on $`\mathrm{\Omega }`$. Were $`\phi `$ not a projection, we could choose $`0<\delta <\frac{1}{2}`$ so that setting $`F=\{\omega \mathrm{\Omega }:\delta \phi (\omega )1\delta \}`$, then $`\mu (F)>0`$. Since $`\mu `$ is atomless, choose a measurable $`EF`$ with $`\mu (F)=\frac{1}{2}\mu (E)`$. Now define $`g`$ by
(2.34)
$$g=\frac{\delta }{2}\chi _E\frac{\delta }{2}\chi _{FE}.$$
Then $`g0`$, $`\tau (g)=0`$, and $`0\phi \pm g1`$. But then $`\tau (\phi \pm g)\epsilon `$, hence $`\phi \pm gK`$ and $`\phi =\frac{(\phi +g)+(\phi g)}{2}`$, contradicting the fact that $`\phi \mathrm{Ext}K`$. (For a proof of this claim in a more general setting, see \[CKS\].)
We finally observe that the supremum in (2.29) is actually attained, thanks to the $`\omega ^{}`$-compactness of $`K`$. But it then follows that this is attained at an extreme point of $`K`$, i.e., there indeed exists a $`G𝒫(_0)`$ with $`\tau (G)=x`$, satisfying (2.22).
We may now also easily obtain (2.8). Letting $`f=f^+f^{}`$ where $`f^+f^{}=0`$ and $`f^+,f^{}0`$, we have (by the proof of (2.7))
(2.35)
$$\begin{array}{cc}\hfill \omega (f,\epsilon )& =sup\{\tau (|f|P):P𝒫(_0),\tau (P)\epsilon \}\hfill \\ & =sup\{\tau (f^+P)+\tau (f^{}P):P𝒫(_0),\tau (P)\epsilon \}\hfill \\ & 2sup\{|\tau (fP)|:P𝒫(_0),\tau (P)\epsilon \}\hfill \\ & 2\underset{¯}{\omega }(f,\epsilon )\hfill \end{array}$$
The first equality in (2.9) follows from the fact that for a general $`f`$ affiliated with $`𝒩`$, there exists a unitary $`U`$ in $`𝒩`$ with $`f=U|f|`$ (thanks to the finiteness of $`𝒩`$). But then $`|f|`$ and $`|f^{}|`$ are unitarily equivalent, which yields that $`\mu (f,t)=\mu (f^{},t)`$ for all $`t`$, and hence the desired equality follows by the final equality in (2.7).
It remains to prove the last inequalities in (2.9) and (2.10), and the final statement of the lemma. Let $`f=g+ih`$ with $`g`$ and $`h`$ self-adjoint (and so in $`L^p(\tau )`$). Then
(2.36)
$$\begin{array}{cc}\hfill \omega _p(f,\epsilon )& \omega _p(g,\epsilon )+\omega _p(h,\epsilon )\text{ by (}\text{2.6}\text{)}\hfill \\ & =\omega _p^s(g,\epsilon )+\omega _p^s(h,\epsilon )\text{ by (}\text{2.7}\text{)}.\hfill \end{array}$$
But if $`\phi =g`$ or $`h`$, then
(2.37)
$$\omega _p^s(\phi ,\epsilon )\omega _p^s(f,\epsilon ).$$
Indeed, if $`P𝒫`$, $`\tau (P)\epsilon `$, then $`PfP=PgP+iPhP`$. But $`PgP`$ and $`PhP`$ are both self adjoint, hence $`P\phi P_pPfP_p`$, yielding (2.37). Of course (2.36) and (2.37) yield the final inequality in (2.9). Similarly, in case $`p=1`$,
(2.38)
$$\begin{array}{cc}\hfill \omega (f,\epsilon )& \omega (g,\epsilon )+\omega (h,\epsilon )\text{ by (}\text{2.6}\text{)}\hfill \\ & 2\underset{¯}{\omega }(g,\epsilon )+2\underset{¯}{\omega }(h,\epsilon )\text{ by (}\text{2.8}\text{)}\hfill \\ & 4\underset{¯}{\omega }(f,\epsilon )\hfill \end{array}$$
since we also have for $`\phi =g`$ or $`h`$, that $`\underset{¯}{\omega }(\phi ,\epsilon )\underset{¯}{\omega }(f,\epsilon )`$ (by an argument similar to that for (2.37)).
To obtain the final assertion of the lemma, let $`r=\mu (f,\epsilon )`$, and let $`E=E_{|f|}`$. Now if $`\overline{\epsilon }=\tau (E[r,\mathrm{}))`$ then since
(2.39)
$$E([r,\mathrm{}))=\{E([s,\mathrm{})):s<r\},$$
we have $`\epsilon \overline{\epsilon }`$. Thus
(2.40)
$$r^p\epsilon r^p\overline{\epsilon }_{[r,\mathrm{})}t^p𝑑\tau E(t)_{[0,\mathrm{})}t^p𝑑\tau E(t)=f_p^p.$$
Hence
(2.41)
$$r\epsilon ^{1/p}f_p.$$
Now also by the definition of $`r`$, $`\tau (E(r,\mathrm{}))\epsilon `$, and so
(2.42)
$$\tau (|f|^pE_{(r,\mathrm{})})=_{(r,\mathrm{})}t^p𝑑\tau E(t)\stackrel{~}{\omega }_p(f,\epsilon )^p.$$
Finally, let $`f=U|f|`$ be the polar decomposition of $`f`$. In particular, $`U`$ is a partial isometry belonging to $`𝒩`$. Then $`P=E([0,r])`$ satisfies (2.11). Indeed, $`fP=U|f|P`$ and $`|f|P_{\mathrm{}}r`$, so also $`U|f|P_{\mathrm{}}r`$, and
$$\begin{array}{cc}\hfill U|f|(IP)_p& |f|(IP)_p=(\tau (|f|^pE_{(r,\mathrm{})})^{1/p}\hfill \\ & \stackrel{~}{\omega }_p(f,\epsilon )\text{ by (}\text{2.42}\text{). }\mathit{}\hfill \end{array}$$
###### Remarks.
1. We have given a self-contained proof of the basic inequality (2.27) for the sake of completeness. An alternate deduction may be obtained as follows. The remarks preceding (2.20) actually yield that for any $`gL^p(\tau )`$, $`g_p=\mu (g,)_p`$. Let $`f`$ be as in the proof of (2.27) and fix a $`P𝒫`$ with $`\tau (P)=\epsilon `$. We apply this observation to $`g=fP`$. First, Proposition 1.1 of \[CS\] yields that for any $`0<x1`$,
$$_0^x\mu (fP,t)𝑑t_0^x\mu (f,t)\mu (P,t)𝑑t.$$
Hence applying the Sublemma and the observation,
$$\begin{array}{cc}\hfill fP_p^p=_0^1\mu (fP,t)^p𝑑t& _0^1(\mu (f,t)\mu (P,t))^p𝑑t\hfill \\ & =_0^\epsilon \mu ^p(f,t)𝑑t\hfill \end{array}$$
which of course yields (2.26) and hence (2.27).
2. Rather than making use of the measure isomorphism of $`(\mathrm{\Omega },𝒮_0,\nu |𝒮_0)`$ with $`([0,1],,m)`$, one can use the following more elementary procedure, in demonstrating (2.20). Let $`r=\mu (f,x)`$. Then it follows that setting $`P=E_{|f|}((r,\mathrm{}))`$, $`\tau (P)x`$ and $`\tau (E_{|f|}([r,\mathrm{})))x`$. Using that $`\tau |`$ is atomless, choose $`Q𝒫()`$ with $`QE_{|f|}(\{r\})`$ so that $`\tau (Q)+\tau (P)=x`$. Then
$$\begin{array}{cc}\hfill \tau (|f(P+Q)|^p)& =\tau (|f|^p(P+Q))\hfill \\ & =r\tau (Q)+_{(r,\mathrm{})}t^p𝑑\tau E_{|f|}(t)\hfill \\ & =_0^x\mu ^p(f,t)𝑑t.\hfill \end{array}$$
Here, the first two equalities are trivial; however the third one follows by a direct elementary (but somewhat involved) argument. (We are indebted to Ken Davidson for this Remark.)
We next use the modulus of uniform integrability to establish a criterion for relative weak compactness.
###### Definition 2.4.
A subset $`W`$ of $`L^1(\tau )`$ is called uniformly integrable if
$$\underset{\epsilon 0}{lim}\underset{fW}{sup}\omega (f,\epsilon )=0.$$
###### Comment.
The assumption that $`\tau `$ is atomless implies uniformly integrable subsets are bounded in $`L^1(\tau )`$. In fact, it then follows that if $`W`$ satisfies that $`sup_{fW}\omega (f,\epsilon _0)<\mathrm{}`$ for some $`\epsilon _0>0`$, $`W`$ is bounded.
###### Proposition 2.5.
Let $`(f_n)`$ be a given sequence in $`L^1(\tau )`$. The following are equivalent
* $`(f_n)`$ is relatively weakly compact in $`L^1(\tau )`$.
* $`(f_n)`$ is uniformly integrable.
* $`(|f_n|)`$ is relatively weakly compact.
* $`(f_n)`$ is bounded in $`L^1(\tau )`$ and $`lim_{\epsilon 0}sup_n\stackrel{~}{\omega }_1(f_n,\epsilon )=0`$.
* For all $`\epsilon >0`$, there exists an $`r<\mathrm{}`$ so that for all $`n`$,
$$d_{L^1(\tau )}(f_n,r_a(𝒩))<\epsilon .$$
Moreover if $`(f_n)`$ is bounded in $`L^1(\tau )`$ and
(2.43)
$$\eta =\underset{\epsilon 0}{lim}\underset{n}{sup}\omega (f_n,\epsilon )>0,$$
there exists a sequence $`P_1,P_2,\mathrm{}`$ of pairwise orthogonal projections in $`𝒫`$ and $`n_1<n_2<\mathrm{}`$ so that
(2.44)
$$|\tau (f_{n_k}P_k)|>\frac{\eta }{5}\text{ for all }k.$$
###### Remark.
$`_a(𝒩)`$ denotes the closed unit ball of $`𝒩`$; thus $`r_a(𝒩)=\{f𝒩:f_{\mathrm{}}r\}`$. For $`WL^1(\tau )`$ and $`fL^1(\tau )`$, $`d_{L^1(\tau )}(f,W)=inf\{fw_1:wW\}`$ by definition. Our proof of (iv) $``$ (v) reduces, via the proof of Lemma 2.3, to a standard truncation argument in the case of commutative $`𝒩`$.
###### Proof.
Once (i) $``$ (ii) is established, the other equivalences in this Proposition follow easily from 2.3. Indeed, we have by the equalities in (2.9) that
$$\underset{\epsilon 0}{lim}\underset{n}{sup}\omega (f_n,t)=\underset{\epsilon 0}{lim}\underset{n}{sup}\omega (|f_n|,\epsilon ),$$
whence we have the equivalence of (i)–(iii). Now trivially (ii) $``$ (iv) since $`\stackrel{~}{\omega }_1(f,\epsilon )\omega (f,\epsilon )`$ for any $`fL^1(\tau )`$ and $`\epsilon >0`$ (see (2.11)). Suppose first that $`(f_n)`$ satisfies (v). Then given $`\epsilon >0`$, for each $`n`$ we may choose $`\psi _n𝒩`$, $`\psi _n_{\mathrm{}}r`$, with
(2.45)
$$f_n\psi _n_{L^1(\tau )}<\epsilon .$$
But then for any $`\delta <\epsilon `$,
(2.46)
$$\omega (f_n,\delta )\omega (f_n\psi _n,\delta )+\omega (\psi _n,\delta )<\epsilon +r\delta .$$
Hence $`\overline{\mathrm{lim}}_{\delta 0}sup_n\omega (f_n,\delta )\epsilon `$, proving (ii). On the other hand, suppose (iv) holds. Let $`\epsilon >0`$, and choose $`\delta >0`$ so that
(2.47)
$$\stackrel{~}{\omega }_1(f_n,\delta )<\epsilon \text{ for all }n.$$
Also, let $`M=supf_n_{L^1(\tau )}`$. Then setting $`r=\delta ^1M`$, it follows by the final statement of Lemma 2.3 that for each $`n`$, we may choose $`\psi _nr_a𝒩`$ with
$$\psi _nf_n_{L^1(\tau )}\stackrel{~}{\omega }_1(f,\delta )<\epsilon ,$$
proving (iv) $``$ (v).
To prove the equivalences of (i) and (ii), we use the following classical criterion due to C. Akemann \[A\]: A bounded set $`W`$ in the predual of a von-Neumann algebra $``$ is relatively compact if and only if for any sequence $`P_1,P_2,\mathrm{}`$ of disjoint projections in $``$,
(2.48)
$$\underset{j\mathrm{}}{lim}\underset{wW}{sup}|P_j(w)|=0.$$
Now suppose first that $`(f_n)`$ is not relatively weakly compact; then choosing disjoint $`P_j`$’s as in the above criteria, we obtain that
(2.49)
$$\underset{j\mathrm{}}{\overline{\mathrm{lim}}}\underset{n}{sup}|\tau (P_jf_n)|=\delta >0.$$
But $`lim\tau (P_j)=0`$, since the $`P_j`$’s are disjoint. It follows immediately that
(2.50)
$$\underset{\epsilon 0}{lim}\underset{n}{sup}\underset{¯}{\omega }(f_n,\epsilon )\delta ,$$
which together with (2.10), proves that (ii) $``$ (i).
Finally, to show that (i) $``$ (ii), assume instead that $`\eta >0`$, where $`\eta `$ is given in (2.43). It now suffices to demonstrate the final assertion of 2.5, for then $`(f_n)`$ is not relatively weakly compact by Akemann’s criterion. Let $`0<\epsilon <\eta `$ with $`\frac{\eta }{4}\epsilon >\frac{\eta }{5}`$. By (2.43), choose $`n_1`$ with
(2.51)
$$\omega (f_{n_1},\frac{1}{2})>\eta \epsilon .$$
Then choose (by (2.10) of Lemma 2.3), $`Q_1𝒫`$ with $`\tau (Q_1)1/2`$ and
(2.52)
$$|\tau (f_{n_1}Q_1)|>\frac{\eta \epsilon }{4}.$$
Since $`f_{n_1}`$ is integrable, $`\{f_{n_1}\}`$ is uniformly integrable, so we may choose $`0<\epsilon _2<1`$ so that
(2.53)
$$\omega (f_{n_1},\epsilon _2)<\frac{\epsilon }{2}.$$
Next by (2.43), choose $`n_2>n_1`$ with
(2.54)
$$\omega (f_{n_2},\epsilon _2)>\eta \epsilon .$$
(It is easily seen, thanks to the uniform integrability of finite sets in $`L^1(\tau )`$, that in fact $`\eta =lim_{\epsilon 0}\overline{\mathrm{lim}}_n\mathrm{}\omega (f_n,\epsilon )`$; thus we may insure that $`n_2`$ may be chosen larger than $`n_1`$.) Again using (2.54) and (2.10), choose $`Q_2𝒫`$ with $`\tau (Q_2)\frac{\epsilon _2}{2^2}`$ and
(2.55)
$$|\tau (f_{n_2}Q_2)|>\frac{\eta \epsilon }{4}.$$
Then choose $`\epsilon _3<\epsilon _2`$ so that
(2.56)
$$\omega (f_{n_2},\epsilon _3)<\frac{\epsilon }{2}.$$
Continuing by induction, we obtain $`n_1<n_2<\mathrm{}`$, $`1=\epsilon _1>\epsilon _2>\mathrm{}`$, and projections $`Q_1,Q_2,\mathrm{}`$ in $`𝒫`$ so that for all $`k`$,
(2.57)
$$\tau (Q_k)\frac{\epsilon _k}{2^k}$$
(2.58)
$$\omega (f_{n_k},\epsilon _{k+1})<\frac{\epsilon }{2}$$
and
(2.59)
$$|\tau (f_{n_k}Q_k)|>\frac{\eta \epsilon }{4}.$$
Now set $`P_k=Q_k(_{j>k}(1Q_j))`$, for $`k=1,2,\mathrm{}`$. Evidently the $`P_k`$’s are pairwise orthogonal. For each $`i`$, let $`\stackrel{~}{Q}_i=Q_iP_i`$. Now by subadditivity of $`\tau `$,
$$\begin{array}{cc}\hfill \tau (P_i)& \tau (Q_i)\left(1\tau \underset{j>i}{}(1Q_j)\right)\hfill \\ & \tau (Q_i)\underset{j>i}{}\tau (Q_j).\hfill \end{array}$$
But
$$\begin{array}{cc}\hfill \underset{j>i}{}\tau (Q_j)\underset{j>i}{}\frac{\epsilon _j}{2^j}& <\epsilon _{i+1}\underset{j>i}{}\frac{1}{2^j}\text{ by (2.57)}\hfill \\ & <\epsilon _{i+1}.\hfill \end{array}$$
Hence we have
(2.60)
$$\tau (\stackrel{~}{Q}_i)\underset{j>i}{}\tau (Q_j)<\epsilon _{i+1}.$$
Thus by (2.58),
(2.61)
$$f_{n_i}\stackrel{~}{Q}_i_1\omega (f_{n_i},\epsilon _{i+1})<\frac{\epsilon }{2}.$$
Hence
$$\begin{array}{cc}\hfill |\tau (f_{n_i}P_i)|& =|\tau (f_{n_i}Q_if_{n_i}\stackrel{~}{Q}_i)|\hfill \\ & \frac{\eta \epsilon }{4}\frac{\epsilon }{2}\text{ by (}\text{2.61}\text{)}\hfill \\ & \frac{\eta }{5}.\hfill \end{array}$$
###### Remark.
The proof of the implication (i) $``$ (ii) itself, may quickly be achieved, using instead Theorem 3.5 of \[DSS\].
The following result is an immediate consequence of 2.5.
###### Corollary 2.6.
A subset of $`L^1(\tau )`$ is relatively weakly compact if and only if it is uniformly integrable.
###### Proof.
Let $`W`$ be the subset, and suppose first $`W`$ is relatively weakly compact, yet $`lim_{\epsilon 0}sup_{fW}\omega (f,\epsilon )\stackrel{\text{def}}{=}\eta >0`$. Then for each $`n`$, choose $`f_nW`$ with $`\omega (f_n,\frac{1}{2^n})>\eta \frac{1}{2^n}`$. It follows immediately that also $`lim_{\epsilon 0}sup_n\omega (f_n,\epsilon )=\eta `$, hence $`(f_n)`$ is not relatively weakly compact by Proposition 2.5. If $`W`$ is uniformly integrable, then $`W`$ is bounded, and then $`W`$ is relatively weakly compact by Akemann’s criterion, (stated preceding (2.48)). ∎
###### Remark.
Suppose $`f_i_11`$ for all $`i`$, and $`(f_i)`$ satisfies (2.43). Letting the $`n_1<n_2<\mathrm{}`$ be as in the proof of 2.5, we show in Section 3, using arguments in \[R1\], that there exists a subsequence $`(f_i^{})`$ of $`(f_{n_i})`$ so that $`(f_i^{})`$ is $`\frac{5}{\eta }`$-equivalent to the usual $`\mathrm{}^1`$-basis, with also $`[f_i^{}]`$ $`\frac{5}{\eta }`$-complemented in $`L^1(\tau )`$. Hence $`(f_i)`$ has a subsequence equivalent to the $`\mathrm{}^1`$-basis, so of course $`(f_i)`$ is not relatively weakly compact.
We note finally a consequence of the proof of 2.5, valid for all $`1p<\mathrm{}`$ and arbitrary (not necessarily atomic) finite von Neumann algebras.
###### Corollary 2.7.
Let $`1p<\mathrm{}`$, let $``$ be a finite von Neumann algebra endowed with a faithful normal tracial state $`\tau `$, and let $`W`$ be a bounded subset of $`L^p(\tau )`$. Then the following are equivalent.
* $`\{|w|^p:wW\}`$ is uniformly integrable.
* $`lim_{\epsilon 0}sup_{fW}\stackrel{~}{\omega }_p(f,\epsilon )=0`$.
* $`lim_r\mathrm{}g_W(r)=0`$,
where the function $`g_W`$ is defined by
(2.62)
$$g_W(r)=\underset{wW}{sup}d_{L^p(\tau )}(w,r_a())\text{ for }r>0.$$
###### Proof.
(i) $``$ (ii) follows immediately from the (obvious) inequality $`\stackrel{~}{\omega }_p(f,\epsilon )\omega _p(f,\epsilon )`$ (stated as part of (2.11) in Lemma 2.3).
(ii) $``$ (iii). Assume that $`w_pM`$ for all $`wW`$. For $`r`$ sufficiently large, define $`\epsilon (r)=\epsilon >0`$ by
(2.63)
$$r=\epsilon ^{1/p}M.$$
Let $`fW`$. Since $`\epsilon ^{1/p}f_pr`$, by the final assertion of Lemma 2.3, we may choose $`P`$ a spectral projection for $`|f|`$ so that
(2.64)
$$fPr_a()\text{ and }f(IP)_p\stackrel{~}{\omega }_p(f,\epsilon ).$$
It follows immediately that
(2.65)
$$g_W(r)\underset{fW}{sup}\stackrel{~}{\omega }_p(f,\epsilon ).$$
Thus (iii) holds by (ii), since $`\epsilon (r)0`$ as $`r\mathrm{}`$. (Note also that the final assertion of 2.3 does not involve the “atomless” hypothesis, since $`\stackrel{~}{\omega }_p(f,\epsilon )`$ is defined in terms of the spectral measure for $`|f|`$.)
(iii) $``$ (i). Given $`fW`$ and $`\epsilon >0`$, choose $`\psi r_a()`$ with
(2.66)
$$f\psi _{L^p(\tau )}<\epsilon .$$
Then for any $`\delta <\epsilon `$,
(2.67)
$$\omega _p(f,\delta )\omega _p(f\psi ,\delta )+\omega _p(\psi ,\delta )<\epsilon +r\delta .$$
Hence $`\overline{\mathrm{lim}}_{\delta 0}sup_{fW}\omega _p(f,\delta )\epsilon `$, proving that (i) holds, since $`\epsilon >0`$ is arbitrary and $`\omega _p(f,t)=(\omega (|f|^p,t))^{1/p}`$ for any $`f`$ and $`t`$, by (2.9) of Lemma 2.3. ∎
## 3. Proof of the Main Theorem
We first assemble some preliminary lemmas, perhaps useful in a wider context. $`𝒩`$ and $`\tau `$ are assumed to be as in Section 2. Let $`r_1,r_2,\mathrm{}`$ denote the Rademacher functions on $`[0,1]`$; equivalently, an independent sequence of $`\{1,1\}`$-valued random variables $`(r_j)`$ with $`P(r_j=1)=P(r_j=1)=\frac{1}{2}`$ for all $`j`$.
###### Lemma 3.1.
Let $`1p<2`$ and $`(f_n)`$ be a bounded unconditional basic sequence in $`L^p(\tau )`$, so that $`(|f_i|^p)_{i=1}^{\mathrm{}}`$ is uniformly integrable. Then $`lim_n\mathrm{}n^{1/p}f_1+\mathrm{}+f_n_{L^p(\tau )}=0`$.
###### Remark.
Recall from the introduction that a sequence $`(x_n)`$ in a Banach space is called unconditional if there is a constant $`u`$ so that
(3.1)
$$\begin{array}{cc}& \left\{\underset{i=1}{\overset{n}{}}\alpha _ic_ix_iu\underset{i=1}{\overset{n}{}}c_ix_i\right\}\text{ for all }n\text{ and scalars}\hfill \\ & c_1,\mathrm{},c_n\text{ and }\alpha _1,\mathrm{},\alpha _n\text{ with }|\alpha _i|=1\text{ for all }i.\hfill \end{array}$$
$`(x_n)`$ is called $`u`$-unconditional if (3.1) holds.
###### Proof of 3.1.
Suppose $`(f_n)`$ is $`u`$-unconditional. Then $`(f_n)`$ is $`u`$-equivalent to $`(f_nr_n)`$ in $`L^p(𝒩\overline{}L^{\mathrm{}})`$, so it suffices to prove the same conclusion for $`(f_nr_n)`$ instead. Let $`\beta =\tau m`$, where $`m`$ is Lebesgue measure on $`[0,1)`$. We may also assume without loss of generality that $`f_n_{L^p(\tau )}1`$ for all $`n`$. Now let $`\epsilon >0`$, and choose $`\delta >0`$ so that
(3.2)
$$\omega (|f_n|^p,\delta )\epsilon \text{ for all }n$$
(using that $`(|f_n|^p)`$ is uniformly integrable). By the final statement of Lemma 2.3, we may by (3.2) choose for each $`j`$ a $`P_j𝒫=𝒫(𝒩)`$ so that $`f_jP_j𝒩`$ with
(3.3)
$$f_jP_j_{\mathrm{}}\frac{1}{\delta }\text{ and }f_j(IP_j)_p^p\epsilon .$$
Then fixing $`n`$,
(3.4)
$$\underset{i=1}{\overset{n}{}}f_ir_i_{L^p(\beta )}\underset{i=1}{\overset{n}{}}f_iP_ir_i_{L^p(\beta )}+\underset{i=1}{\overset{n}{}}f_i(IP_i)r_i_{L^p(\beta )}.$$
But
(3.5)
$$\underset{i=1}{\overset{n}{}}f_iP_ir_i_{L^p(\beta )}\underset{i=1}{\overset{n}{}}f_iP_ir_i_{L^2(\beta )}\frac{\sqrt{n}}{\delta }$$
since $`f_iP_i_{\mathrm{}}\frac{1}{\delta }`$ for all $`i`$.
On the other hand, since $`L^p()`$ is type $`p`$ with type $`p`$ constant 1 for any von-Neumann algebra $``$,
(3.6)
$$\begin{array}{cc}\hfill \underset{i=1}{\overset{n}{}}f_i(IP_i)r_i_{L^p(\beta )}& \left(\underset{i=1}{\overset{n}{}}f_i(IP_i)_{L^p(\tau )}^p\right)^{1/p}\hfill \\ & \epsilon n^{1/p}\text{ by (}\text{3.3}\text{).}\hfill \end{array}$$
(This fact follows by Clarkson’s inequalities — see the discussion in the proof of the next lemma.) We thus have that
(3.7)
$$\underset{n\mathrm{}}{\overline{\mathrm{lim}}}n^{1/p}\underset{i=1}{\overset{n}{}}f_ir_i_{L^p(\beta )}\underset{n\mathrm{}}{lim}\frac{n^{1/2}}{\delta n^{1/p}}+\epsilon =\epsilon $$
by (3.5) and (3.6). Since $`\epsilon >0`$ is arbitrary, the conclusion of the lemma follows. ∎
###### Remarks.
1. It follows easily from the above proof that in fact if $`(f_n)`$ satisfies the hypothesis of 3.1, then $`lim_n\mathrm{}n^{1/p}f_1^{}+\mathrm{}+f_n^{}_p=0`$ uniformly over all subsequences $`(f_n^{})`$ of $`f_n`$.
2. The proof of Lemma 3.1 yields the following quantitative result. Fix $`\epsilon >0`$, and let $`(f_j)`$ be a bounded sequence in $`L^p(\tau )`$ so that there exists an $`r<\mathrm{}`$ with $`d_{L^p(\tau )}(f_j,r_a𝒩)<\epsilon `$ for all $`j`$. Then $`\overline{\mathrm{lim}}_n\mathrm{}𝔼_\omega n^{1/p}_{j=1}^nr_j(w)f_j_{L^p(\tau )}\epsilon `$. Indeed, for each $`j`$, choose $`\phi _jr_a𝒩`$ with $`f_j\phi _j_{L^p(\tau )}\epsilon `$. Then fixing $`n`$, (3.4)–(3.6) yield
$$\begin{array}{cc}\hfill \underset{i=1}{\overset{n}{}}f_ir_i_{L^p(\beta )}& \underset{i=1}{\overset{n}{}}\phi _ir_i_{L^p(\beta )}+\underset{i=1}{\overset{n}{}}(f_i\phi _i)r_i_{L^p(\beta )}\hfill \\ & r\sqrt{n}+\epsilon n^{1/p}.\hfill \end{array}$$
Hence $`\overline{\mathrm{lim}}_n\mathrm{}n^{1/p}_{i=1}^nf_ir_i_{L^p(\beta )}\epsilon `$ as desired.∎
We next give a criterion for a finite or infinite sequence in $`L^p(\tau )`$ to be equivalent to the usual $`\mathrm{}^p`$ basis.
###### Lemma 3.2.
Let $`u1`$, $`\delta >0`$, $`1p<2`$, and $`f_1,\mathrm{},f_n`$ elements of $`_a(L^p(𝒩))`$ be given so that $`(f_i)_{i=1}^n`$ is $`u`$-unconditional. Assume there exist pairwise orthogonal projections $`P_1,\mathrm{},P_n`$ in $`𝒫`$ so that
(3.8)
$$\tau (|P_jf_jP_j|^p)\delta ^p\text{ for all }1jn.$$
Then $`(f_i)_{i=1}^n`$ is $`C`$-equivalent to the usual $`\mathrm{}_n^p`$ basis, where $`C=u\sqrt{3}\delta ^1`$.
###### Proof.
We first note that (using interpolation), $`L^p(\tau )`$ satisfies Clarkson’s inequalities: for all $`x,yL^p(\tau )`$,
(3.9)
$$x+y_p^p+xy_p^p2(x_p^p+y_p^p).$$
It follows immediately by induction on $`n`$ that $`L^p(\tau )`$ is type $`p`$ with constant one; that is, for any $`x_1,\mathrm{},x_n`$ in $`L^p(\tau )`$,
(3.10)
$$\begin{array}{cc}\hfill \underset{A\pm }{}\pm x_1\pm \mathrm{}\pm x_n_p^p& =_0^1\underset{i=1}{\overset{n}{}}r_i(\omega )x_i_p^p𝑑\omega \hfill \\ & \left(\underset{i=1}{\overset{n}{}}x_i_p^p\right).\hfill \end{array}$$
Now let scalars $`a_1,\mathrm{},a_n`$ be given, and let $`f=_{i=1}^na_if_i`$. We obtain from (3.10) that since $`(f_i)`$ is $`u`$-unconditional,
(3.11)
$$f_pu\left(\underset{i=1}{\overset{n}{}}|a_i|^p\right)^{1/p}.$$
Now fix $`\omega `$ and set $`f_\omega =_{i=1}^na_ir_i(\omega )f_i`$. Then
(3.12)
$$f_\omega _p^p\underset{j=1}{\overset{n}{}}P_jf_\omega P_j_p^p.$$
Thus integrating over $`\omega `$ and again using unconditionality,
(3.13)
$$\begin{array}{cc}\hfill f_p^p& \frac{1}{u^p}_0^1f_\omega _p^p𝑑\omega \hfill \\ & \frac{1}{u^p}\underset{j=1}{\overset{n}{}}_0^1P_jf_\omega P_j_p^p𝑑\omega \text{ by (}\text{3.12}\text{).}\hfill \end{array}$$
But fixing $`j`$, since $`L^p(\tau )`$ is cotype 2 with constant at most $`3^{1/2}`$,
(3.14)
$$\begin{array}{cc}\hfill _0^1P_jf_\omega P_j_p^p𝑑\omega & \frac{1}{3^{p/2}}\left(\underset{i}{}P_ja_if_iP_j_p^2\right)^{p/2}\hfill \\ & \frac{1}{3^{p/2}}P_ja_jf_jP_j_p^p\hfill \\ & \frac{1}{3^{p/2}}|a_j|^p\delta ^p\text{ by (}\text{3.8}\text{).}\hfill \end{array}$$
Thus in view of (3.13),
(3.15)
$$f_p^p\frac{\delta ^p}{u^p3^{p/2}}\left(\underset{j=1}{\overset{n}{}}|a_j|^p\right),$$
so (3.11) and (3.15) now imply the conclusion of Lemma 3.2. ∎
Our last preliminary result yields an estimate for equivalence to the $`\mathrm{}_n^p`$ basis in terms of $`p`$-moduli.
###### Lemma 3.3.
Let $`0<\epsilon <\eta /2`$, $`n1`$, and $`f_1,\mathrm{},f_n_aL^p(\tau )`$ be such that $`(f_1,\mathrm{},f_n)`$ is $`u`$-unconditional and there are $`\delta _1\delta _2\mathrm{}\delta _n>0`$ so that for all $`1jn`$ and all $`k`$ with $`j<k`$ (if $`j<n`$)
(3.16)
$$\omega _p(f_j,\delta _j)>\eta \text{ and }\omega _p(f_j,\delta _k+\delta _{k+1}+\mathrm{}+\delta _n)<\frac{\epsilon }{2}.$$
Then $`(f_1,\mathrm{},f_n)`$ is $`C`$-equivalent to the $`\mathrm{}_n^p`$ basis where
$$Cu\sqrt{3}\left(\frac{\eta }{2}\epsilon \right)^1.$$
###### Proof.
By Lemma 2.3, (see (2.9)), we have, fixing $`1jn`$, that
(3.17)
$$\omega _p^s(f_j,\delta _j)>\frac{\eta }{2}.$$
Hence we may choose $`Q_j𝒫`$ with
(3.18)
$$Q_jf_jQ_j_p>\frac{\eta }{2}\text{ and }\tau (Q_j)\delta _j.$$
Define projections $`P_j`$ and $`\stackrel{~}{Q}_j`$ by
(3.19)
$$P_j=Q_j\underset{k>j}{}(1Q_k)\text{ and }\stackrel{~}{Q}_j=Q_jP_j.$$
Then
(3.20)
$$Q_jf_jQ_j=P_jf_jP_j+\stackrel{~}{Q}_jf_jP_j+Q_jf_j\stackrel{~}{Q}_j.$$
Now we have by subadditivity of $`\tau `$ that $`\tau (_{k>j}(1Q_k))1_{k>j}\delta _k`$, and so again by subadditivity,
$$\begin{array}{cc}\hfill \tau (P_j)& \tau (Q_j)\left(1\tau \left(\underset{k>j}{}1Q_k\right)\right)\hfill \\ & \tau (Q_j)\underset{k>j}{}\delta _k.\hfill \end{array}$$
Thus $`\tau (\stackrel{~}{Q}_j)<_{k>j}\delta _k`$. Hence we have
(3.21)
$$\begin{array}{cc}\hfill \stackrel{~}{Q}_jf_jP_j_p\stackrel{~}{Q}_jf_j_p& \omega _p(f_j^{},\underset{k>j}{}\delta _k)\hfill \\ & =\omega _p(f_j,\underset{k>j}{}\delta _k)\frac{\epsilon }{2}\text{ by (}\text{3.16}\text{)).}\hfill \end{array}$$
By the same argument,
(3.22)
$$Q_jf_j\stackrel{~}{Q}_j_p\frac{\epsilon }{2}.$$
Thus from (3.18), (3.20), (3.21) and (3.22), we obtain
(3.23)
$$P_jf_jP_j_p\frac{\eta }{2}\epsilon .$$
Of course $`P_1,\mathrm{},P_n`$ are pairwise orthogonal; hence Lemma 3.2 now immediately yields the conclusion of 3.3. ∎
Lemma 3.3 immediately yields an infinite dimensional conclusion as well. Combining this and Lemma 3.1 we obtain the following definitive result.
###### Corollary 3.4.
Let $`(f_n)`$ be a bounded unconditional sequence in $`L^p(\tau )`$, $`1p<2`$. The following are equivalent:
* $`(f_n)`$ has a subsequence equivalent to the usual $`\mathrm{}^p`$ basis.
* $`(|f_n|^p)`$ is not uniformly integrable.
###### Proof.
(a) $``$ (b) follows immediately from Lemma 3.1. Assume that (b) holds and also assume without loss of generality that $`f_n_p1`$ for all $`n`$. Then by Lemma 3.1,
(3.24)
$$\eta \stackrel{\text{def}}{=}\underset{\epsilon 0}{lim}\underset{n}{sup}\omega _p(f_n,\epsilon )>0.$$
Now Lemma 3.3 yields that there is a subsequence $`(f_n^{})`$ of $`(f_n)`$ so that
(3.25)
$$(f_n^{})\text{ is }\frac{cu}{\eta }\text{-equivalent to the }\mathrm{}^p\text{ basis,}$$
where $`c`$ is an absolute constant.
Indeed, fix $`0<\epsilon <\frac{\eta }{2}`$. Choose $`\delta _11`$ and $`n_1`$ so that
(3.26)
$$\omega _p(f_{n_1},\delta _1)>\eta \epsilon .$$
Suppose $`n_1<\mathrm{}<n_j`$ and $`\delta _1>\delta _2\mathrm{}>\delta _j`$ chosen so that
$$\omega _p(f_{n_i},\delta _{i+1}+\mathrm{}+\delta _j)<\frac{\epsilon }{2}\text{ for all }1i<j.$$
By continuity of the functions $`t\omega _p(f_{n_i},t)`$ for $`i<j`$ and the fact that $`f_{n_j}L^p(\tau )`$, choose $`\overline{\delta }_{j+1}<\delta _j`$ so that
(3.27)
$$\omega _p(f_{n_i},\delta _{i+1}+\mathrm{}+\delta _j+\overline{\delta }_{j+1})<\frac{\epsilon }{2}\text{ for all }1ij.$$
Then choose $`\delta _{j+1}\overline{\delta }_{j+1}`$ and $`n_{j+1}>n_j`$ so that
(3.28)
$$\omega _p(f_{n_{j+1}},\delta _{j+1})>\eta \epsilon .$$
This completes the inductive choice of $`n_1<n_2<\mathrm{}`$.
Setting $`f_k^{}=f_{n_k}`$, then $`(f_1^{},\mathrm{},f_n^{})`$ satisfies the hypotheses of Lemma 3.3 for all $`n`$, and hence $`(f_n^{})`$ is $`u\sqrt{3}(\frac{\eta }{2}\epsilon )^1`$-equivalent to the $`\mathrm{}^p`$ basis by 3.3. By taking $`\epsilon `$ small enough, we obtain $`c7`$ in (3.25). ∎
###### Remark.
The hypothesis that $`(f_n)`$ is unconditional may be omitted when $`p=1`$, as pointed out in the remark following the proof of Proposition 2.5. Also, it’s not hard to show that the sequence $`(f_n^{})`$ constructed above has its closed linear span complemented in $`L^p(\tau )`$. Finally, it follows from known (rather non-trivial) results that if $`1<p<\mathrm{}`$ and $`𝒩`$ is hyperfinite, then every semi-normalized weakly null sequence in $`L^p(𝒩)`$ has an unconditional subsequence. Indeed, assuming (as we may) that $`𝒩`$ acts on a separable Hilbert space, $`L^p(𝒩)`$ has an unconditional finite dimensional decomposition (see \[SF\], \[PX1\]), which yields the above statement. Thus also in the hyperfinite case, the hypothesis that $`(f_n)`$ is unconditional may be omitted. We do not know, however, if this is so for general $`𝒩`$.
###### Corollary 3.5.
Let $`(f_n)`$ be a bounded unconditional sequence in $`L^p(\tau )`$, $`1p<2`$. The following are equivalent.
* For every subsequence $`(f_n^{})`$ of $`(f_n)`$
$$\underset{n\mathrm{}}{lim}n^{1/p}\underset{i=1}{\overset{n}{}}f_i^{}_{L^p(\tau )}=0.$$
* $`(|f_n|^p)`$ is uniformly integrable.
###### Proof.
(a) $``$ (b): Assume (b) is false. Then by Corollary 3.4 there exists a subsequence $`(f_n^{})`$ equivalent to the usual $`\mathrm{}^p`$-basis. In particular
$$\underset{n\mathrm{}}{lim\; inf}n^{1/p}\underset{i=1}{\overset{n}{}}f_i^{}_{L^p(\tau )}>0.$$
which contradicts (a).
(b) $``$ (a). This follows from Lemma 3.1, since condition (b) implies that $`(|f_n^{}|)^p`$ is uniformly integrable for any subsequence $`(f_n^{})`$ of $`(f_n)`$. ∎
We now turn to the proof of the Main Theorem. First we give some preliminary results concerning ultrapowers of Banach spaces and the standard construction of the ultrapower of a finite von Neumann algebra (cf. \[McD\], \[V\]).
Fix $`U`$ a free ultrafilter on $``$. For a given Banach space $`X`$, let $`\mathrm{}^{\mathrm{}}(X)`$ denote the set of bounded sequences in $`X`$, under the norm $`(x_n)=sup_nx_n`$, and set
(3.29)
$$E_U=\{(x_n)\mathrm{}^{\mathrm{}}(X):\underset{nU}{lim}x_n=0\}.$$
Then $`X_U`$, the ultrapower of $`X`$ with respect to $`U`$, is given by
(3.30)
$$X_U=\mathrm{}^{\mathrm{}}(X)/E_U.$$
Now fix $`𝒩`$ a finite von Neumann algebra with a normal faithful tracial state $`\tau `$, and define $`I_U`$ by
(3.31)
$$I_U=\{(x_n)\mathrm{}^{\mathrm{}}(𝒩):\underset{nU}{lim}\tau (x_n^{}x_n)=0\}.$$
Then $`I_U`$ is a norm-closed two-sided ideal in $`\mathrm{}^{\mathrm{}}(X)`$; we define $`𝒩^U`$ (a different object than $`𝒩_U`$!) by
(3.32)
$$𝒩^U=\mathrm{}^{\mathrm{}}(𝒩)/I_U.$$
Then by the references cited above, $`𝒩^U`$ is a $`W^{}`$-algebra (i.e., an abstract von Neumann algebra) with a normal faithful tracial state $`\tau _U`$ given by
(3.33)
$$\tau _U(\pi (x_n))=\underset{nU}{lim}\tau (x_n)$$
where $`\pi :\mathrm{}^{\mathrm{}}(𝒩)𝒩^U`$ is the quotient map.
The next result yields that $`L^p(𝒩^U)`$ may be regarded as a subspace of the Banach space ultrapower $`L^p(𝒩)^U`$.
###### Lemma 3.6.
Let $`1p<\mathrm{}`$ and let $`Y_p`$ denote the closure of $`\mathrm{}^{\mathrm{}}(𝒩)`$ in the Banach space $`\mathrm{}^{\mathrm{}}(L^p(𝒩))`$. Then $`\pi `$ has a unique extension to a bounded linear map $`\stackrel{~}{\pi }:Y_pL^p(𝒩^U)`$. Moreover, for $`(x_n)Y_p`$,
(3.34)
$$\stackrel{~}{\pi }((x_n))_{L^p(\tau _U)}=\underset{nU}{lim}x_n_{L^p(\tau )}.$$
Fixing $`p`$ as in 3.6 and letting $`\rho :\mathrm{}^{\mathrm{}}(L^p(𝒩))L^p(𝒩)^U`$ be the quotient map, Lemma 3.6 yields there is a unique isometric embedding $`i:L^p(𝒩^U)L^p(𝒩)^U`$ so that the following diagram commutes:
(3.35) .
###### Proof.
Since $`\pi `$ is a $``$-homomorphism of $`\mathrm{}^{\mathrm{}}(𝒩)`$ onto $`𝒩^U`$, we have for any continuous function $`f:[0,\mathrm{})`$ and any $`x=(x_n)\mathrm{}^{\mathrm{}}(𝒩)`$,
(3.36)
$$\pi \left((f(x_n^{}x_n))_{n=1}^{\mathrm{}}\right)=f(\pi (x^{})\pi (x)).$$
Applying this to $`f(t)=|t|^{p/2}`$, we get by the trace formula (3.33) that
(3.37)
$$\pi (x)_{L^p(\tau _U)}=\underset{nU}{lim}x_n_{L^p(\tau )}.$$
In particular,
(3.38)
$$\begin{array}{cc}\hfill \pi (x)_{L^p(\tau _U)}& \underset{n}{sup}x_n_{L^p(\tau )}\hfill \\ & =x_{\mathrm{}^{\mathrm{}}(L^p(𝒩))}.\hfill \end{array}$$
Thus $`\pi `$ extends by continuity to a contraction $`\stackrel{~}{\pi }:Y_pL^p(𝒩^U)`$. Now let $`x=(x_n)`$ belong to $`Y_p`$, and let $`\epsilon >0`$. Then choose $`y=(y_n)`$ in $`\mathrm{}^{\mathrm{}}(𝒩)`$ so that
(3.39)
$$xy_{\mathrm{}^{\mathrm{}}(L^p(𝒩))}<\epsilon .$$
It follows from (3.39) that
(3.40)
$$\left|\pi (x)_{L^p(\tau _U)}\pi (y)_{L^p(\tau _U)}\right|<\epsilon $$
and
(3.41)
$$\left|\underset{nU}{lim}x_n_{L^p(\tau )}\underset{nU}{lim}y_n_{L^p(\tau )}\right|<\epsilon .$$
Since (3.37) holds, replacing “$`x`$” by “$`y`$” in its statement, we have from (3.40) and (3.41) that
(3.42)
$$\left|\pi (x)_{L^p(\tau _U)}\underset{nU}{lim}x_n_{L^p(\tau )}\right|<2\epsilon .$$
Since $`\epsilon >0`$ is arbitrary, (3.34) holds for all $`x=(x_n)`$ in $`Y_p`$. ∎
###### Lemma 3.7.
Let $`1p<2`$, and let $`(x_{ij})`$ be an infinite matrix in $`L^p(𝒩)`$ so that for some $`C1`$, each row and each column of $`(x_{ij})`$ is $`C`$-equivalent to the usual $`\mathrm{}^2`$-basis. Then for every free ultrafilter $`U`$ on $``$
(3.43)
$$\underset{j}{sup}\underset{iU}{lim}d_{L^p(\tau )}(x_{ij},r_a(𝒩))0\text{ as }r\mathrm{}$$
###### Proof.
Define for each $`j`$ a function $`g_j:^+^+`$ by
$$g_j(r)=\underset{i}{sup}d_{L^p(\tau )}(x_{ij},r_a(𝒩)).$$
For fixed $`j`$, $`(x_{ij})_{i=1}^{\mathrm{}}`$ is $`C`$-equivalent to the usual $`\mathrm{}^2`$-basis, so by Corollary 3.4 and Corollary 2.7, $`(|x_{ij}|^p)_{i=1}^{\mathrm{}}`$ is uniformly integrable and
(3.44)
$$\underset{r\mathrm{}}{lim}g_j(r)=0.$$
Now (3.44) implies that $`(x_{ij})_{i=1}^{\mathrm{}}`$ belongs to $`Y_p`$. Let $`\stackrel{~}{\pi }`$ be as in the statement of Lemma 3.6 and define $`x_j`$ by
$$x_j=\stackrel{~}{\pi }\left((x_{ij})_{i=1}^{\mathrm{}}\right)L^p(𝒩^U).$$
Now we claim that
(3.45)
$$(x_j)\text{ is }C\text{-equivalent to the }\mathrm{}^2\text{-basis.}$$
Indeed, using the hypotheses of Theorem 1.1 and Lemma 3.6, we have for any $`n`$ and scalars $`c_1,\mathrm{},c_n`$, that
$$\begin{array}{cc}\hfill \underset{j=1}{\overset{n}{}}c_jx_j_{L^p(\tau _U)}& =\stackrel{~}{\pi }\left(\left(\underset{j=1}{\overset{n}{}}c_jx_{ij}\right)_{i=1}^{\mathrm{}}\right)_{L^p(\tau _U)}\hfill \\ & =\underset{iU}{lim}\underset{j=1}{\overset{n}{}}c_jx_{ij}_{L^p(\tau )}\text{ by (}\text{3.34}\text{)}\hfill \\ & \stackrel{C}{}\left(|c_j|^2\right)^{1/2}.\hfill \end{array}$$
Now define $`g:^+^+`$ by
$$g(r)=\underset{j}{sup}d_{L^p(\tau _U)}(x_j,r_a(𝒩^U)).$$
Again by (3.45) and Corollary 3.4, $`(|x_j|^p)_{j=1}^{\mathrm{}}`$ is uniformly integrable in $`L^p(\tau _U)`$, so by Corollary 2.7 we have that
(3.46)
$$\underset{r\mathrm{}}{lim}g(r)=0.$$
Now let $`\epsilon >0`$. Since $`\pi `$ is a quotient map of $`\mathrm{}^{\mathrm{}}(𝒩)`$ onto $`𝒩^U`$, it follows that fixing $`j`$, there exists for every $`r>0`$, $`(y_{ij})_{i=1}^{\mathrm{}}r_a(𝒩)`$ so that
$$x_j\pi ((y_{ij})_{i=1}^{\mathrm{}})_{L^p(\tau _U)}<g(r)+\epsilon .$$
Hence by Lemma 3.6,
$$\underset{iU}{lim}x_{ij}y_{ij}_{L^p(\tau )}<g(r)+\epsilon ,$$
which implies that
$$\underset{iU}{lim}d_{L^p(\tau )}(x_{ij},r_a(𝒩))<g(r)+\epsilon .$$
Hence by (3.46)
$$\underset{r\mathrm{}}{lim\; sup}\left(\underset{j}{sup}\underset{iU}{lim}d_{L^p(\tau )}(x_{ij},r_a(𝒩))\right)\epsilon .$$
Since $`\epsilon >0`$ was arbitrary, we get (3.43). ∎
###### Proof of Theorem 1.1.
Let $`1p<2`$, and let $`(x_{ij})`$ be as in Theorem 1.1, and let $`U`$ be a free ultrafilter on $``$. Put
(3.47)
$$h(r)=\underset{j}{sup}\underset{iU}{lim}d_{L^p(\tau )}(x_{ij},r_a(𝒩)),r_+.$$
Then $`h:^+^+`$ is a decreasing function and by (3.43)
(3.48)
$$\underset{r\mathrm{}}{lim}h(r)=0.$$
We claim that (3.47) and (3.48) imply that for a suitable choice of natural numbers $`i_1<i_2<\mathrm{}`$ one has
(3.49)
$$(|x_{i_j,j}|^p)_{j=1}^{\mathrm{}}\text{ is uniformly integrable.}$$
To prove (3.49) put for $`j`$
(3.50)
$$G_j=\underset{r=1}{\overset{j}{}}G_{j,r}$$
where for $`j,r`$,
(3.51)
$$G_{j,r}=\{id_{L^p(\tau )}(x_{ij},r_a(𝒩))<h(r)+\frac{1}{r}\}.$$
By (3.47) each $`G_{j,r}U`$, and hence also $`G_jU`$ for all $`j`$. Since $`U`$ is a free ultrafilter, each $`G_j`$ is infinite, so we can choose successively $`i_1<i_2<\mathrm{}`$ such that $`i_jG_j`$ for all $`j`$. Put $`y_j=x_{i_j,j}`$, $`j`$ and $`W=\{y_j,j\}`$, and put as in Corollary 2.7
(3.52)
$$g_W(r)=\underset{j}{sup}d_{L^p(\tau )}(y_j,r_a(𝒩)),r^+.$$
To prove (3.49) we just have to show that $`g_W(r)0`$ when $`r\mathrm{}`$ (cf. Corollary 2.7). Let $`\epsilon >0`$. By (3.48) we can choose $`r_0`$ such that
(3.53)
$$h(r_0)+\frac{1}{r_0}<\epsilon .$$
When $`jr_0`$, $`i_jG_jG_{j,r_0}`$. Hence by (3.51) and (3.53)
(3.54)
$$d_{L^p(\tau )}(y_j,r_0_a(𝒩))<\epsilon ,jr_0.$$
Since $`𝒩=_{r>0}r_a(𝒩)`$ is dense in $`L^p(\tau )`$ we have for every $`j`$,
$$\underset{r\mathrm{}}{lim}d_{L^p(\tau )}(y_j,r_a(𝒩))=0.$$
Hence, we may choose $`r_1r_0`$, such that
(3.55)
$$d_{L^p(\tau )}(y_j,r_1_a(𝒩))<\epsilon ,j=1,\mathrm{},r_01.$$
By (3.54) and (3.55), $`g_W(r)<\epsilon `$ for all $`rr_1`$. This shows that $`lim_r\mathrm{}g_W(r)=0`$ and hence by Corollary 2.7, $`(|y_j|^p)_{j=1}^{\mathrm{}}`$ is uniformly integrable, i.e., (3.49) holds. Thus by the assumption that $`(y_j)`$ is unconditional, Corollary 3.5 yields that for any subsequence $`(y_j^{})`$ of $`(y_j)`$,
(3.56)
$$\underset{n\mathrm{}}{lim}n^{1/p}\underset{j=1}{\overset{n}{}}y_j^{}_{L^p(\tau )}=0.$$
Putting now $`j_k=k`$, we have $`y_k=x_{i_k,j_k}`$ and Theorem 1.1 follows. ∎
## 4. Improvements to the Main Theorem
We obtain here results that are stronger than the Main Theorem. In particular, Theorem 4.1 is also needed in Section 6 (specifically, for the proof of Theorem 6.9). The arguments in this section do not use the ultraproduct construction and technique of Section 3. They are in a sense more elementary, and also more delicate, than those of Section 3.
We use the following terminology: given a matrix $`(x_{ij})`$, a sequence $`(x_{i_k,j_k})`$ of elements of the matrix is called a generalized diagonal if $`i_1<i_2<\mathrm{}`$ and $`j_1<j_2<\mathrm{}`$. A set $`W`$ (or matrix $`(x_{ij})`$) in a Banach space is called semi-normalized if there are $`0<\delta K<\mathrm{}`$ with $`\delta wK`$ for all $`wW`$. The main theorem follows also immediately from the following result.
###### Theorem 4.1.
Let $`𝒩`$ be a finite von-Neumann algebra, $`1p<2`$, and $`(x_{ij})`$ be an infinite semi-normalized matrix in $`L^p(𝒩)`$. Assume that every column and generalized diagonal is unconditional, and there is a $`u1`$ so that every row is $`u`$-unconditional. Then one of the following three alternatives holds.
* Some column has a subsequence equivalent to the usual $`\mathrm{}^p`$ basis.
* There is a $`C1`$ so that for all $`n`$, there exists a row which contains $`n`$ elements $`C`$-equivalent to the usual $`\mathrm{}_n^p`$ basis.
* There is a generalized diagonal $`(y_k)`$ so that
$$n^{1/p}\underset{i=1}{\overset{n}{}}y_i^{}_p0\text{ as }n\mathrm{}$$
for all subsequences $`(y_i^{})`$ of $`(y_i)`$.
To recover the Main Theorem from 4.1, let $`(x_{ij})`$ be as in the hypotheses of the Main Theorem, and simply note that Cases I and II of 4.1 are impossible, since otherwise one would obtain a constant $`\lambda `$ so that the $`\mathrm{}_n^p`$ and $`\mathrm{}_n^2`$ bases are $`\lambda `$-equivalent for all $`n`$. Case III now yields the conclusion of the Main Theorem.
###### Remark.
Let us say that the rows of $`(x_{ij})`$ contain $`\mathrm{}_n^p`$-sequences if condition II of 4.1 holds, with a similar definition for the columns. Since obviously we can interchange rows and columns in the statement of 4.1, we then obtain the following immediate consequence: Let $`𝒩`$, $`p`$ and $`(x_{ij})`$ be as in the first sentence of Theorem 4.1. Assume that every generalized diagonal is unconditional and there is a $`u1`$ so that every row and column are $`u`$-unconditional. Then one of the following holds.
* Some column or some row has a subsequence equivalent to the usual $`\mathrm{}^p`$ basis.
* Both the rows and the columns contain $`\mathrm{}_n^p`$-sequences.
* Condition III of 4.1 holds.
## Proof of Theorem 4.1
We may assume without loss of generality that $`x_{ij}_p1`$ for all $`i`$ and $`j`$. We introduce the following notation, for all $`\epsilon >0`$ and all $`i,j=1,2,\mathrm{}`$.
(4.1) $`\omega _{ij}(\epsilon )`$ $`=`$ $`\omega _p(x_{ij},\epsilon )`$
(4.2) $`\omega _j(\epsilon )`$ $`=`$ $`\underset{i}{sup}\omega _{ij}(\epsilon ).`$
Now assume that Case I of Theorem 4.1 does not occur. We then have by Corollary 3.4 (and Lemma 2.3) that $`(|x_{ij}|^p)_{i=1}^{\mathrm{}}`$ is uniformly integrable for all $`j`$, and hence
(4.3)
$$\underset{\epsilon 0}{lim}\omega _j(\epsilon )=0\text{ for all }j.$$
We now use the following (hopefully intuitive) convention. A set of rows $``$ of $`(x_{ij})`$ is identified with a set $`𝒥\{1,2,\mathrm{}\}`$ via $`=\{R_i:i𝒥\}`$ where $`R_i=\{x_{ij}:j=1,2,\mathrm{}\}`$ for all $`i𝒥`$. Columns are just identified with $`j`$; i.e., $`jC_j=\{x_{ij}:i=1,2,\mathrm{}\}`$.
### Case II
There is an $`\eta >0`$ and an infinite set of rows $`𝒥`$ so that for all further infinite sets of rows $`𝒥^{}𝒥`$, all $`\delta >0`$, and all columns $`j_0`$, there is a column $`j>j_0`$ so that
(4.4)
$$\{i𝒥^{}:\omega _{i,j}(\delta )>\eta \}\text{ is infinite.}$$
Intuitively, the final statement means that looking down the $`j^{th}`$ column of the submatrix with rows $`𝒥^{}`$, then infinitely many of the moduli $`\omega _{i,j}(\delta )`$ are bigger than $`\eta `$.
We shall show that Case II yields II of Theorem 4.1. In fact, we shall show that then, via Lemma 3.3,
(4.5)
$$\{\begin{array}{cc}& \text{for every }n\text{, there exists a row }R_i\text{ and elements }x_{ij_1},\mathrm{},x_{ij_n}\text{ in}\hfill \\ & R_i\text{}j_1<\mathrm{}<j_n\text{, with }(x_{ij_k})_{k=1}^n\text{ }\frac{7u}{\eta }\text{-equivalent to the }\mathrm{}_n^p\text{ basis.}\hfill \end{array}$$
Let $`𝒥_0`$ be the initial set of rows satisfying Case II. Let $`\delta _1=1/2`$, and choose $`j_1`$ so that
(4.6)
$$𝒥_1\stackrel{\text{def}}{=}\{i𝒥_0:\omega _{ij_1}(\delta _1)>\eta \}\text{ if infinite.}$$
Next, using (4.3), choose $`\overline{\delta }_2<\delta _1`$ so that
(4.7)
$$\omega _{j_1}(\overline{\delta }_2)<\frac{\epsilon }{2},$$
and choose $`\delta _2<\overline{\delta }_2`$. Now using the assumptions of Case II, choose $`j_2>j_1`$ so that
(4.8)
$$𝒥_2\stackrel{\text{def}}{=}\{i𝒥_1:\omega _{ij_2}(\delta _2)>\eta \}\text{ is infinite.}$$
For the general inductive step, suppose $`n>1`$, infinite $`𝒥_1\mathrm{}𝒥_{n1}`$ and $`j_1<\mathrm{}<j_{n1}`$, $`\delta _1>\overline{\delta }_2>\delta _2>\mathrm{}>\overline{\delta }_{n1}>\delta _{n1}>0`$ have been chosen so that for all $`1\mathrm{}<n1`$, $`\omega _j_{\mathrm{}}(\overline{\delta }_{\mathrm{}+1})<\frac{\epsilon }{2}`$ and $`\delta _{\mathrm{}+1}+\mathrm{}+\delta _{n1}<\overline{\delta }_{\mathrm{}+1}`$. Using (4.3), choose $`0<\overline{\delta }_n<\delta _{n1}`$ so that $`\omega _{j_{n1}}(\overline{\delta }_n)<\frac{\epsilon }{2}`$; then choose $`0<\delta _n<\overline{\delta }_n`$ so that also $`\delta _{\mathrm{}+1}+\mathrm{}+\delta _n<\overline{\delta }_{\mathrm{}+1}`$ for all $`1\mathrm{}<n1`$. We thus have that
(4.9)
$$\omega _j_{\mathrm{}}(\delta _{\mathrm{}+1}+\mathrm{}+\delta _n)<\frac{\epsilon }{2}\text{ for all }1\mathrm{}n1.$$
Then choose $`j_n>j_{n1}`$ so that
(4.10)
$$𝒥_n\stackrel{\text{def}}{=}\{i𝒥_{n1}:\omega _{ij_n}(\delta _n)>\eta \}\text{ is infinite.}$$
This completes the inductive construction. Now fix $`n`$, let $`i𝒥_n`$, and let $`f_k=x_{ij_k}`$ for $`1kn`$. Then $`(f_1,\mathrm{},f_n)`$ satisfies the assumption of Lemma 3.3. Indeed, the $`f_i`$’s are $`u`$-unconditional by hypothesis, and for each $`k`$, $`1kn`$
(4.11)
$$\omega _{ij_k}(\delta _k)=\omega _p(f_k,\delta _k)>\eta $$
and
(4.12)
$$\omega _p(f_k,\delta _m+\delta _{m+1}+\mathrm{}+\delta _n)\omega _{j_k}(\delta _m+\delta _{m+1}+\mathrm{}+\delta _n)<\frac{\epsilon }{2}\text{ for }k<mn.$$
Thus $`(x_{ij_k})_{k=1}^n`$ satisfies the conclusion of (4.5) in view of Lemma 3.3, proving Case II of 4.1 holds.
We now suppose that Case II does not hold, i.e., we have
### Case III
For all $`\eta >0`$ and infinite sets of rows $`𝒥`$, there exists an infinite set of rows $`𝒥^{}𝒥`$, a $`\delta >0`$, and a column $`𝐣`$ so that for all columns $`j𝐣`$,
(4.13)
$$\omega _{i^{}j}(\delta )\eta \text{ for all but finitely many }i^{}𝒥^{}.$$
(Note that we get $`j𝐣`$ instead of $`j>𝐣`$ by just replacing $`𝐣`$ by $`𝐣+1`$).
Intuitively, the final statement means that now, looking down the $`j^{th}`$ column of the submatrix with rows $`𝒥^{}`$, then all but finitely many of the moduli $`\omega _{i^{},j}(\delta )`$ are no bigger than $`\eta `$.
We shall now construct $`i_1<i_2<\mathrm{}`$ and $`j_1<j_2<\mathrm{}`$ so that
(4.14)
$$\underset{\epsilon 0}{lim}\underset{k}{sup}\omega _{i_k,j_k}(\epsilon )=0.$$
Thus we obtain that $`(|x_{i_kj_k}|^p)_{k=1}^{\mathrm{}}`$ is uniformly integrable, and hence Case III of Theorem 4.1 holds by Corollary 3.5.
We first claim that we may choose infinite sets of rows $`𝒥_1𝒥_2\mathrm{}`$, columns $`j_1<j_2<\mathrm{}`$, and numbers $`1\delta _1`$, $`\frac{1}{2}\delta _2`$, $`\frac{1}{3}\delta _3\mathrm{}`$ so that for all $`k`$,
(4.15)
$$\text{ for all }jj_k\text{}\omega _{ij}(\delta _k)\frac{1}{2^k}\text{ for all but finitely many }i𝒥_k.$$
Indeed, first choose $`𝒥_1`$ an infinite set of rows, $`j_1`$ and $`\delta _1>0`$ so that for all $`jj_1`$, (4.13) holds, where $`𝒥^{}=𝒥`$, $`\eta =1/2`$, and $`\delta _1=\delta `$.
Now suppose $`𝒥_k`$, $`j_k`$, and $`\delta _k`$ have been chosen. Setting $`\eta =1/2^{k+1}`$, choose an infinite $`𝒥_{k+1}𝒥_k`$, $`𝐣>j_k`$ and a $`\delta >0`$ so that for all $`j𝐣`$, (4.13) holds for $`𝒥^{}=𝒥_{k+1}`$. Now simply let $`\delta _{k+1}=\mathrm{min}\{\delta ,2^1\delta _k,\frac{1}{k+1}\}`$. Since the functions $`\omega _i\mathrm{}`$ are non-decreasing, we have that also for all $`j>𝐣`$, $`\omega _{ij}(\delta _{k+1})1/2^{k+1}`$ for all but finitely many $`i𝒥_{k+1}`$. This completes the inductive construction, with (4.15) holding for all $`k`$.
Now choose $`i_1𝒥_1`$ with $`\omega _{i_1,j_1}(\delta _1)1/2`$. Then also for all but finitely many $`i𝒥_2`$, $`\omega _{i,j_2}(\delta _1)1/2`$ and $`\omega _{i,j_2}(\delta _2)1/4`$. Hence we can choose $`i_2>i_1`$ $`(i_2𝒥_2)`$, with
(4.16)
$$\omega _{i_2,j_2}(\delta _1)\frac{1}{2}\text{ and }\omega _{i_2,j_2}(\delta _2)\frac{1}{4}.$$
But we can also choose $`0<\epsilon _2\delta _2`$ so that
(4.17)
$$\omega _{i_1,j_1}(\epsilon _2)\frac{1}{4}.$$
Thus also
(4.18)
$$\omega _{i_2,j_2}(\epsilon _2)\frac{1}{4}.$$
Now suppose $`i_1<\mathrm{}<i_n`$ and $`\delta _1=\epsilon _1,\mathrm{},\epsilon _n`$ have been chosen so that $`\epsilon _j\delta _j`$ for all $`jn`$ and
(4.19)
$$\omega _{i_k,j_k}(\epsilon _i)\frac{1}{2^i}\text{ for all }1kn,1in.$$
Now by (4.15), choose $`i_{n+1}>i_n`$ ($`i_{n+1}𝒥_{n+1}`$) so that
(4.20)
$$\omega _{i_{n+1},j_{n+1}}(\delta _{\mathrm{}})\frac{1}{2^{\mathrm{}}}\text{ for all }1\mathrm{}n+1.$$
This is possible, since for each $`\mathrm{}`$, $`\omega _{i,j_{n+1}}(\delta _{\mathrm{}})1/2^{\mathrm{}}`$ for all but finitely many $`i𝒥_{n+1}`$.
Again, since the $`\epsilon _{\mathrm{}}`$’s are smaller than the $`\delta _{\mathrm{}}`$’s,
(4.21)
$$\omega _{i_{n+1},j_{n+1}}(\epsilon _{\mathrm{}})\frac{1}{2^{\mathrm{}}}\text{ for all }1\mathrm{}n.$$
Finally, choose $`\epsilon _{n+1}\delta _{n+1}`$ so that
(4.22)
$$\omega _{i_{\mathrm{}},j_{\mathrm{}}}(\epsilon _{n+1})\frac{1}{2^{n+1}}\text{ for all }1\mathrm{}n.$$
Again, we also have
(4.23)
$$\omega _{i_{n+1},j_{n+1}}(\epsilon _{n+1})\frac{1}{2^{n+1}}.$$
This completes the inductive construction of $`i_1<i_2<\mathrm{}`$ and $`\epsilon _1,\epsilon _2,\mathrm{}`$. Then for each $`i`$, we have
(4.24)
$$\underset{k}{sup}\omega _{i_k,j_k}(\epsilon _i)\frac{1}{2^i}.$$
It then follows immediately that (4.14) holds, since if $`\epsilon \epsilon _i`$, then also
(4.25)
$$\underset{k}{sup}\omega _{i_k,j_k}(\epsilon )\frac{1}{2^i}.$$
This completes the proof of Theorem 4.1, in view of the comment after (4.14).∎
Using theorems from Banach space theory, we next obtain a stronger version of 4.1.
###### Theorem 4.2.
Let $`𝒩`$, $`p`$ and $`(x_{ij})`$ be as in 4.1. The conclusion of 4.1 holds under the following assumptions:
* $`1<p`$, and every column is an unconditional basic sequence, every generalized diagonal is a basic sequence, and there is a $`\lambda 1`$ so that every row is a $`\lambda `$-basic sequence
or
* $`p=1`$ and every generalized diagonal is a basic sequence.
Moreover the unconditional assumption in (a) may be dropped if $`𝒩`$ is hyperfinite.
###### Remark.
We do not know if the unconditional assumption in (a) may be dropped in general. However our proof of 4.2 yields the following result, for arbitrary finite $`𝒩`$ and $`1<p<2`$. Assume (a) with “unconditional” deleted. Then the following three alternatives hold: II or III of Theorem 4.1, or
I. There is a $`C1`$ and a column so that for all $`n`$, the column contains $`n`$ elements $`C`$-equivalent to the usual $`\mathrm{}_n^p`$ basis.
To obtain the case $`p>1`$, we require the following remarkable result, due to Brunel and Sucheston (\[BrS1\], \[BrS2\]; see also \[G\]). (A sequence $`(x_j)`$ of non-zero elements in a Banach space is called $`\beta `$-suppression unconditional if for all $`n`$, scalars $`c_1,\mathrm{},c_n`$, and $`F\{1,\mathrm{},n\}`$, $`_{j=F}c_jx_j\beta _{j=1}^nc_jx_j`$. It is easily seen that if $`(x_j)`$ is $`\lambda `$-suppression unconditional, it is $`2\lambda `$-unconditional over real scalars and $`4\lambda `$-unconditional over complex scalars. Actually, a neat result of Kaufman-Rickert yields that such a sequence is $`\pi \lambda `$-unconditional (over complex scalars) \[KR\].)
###### Lemma 4.3.
Let $`(x_n)`$ be a semi-normalized weakly null sequence in a Banach space $`X`$, and let $`\epsilon >0`$. Then there exists a subsequence $`(y_j)`$ of $`(x_j)`$ so that for any $`kj_1<j_2<\mathrm{}<j_{2^k}`$, $`(y_{j_i})_{i=1}^{2^k}`$ is $`(1+\epsilon )`$-suppression unconditional (and hence $`\pi (1+\epsilon )`$-unconditional).
###### Remarks.
1. Actually, the results of Brunel-Sucheston yield much more than this. They obtain that under the hypotheses of Lemma 4.3, there exists a Banach space $`E`$ with a suppression 1-unconditional semi-normalized basis $`(e_j)`$ and a basic subsequence $`(y_j)`$ of $`(x_j)`$ so that:
* $`(e_j)`$ is isometrically equivalent to all of its subsequences and
* for all $`\epsilon >0`$ and $`k`$ large enough, and any $`kj_1<\mathrm{}<j_{2^k}`$, $`(y_{j_i})_{i=1}^{2^k}`$ is $`(1+\epsilon )`$-equivalent to $`(e_1,\mathrm{},e_{2^k})`$.
In the standard Banach space terminology, $`(e_j)`$ is called a subsymmetric basis for $`E`$, and a spreading model for $`(x_j)`$.
2. A classical result of Bessaga-Pełczyński yields that any seminormalized weakly null sequence in a Banach space has a basic subsequence (in fact, for every $`\epsilon >0`$, a subsequence which is $`(1+\epsilon )`$-basic). However it is obtained in \[MR\] that there exists a normalized weakly null sequence in a certain Banach space with no unconditional subsequence, and in \[GM\] that there exists an (infinite dimensional) reflexive Banach space with no (infinite) unconditional basic sequences at all. Thus in a sense, Lemma 4.3 is the best possible positive result in this direction.
We now give consequences of this lemma that are needed for Theorem 4.2. The first one follows from Lemma 3.1 and Lemma 4.3.
###### Corollary 4.4.
Let $`1p<2`$ and $`(f_n)`$ be a weakly null sequence in $`L^p(\tau )`$ so that $`(|f_i|^p)_{i=1}^{\mathrm{}}`$ is uniformly integrable. Then there is a subsequence $`(f_i^{})`$ of $`(f_i)`$ so that
$$\underset{n\mathrm{}}{lim}n^{1/p}\epsilon _1y_1+\mathrm{}+\epsilon _ny_n_{L^p(\tau )}=0$$
uniformly over all subsequences $`(y_i)`$ of $`(f_i^{})`$ and all choices $`(\epsilon _j)`$ of scalars with $`|\epsilon _j|1`$ for all $`j`$.
###### Remark.
The result shows (and also follows from): any spreading model for $`(f_j)`$ is not equivalent to the $`\mathrm{}^p`$-basis.
###### Proof of 4.4.
We may assume without loss of generality that $`f_j_p1`$ for all $`j`$. Let $`\epsilon >0`$ be such that $`\pi (1+\epsilon )4`$, and choose $`(y_j)`$ a subsequence of $`(f_j)`$ satisfying the conclusion of Lemma 4.3. Let $`(r_j)`$ denote the Rademacher functions on $`[0,1]`$ (as defined in Section 3), set $`\stackrel{~}{𝒩}=𝒩\overline{}L^{\mathrm{}}`$, and let $`g_j=y_jr_j`$ for all $`j`$. Then $`(g_j)`$ is 2-unconditional (over complex scalars) and of course $`(|g_j|^p)`$ is also uniformly integrable in $`L^1(\stackrel{~}{𝒩})`$, whence by Lemma 3.1,
(4.26)
$$\underset{n\mathrm{}}{lim}n^{1/p}g_1+\mathrm{}+g_n_{L^p(\stackrel{~}{𝒩})}=0.$$
Let $`\epsilon >0`$, and choose $`N`$ so that if $`nN`$, then
(4.27)
$$n^{1/p}g_1+\mathrm{}+g_n_{L^p(\stackrel{~}{𝒩})}<\frac{\epsilon }{16}$$
and
(4.28)
$$n^{1/p}(1+\mathrm{log}_2n)<\frac{\epsilon }{2}.$$
Now fix $`n`$, and choose $`k`$ with
(4.29)
$$2^{k1}n<2^k.$$
Of course then
(4.30)
$$k1+\mathrm{log}_2n.$$
Now if $`\epsilon _1,\mathrm{},\epsilon _n`$ are given scalars of modulus at most one, then
(4.31)
$$\underset{j=k+1}{\overset{n}{}}\epsilon _jy_j_{L^p(𝒩)}16\underset{j=k+1}{\overset{n}{}}g_j_{L^p(\stackrel{~}{𝒩})}.$$
Indeed, $`y_{k+1},\mathrm{},y_n`$ is 4-unconditional by the conclusion of Lemma 4.3 (since $`nk<n<2^k`$), yielding (4.31). On the other hand,
(4.32)
$$\underset{j=1}{\overset{k}{}}\epsilon _jy_j_{L^p(𝒩)}k1+\mathrm{log}_2n\text{ by (}\text{4.30}\text{).}$$
Thus we have
(4.33)
$$\begin{array}{cc}\hfill n^{1/p}\underset{j=1}{\overset{n}{}}\epsilon _jy_j_p& n^{1/p}\underset{j=1}{\overset{k}{}}\epsilon _jy_j_p+n^{1/p}\underset{j=k+1}{\overset{n}{}}\epsilon _jy_j_p\hfill \\ \\ & n^{1/p}(1+\mathrm{log}_2n)+8n^{1/p}\underset{j=k+1}{\overset{n}{}}g_j_{L^p(\stackrel{~}{𝒩})}\hfill \\ \\ & \frac{\epsilon }{2}+8n^{1/p}\underset{j=1}{\overset{n}{}}g_j_{L^p(\stackrel{~}{𝒩})}\hfill \\ \\ & \frac{\epsilon }{2}+\frac{\epsilon }{2}=\epsilon .\hfill \end{array}$$
(The last inequality holds by (4.27); the next to the last by (4.28) and the fact that $`(g_j)`$ is 1-unconditional over real scalars.) The uniformity of the limit over all subsequences of $`(y_i)`$ follows from the fact that the limit in (4.26) is uniform over all subsequences of $`(g_i)`$, thanks to the proof of Lemma 3.1. ∎
We next note a general consequence of Lemma 4.3, which follows from ultraproducts.
###### Corollary 4.5.
Let $`X`$ be a uniformly convex Banach space and let $`\lambda 1`$, $`\epsilon >0`$, and $`k`$ be given. Then there is an $`nk`$ so that for any $`\lambda `$-basic sequence $`(x_1,\mathrm{},x_n)`$ in $`X`$, there exist $`1j_1<j_2<\mathrm{}<j_k`$ so that $`(x_{j_1},\mathrm{},x_{j_k})`$ is suppression $`(1+\epsilon )`$-unconditional (and hence $`\pi (1+\epsilon )`$-unconditional).
###### Proof.
Suppose the conclusion were false. Then we could find for every $`nk`$, an $`n`$-tuple $`(x_1^n,\mathrm{},x_n^n)`$ of elements in $`X`$ so that $`(x_1^n,\mathrm{},x_n^n)`$ is $`\lambda `$-basic, yet no $`k`$ terms are suppression $`(1+\epsilon )`$-unconditional. By homogeneity, we may assume that $`x_i^n=1`$ for all $`n`$ and $`in`$. Now let $`𝒰`$ be a non-trivial ultrafilter on $``$ and let $`X_𝒰`$ denote the ultrapower of $`X`$ with respect to $`𝒰`$. (That is, we let $`E_𝒰`$ denote the subspace of $`\mathrm{}^{\mathrm{}}(X)`$ consisting of all bounded sequences $`(x_j)`$ in $`X`$ with $`lim_{j𝒰}x_j=0`$, and then set $`X_𝒰=\mathrm{}^{\mathrm{}}(X)/E_𝒰`$.) Since $`X`$ is uniformly convex, so is $`X_𝒰`$. Now define a sequence $`(\stackrel{~}{x}_j)`$ in $`X_𝒰`$ by $`\stackrel{~}{x}_j=\pi (x_j^n)_{n=1}^{\mathrm{}}`$, for all $`j`$, where $`\pi :\mathrm{}^{\mathrm{}}(x)X_𝒰`$ is the quotient map and we set $`x_j^n=0`$ if $`n<j`$. It then follows that $`(\stackrel{~}{x}_j)`$ is also $`\lambda `$-basic and normalized; since $`X_𝒰`$ is reflexive, $`(\stackrel{~}{x}_j)`$ is weakly null. But then by Lemma 4.3, there exist $`k`$ terms $`\stackrel{~}{x}_{j_1},\mathrm{},\stackrel{~}{x}_{j_k}`$ of this sequence with $`(\stackrel{~}{x}_{j_i})_{i=1}^k`$ $`(1+\frac{\epsilon }{2})`$-suppression unconditional. Standard ultraproduct techniques yield that $`\eta >0`$ given, there exists an $`n>j_k`$ so that $`(\stackrel{~}{x}_{j_1},\mathrm{},\stackrel{~}{x}_{j_k})`$ is $`(1+\eta )`$-equivalent to $`(x_{j_1}^n,\mathrm{},x_{j_k}^n)`$ and hence the latter is $`(1+\eta )`$ $`(1+\frac{\epsilon }{2})`$-suppression unconditional. Of course we have a contradiction if $`(1+\eta )(1+\frac{\epsilon }{2})<1+\epsilon `$. ∎
###### Proof of Theorem 4.2(a).
We use the same notations and assumptions as in the proof of Theorem 4.1 (e.g., we assume that $`x_{ij}_p1`$ for all $`i`$ and $`j`$). Assume that Case I of 4.1 does not occur. Then again we have that $`(|x_{ij}|^p)_{i=1}^{\mathrm{}}`$ is uniformly integrable for all $`j`$, and hence (3.23) holds. This is also the case if $`𝒩`$ is hyperfinite and the unconditional assumption in (a) is dropped. For suppose to the contrary that for some $`i`$, $`(f_j)\stackrel{\text{def}}{=}(x_{ij})`$ has the property that $`(|f_j|^p)`$ is not uniformly integrable. Then setting $`g_j=f_jr_j`$ in $`L^p(\stackrel{~}{𝒩})`$ (as defined in the proof of Corollary 4.4), $`(g_j)`$ is unconditional and again $`(|g_j|^p)`$ is not uniformly integrable, hence there exist $`n_1<n_2<\mathrm{}`$ with $`(g_{n_j})`$ equivalent to the usual $`\mathrm{}^p`$-basis, by Corollary 3.4). But $`(f_{n_j})`$ has an unconditional subsequence $`(f_j^{})`$ by \[SF\], \[PX1\]. Of course then $`(f_j^{})`$ is equivalent to $`(g_j^{})\stackrel{\text{def}}{=}(f_j^{}r_j)`$, a subsequence of $`(g_{n_j})`$, whence $`(f_j^{})`$ is equivalent to the $`\mathrm{}^p`$ basis.
Now replace the entire matrix $`(x_{ij})`$ by $`(\stackrel{~}{x}_{ij})\stackrel{\text{def}}{=}(x_{ij}r_{ij})`$ in $`L^p(\stackrel{~}{𝒩})`$ (where $`\stackrel{~}{𝒩}=𝒩\overline{}L^{\mathrm{}}`$), where $`r_{ij}`$ is just a “renumbering” of $`(r_j)`$ via $`\times `$ (precisely, let $`\phi :\times `$ be a bijection, and set $`r_{ij}=r_{\phi (i,j)}`$). Now $`\omega _p(x_{ij},\epsilon )=\omega _p(\stackrel{~}{x}_{ij},\epsilon )`$ for all $`i,j`$, and $`\epsilon `$; hence assuming Case II in the proof of Theorem 4.1 occurs, we have that II of 4.1 holds for the matrix $`(\stackrel{~}{x}_{ij})`$. But then since $`L^p(𝒩)`$ is uniformly convex, II holds for $`(x_{ij})`$ itself, by Corollary 4.5. Indeed, let $`C`$ be as in II of 4.1, let $`k`$ be given. Choose $`nk`$ satisfying the conclusion of 4.5 for $`X=L^p(𝒩)`$ (with $`\pi (1+\epsilon )4`$, say). Choose $`i`$ and $`m_1<\mathrm{}<m_n`$ with $`(\stackrel{~}{x}_j)_{j=1}^n`$ $`C`$-equivalent to the $`\mathrm{}_n^p`$ basis where we set $`x_j=x_{im_j}`$ and $`\stackrel{~}{x}_j=\stackrel{~}{x}_{im_j}`$ for all $`j`$. Then choose $`j_1<\mathrm{}j_k`$ with $`(x_{j_i})`$ 4-unconditional. But then $`(x_{j_i})`$ is 8-equivalent to $`(\stackrel{~}{x}_{j_i})`$, and is hence $`8C`$-equivalent to the $`\mathrm{}_k^p`$ basis.
If Case II in the proof of 4.1 does not occur, we have by Case III that there exists a generalized diagonal $`(\stackrel{~}{x}_{i_n,j_n})_{n=1}^{\mathrm{}}`$ of $`(x_{ij})`$ so that $`(|\stackrel{~}{x}_{i_n,j_n}|^p)_{n=1}^{\mathrm{}}`$ is uniformly integrable. Hence immediately, $`(|x_{i_n,j_n}|^p)_{n=1}^{\mathrm{}}`$ is uniformly integrable, and so by Corollary 4.4, $`(x_{i_n,j_n})`$ has a subsequence $`(y_k)`$ (which is of course also a generalized diagonal) satisfying III of 4.1. This completes the proof of Theorem 4.2(a).∎
To obtain 4.2(b), we need two further “Banach” properties of preduals of von Neumann algebras. The first one holds in complete generality.
###### Lemma 4.6.
Let $``$ be a von-Neumann algebra, and let $`(f_n)`$ be a bounded sequence in $`_{}`$ such that $`(f_n)`$ is not relatively weakly compact. Then $`(f_n)`$ has a subsequence equivalent to the $`\mathrm{}^1`$-basis.
We give a “quantitative” proof of this result at the end of this section, using the case for commutative $`𝒩`$ established in \[R1\]. In fact, Lemma 4.6 is due to H. Pfitzner \[Pf\]. However, the second result we need is a “localization” of our proof, which does not seem to follow directly from previously known material. This result yields that if $`n`$ elements of $`_a(𝒩_{})`$ ($`𝒩`$ finite) have mass at least $`\eta `$ on pairwise orthogonal projections, then $`k`$ of these are $`C`$-equivalent to the $`\mathrm{}_k^1`$-basis. Here, $`C`$ depends only on $`\eta `$, $`n`$ on $`k`$ and $`\eta `$. To make this more manageable, let us simply say that $`n`$ elements $`f_1,\mathrm{},f_n`$ of the predual of a von-Neumann algebra $``$ are $`\eta `$-disjoint provided there exist pairwise orthogonal projections $`P_1,\mathrm{},P_n`$ in $``$ such that
(4.34)
$$P_if_iP_i_1\eta \text{ for all }i.$$
(Here, if $`P`$ and $`f_{}`$, $`PfP`$ is defined by: $`T,PfP=PTP,f`$ for all $`T`$. Also, $`_1`$ denotes the predual norm on $`_{}`$.) (We shall also say $`f_1,\mathrm{},f_n`$ are disjoint provided there are pairwise orthogonal projections $`P_1,\mathrm{},P_n`$ in $``$ with $`f_i=P_if_iP_i`$ for all $`i`$. Evidently if the $`f_i`$’s are normalized, they are disjoint iff they are 1-disjoint.)
###### Lemma 4.7.
Given $`\eta >0`$, then if $`C>\frac{1}{\eta }`$, then for all $`k1`$, there is an $`nk`$ so that for any von-Neumann algebra $`𝒩`$ and $`\eta `$-disjoint elements $`f_1,\mathrm{},f_n`$ in $`_a(𝒩_{})`$, there exist $`j_1<\mathrm{}<j_k`$ with $`(f_{j_i})_{i=1}^k`$ $`C`$-equivalent to the $`\mathrm{}_k^1`$ basis.
We delay the proof of this result, and complete the proof of Theorem 4.2, i.e., the case $`p=1`$. Again we make the same assumptions and use the same notation as in the proof of 4.1(a). Now suppose that Case I of Theorem 4.1 does not occur. We now have, immediately from Proposition 2.5 and Lemma 4.6, that $`(x_{ij})_{j=1}^{\mathrm{}}`$ is uniformly integrable for all $`j`$, and hence again (3.23) holds. Now again assume Case II of the proof 4.1 holds. Then the proof of 4.1II yields that for all $`n`$, there exists a row $`i`$ and $`j_1<\mathrm{}<j_n`$ so that $`(f_k)_{k=1}^n`$ is $`\frac{\eta }{3}`$-disjoint, where $`f_k=x_{ij_k}`$ for all $`k`$.
Indeed, we obtain there (following formula (4.3)), that for all $`n`$, there is a sequence $`(f_1,\mathrm{},f_n)`$ satisfying the assumptions of Lemma 3.3 (for $`\eta >0`$ and $`0<\epsilon <\frac{\eta }{2}`$) except for the $`u`$-unconditionality assumption. But the proof of Lemma 3.3 yields precisely that $`(f_1,\mathrm{},f_n)`$ is $`\frac{\eta }{2}\epsilon `$ disjoint; the unconditionality assumption was only used, in invoking Lemma 3.2. Of course we may choose $`\epsilon =\frac{\eta }{6}`$, and so $`(f_1,\mathrm{},f_n)`$ is then $`\frac{\eta }{3}`$-disjoint.
Then Lemma 4.7 immediately yields Case II of Theorem 4.1. Finally, assuming Case II of the proof of 4.1 does not occur, we obtain again from the proof of Case III that there exists a generalized diagonal $`(g_k)`$ of $`(x_{ij})`$ with $`(g_k)`$ uniformly integrable. Hence there exists a weakly convergent subsequence $`(f_j)`$ of $`(g_k)`$, by Proposition 2.5. But since we assume the generalized diagonals of $`(x_{ij})`$ are basic sequences, $`(f_j)`$ must be weakly null. Now Corollary 4.4 immediately yields Case III of Theorem 4.1. ∎
###### Remark.
The case $`p=1`$ of Theorem 4.2 may be alternatively formulated as follows (with essentially no assumptions at all on the matrix $`(x_{ij})`$).
###### Theorem 4.2(b).
Let $`𝒩`$ be a finite von-Neumann algebra and let $`(x_{ij})`$ be an infinite semi-normalized matrix in $`𝒩_{}`$. Then one of the following holds.
* Some column has a subsequence equivalent to the usual $`\mathrm{}^1`$ basis.
* There is a $`C1`$ so that for all $`n`$, there exists a row with $`n`$ elements $`C`$-equivalent to the usual $`\mathrm{}_n^1`$ basis.
* Some generalized diagonal of $`(x_{ij})`$ is weakly convergent.
It remains to prove Lemma 4.7. This is an immediate consequence of the following two results, which in turn follow from the techniques in \[R1\]. (We denote the “predual norm” of a general von-Neumann algebra by $`_1`$.)
###### Lemma 4.8.
Let $`𝒩`$ be an arbitrary von-Neumann algebra, and $`f_1,f_2,\mathrm{}`$ be a finite or infinite sequence in $`𝒩_{}`$ with $`f_i_11`$ for all $`i`$. Assume there are pairwise orthogonal projections $`P_1,P_2,\mathrm{}`$ in $`𝒩`$ and $`0<\epsilon <\delta 1`$ so that for all $`i`$,
(4.35)
$$P_if_iP_i_1\delta \text{ and }\underset{ji}{}P_jf_iP_j_1\epsilon .$$
Then $`f_1,f_2,\mathrm{}`$ is $`\frac{1}{\delta \epsilon }`$ equivalent to the usual basis of $`\mathrm{}^1`$ (resp. $`\mathrm{}_n^1`$ if the sequence has $`n`$ terms).
###### Lemma 4.9.
Let $`k1`$ and $`0<\epsilon <1`$ be given. There is an $`nk`$ so that given any von Neumann algebra $`𝒩`$, $`f_1,\mathrm{},f_n_a(𝒩_{})`$, and pairwise orthogonal projections $`P_1,\mathrm{},P_n`$ in $`𝒩`$, there exist $`j_1<j_2<\mathrm{}<j_k`$ so that for all $`1ik`$,
(4.36)
$$\underset{ri}{}P_{j_r}f_{j_i}P_{j_r}_1<\epsilon .$$
###### Remark.
We obtain that we may choose $`n=k^{\mathrm{}}`$ where $`\mathrm{}=[1/\epsilon ]+1`$.
###### Proof of Lemma 4.7.
Let $`C>\frac{1}{\eta }`$ and choose $`0<\epsilon <\eta `$ with $`\frac{1}{\eta \epsilon }<C`$. Let $`n`$ be as in Lemma 4.9, $`f_1,\mathrm{},f_n`$ as in the hypotheses of 4.7, and choose $`j_1,\mathrm{},j_k`$ satisfying the conclusion of 4.9. Then $`(f_{j_i})_{i=1}^k`$ is $`C`$-equivalent to the $`\mathrm{}_k^1`$ basis by Lemma 4.8. ∎
###### Proof of Lemma 4.8.
Let $`n<\mathrm{}`$ be less than or equal to the number of terms in the sequence, and let $`c_1,\mathrm{},c_n`$ be given scalars with
(4.37)
$$\underset{i=1}{\overset{n}{}}|c_i|=1.$$
Let $`g=_{i=1}^nc_if_i`$. Since the $`P_j`$’s are pairwise orthogonal, we have that
(4.38)
$$g_1\underset{j=1}{\overset{n}{}}P_jgP_j_1.$$
Now fixing $`j`$,
(4.39)
$$\begin{array}{cc}\hfill P_jgP_j_1& P_jc_jf_jP_j+P_j\underset{ij}{}c_if_iP_j_1\hfill \\ & |c_j|\delta \underset{ij}{}|c_i|P_jf_iP_j_1\hfill \end{array}$$
by (4.35) and the triangle inequality. Hence using (4.38) and (4.39),
(4.40)
$$\begin{array}{cc}\hfill g_1& \underset{j=1}{\overset{n}{}}|c_j|\delta \underset{j=1}{\overset{n}{}}\underset{ij}{}|c_i|P_jf_iP_j_1\hfill \\ & =\delta \underset{i=1}{\overset{n}{}}|c_i|\underset{ji}{}P_jf_iP_j_1\text{ by (}\text{4.37}\text{)}\hfill \\ & \delta \epsilon \text{ by (}\text{4.37}\text{) and (}\text{4.35}\text{).}\hfill \end{array}$$
This completes the proof. ∎
We finally deal with Lemma 4.9. This result follows from the simplest possible setting: $`𝒩=\mathrm{}_n^{\mathrm{}}`$, the $`f_i`$’s are in $`\mathrm{}_n^{1+}`$ (i.e., the positive part of $`𝒩_{}=\mathrm{}_n^1`$), and the orthogonal projections $`P_i`$ correspond to multiplication by $`\chi _{\{i\}}`$ for all $`i`$. That is, we finally have the following elementary disjointness result.
###### Lemma 4.10.
A. Let $`f_1,f_2,\mathrm{}`$ be a bounded infinite subset of $`\mathrm{}^{1+}`$, and let $`\epsilon >0`$. There exist $`n_1<n_2<\mathrm{}`$ so that for all $`i`$,
(4.41)
$$\underset{ji}{}f_{n_i}(n_j)<\epsilon .$$
B. Let $`k`$ and $`\epsilon >0`$ be given. There exists an $`Nk`$ so that given $`f_1,\mathrm{},f_N_a\mathrm{}_N^{1+}`$, there exist $`n_1<n_2<\mathrm{}<n_k`$ so that for all $`1ik`$, (4.41) holds.
###### Remark.
Part A is a special case of Lemma 1.1 of \[R1\]. Part B appears to be new. We obtain in fact that we may let $`N=k^{\mathrm{}}`$ where $`\mathrm{}=[1/\epsilon ]+1`$.
###### Proof of Lemma 4.9.
Let $`\epsilon >0`$ and $`N`$ be as in the conclusion of 4.10B. Let the $`f_i`$’s and $`P_i`$’s be as in the statement of 4.9. For each $`i`$, define $`\stackrel{~}{f}_i`$ in $`\mathrm{}^{1+}`$ by $`\stackrel{~}{f}_i(j)=P_jf_iP_j_1`$ for all $`1jN`$. Then
(4.42)
$$\underset{j=1}{\overset{N}{}}P_jf_iP_j_1=\stackrel{~}{f}_i_1f_i_11$$
for all $`i`$. Now the conclusion of B yields $`j_1<\mathrm{}<j_k`$ so that
(4.43)
$$\underset{ri}{}\stackrel{~}{f}_{j_i}(j_r)<\epsilon \text{ for all }1ik.$$
Then $`f_{j_1},\mathrm{},f_{j_k}`$ satisfies the conclusion of Lemma 4.9. ∎
At last, we give the proof of Lemma 4.10.
We first prove A, using an argument due to J. Kupka \[Ku\]. We then adapt this argument to obtain Part B. We regard elements of $`\mathrm{}^{1+}`$ as finite measures on subsets of $``$ and use the notation: $`f(E)=_{jE}f(j)`$ for $`f\mathrm{}^{1+}`$ and $`E`$. Thus, the conclusion of A may be restated: There exists an infinite $`M`$ so that
(4.44)
$$f_i(M\{i\})<\epsilon \text{ for all }iM.$$
Let $`N_1,N_2,\mathrm{}`$ be pairwise disjoint infinite subsets of $``$ with $`=_{j=1}^{\mathrm{}}N_j`$.
### Case I
For each $`i`$, there exists $`n_iN_i`$ so that
(4.45)
$$f_{n_i}(N_i)<\epsilon .$$
It then follows that $`M=\{n_1,n_2,\mathrm{}\}`$ satisfies (4.44). Indeed, for all $`i`$,
(4.46)
$$\{n_1,n_2,\mathrm{},n_{i1},n_{i+1},\mathrm{}\}N_i$$
since the $`N_i`$ are disjoint, so (4.44) follows from (4.45) and (4.46).
### Case II
Case I fails. Thus we may choose $`i_1`$ so that
(4.47)
$$f_j(N_{i_1})\epsilon \text{ for all }jN_{i_1}.$$
Now repeat the same procedure; let $`M_1=N_{i_1}`$, and choose $`M_1^1,M_1^2,\mathrm{}`$ disjoint infinite subsets of $`M_1`$ with $`M_1=_{j=1}^{\mathrm{}}M_1^j`$. If Case I fails for $`M_1`$, we will obtain $`M_2\stackrel{\text{def}}{=}M_1^j`$ (for some $`j`$) so that
(4.48)
$$f_j(M_1M_2)\epsilon \text{ for all }jM_2.$$
Again divide up $`M_2`$. This “failure of Case I” must terminate before $`\mathrm{}`$ steps, where $`f_j_1<\mathrm{}\epsilon `$ for all $`j`$. Indeed, otherwise, we finally obtain $`=M_0M_1M_2\mathrm{}M_{\mathrm{}}`$ and a $`jM_{\mathrm{}}`$ with
(4.49)
$$f_j(M_{i1}M_i)\epsilon \text{ for all }i,$$
whence $`f_j\mathrm{}\epsilon `$, a contradiction.
###### Proof of Part B.
Let $`\mathrm{}=[1/\epsilon ]+1`$ and let $`N=k^{\mathrm{}}`$. Let then $`f_1,\mathrm{},f_N_a(\mathrm{}_N^{1+})`$ be given. Of course the conclusion of Part B may be restated: There exists an $`M\{1,\mathrm{},N\}`$ with $`\mathrm{\#}M=k`$ so that (4.44) holds.
Let $`N_1,\mathrm{},N_k`$ be disjoint subsets of $`\{1,\mathrm{},N\}`$, each of cardinality $`k^\mathrm{}1`$, and just repeat the argument for Part A, Case I. If Case I fails, we repeat again the rest of the argument: that is, we find $`i_1`$ satisfying (4.47) and set $`M_1=N_{i_1}`$. Now we just choose $`M_1^1,\mathrm{},M_1^k`$ disjoint subsets of $`M_1`$, each of cardinality $`k^\mathrm{}2`$; if Case I fails for $`M_1`$, we continue as before, with $`M_2`$ satisfying (4.48) and $`M_2M_1`$, $`\mathrm{\#}M_2=k^\mathrm{}2`$. If Case I fails for $`\mathrm{}`$ steps, we obtain finally $`\{1,\mathrm{},N\}=M_0M_1\mathrm{}M_{\mathrm{}}`$ with $`\mathrm{\#}M_i=k^\mathrm{}i`$ for all $`i`$, so $`\mathrm{\#}M_{\mathrm{}}=1`$; and for $`j`$ the unique number of $`M_{\mathrm{}}`$, (4.49) holds, whence again $`f_j\mathrm{}\epsilon >1`$, a contradiction. ∎
Let us say that a finite or infinite sequence $`(f_i)`$ satisfying the hypotheses of Lemma 4.8 is $`(\delta ,\epsilon )`$-relatively disjoint. It then follows from arguments in \[R1\] that the closed linear span of such a sequence is $`K`$-complemented in $`𝒩_{}`$, where $`K`$ depends only on $`\delta `$ and $`\epsilon `$. Indeed, let $`W`$ denote the closed linear span of the $`f_i`$’s; let $`P_1,P_2,\mathrm{}`$ be as in the statement of 4.8, and let $`g_j=P_jf_jP_j`$ for all $`j`$, then let $`Z`$ denote the closed linear span of the $`g_j`$’s. Of course then $`Z`$ is isometric to $`\mathrm{}^1`$ (or $`\mathrm{}_n^1`$ if the sequence has $`n`$ terms). We may easily define a contractive projection $`R:𝒩_{}Z`$ as follows. For each $`j`$, choose by duality an element $`\phi _j𝒩`$ of norm one with $`\phi _j=P_j\phi _jP_j`$ and
(4.50)
$$\phi _j,g_j=g_j_1.$$
(Note that $`1g_j_1\delta `$ for all $`j`$.) Then define
(4.51)
$$R(f)=\phi _j,fg_j_1^1g_j$$
for $`f𝒩_{}`$. Next, define an operator $`U:WZ`$ by
(4.52)
$$U(c_jf_j)=c_jg_j$$
for all $`c_j`$’s with $`|c_j|<\mathrm{}`$. Then Lemma 4.8 yields that $`U`$ is invertible with
(4.53)
$$U^1(\delta \epsilon )^1.$$
Now a simple computation yields that
(4.54)
$$U(w)R(w)\frac{\epsilon }{\delta }U(w)\text{ for all }wW.$$
It then follows that $`R|W`$ is an isomorphism mapping $`W`$ onto $`Z`$, with
(4.55)
$$(R|W)^1\left[(1\frac{\epsilon }{\delta })(\delta \epsilon )\right]^1\stackrel{\text{def}}{=}K.$$
Finally, $`Q\stackrel{\text{def}}{=}(R|W)^1R`$ is thus a projection from $`𝒩`$ onto $`W`$, with $`QK`$. It then follows that the elements satisfying the conclusion of Lemma 4.7 have a “well-complemented” linear span.
We also obtain finally, a quantitative proof of Lemma 4.6, yielding also the result of H. Pfitzner \[Pf\] that the preduals of von Neumann algebras have Pełczyński’s property $`(V^{})`$.
###### Lemma 4.6.
Let $`𝒩`$ be an arbitrary von Neumann algebra, and $`W`$ be a subset of $`_a𝒩_{}`$ so that there exists a sequence $`P_1,P_2,\mathrm{}`$ of orthogonal projections in $`𝒩`$ with
(4.56)
$$\overline{\mathrm{lim}}_j\underset{wW}{sup}|P_j,w|\stackrel{\text{def}}{=}\eta >0.$$
Then given $`C>\frac{1}{\eta }`$, there exists a sequence $`w_1,w_2,\mathrm{}`$ in $`W`$ which is $`C`$-equivalent to the usual $`\mathrm{}^1`$-basis, with closed linear span $`C`$-complemented in $`𝒩_{}`$.
###### Remark.
By Akemann’s criterion \[A\], it thus follows that any bounded non-relatively weakly compact subset of $`𝒩_{}`$ contains a sequence equivalent to the $`\mathrm{}^1`$-basis, with complemented span. This is an equivalent formulation of property $`(V^{})`$.
###### Proof.
It follows easily that we may choose $`(f_i)`$ a sequence in $`W`$ and $`n_1<n_2<\mathrm{}`$ so that
(4.57)
$$\underset{¯}{\mathrm{lim}}|P_{n_j},f_j|\eta .$$
Then given $`0<\epsilon <\eta ^{}<\eta `$, Lemma 4.10A yields a subsequence $`(f_j^{})`$ of $`(f_j)`$ so that $`(f_j^{})`$ is $`(\eta ^{},\epsilon )`$-relatively disjoint. Finally, since $`\eta ^{}`$ may be arbitrarily close to $`\eta `$ and $`\epsilon `$ arbitrarily small, we deduce from Lemma 4.8 and (4.55) that given $`C>\frac{1}{\eta }`$, $`(f_i^{})`$ may be chosen $`C`$-equivalent to the $`\mathrm{}^1`$-basis with span $`C`$-complemented in $`𝒩_{}`$. ∎
## 5. Complements on the Banach/operator space structure of $`L^p(𝒩)`$-spaces
We give here several applications of our main result, and the techniques used in its proof. For the first one, we let $`\mathrm{Row}`$ (resp. $`\mathrm{Col}`$) denote the operator row (resp. column) space. We also follow the notation in \[Pi2\]: for a given operator space $`X`$, $`X^{\mathrm{op}}`$ (the “opposite” of $`X`$) denotes the following operator space: if $`XB(H)`$ and $`(x_{ij})`$ is an element of $`𝒦_{\mathrm{op}}X`$, regarded as a matrix, then $`X^{\mathrm{op}}\stackrel{\text{def}}{=}\{(x_{ji}):(x_{ij})𝒦_{\mathrm{op}}X\}`$, where $`𝒦`$ denotes the space of compact opertors on $`\mathrm{}^2`$. One then has that $`\mathrm{Row}^{}=\mathrm{Row}^{\mathrm{op}}=\mathrm{Col}`$.
###### Proposition 5.1.
Let $`𝒩`$ be a finite von Neumann algebra. Then neither $`\mathrm{Row}`$ nor $`\mathrm{Col}`$ is completely isomorphic to a subspace of $`L^1(𝒩)`$.
###### Proof.
Suppose to the contrary that there exists an $`XL^1(𝒩)`$ with $`X`$ completely isomorphic to $`\mathrm{Row}`$. But then $`X^{\mathrm{op}}L^1(𝒩^{\mathrm{op}})`$ is completely isomorphic to $`\mathrm{Col}`$. Let then $`=𝒩^{\mathrm{op}}\overline{}𝒩`$. $``$ is again a finite von-Neumann algebra, and now $`X^{\mathrm{op}}\widehat{}X`$ is a subspace of $`L^1()`$; that is, $`\mathrm{Col}\widehat{}\mathrm{Row}`$ is completely isomorphic to a subspace of $`L^1()`$. But $`\mathrm{Col}\widehat{}\mathrm{Row}`$ is (completely isometric to) $`C_1`$; this contradicts our main result. ∎
###### Remark.
An operator space $`X`$ is called homogeneous if every bounded linear operator on $`X`$ is completely bounded; $`X`$ is called Hilbertian if it is Banach isomorphic to a Hilbert space. The above argument then yields the following generalization (since $`\mathrm{Row}`$ is indeed a homogeneous Hilbertian operator space).
###### Proposition.
Let $`X`$ be an infinite dimensional Hilbertian homogeneous operator space so that $`X^{}`$ is completely isomorphic to $`X^{\mathrm{op}}`$, and let $`𝒩`$ be a finite von Neumann algebra. Then $`X`$ is not completely isomorphic to a subspace of $`L^1(𝒩)`$.
To obtain this, first observe that the hypotheses yield that $`X^{}_{\mathrm{op}}X`$ is Banach isomorphic to $`𝒦`$. Hence $`X^{}\widehat{}X`$ is Banach isomorphic to $`C_1`$. But $`X^{}\widehat{}X`$ is completely isomorphic to $`X^{\mathrm{op}}\widehat{}X`$ by hypothesis; as above, if we then assume that $`XL^1(𝒩)`$, we obtain that $`C_1`$ Banach embeds in $`L^1()`$, again contradicting our main result.∎
Our next result yields characterizations of those subspaces of $`L^p(𝒩)`$ which contain $`\mathrm{}^p`$ isomorphically ($`1p<2`$, $`𝒩`$ finite). We have need of the following concept. (For isomorphic Banach spaces $`X`$ and $`Y`$, $`d(X,Y)=inf\{TT^1:T:XY`$ is a surjective isomorphism).
###### Definition 5.2.
Let $`1p\mathrm{}`$. A Banach space $`X`$ is said to contain $`\mathrm{}_n^p`$’s if there is a $`C1`$ so that for all $`n`$, there exists a subspace $`X_n`$ of $`X`$ with $`d(X_n,\mathrm{}_n^p)C`$.
A remarkable result of J.L. Krivine yields that if a Banach space contains $`\mathrm{}_n^p`$’s, it contains them almost isometrically (\[Kr\]; cf. also \[R3\], \[L\]). That is, then for every $`\epsilon `$ and $`n`$, one can choose $`X_nX`$ with $`d(X_n,\mathrm{}_n^p)<1+\epsilon `$. (Of course the famous Dvoretzky theorem yields that every infinite dimensional Banach space contains $`\mathrm{}_n^2`$’s almost isometrically; also the case $`p=1`$ or $`\mathrm{}`$ in Krivine’s Theorem was established previously by Giesy-James \[GJ\].)
We also need the following natural notion.
###### Definition 5.3.
Let $`𝒩`$ be a von Neumann algebra and $`1p<\mathrm{}`$. A sequence $`(g_n)`$ in $`L^p(𝒩)`$ is called disjointly supported provided there exists a sequence $`P_1,P_2,\mathrm{}`$ of pairwise orthogonal projections in $`𝒩`$ so that $`g_j=P_jg_jP_j`$ for all $`j`$. A semi-normalized sequence $`(f_n)`$ in $`L^p(𝒩)`$ is called almost disjointly supported if there exists a disjointly supported sequence $`(g_j)`$ in $`L^p(𝒩)`$ so that $`lim_n\mathrm{}f_ng_n_{L^p(𝒩)}=0`$.
Of course a disjointly supported sequence of non-zero elements spans a subspace isometric to $`\mathrm{}^p`$. A standard elementary perturbation argument then yields that an almost disjointly supported sequence in $`L^p(𝒩)`$ has, for every $`\epsilon >0`$, a subsequence spanning a subspace $`(1+\epsilon )`$-isomorphic to $`\mathrm{}^p`$. The next result yields in particular that for $`𝒩`$ finite, and $`1p<2`$, subspaces of $`L^p(𝒩)`$ which are isomorphic to $`\mathrm{}^p`$ always contain almost disjointly supported sequences.
###### Theorem 5.4.
Let $`1p<2`$ and $`𝒩`$ be a finite von Neumann algebra; let $`\tau `$ be a faithful normal tracial state on $`𝒩`$. Let $`X`$ be a closed linear subspace of $`L^p(𝒩)`$. The following assertions are equivalent.
* $`X`$ contains a subspace isomorphic to $`\mathrm{}^p`$.
* $`X`$ contains $`\mathrm{}_n^p`$’s.
* $`\{|x|^p:x_a(X)\}`$ is not uniformly integrable.
* $`sup_{f_a(X)}\omega _p(f,\epsilon )=sup_{f_a(X)}\stackrel{~}{\omega }_p(f,\epsilon )=1`$ for all $`\epsilon >0`$.
* The $`p`$ and $`1`$ norms on $`X`$ are not equivalent (in case $`p>1`$).
* $`X`$ contains an almost disjointly supported sequence.
* For all $`\epsilon >0`$, $`X`$ contains a subspace $`(1+\epsilon )`$-isomorphic to $`\mathrm{}^p`$.
###### Remarks.
1. This result is established for the commutative case in \[R2\]; the case $`p>2`$ is also valid, and follows (with some extra work for assertion 5) from the results in \[S1\]. Again, the commutative case for $`p>2`$ is immediate from the classical work of Kadec-Pełczyński \[KP\]. Also, condition 5 may be replaced by the following one, valid also for $`p=1`$:
* The $`p`$ and $`q`$ quasi-norms are not equivalent on $`X`$ for all $`0<q<p`$.
2. The equivalences of 1, 5, 6 and 7 of Theorem 5.4 follow also from recent work of N. Randrianantoanina, which establishes these also for semi-finite von-Neumann algebras $`𝒩`$ and $`1p<\mathrm{}`$, $`p2`$ (\[Ra1\] and \[Ra2\]).
###### Proof.
We show $`124671`$, $`432`$, and $`453`$ in case $`p>1`$. Of course $`12`$ and $`71`$ are trivial. So is $`43`$, in virtue of Lemma 2.3.
$`24`$. Fix $`\delta >0`$. Choosing an “almost isometric” copy of $`\mathrm{}_n^p`$ in $`X`$ by Krivine’s theorem, we shall show that for $`n`$ large enough, one of the natural basis elements $`f_i`$ of this copy is such that $`\stackrel{~}{\omega }_p(f_i,\delta )`$ is almost equal to 1.
Define $`\lambda `$ by
(5.1)
$$\lambda =sup\{\stackrel{~}{\omega }_p(x,\delta ):xX,x1\}.$$
Let $`C>1`$, and using Krivine’s theorem, choose $`f_1,\mathrm{},f_n_a(X)`$ with $`(f_1,\mathrm{},f_n)`$ $`C`$-equivalent to the $`\mathrm{}_n^p`$ basis. In particular, we have that
(5.2)
$$\underset{i=1}{\overset{n}{}}\pm f_i_p\frac{1}{C}n^{1/p}\text{ for all choices of }\pm .$$
Again by the final assertion of Lemma 2.3, we may choose for each $`i`$ a $`\psi _i𝒩`$ so that
(5.3)
$$\psi _i_{\mathrm{}}\delta ^{1/p}\text{ and }f_i\psi _i\stackrel{~}{\omega }_p(f_i,\delta )\lambda .$$
Thus letting $`\beta `$ be as in the proof of Lemma 3.1, again we have
(5.4)
$$\begin{array}{cc}\hfill \frac{1}{C}n^{1/p}& f_ir_i_{L^p(\beta )}\text{ by (}\text{5.2}\text{)}\hfill \\ & \psi _ir_i_{L^2(\beta )}+(f_i\psi _i)r_i_{L^p(\beta )}\hfill \\ & \delta ^{1/p}\sqrt{n}+\lambda n^{1/p}\hfill \end{array}$$
by (5.3) and the fact that $`L^p(\beta )`$ is type $`p`$ with constant one.
Thus
(5.5)
$$\frac{1}{C}\frac{1}{\delta ^{1/p}n^{\frac{1}{p}\frac{1}{2}}}\lambda .$$
Since $`C>1`$ and $`n`$ are arbitrary, we obtain that $`\lambda =1`$, proving $`24`$.
$`46`$. We first note that assuming 4, then given $`1>\epsilon >0`$, we may choose $`fX`$ with $`f_p=1`$ and $`P𝒫(𝒩)`$ with $`\tau (P)<\epsilon `$ so that
(5.6)
$$fP_p>1\epsilon \text{ and }f(IP)_p<\epsilon .$$
Indeed, choose $`f`$ in $`X`$ of norm one so that $`\stackrel{~}{\omega }_p(f,\epsilon )>1\epsilon `$. Then choose $`P`$ a spectral projection for $`|f|`$ with $`fP_p>(1\epsilon ^p)^{1/p}`$. But then since $`P`$ commutes with $`|f|`$,
(5.7)
$$fP_p^p=\tau (|f|^pP)\text{ and }f(IP)_p^p=\tau (|f|^p(IP)),$$
whence
(5.8) $`1`$ $`\tau (|f|^pP)+\tau (|f|^p(IP))(1\epsilon )+f(IP)_p^p`$
(5.9) $`1`$ $`\tau (|f|^pP)+\tau (|f|^p(IP))`$
$`1\epsilon ^p+f(IP)_p^p,`$
so $`f(IP)_p<\epsilon `$ as desired. Now since $`|f|`$ and $`|f^{}|`$ are unitarily equivalent in $`𝒩`$, we also obtain the existence of a $`QP(𝒩)`$ with $`\tau (Q)<\epsilon `$ so that
(5.10)
$$Qf_p>1\epsilon \text{ and }f(IQ)_p<\epsilon .$$
Then let $`R=PQ`$. We have
(5.11)
$$\tau (R)<2\epsilon \text{ and }fRfR<2\epsilon .$$
Indeed, the first estimate is trivial; but
$$fRfR=f(IR)+(IR)fR=f(IP)(IR)+(IR)(IQ)fR$$
and so (5.11) follows from (5.6) and (5.10).
Now using that for $`\epsilon >0`$, $`f`$ of norm 1 in $`X`$ and $`R`$ may be chosen satisfying (5.11) we choose inductively $`f_1,f_2,\mathrm{}`$ in $`X`$ of norm one, $`1>\delta _1>\delta _2>\mathrm{}>0`$, and $`Q_1,Q_2,\mathrm{}`$ in $`𝒫(𝒩)`$ so that for all $`j`$,
(5.12)
$$f_jQ_jf_jQ_j_p<\frac{1}{2^j}\text{ and }\tau (Q_j)\frac{\delta _j}{2^j}$$
(5.13)
$$\omega _p(f_j,\delta _{j+1})<\frac{1}{2^j}.$$
To see this is possible, just choose $`\delta _1=1/2`$, then choose $`f_1`$ and $`Q_1`$ thanks to (5.11). Suppose $`f_1,\mathrm{},f_n`$, and $`\delta _n`$ chosen. By uniform integrability of $`\{|f_n|^p\}`$, choose $`\delta _{n+1}<\delta _n`$ so that $`\omega _p(f_n,\delta _{n+1})<1/2^{n+1}`$. Then choose $`f_{n+1}`$ and $`Q_{n+1}`$ satisfying (5.12) for $`j=n+1`$.
Now define projections $`P_j`$ and $`\stackrel{~}{Q}_j`$ by (3.19). The $`P_j`$’s are orthogonal and by the argument for the last part of Proposition 2.5, fixing $`j`$, we have
(5.14)
$$\begin{array}{cc}\hfill \tau (\stackrel{~}{Q}_j)\underset{k>j}{}\tau (Q_k)& \delta _{j+1}\underset{k>j}{}\frac{1}{2^k}\text{ by (}\text{5.12}\text{)}\hfill \\ & <\delta _{j+1}.\hfill \end{array}$$
Hence
$$\stackrel{~}{Q}_jf_j_p\omega _p(f_j^{},\delta _{j+1})=\omega _p(f_j,\delta _{j+1})<\frac{1}{2^j}$$
(by (5.13)) and also
$$f_j\stackrel{~}{Q}_j_p\omega _p(f_j,\delta _{j+1})<\frac{1}{2^j}.$$
Hence
(5.15)
$$\stackrel{~}{Q}_jf_jQ_j_p<\frac{1}{2^j}\text{ and }Q_jf_j\stackrel{~}{Q}_j_p<\frac{1}{2^j}.$$
Hence finally we have by (5.12) and (5.15),
(5.16)
$$f_jP_jf_jP_j\frac{3}{2^j}\text{ for all }j.$$
Thus $`(f_j)`$ is almost disjointly supported, proving that 6 holds.
$`67`$ is a standard perturbation argument in Banach space theory. Assuming 6 holds, we may choose a normalized disjointly supported sequence $`(g_n)`$ in $`L^p(𝒩)`$ and a sequence $`(f_n)`$ in $`X`$ so that
(5.17)
$$g_nf_n_p<\mathrm{}.$$
But now $`(g_n)`$ is 1-equivalent to the $`\mathrm{}^p`$-basis, and a simple perturbation argument gives that given $`\epsilon >0`$, there is an $`N`$ so that $`(f_n)_{nN}`$ is $`(1+\epsilon )`$-equivalent to the $`\mathrm{}^p`$ basis. (Thus $`(f_n)`$ is “almost isometrically equivalent” to the $`\mathrm{}^p`$ basis.)
$`32`$. We have that if $`p=1`$, $`X`$ contains a subspace isomorphic to $`\mathrm{}^1`$ by Lemma 4.6, so assume $`p>1`$. We may choose a sequence $`(f_n)`$ of norm-1 elements of $`X`$, $`\delta _1>\delta _2>\mathrm{}`$ with $`\delta _n0`$ and $`\eta >0`$ so that
(5.18)
$$\omega _p(f_n,\delta _n)>\eta \text{ for all }n.$$
By passing to a subsequence, we may assume without loss of generality that $`(f_n)`$ is weakly convergent, with weak limit $`f`$, say. But
(5.19)
$$\omega _p(f_nf,\delta _n)\omega _p(f_n,\delta _n)\omega _p(f,\delta _n)$$
and hence
(5.20)
$$\underset{¯}{\mathrm{lim}}_n\mathrm{}\omega _p(f_nf,\delta _n)\eta .$$
That is, we have now obtained a weakly null sequence $`(g_n)`$ in $`X`$ so that
(5.21)
$$(|g_n|^p)\text{ is not uniformly integrable.}$$
By Corollary 3.4, after passing to a subsequence of $`(g_n)`$, we may assume
(5.22)
$$(g_nr_n)\text{ is }C\text{-equivalent to the usual }\mathrm{}^p\text{-basis in }L^p(\beta )\text{ for some }C\text{.}$$
Now Lemma 4.3 yields that for all $`n`$, there exist $`m_1<m_2<\mathrm{}<m_n`$ so that $`g_{m_1},\mathrm{},g_{m_n}`$ is 4-unconditional, and hence
(5.23)
$$(g_{m_i})_{i=1}^n\text{ is }4C\text{-equivalent to the }\mathrm{}_n^p\text{-basis.}$$
This proves that 2 holds. Now assume $`p>1`$.
$`45`$. Let $`\epsilon >0`$ and choose $`fX`$ with $`f_p=1`$ and $`P𝒫(𝒩)`$ with $`\tau (P)<\epsilon `$ so that (5.6) holds. Then of course
(5.24)
$$f(IP)_1<\epsilon .$$
Now letting $`\frac{1}{p}+\frac{1}{q}=1`$,
(5.25)
$$fP_1f_pP_q\epsilon ^{1/q}\text{ by Hölder’s inequality.}$$
Thus
(5.26)
$$f_1<\epsilon +\epsilon ^{1/q}.$$
Since $`f_p=1`$ and $`\epsilon >0`$ is arbitrary, 5 holds.
$`53`$. Suppose 5 holds, yet 3 were false. Choose $`0<\delta `$ so that
(5.27)
$$\stackrel{~}{\omega }_p(f,\delta )\frac{1}{2}\text{ for all }f_a(X).$$
Let $`fX`$, $`f_p=1`$. By the last statement of Lemma 2.3, choose $`P`$ a spectral projection for $`|f|`$ so that $`fP𝒩`$ with
(5.28)
$$f(IP)_p\frac{1}{2}\text{ and }fP_{\mathrm{}}\delta ^{1/p}.$$
Then
(5.29)
$$\begin{array}{cc}\hfill \frac{1}{2^p}fP_p^p& =\tau (|f|^pP)\text{ (since }P|f|\text{)}\hfill \\ & =\tau (|f||f|^{p1}P)\hfill \\ & f_1\delta ^{1\frac{1}{p}}.\hfill \end{array}$$
That is,
(5.30)
$$f_12^{1/p}\delta ^{\frac{1}{p}1}\stackrel{\text{def}}{=}C.$$
(5.30) yields that $`g_pCg_1`$ for all $`gX`$; i.e., 5 does not hold, a contradiction. This completes the proof of the theorem. ∎
The final result of this section deals with the Banach-Saks property.
###### Definition 5.5.
Let $`X`$ be a Banach space, and $`1<p<\mathrm{}`$.
(a) Let $`(x_n)`$ be a weakly null sequence in $`X`$. $`(x_n)`$ is called
* a Banach-Saks sequence if
(5.31)
$$\underset{n\mathrm{}}{lim}n^1\underset{j=1}{\overset{n}{}}y_j=0\text{ for all subsequences }(y_j)\text{ of }(x_j).$$
* a $`p`$-Banach-Saks sequence if
(5.32)
$$\text{ there is a }C<\mathrm{}\text{ so that }\overline{\mathrm{lim}}_n\mathrm{}n^{1/p}\underset{j=1}{\overset{n}{}}y_jC\text{ for all subsequences }(y_j)\text{ of }(x_j)\text{.}$$
* a strong $`p`$-Banach-Saks sequence if
(5.33)
$$\underset{n\mathrm{}}{lim}n^{1/p}\underset{j=1}{\overset{n}{}}y_j=0\text{ for all subsequences }(y_j)\text{ of }(x_j)\text{.}$$
(b) $`X`$ is said to have the Banach-Saks property (resp. the $`p`$-Banach-Saks property) (resp. the strong $`p`$-Banach-Saks property) if every weakly null sequence in $`X`$ has a Banach-Saks (resp. $`p`$-Banach-Saks) (resp. strong $`p`$-Banach-Saks) subsequence.
The classical paper of Banach-Saks \[BS\] yields that commutative $`L^p`$ spaces have the $`p`$-Banach-Saks property, for $`1<p2`$; the fact that $`L^1`$-spaces have the Banach-Saks property was proved later by Szlenk \[Sz\]. Our last result yields in particular a generalization to the spaces $`L^p(𝒩)`$, $`𝒩`$ finite. Most of its assertions follow very quickly from our previous results.
###### Proposition 5.6.
Let $`𝒩`$ be a finite von-Neumann algebra and $`1<p<2`$.
* $`L^1(𝒩)`$ has the Banach-Saks property and $`L^p(𝒩)`$ has the $`p`$-Banach-Saks property.
* A weakly null sequence $`(f_n)`$ in $`L^p(𝒩)`$ has a strong $`p`$-Banach-Saks subsequence if $`(|f_n|^p)`$ is uniformly integrable. If $`(|f_n|^p)`$ is not uniformly integrable, $`(f_n)`$ has a subsequence $`(f_n^{})`$ so that for some $`c>0`$ and all subsequences $`(y_j)`$ of $`(f_j^{})`$,
(5.34)
$$\underset{¯}{\mathrm{lim}}n^{1/p}\underset{j=1}{\overset{n}{}}y_jc.$$
* A closed linear subspace $`X`$ of $`L^p(𝒩)`$ has the strong $`p`$-Banach-Saks property if and only if $`X`$ has no subspace isomorphic to $`\mathrm{}^p`$.
###### Proof.
Corollary 4.4 together with Proposition 2.5 yields that $`L^1(𝒩)`$ has the Banach-Saks property. It also yields the first assertion in 2. Suppose $`(|f_n|^p)`$ is not uniformly integrable and assume (without loss of generality) that $`f_n_p1`$ for all $`n`$. Applying Corollary 3.4 and Lemma 4.3, we may choose a subsequence $`(f_n^{})`$ of $`(f_n)`$ so that for some $`C1`$,
(5.35)
$$(f_n^{}r_n)\text{ is }C\text{-equivalent to the usual }\mathrm{}^p\text{-basis.}$$
and
(5.36)
$$(f_{n_1}^{},\mathrm{},f_{n_{2^k}}^{})\text{ is 4-unconditional for all }kn_1<n_2<\mathrm{}<n_{2^k}.$$
Suppose $`(y_j)`$ is a subsequence of $`(f_j^{})`$. Then it follows that for all $`k`$,
(5.37)
$$(y_{k+1},\mathrm{},y_{k+2^k})\text{ is }(4C)\text{-equivalent to the }\mathrm{}_{2^k}^p\text{-basis.}$$
Let $`n`$ be a “large” integer and choose $`k`$ with
(5.38)
$$2^{k1}n<2^k.$$
Then
(5.39)
$$\underset{j=k+1}{\overset{n}{}}y_j\frac{(nk)^{1/p}}{4C}\text{ by (}\text{5.37}\text{)}.$$
Thus
(5.40)
$$\underset{j=1}{\overset{n}{}}y_j_p\frac{(nk)^{1/p}}{4C}k\frac{(n\mathrm{log}_2n1)^{1/p}}{4C}\mathrm{log}_2n1.$$
Hence
(5.41)
$$\underset{¯}{\mathrm{lim}}_n\mathrm{}n^{1/p}\underset{j=1}{\overset{n}{}}y_j_p\frac{1}{4C}.$$
This completes the proof of assertion 2 of the Proposition. But we also have that
(5.42)
$$\underset{j=k+1}{\overset{n}{}}y_j_p4C(nk)^{1/p}\text{ by (}\text{5.37}\text{)},$$
and so
(5.43)
$$\underset{j=1}{\overset{n}{}}y_j_p4C(n\mathrm{log}_2n)^{1/p}+\mathrm{log}_2n+1,$$
thus
(5.44)
$$\overline{\mathrm{lim}}_n\mathrm{}n^{1/p}\underset{j=1}{\overset{n}{}}y_j_p4C.$$
This proves that $`L^p(𝒩)`$ has the $`p`$-Banach-Saks property, for any weakly null sequence $`(f_n)`$ in $`L^p(𝒩)`$ either has $`(|f_n|^p)`$ uniformly integrable (and hence a strong $`p`$-Banach-Saks subsequence), or a subsequence $`(f_n^{})`$ as above.
The final assertion of the Proposition follows immediately from Theorem 5.4 and assertion 2. ∎
###### Remark.
Of course Hilbert space has the 2-Banach Saks property. Actually, it can be shown that $`L^p(𝒩)`$ has the 2-Banach Saks property for $`p>2`$ and $`𝒩`$ finite, and this is best possible (in general). Indeed, if $`(f_j)`$ is a weakly null sequence in $`L^p(𝒩)`$, then if $`f_j_p0`$, $`(f_j)`$ trivially has a $`p`$-Banach Saks subsequence; the same is true if $`(f_j)`$ has a subsequence equivalent to the $`\mathrm{}^p`$-basis (and of course a $`p`$-Banach Saks sequence is a 2-Banach Saks sequence). Otherwise, combining arguments in \[S1\] Theorem 2.4 with the arguments in Proposition 5.6, we see that there exists a subsequence $`(f_j^{})`$ of $`(f_j)`$ such that its all subsequences $`(y_n)`$ are 2-Banach Saks.
We conclude this section with a brief discussion of the following open
###### Problem.
Let $`1<p<2`$ and $`(f_n)`$ be a seminormalized weakly null sequence in $`L^p(𝒩)`$ ($`𝒩`$ a finite von Neumann algebra) such that $`(|f_n|^p)`$ is not uniformly integrable. Does $`(f_n)`$ have a subsequence equivalent to the usual $`\mathrm{}^p`$ basis?
As pointed out previously, the answer is affirmative if $`(f_n)`$ has an unconditional subsequence. Actually, it can be proved that if $`(f_n)`$ satisfies the hypotheses of this Problem, it has a subsequence $`(f_n^{})`$ which dominates the $`\mathrm{}^p`$-basis and moreover has spreading model equivalent to the $`\mathrm{}^p`$-basis. (The last assertion follows immediately from our proof of Proposition 5.6.) It may then be shown that the above Problem is equivalent to the following one (in which the hypothesis concerning $`(|f_n|^p)`$ no longer enters).
###### Problem.
Let $`(f_n)`$ be a seminormalized basic sequence in $`L^p(𝒩)`$, $`p`$ and $`𝒩`$ as above. Does $`(f_n)`$ have a subsequence $`(f_n^{})`$ which is dominated by the $`\mathrm{}^p`$-basis? i.e., such that $`c_jf_j^{}`$ converges in $`L^p(𝒩)`$ whenever $`|c_j|^p<\mathrm{}`$?
## 6. The Banach isomorphic classification of the spaces $`L^p(𝒩)`$ for $`𝒩`$ hyperfinite semi-finite
We first fix some notation. Let $`1p<\mathrm{}`$. We let $`S_p=(_{n=1}^{\mathrm{}}C_p^n)_p`$ ($`=L^p(M_n)_{\mathrm{}}`$). To avoid confusion, we denote by $`L_p_pX`$ the Bochner space $`L_p(X,m)`$, where $`m`$ is Lebesgue measure and $`X`$ is a Banach space. Thus e.g., $`L_p_pC_p=L_p(C_p)=L^p(L^{\mathrm{}}(m)\overline{}B(\mathrm{}^2))`$. $``$ denotes the hyperfinite type II factor, and $`L^p()_pC_p`$ denotes $`L^p(\overline{}B(\mathrm{}^2))`$ (so $`\overline{}B(\mathrm{}^2)`$ is the hyperfinite type II factor).
The main motivating result of this section is as follows.
###### Theorem 6.1.
Let $`𝒩`$ be a hyperfinite semi-finite infinite dimensional von-Neumann algebra, and let $`1p<\mathrm{}`$, $`p2`$. Then $`L^p(𝒩)`$ is (completely) isomorphic to precisely one of the following thirteen Banach spaces.
$$\mathrm{}_p,S_p,L_p,C_p,S_pL_p,C_pL_p,L_p_pS_p,C_p(L_p_pS_p)$$
$$L^p(),L_p_pC_p,C_pL^p(),L^p()(L_p_pC_p),L^p()_pC_p.$$
Theorem 6.1 is an immediate consequence of the following finer result concerning embeddings.
###### Theorem 6.2.
Let $`1p<2`$. If $`𝒩`$ is as in 6.1, then $`L^p(𝒩)`$ is (completely) isomorphic to one of the thirteen spaces in the tree in Figure 1. If $`XY`$ are listed in the tree, then $`X`$ is Banach isomorphic to a subspace of $`Y`$ if and only if $`X`$ can be joined to $`Y`$ through a descending branch (in which case $`X`$ is completely isometric to a subspace of $`Y`$).
###### Remark.
In the language of graph theory, Theorem 6.2 asserts that the tree in Figure 1 is the Hasse diagram for the partially ordered set consisting of the equivalence classes of $`L^p(𝒩)`$ under Banach isomorphism (over $`𝒩`$ as in 6.1), with the order relation: $`[X][Y]`$ provided $`X`$ is isomorphic to a subspace of $`Y`$.
Parts of Theorem 6.2 require previously known results, some of which are very recent. It is established in \[S2\] that the first nine spaces in the list in Theorem 6.1 are isomorphically distinct when $`p=1`$, and exhaust the list of the possible Banach isomorphism types of $`L^p(𝒩)`$ for $`𝒩`$ type I ($`𝒩`$ as in 6.1), $`p2`$.
Theorem 6.2 yields the new result in the type I case: $`L_p_pC_p`$ does not embed in $`C_p(L_p_pS_p)`$ for $`1p<2`$; (another new result in this case, that $`C_p`$ does not embed in $`L_p_pS_p`$, follows immediately from Corollary 1.2); the other embedding results stated in 6.2 for the type I case are given in \[S2\]. We give here a new proof of the particular case that $`L_p_pS_p`$ does not embed in $`L_pC_p`$, using the Main Result of this paper.
We first proceed with the non-embedding results required for Theorem 6.2. The following theorem is crucial.
###### Theorem 6.3.
Let $`𝒩`$ be a finite von Neumann algebra and $`1p<2`$. Then $`L_p_pC_p`$ is not isomorphic to a subspace of $`C_pL^p(𝒩)`$.
We now fix $`1p<2`$ for the remainder of this section.
We first require
###### Lemma 6.4.
Let $`T:L_pC_p`$ be a given bounded linear operator, and let $`\epsilon >0`$. Then there exists an $`fL_p`$ with $`f`$ $`\{1,1\}`$-valued so that $`Tf<\epsilon `$.
###### Sublemma.
The conclusion of 6.4 holds, replacing $`C_p`$ by $`\mathrm{}^2`$ in its hypotheses.
###### Proof.
Fix $`n`$ a positive integer. Using the generalized parallelogram identity,
(6.1)
$$\begin{array}{cc}\hfill av_\pm T\underset{j=1}{\overset{n}{}}\pm \chi _{[\frac{j1}{n},\frac{j}{n})}_2^2& =\underset{j=1}{\overset{n}{}}T(\chi _{[\frac{j1}{n},\frac{j}{n})})_2^2\hfill \\ \\ & T^2\underset{j=1}{\overset{n}{}}\chi _{[\frac{j1}{n},\frac{j}{n})}_p^2\hfill \\ \\ & =T^2\frac{n}{n^{2/p}}=T^2\frac{1}{n^{2/p1}}.\hfill \end{array}$$
It follows that we may choose $`\eta _j=\pm 1`$ for all $`j`$ so that
(6.2)
$$T\left(\underset{j=1}{\overset{n}{}}\eta _j\chi _{[\frac{j1}{n},\frac{j}{n})}\right)_2\frac{T}{n^{\frac{1}{p}\frac{1}{2}}}.$$
Now simply choose $`n`$ so that $`\frac{T}{n^{\frac{1}{p}\frac{1}{2}}}<\epsilon `$ and let $`f=_{j=1}^n\eta _j\chi _{[\frac{j1}{n},\frac{j}{n})}`$. ∎
###### Proof of Lemma 6.4.
Let $`(e_{ij})`$ be the matrix units basis for $`C_p`$, and define for each $`n`$,
(6.3)
$$H_n=[e_{ij}:1in\text{ and }1j<\mathrm{}\text{ or }1i<\mathrm{}\text{ and }1jn].$$
Let $`P_n`$ be the natural basis projection onto $`H_n`$; i.e., $`P_n:C_pC_p`$ is the projection with $`P_n(e_{ij})=0`$ if $`e_{ij}H_n`$, $`P_n(e_{ij})=e_{ij}`$ if $`e_{ij}H_n`$ (so $`P_n2`$). Then $`H_n`$ is isomorphic to $`\mathrm{}^2`$, so by the sub-lemma we may choose $`f_n`$ in $`L^p`$ with $`f_n`$ $`\{1,1\}`$-valued and
(6.4)
$$P_nTf_n\frac{1}{2^n}.$$
We claim that
(6.5)
$$\underset{n\mathrm{}}{lim}Tf_n=0.$$
Of course (6.5) yields the conclusion of the Lemma. Suppose (6.5) were false. It follows that $`(f_n)`$ has a subsequence $`(f_n^{})`$ so that
(6.6)
$$(Tf_n^{})\text{ is equivalent to the usual }\mathrm{}^p\text{-basis}$$
and
(6.7)
$$(f_n^{})\text{ converges weakly in }L^2.$$
((6.6) follows because $`(f_n^{})`$ may be chosen to be a small perturbation of a “block-off-diagonal sequence”, by 6.4).
Of course $`(f_n^{})`$ converges weakly in $`L^p`$ as well, hence $`(Tf_n^{})`$ also converges weakly, a contradiction when $`p=1`$ since then $`(Tf_n^{})`$ is equivalent to the $`\mathrm{}^1`$-basis.
When $`p>1`$, letting $`f`$ be the weak limit of $`(f_n)`$, we have that $`Tf=0`$ since $`Tf_n^{}0`$ weakly. Moreover $`f_{\mathrm{}}2`$, so letting $`f_n^{\prime \prime }=f_n^{}f`$ for all $`n`$, $`(f_n^{\prime \prime })`$ is a uniformly bounded weakly null sequence in $`L^p`$ with $`(Tf_n^{\prime \prime })=(Tf_n^{})`$ equivalent to the $`\mathrm{}^p`$-basis. Finally, since $`(f_n^{\prime \prime })`$ is also semi-normalized in $`L^p`$, $`(f_n^{\prime \prime })`$ has a subsequence $`(g_n)`$ equivalent to the usual $`\mathrm{}^2`$-basis. (Indeed, we may choose $`(g_n)`$ equivalent to the $`\mathrm{}^2`$-basis in $`L^2`$-norm, and unconditional. But then since $`L^p`$ has cotype 2, $`(g_n)`$ is equivalent to the $`\mathrm{}^2`$-basis in the $`L^p`$-norm). Still, $`(Tg_n)`$ is equivalent to the $`\mathrm{}^p`$-basis; this is impossible since $`p<2`$. ∎
We now apply our Main Result and Lemma 6.4, to give the
###### Proof of Theorem 6.3.
Suppose to the contrary that $`𝒩`$ is a finite von Neumann algebra and $`T:L_p_pC_pC_pL^p(𝒩)`$ is an isomorphic embedding. Of course we may assume that $`T=1`$; let $`\epsilon =T^1^1`$. Thus we have
(6.8)
$$Tf\epsilon f\text{ for all }fL_p_pC_p.$$
Let $`P`$ be the projection of $`C_pL^p(𝒩)`$ onto $`C_p`$ with kernel $`L^p(𝒩)`$, and set $`Q=IP`$. Also, for each $`i`$ and $`j`$, let $`Q_{ij}`$ be the natural projection of $`L_p_pC_p`$ onto the space
(6.9)
$$E_{ij}\stackrel{\text{def}}{=}\{fe_{ij}:fL_p\}.$$
(As before, $`e_{ij}`$ denotes the $`i,j^{\text{th}}`$ matrix unit for $`C_p`$. Visualizing $`C_p`$ as matrices of scalars and $`L_p_pC_p`$ as all matrices $`(f_{ij})`$ of functions in $`L_p`$ with
$$(f_{ij})=\left((f_{ij}(w))_{C_p}^p𝑑w\right)^{1/p}<\mathrm{},$$
then $`Q_{ij}((f_k\mathrm{}))=f_{ij}e_{ij}`$. $`E_{ij}`$ is just the space of matrices with all entries zero except in the $`ij^{\text{th}}`$ slot). Now fix $`i`$ and $`j`$. Of course $`E_{ij}`$ is isometric to $`L_p`$.
Thus by Lemma 6.4, we may choose $`f_{ij}L_p`$ with $`f_{ij}`$ $`\{1,1\}`$-valued so that
(6.10)
$$PTf_{ij}e_{ij}<\frac{\epsilon }{2^{i+j+2}}.$$
Now letting $`X=[f_{ij}e_{ij}:i,j=1,2,\mathrm{}]`$, then $`X`$ is a 1-$`GC_p`$ space, in the terminology of the Introduction. That is, every row and column of $`(f_{ij}e_{ij})`$ is 1-equivalent to the $`\mathrm{}^2`$ basis, while every generalized diagonal is 1-equivalent to the $`\mathrm{}^p`$ basis. Hence $`X`$ is not isomorphic to a subspace of $`L^p(𝒩)`$ by our Main Theorem (i.e. Corollary 1.2). However
(6.11)
$$QT|X\text{ is an isomorphic embedding.}$$
Indeed, if $`x=c_{ij}(f_{ij}e_{ij})`$ with only finitely many $`c_{ij}`$’s non zero, and $`x=1`$, then $`|c_{ij}|1`$ for all $`i`$ and $`j`$ (since the $`Q_{ij}`$’s are contractive and $`f_{ij}=1`$ for all $`i`$ and $`j`$), and so
(6.12)
$$\begin{array}{cc}\hfill PTx& \underset{i,j}{\mathrm{max}}|c_{ij}|\underset{i,j}{}T(f_{ij}e_{ij})\hfill \\ & \underset{i=1}{\overset{\mathrm{}}{}}\underset{j=1}{\overset{\mathrm{}}{}}\frac{\epsilon }{2^{i+j+2}}=\frac{\epsilon }{2}\hfill \end{array}$$
using (6.10) and our assumption that $`T`$ is a contraction. Hence
(6.13)
$$QTx\frac{\epsilon }{2}\text{ by (}\text{6.8}\text{).}$$
This proves (6.11), and completes the proof by contradiction. ∎
Our localization result, Corollary 1.4, and the preceding proof, yield an alternate proof of the following result, obtained in \[S2\].
###### Proposition 6.5.
$`L^p_pS_p`$ is not isomorphic to a subspace of $`C_pL_p`$.
###### Proof.
We have that $`L^p_pS_p`$ is (linearly isometric to) $`(_{n=1}^{\mathrm{}}L_p_pC_p^n)_p`$. Thus it suffices to prove that
(6.14)
$$\underset{n\mathrm{}}{lim}\lambda _n=\mathrm{}$$
where
(6.15)
$$\lambda _n=inf\{d(L_p_pC_p^n,Y):Y\text{ is a subspace of }C_pL_p\}$$
and “$`d`$” denotes the Banach Mazur distance-coefficient (defined just preceding Corollary 1.4).
Now fix $`n`$, and let $`T:L_p_pC_p^nYC_pL_p`$ be an isomorphic embedding onto $`Y`$, with
(6.16)
$$T=1\text{ and }T^12\lambda _n.$$
Using the notation and reasoning in the proof of Theorem 6.3, and setting $`\epsilon =1/(2\lambda _n)`$, we may choose for each $`i`$ and $`j`$ with $`1i,jn`$, a $`\{1,1\}`$-valued $`f_{ij}L^p`$ satisfying (6.10). We thus obtain that $`PT|X\epsilon /2`$ by (6.12). Hence for all $`xX`$,
(6.17)
$$QT(x)\left(\frac{1}{2\lambda _n}\frac{\epsilon }{2}\right)x=\frac{1}{4\lambda _n}x$$
using also (6.16). That is, setting $`Z=QT(X)`$, we have that
(6.18)
$$d(X,Z)4\lambda _n.$$
Now $`X`$ is a 1-$`GC_p^n`$-space; thus
(6.19)
$$4\lambda _n\beta _{n,1}\text{ for all }n$$
(in the notation of Corollary 1.4), so (6.14) holds by Corollary 1.4. ∎
We also require the following rather deep result, due to M. Junge \[J\].
###### Theorem 6.6.
$`C_q`$ is isomorphic to a subspace of $`L^p()`$ for all $`p<q<2`$.
Finally, we require the following (unpublished) result, due to G. Pisier and Q. Xu \[PX2\].
###### Lemma 6.7.
Let $`X`$ be a (closed linear) subspace of $`L_p_pC_p`$. Then either $`X`$ embeds in $`L_p`$ or $`\mathrm{}^p`$ embeds in $`X`$.
For the sake of completeness, we sketch a proof. First, we give an important, quick consequence of these last two results.
###### Corollary 6.8.
$`L^p()`$ is not isomorphic to a subspace of $`L_p_pC_p`$.
###### Proof.
By Theorem 6.6, it suffices to prove that $`C_q`$ does not embed in $`L_p_pC_p`$ if $`p<q<2`$. If $`C_q`$ did embed, then since it does not embed in $`L_p`$, it would have a subspace isomorphic to $`\mathrm{}^p`$, by Lemma 6.7. However it is a standard fact that every infinite-dimensional subspace of $`C_p`$ is either isomorphic to $`\mathrm{}^2`$ or contains a subspace isomorphic to $`\mathrm{}^p`$, a contradiction.∎
We next sketch the proof of Lemma 6.7 (which also yields the above mentioned standard fact).
Let $`(x_{ij})`$ be a given matrix in a linear space $`X`$. Call a sequence $`(f_k)`$ in $`X`$ a generalized block diagonal of $`(x_{ij})`$ if there exist $`i_1<i_2<\mathrm{}`$ and $`j_1<j_2<\mathrm{}`$ so that for all $`k`$,
(6.20)
$$f_k[x_{ij}:i_ki<i_{k+1}\text{ and }j_kj<j_{k+1}].$$
Now if $`(f_k)`$ is a generalized block diagonal of the matrix $`(e_{ij})`$ consisting of non-zero terms, $`e_{ij}`$ the matrix units for $`C_p`$ (as above), then $`(f_k/f_k)`$ is isometrically equivalent to the $`\mathrm{}^p`$-basis. But then it follows immediately that if $`(f_k)`$ is a normalized generalized block diagonal of $`(\text{1}e_{ij})`$ (in $`L^p_pC_p`$) consisting of non-zero terms, $`(f_k)`$ is also isometrically equivalent to the $`\mathrm{}^p`$-basis. Indeed, for any scalars $`c_1,c_2,\mathrm{}`$ with only finitely many non-zero terms, and any $`0c_j1`$,
(6.21)
$$c_jf_j(w)_{C_p}^p=|c_j|^p|f_j(w)|^p.$$
Hence
(6.22)
$$c_jf_j^p=c_jf_j(w)_{C_p}^p𝑑w=|c_j|^p.$$
Now fix $`n`$, and let $`H_n`$ be the subspace of $`C_p`$ defined in the proof of Lemma 6.4 (specifically, in (6.3)). Standard results yield that $`L^p_pH_n`$ embeds in $`L^p`$ (actually, $`L^p_pH_n`$ is isomorphic to $`L^p`$ if $`p>1`$), and of course $`IP_n`$ is a projection onto $`L^p_pH_n`$ with $`IP_n2`$ ($`P_n`$ as defined in the proof of 6.4). Now let $`X`$ be as in Lemma 6.7, and suppose $`X`$ does not embed in $`L_p`$. Then for each $`n`$, we may choose an $`x_nX`$ with
(6.23)
$$x_n=1\text{ and }(IP_n)x_n<\frac{1}{2^n}.$$
But it follows that for any $`fL_p_pC_p`$,
(6.24)
$$(IP_n)(f)f\text{ as }n\mathrm{}.$$
A standard travelling hump argument now yields a normalized generalized block diagonal $`(f_k)`$ of $`(\text{1}e_{ij})`$ and a subsequence $`(x_j^{})`$ of $`(x_j)`$ so that
(6.25)
$$x_k^{}f_k<\frac{1}{2^k}\text{ for all }k.$$
It follows immediately that $`(x_k^{})`$ is equivalent to the $`\mathrm{}^p`$-basis. ∎
###### Remark.
The last part of this proof also yields the fact (due to Y. Friedman \[F\]) that if $`X`$ is an infinite-dimensional subspace of $`C_p`$, then $`X`$ is isomorphic to $`\mathrm{}^2`$ or $`\mathrm{}^p`$ embeds in $`X`$. Indeed, assuming $`X`$ is not isomorphic to $`\mathrm{}^2`$, then since $`H_n`$ is isomorphic to $`\mathrm{}^2`$, we obtain for each $`n`$ and $`x_nX`$ with $`x_n=1`$ and $`P_nx_n<\frac{1}{2^n}`$. Again we then obtain a normalized block diagonal $`(f_k)`$ of $`(e_{ij})`$ and a subsequence $`(x_j^{})`$ of $`(x_j)`$ satisfying (6.25), and then $`(x_k^{})`$ is equivalent to the $`\mathrm{}^p`$ basis.
We now give the last and perhaps most delicate of the needed non-embedding results; its proof requires Theorem 4.1, the “fine” version of our Main Result.
###### Theorem 6.9.
Let $`𝒩`$ be a finite von Neumann algebra. Then $`L^p()_pC_p`$ is not isomorphic to a subspace of $`L^p(𝒩)(L_p_pC_p)`$.
We first give some notation used in the proof. As always, $`e_{ij}`$’s denote the matrix units for $`C_p`$. Thus $`L^p()_pC_p=L^p(\overline{}B(\mathrm{}^2))=`$ the closed linear span of the elementary tensors $`fe_{ij}`$, $`fL^p()`$, $`i`$ and $`j`$ arbitrary. We denote also the norm on $`L^p()_pC_p`$ as $`_p`$. If $`X`$ is a closed linear subspace of $`L^p()`$,
(6.26)
$$X_pC_p\stackrel{\text{def}}{=}[xe_{ij}:xX,i,j]$$
(where the closed linear span above is taken in $`L^p()_pC_p`$). Next, we need expressions for the norm on $`L^p()\mathrm{Row}`$, $`L^p()`$ Column. We easily see that given $`x_1,\mathrm{},x_n`$ in $`L^p()`$, then for any $`i`$,
(6.27)
$$\underset{j=1}{\overset{n}{}}x_je_{ij}_p=\left(\underset{j=1}{\overset{n}{}}x_jx_j^{}\right)^{1/2}_p$$
and
(6.28)
$$\underset{j=1}{\overset{n}{}}x_je_{ji}_p=\left(\underset{j=1}{\overset{n}{}}x_j^{}x_j\right)^{1/2}_p.$$
Evidently (6.27) and (6.28) show that if we consider a matrix of the form $`(x_{ij}e_{ij})`$ with $`x_{ij}`$ non-zero elements of $`L^p()`$ for all $`i`$ and $`j`$, then all rows and columns of this matrix are 1-unconditional sequences.
The next result is really a “localization” of Lemma 3.1 (and could be formulated instead for $`L^p(𝒩)`$, $`𝒩`$ any finite von Neumann algebra).
###### Lemma 6.10.
Let $`X`$ be a closed linear subspace of $`L^p()`$ containing no subspace isomorphic to $`\mathrm{}^p`$. Then given $`\epsilon >0`$, there is an $`N`$ so that given any $`nN`$ and $`x_1,\mathrm{},x_n`$ in $`_a(X)`$,
(6.29)
$$n^{1/p}\left(\underset{i=1}{\overset{n}{}}x_ix_i^{}\right)^{1/2}_p\epsilon \text{ and }n^{1/p}\left(\underset{i=1}{\overset{n}{}}x_i^{}x_i\right)^{1/2}_p\epsilon .$$
###### Proof.
Let $`\tau `$ be the normal faithful tracial state in $``$. By Theorem 5.4, $`\{|x|^p:x_a(X)\}`$ is uniformly integrable. Let $`\eta >0`$, to be decided later. Choose $`\delta >0`$ so that
(6.30)
$$\omega (|x|^p,\delta )\eta ^p\text{ for all }x_a(X).$$
Let $`x_1,\mathrm{},x_n`$ be elements of $`_a(X)`$. By the final statement of Lemma 2.3, we may choose for each $`j`$ a $`P_j𝒫()`$ so that $`x_jP_j`$ with
(6.31)
$$x_jP_j_{\mathrm{}}\delta ^{1/p}\text{ and }x_j(IP_j)_p\eta .$$
Then
(6.32)
$$\begin{array}{cc}\hfill \left(\underset{j=1}{\overset{n}{}}x_jx_j^{}\right)^{1/2}_p& =\underset{j=1}{\overset{n}{}}x_je_{ij}_p\text{ by (}\text{6.27}\text{)}\hfill \\ & \underset{j=1}{\overset{n}{}}x_jP_je_{1j}_p+\underset{j=1}{\overset{n}{}}x_j(IP_j)e_{1j}_p.\hfill \end{array}$$
Since $`(x_j(IP_j)e_{1j})_{j=1}^n`$ is 1-unconditional and $`L^p()_pC_p`$ is type $`p`$ with constant one,
(6.33)
$$\begin{array}{cc}\hfill \underset{j=1}{\overset{n}{}}x_j(IP_j)e_{1j}_p& \left(\underset{j=1}{\overset{n}{}}x_j(IP_j)_p^p\right)^{1/p}\hfill \\ & \eta n^{1/p}\text{ by (}\text{6.31}\text{)}.\hfill \end{array}$$
Now
(6.34)
$$\begin{array}{cc}\hfill \underset{j=1}{\overset{n}{}}x_jP_je_{1j}_p& =\left[\tau \left(\underset{j=1}{\overset{n}{}}x_jP_jx_j^{}\right)^{p/2}\right]^{1/p}\hfill \\ & \left[\tau \left(\underset{j=1}{\overset{n}{}}x_jP_jx_j^{}\right)\right]^{1/2}\text{ (since }p<2\text{)}\hfill \\ & n^{1/2}\delta ^{1/p}\text{ by (}\text{6.31}\text{).}\hfill \end{array}$$
Thus (6.32)–(6.34) yield that
(6.35)
$$n^{1/p}\left(\underset{j=1}{\overset{n}{}}x_jx_j^{}\right)^{1/2}_p\eta +\frac{1}{n^{\frac{1}{p}\frac{1}{2}}}\delta ^{1/p}.$$
Evidently we now need only take $`\eta \frac{\epsilon }{2}`$; then choose $`N`$ so that $`N^{(\frac{1}{p}\frac{1}{2})}\delta ^{1/p}\frac{\epsilon }{2}`$; the identical argument for $`(x_i^{}x_i)_{i=1}^n`$ now yields that (6.29) holds for all $`nN`$. ∎
We may now easily obtain our final needed preliminary result. (See the Remark following Theorem 4.1 for the definition of: the rows or columns of a matrix contain $`\mathrm{}_n^p`$-sequences.)
###### Corollary 6.11.
Let $`X`$ be a closed linear subspace of $`L^p()`$ containing no subspace isomorphic to $`\mathrm{}^p`$, and let $`(x_{ij})`$ be a seminormalized matrix whose terms lie in $`X`$. Then the matrix $`(x_{ij}e_{ij})`$ in $`X_pC_p`$ has the following properties:
* Neither the rows nor the columns contain $`\mathrm{}_n^p`$-sequences.
* Every row and column is 1-unconditional.
* Every generalized diagonal is equivalent to the usual $`\mathrm{}^p`$ basis.
###### Proof.
(i) follows immediately from Lemma 6.10 and (6.27), and the latter also immediately yields (ii). If $`(f_i)`$ is a generalized diagonal of the matrix, then there exist projections $`P_1,P_2,\mathrm{}`$, $`Q_1,Q_2,\mathrm{}`$ in $`\overline{}B(\mathrm{}^2)`$ so that the $`P_j`$’s and the $`Q_j`$’s are pairwise orthogonal, with $`f_j=P_jf_jQ_j`$ for all $`j`$. (That is, $`(f_j)`$ is “right and left disjointly supported”.) It then follows that for any $`n`$ and scalars $`c_1,\mathrm{},c_n`$,
(6.36)
$$\underset{j=1}{\overset{n}{}}c_jf_j_p=\left(\underset{j=1}{\overset{n}{}}|c_j|^pf_j_p^p\right)^{1/p},$$
which immediately yields (iii) since $`(x_{ij}e_{ij})`$ is semi-normalized. ∎
We are finally prepared for the
###### Proof of Theorem 6.9.
Let $`p<q<2`$ and let $`X`$ be a subspace of $`L^p()`$ so that $`X`$ is isomorphic to $`C_q`$ (using Junge’s result, formulated as Theorem 6.6 above). We claim that $`X_pC_p`$ is not isomorphic to a subspace of $`L^p(𝒩)(L_p_pC_p)`$ (which of course proves Theorem 6.9). Suppose to the contrary that $`T:X_pC_pL^p(𝒩)(L_p_pC_p)`$ is an isomorphic embedding. Assume without loss of generality that $`T=1`$. Let $`\epsilon >0`$ be chosen so that $`Tf\epsilon f`$ for all $`fX_pC_p`$. Let $`P`$ denote the projection of $`L^p(𝒩)(L_p_pC_p)`$ onto $`L^p(𝒩)`$, with kernel $`L_p_pC_p`$; and set $`Q=IP`$. Now fix $`i`$ and $`j`$. Then of course $`Xe_{ij}`$ is isometric to $`X`$. Thus by Lemma 6.7, $`QT|(Xe_{ij})`$ cannot be an isomorphic embedding (that is, $`C_q`$ does not embed in $`L_p_pC_p`$). Hence we may choose $`x_{ij}X`$ with
(6.37)
$$x_{ij}=1\text{ and }QT(x_{ij}e_{ij})<\frac{\epsilon }{2^{i+j+2}}.$$
Now let $`Y=[x_{ij}e_{ij}:i,j=1,2,\mathrm{}]`$. Since $`\mathrm{}^p`$ does not embed in $`X`$, the conclusion of Corollary 6.11 holds for the matrix $`(x_{ij}e_{ij})`$.
It follows from (6.37) that
(6.38)
$$QT|Y<\frac{\epsilon }{2}.$$
Hence we obtain that
(6.39)
$$PT(y)\frac{\epsilon }{2}y\text{ for all }yY.$$
Thus $`Y`$ is isomorphic to a subspace $`Z`$ of $`L^p(𝒩)`$. Let $`z_{ij}=PT(x_{ij}e_{ij})`$ for all $`i`$ and $`j`$. Now since $`PT|Y`$ is an isomorphism, Corollary 6.11 yields that there is a $`u`$ so that every row and column of $`(z_{ij})`$ is $`u`$-conditional, every generalized diagonal of $`(z_{ij})`$ is equivalent to the $`\mathrm{}^p`$-basis, yet neither the rows nor the columns of $`(z_{ij})`$ contain $`\mathrm{}_n^p`$-sequences. This is impossible by Theorem 4.1. ∎
The following result is an immediate consequence of Theorem 6.9 and known structural results for von-Neumann algebras.
###### Corollary 6.12.
Let $`𝒩,`$ be von Neumann algebras so that $``$ has a direct summand of type II or of type III. If $`L^p()`$ is Banach isomorphic to a subspace of $`L^p(𝒩)`$, then also $`𝒩`$ has a direct summand of type II or of type III.
###### Proof.
The hypotheses imply (via known results, cf. \[HS\]) that $`\overline{}B(\mathrm{}^2)`$ is isomorphic to a von Neumann subalgebra of $``$, which is the range of a normal conditional expectation, whence $`L^p()_pC_p`$ is completely isometric to a subspace of $`L^p()`$. Since $`L^p()C_p`$ is separable, we can assume without loss of generality that $`𝒩`$ acts on a separable Hilbert space. Then if $`𝒩`$ fails the conclusion, there exists a finite von Neumann algebra $`\stackrel{~}{𝒩}`$ so that $`𝒩`$ is isomorphic to a subalgebra of $`\stackrel{~}{𝒩}(L^{\mathrm{}}\overline{}B(\mathrm{}^2))`$, and hence $`L^p(𝒩)`$ is completely isometric to a subspace of $`L^p(\stackrel{~}{𝒩})(L_p_pC_p)`$. But then $`L^p()`$ does not Banach embed in $`L^p(𝒩)`$, since $`L^p()_pC_p`$ does not embed in $`L^p(\stackrel{~}{𝒩})(L_p_pC_p)`$ by Theorem 6.9. ∎
###### Remark.
Of course Corollary 6.8 (i.e., the results of Junge and Pisier-Xu cited above) also immediately yields that if $``$ and $`𝒩`$ are von Neumann algebras so that $``$ has a type II<sub>1</sub> summand, and $`L^p()`$ embeds in $`L^p(𝒩)`$, then $`𝒩`$ must have also have a summand of type II or type III. Combining these two results, we have that if $`L^p()`$ is Banach isomorphic to a subspace of $`L^p(𝒩)`$ and $``$ has no type III summand, then $`𝒩`$ has a direct summand of type at least as large as these of the summands of $`𝒩`$. It remains a most intriguing problem to see if one can eliminate the non-type III summand hypothesis in this statement.
We now complete the proof of Theorem 6.2. We shall formulate the “positive” results in the language of operator spaces; the reader unfamiliar with the relevant terms may just ignore the adjective “complete” in all the statements, for of course all positive operator space claims imply the pure Banach space ones. Given operator spaces $`X`$ and $`Y`$, let us say that $`X`$ completely contractively factors through $`Y`$ if $`X`$ is completely isometric to a subspace $`X^{}`$ of $`Y`$ such that there exists a completely contractive projection mapping $`Y`$ onto $`X^{}`$. Equivalently, there exist complete contractions $`U:XY`$ and $`V:YX`$ such that $`VU=I_X`$, $`I_X`$ the identity operator on $`X`$, that is,
(6.40) .
Now we easily see that
(6.41)
$$(L^p()L^p()\mathrm{})_p\text{ completely contractively factors through }L^p().$$
Indeed, simply let $`P_1,P_2,\mathrm{}`$ be pairwise orthogonal non-zero projections in $``$. As is well known, then $`P_iP_i`$ is isomorphic to $``$ and hence $`P_iL^p()P_i`$ is completely isometric to $`L^p()`$ for all $`i`$; then the map on $`L^p()`$ defined by $`fP_ifP_i`$ witnesses (6.41).
Since $`\overline{}`$ is isomorphic to $``$,
(6.42)
$$L^p()_pL^p()\stackrel{\text{def}}{=}L^p(\overline{})\text{ is completely isometric to }L^p().$$
Using (6.41) and (6.42), we may now easily see that if $`Y`$ is immediately below $`X`$ in the tree (and lying on a branch), then $`X`$ completely contractively factors through $`Y`$. Using the notation $`X\stackrel{cc}{}Y`$ to mean that $`X`$ completely contractively factors through $`Y`$, we see, e.g., that $`L_p\stackrel{cc}{}L^p()L_p_pC_n^p\stackrel{cc}{}L^p()_pC_p^n\stackrel{cc}{}L^p()_pL^p()`$, whence
$$L_p_pS_p=\left(\underset{n=1}{\overset{\mathrm{}}{}}(L_p_pC_p^n)\right)_p\stackrel{cc}{}\left(\underset{n=1}{\overset{\mathrm{}}{}}L_p_pL^p()\right)_p\stackrel{cc}{}L^p(),$$
i.e.,
(6.43)
$$L_p_pS_p\stackrel{cc}{}L^p().$$
Writing $`XY`$ to mean: $`X`$ is completely isometric to $`Y`$, we have
(6.44)
$$C_p(L_p_pS_p)\stackrel{cc}{}C_pL_p_pC_p\stackrel{cc}{}(L_pC_p)(L_pC_p)L_pC_p$$
(where we use $`\mathrm{}^p`$-direct sums).
$`X\stackrel{cc}{}Y`$ if $`X`$ is the level 7 space and $`Y`$ is the level 8 space, since the same argument for (6.41) yields also
(6.45)
$$\left((L^p()_pC_p)(L^p()_pC_p)\mathrm{}\right)\stackrel{cc}{}L^p()_pC_p.$$
The reader may now easily check that the remaining “positive” assertions on the tree. For the far deeper negative assertions, let us use the notation: $`X\hookrightarrow ̸Y`$ to mean that the Banach space $`X`$ is not isomorphic to a subspace of $`Y`$.
Now suppose $`XY`$ are on the tree and $`Y`$ cannot be connected to $`X`$ by a descending branch; we claim that $`X\hookrightarrow ̸Y`$.
It suffices to prove this assertion by showing by induction on $`j=2,3,\mathrm{}`$ that $`X`$ lies at level $`j`$ and
(6.46) there is a $`kj`$ so that $`Y`$ is at the $`k^{\text{th}}`$ level, but if $`Z`$ is a higher
level than $`k`$, connected to $`Y`$, $`ZX`$, then $`X`$ is connected to $`Z`$
and moreover there is no $`X^{}`$ connected to $`X`$ but not to $`Y`$ with
level $`X^{}<j`$
or
(6.47) $`Y`$ is at the $`(j1)^{\text{st}}`$ level, but if $`Y`$ is connected to $`Z`$ at level $`kj`$
with $`ZX`$, then $`X`$ is connected to $`Z`$ and moreover if $`Z`$ is connected
to $`X`$ with level $`Z<j`$, then $`Z`$ is connected to $`Y`$.
* $`S_p\hookrightarrow ̸L_p`$ is classical (and also follows from our Corollary 1.4). $`L_p\hookrightarrow ̸C_p`$ since $`\mathrm{}_qL_p`$ if $`p<q<2`$ but $`\mathrm{}_q\hookrightarrow ̸C_p`$.
* $`C_p\hookrightarrow ̸L^p()`$, the main result of the paper.
* $`L_p_pS_p\hookrightarrow ̸C_pL_p`$ by Proposition 6.5.
* $`L^p()\hookrightarrow ̸L_p_pC_p`$ by Corollary 6.8.
* $`L_p_pC_p\hookrightarrow ̸C_pL^p()`$ by Theorem 6.3.
* There is no $`Y`$ satisfying (6.46) or (6.47).
* Theorem 6.9 gives the one required non-embedding result.
This completes the proof of the final statement of Theorem 6.2. It remains to prove the first statement. This follows via the known type-decomposition and structure of hyperfinite von-Neumann algebras, and the following operator space version of the Pełczyński decomposition method (whose proof is exactly as Pełczyński’s proof for the Banach space case \[P\]; see also p.54 of \[LT\] and \[Ar\]).
###### Lemma 6.13.
Let $`X`$ and $`Y`$ be operator spaces so that
* each completely factors through the other
and so that either
* $`X`$ is completely isomorphic to $`XX`$ and $`Y`$ is completely isomorphic to $`YY`$
or
* $`X`$ is completely isomorphic to $`(XX\mathrm{})_q`$ for some $`q[1,\mathrm{}]`$.
Then $`X`$ and $`Y`$ are completely isomorphic.
(We say that $`X`$ completely factors through $`Y`$ if $`X`$ is completely isomorphic to a completely complemented subspace of $`Y`$.)
###### Corollary 6.14.
If $`(XX\mathrm{})_p`$ completely factors through the operator space $`X`$, then $`X`$ is completely isomorphic to $`(XX\mathrm{})_p`$.
End of the proof of Theorem 6.2. $`(XX\mathrm{})_p`$ completely contractively factors through $`X`$ for all of the 13 spaces $`X`$ listed in Theorem 6.2 (applying (6.41), (6.45), and the analogous results for $`C_p`$, $`L_p`$, and $`L_p_pC_p`$). Thus the conclusion of 6.14 applies.
Now let $`𝒩`$ be as in the statement of Theorem 6.2. If $`𝒩`$ is type I, then by the results in \[S2\] $`L^p(𝒩)`$ is completely isomorphic to one of the first nine spaces listed in Theorem 6.1, so assume that $`𝒩`$ is not type I. Then we have that
$$𝒩=𝒩_\text{I}𝒩_{\text{II}_1}𝒩_{\text{II}_{\mathrm{}}},$$
where for each $`i`$, $`𝒩_i=\{0\}`$ or $`𝒩_i`$ is a hyperfinite von Neumann algebra of type $`i`$, so that also $`𝒩_{\text{II}_1}𝒩_{\text{II}_{\mathrm{}}}0`$.
Now suppose that $`𝒩`$ is finite. It then follows from work of A. Connes \[C2\] that
(6.48)
$$𝒩_\text{I}𝒩_{\text{II}_1}\text{ is isomorphic to a von-Neumann subalgebra of }.$$
Indeed, by disintegration and Proposition 6.5 of \[C2\], any finite hyperfinite von Neumann algebra (with separable predual) is a countable $`\mathrm{}^{\mathrm{}}`$-direct sum of von Neumann algebras of the form $`𝒜\overline{}`$, where $`𝒜`$ is abelian and $``$ is either $`M_n`$ for some $`n<\mathrm{}`$ or $``$. But such an algebra $`𝒜\overline{}`$ can be realized as a sub-algebra of $``$; since also $`\overline{}`$ is isomorphic to $``$, and $`(\mathrm{})_{\mathrm{}^{\mathrm{}}}`$ is (isomorphic to) a von Neumann subalgebra of $``$, (6.48) holds. Since $`𝒩_{\text{II}_1}0`$, we have by the above discussion that also
(6.49)
$$\text{ is isomorphic to a von-Neumann subalgebra of }𝒩.$$
Thus, we have that if $`𝒜`$ or $``$ equals $`𝒩`$ or $``$, then
(6.50) $`𝒜`$ is (isomorphic to) a subalgebra of $``$, which is
the range of a normal conditional expectation.
Now if (6.49) holds for any two von Neumann algebras $`𝒜`$ and $``$, then $`L^p(𝒜)`$ completely contractively factors through $`L^p()`$. Thus by Lemma 6.13 and Corollary 6.14 applied to $`X=L^p()`$, we obtain that $`L^p(𝒩)`$ is isomorphic to $`L^p()`$.
Now if $`𝒩_{\text{II}_{\mathrm{}}}0`$, again using the deep results in \[C2\], $`𝒩_{\text{II}_{\mathrm{}}}`$ is (isomorphic to) $`\overline{}B(\mathrm{}^2)`$ where $``$ is a finite hyperfinite von Neumann algebra, whence letting $`𝒜`$ and $``$ equal $`𝒩`$ or $`R\overline{}B(\mathrm{}^2)`$, (6.48) holds, whence $`L^p(𝒩)`$ is completely isomorphic to $`L^p()_pC_p`$ again by Lemma 6.13 and Corollary 6.14 applied to $`L^p()_pC_p`$.
Now assume $`𝒩_{\text{II}_{\mathrm{}}}=\{0\}`$, so $`𝒩_{\text{II}_1}\{0\}`$, and suppose $`𝒩`$ is infinite; since $`𝒩_{\text{II}_{\mathrm{}}}=\{0\}`$, we must have that $`𝒩_I`$ is infinite. But then by the classification of the $`L^p`$ spaces of type I algebras, we have that $`L^p(𝒩_I)`$ is completely isomorphic to either $`C_p`$, $`L_pC_p`$, $`C_pL_p`$, or $`C_p(L_p_pS_p)`$.
But $`C_pL_pL^p()`$ and $`C_p(L_p_pS_p)L^p()`$ are both completely isomorphic to $`C_pL^p()`$, by our analysis of the finite case. Hence $`L^p(𝒩)`$ is completely isomorphic either to $`C_pL^p()`$ or to $`(L_p_pC_p)L^p()`$, completing the entire proof.∎
## 7. $`L^p(𝒩)`$-isomorphism results for $`𝒩`$ type III hyperfinite or a free group von Neumann algebra
We first formulate the results of this section for the case of preduals of von Neumann algebras $`𝒩`$, i.e., $`L^1(𝒩)`$, and then show they hold also for the spaces $`L^p(𝒩)`$ for $`1<p<\mathrm{}`$, as in the preceding sections. The following result is an immediate consequence of Corollary 6.12. We prefer to give a quick proof just using Corollary 1.2.
###### Theorem 7.1.
Let $`𝒩`$ be a factor of type II<sub>1</sub> and let $``$ be a factor of type II or type III. Then the preduals $`𝒩_{}`$ and $`_{}`$ are not Banach space isomorphic.
###### Proof.
By the assumptions $``$ is a properly infinite von Neumann algebra, i.e., $`\overline{}B(\mathrm{}^2)`$ as von Neumann algebras (where $`\overline{}`$ is the standard von Neumann algebra tensor product). In particular $`_{}`$ is isometrically isomorphic to $`_{}_\gamma C_1`$ for some crossnorm $`\gamma `$ on the algebraic tensor product $`_{}C_1`$, and therefore $`C_1`$ imbeds isometrically in $`_{}`$. By Corollary 1.2, $`C_1`$ does not Banach space imbed in $`𝒩_{}`$. ∎
It would be interesting to know, whether a type II-factor and a type III-factor can be distinguished by the Banach space isomorphism classes of their preduals. (As noted in the Introduction, we do not know the answer for the special case of injective factors.) In \[C1\] Connes introduced a subclassification of factors of type III into factors of type III<sub>λ</sub>, where $`\lambda `$ can take any value in the closed interval $`[0,1]`$. Theorem 7.2 below shows that the number $`\lambda `$ in this classification cannot be determined by the Banach space isomorphism class (or even operator space isomorphism class) of the predual. Recall from \[C2\] and \[H\], that for each $`\lambda (0,1]`$, there is up to von Neumann algebra isomorphism only one injective factor of type III<sub>λ</sub> acting on a separable Hilbert space. For $`0<\lambda <1`$ it is the Powers factor
$$R_\lambda =\underset{n=1}{\overset{\mathrm{}}{}}(M_2(),\phi _\lambda )$$
where $`\phi _\lambda `$ is the state on the $`2\times 2`$ complex matrices given by
$$\phi _\lambda \left(\begin{array}{cc}x_{11}& x_{12}\\ x_{21}& x_{22}\end{array}\right)=\frac{\lambda }{1+\lambda }x_{11}+\frac{1}{1+\lambda }x_{22}$$
and for $`\lambda =1`$ it is the Araki-Woods factor $`R_{\mathrm{}}`$, which can be obtained (up to von Neumann-isomorphism) as the tensor product of two Powers factors
$$R_{\mathrm{}}R_{\lambda _1}\overline{}R_{\lambda _2}$$
provided $`\frac{\mathrm{log}\lambda _1}{\mathrm{log}\lambda _2}`$. On the hand there are uncountably many injective factors of type III<sub>0</sub> acting on a separable Hilbert space (cf. \[C1\], \[C2\]). We will consider the predual of a von Neumann algebra as an operator space with the standard dual operator space structure (cf. \[Bl\]).
###### Theorem 7.2.
Let for $`0<\lambda <1`$, $`R_\lambda `$ denote the Powers factor of type III<sub>λ</sub> and let $`R_{\mathrm{}}`$ denote the Araki-Woods factor of type III<sub>1</sub>.
* For every $`\lambda (0,1)`$ the predual $`(R_\lambda )_{}`$ is completely isomorphic to $`(R_{\mathrm{}})_{}`$.
* There is an uncountable family $`(𝒩_i)_{iI}`$ of mutually non-isomorphic (in the von Neumann algebra sense) injective type III<sub>0</sub>-factors on a separable Hilbert space for which $`(𝒩_i)_{}`$ is completely isomorphic to $`(R_{\mathrm{}})_{}`$.
###### Remark.
In \[ChrS\], Christensen and Sinclair proved that all injective infinite dimensional factors acting on separable Hilbert space are completely isomorphic. This does not imply that their preduals are completely isomorphic. Indeed the unique injective type II<sub>1</sub>-factor $``$ and the unique injective type II-factor $`\overline{}B(\mathrm{}^2)`$ have non-isomorphic preduals by Theorem 7.1. Theorem 7.2 as well as the results in \[ChrS\] are based on the completely bounded version of the Pełczyński decomposition method stated as Lemma 6.13 above.
###### Proof of Theorem 7.2.
(a) Let $`0<\lambda <1`$ and put $`𝒩=R_\lambda `$, $`=R_{\mathrm{}}`$. Since $`𝒩`$ is a properly infinite von Neumann algebra, there exists two isometries $`u_1,u_2𝒩`$, such that $`u_1u_1^{}`$ and $`u_2u_2^{}`$ are two orthogonal projections with sum 1. Define now
$$\mathrm{\Phi }:𝒩𝒩𝒩\text{ by }\mathrm{\Phi }(x)=(u_1^{}x,u_2^{}x)$$
and
$$\mathrm{\Psi }:𝒩𝒩𝒩\text{ by }\mathrm{\Psi }(x,y)=(u,x+u_2y)$$
Then $`\mathrm{\Phi }\mathrm{\Psi }=\mathrm{id}_{𝒩𝒩}`$ and $`\mathrm{\Psi }\mathrm{\Phi }=\mathrm{id}_𝒩`$. Since $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ are normal (i.e., continuous) in the $`\omega ^{}`$-topologies on $`𝒩`$ and $`𝒩𝒩`$) and also are completely bounded maps, it follows that $`𝒩_{}_{\mathrm{cb}}𝒩_{}𝒩_{}`$. Similary we have $`_{}_{\mathrm{cb}}_{}_{}`$. Thus the pair $`(_{},𝒩_{})`$ satisfies (ii) in Lemma 6.13. We next check condition (i) in Lemma 6.13.
Since $`R_{\mathrm{}}R_\lambda \overline{}R_{\mathrm{}}`$ as von Neumann algebras (cf. \[C1, Sect.3.6\]), we can without loss of generality assume that $`=𝒩\overline{}𝒫`$ where $`𝒫R_{\mathrm{}}`$. Let $`\phi `$ be a normal faithful state on $`𝒫`$ and define
$$\pi :𝒩𝒩\overline{}𝒫\text{ by }\pi (x)=x\text{1},$$
and let $`\rho :𝒩\overline{}𝒫𝒩`$ be the left slice map given by $`\phi `$, i.e., the unique normal linear map $`𝒩\overline{}P𝒩`$ for which
$$\rho (xy)=\phi (y)x,x𝒩,y𝒫.$$
Thus $`\pi _{\mathrm{cb}}=\rho _{\mathrm{cb}}=1`$ and $`\rho \pi =\mathrm{id}_𝒩`$. Hence $`\mathrm{id}_𝒩_{}`$ has a completely bounded factorization through $`_{}`$, i.e., $`𝒩_{}`$ is $`\mathrm{cb}`$-isomorphic to a $`\mathrm{cb}`$-complemented subspace of $`_{}`$. To prove the converse, we use that if $`\phi `$ is a normal faithful state on the III<sub>1</sub>-factor $`=R_{\mathrm{}}`$ and $`\alpha =\sigma _{t_0}^\phi `$ is the moduluar automorphism associated with $`\phi `$ at $`t_0=\frac{2\pi }{\mathrm{log}\lambda }`$, then the crossed product $`R_{\mathrm{}}_\alpha `$ is a factor of type III<sub>λ</sub> (cf. \[HW, proof of Lemma 2.9\]). Moreover injectivity of $`R_{\mathrm{}}`$ implies that the crossed product is injective (cf. \[C2\]). Hence $`R_{\mathrm{}}_\alpha R_\lambda `$ as von Neumann algebras, so in this part of the proof we may assume that $`_\alpha =𝒩`$. Further, after identifying $``$ with its natural imbedding in the crossed product, we have that $`𝒩`$ is generated as a von Neumann algebra by $``$ and a certain unitary group $`\{u^nn\}`$ coming from the crossed product construction (cf. \[C1\]). Let $`i:_\alpha `$ be the imbedding and let $`\epsilon :_\alpha i()`$ be the unique normal faithful conditional expectation of $`_\alpha `$ onto $`i()`$ for which $`\epsilon (u^n)=0`$, for $`n\{0\}`$ (see again \[C1\]). Then $`i`$ and $`\epsilon `$ are normal maps and $`i^1\epsilon i=\mathrm{id}_{}`$, so as above, we obtain that $`_{}`$ is $`\mathrm{cb}`$-isomorphic to a $`\mathrm{cb}`$-complemented subspace of $`𝒩_{}`$. Hence a) follows from Lemma 6.13.
(b) Put again $`=R_{\mathrm{}}`$ and let $`G`$ be a dense countable subgroup. Let $`\phi `$ be a normal faithful state on $`R_{\mathrm{}}`$ and put $`𝒩=R_{\mathrm{}}_\alpha G`$ where $`\alpha :G\mathrm{Aut}()`$ is the restriction of the modular automorphism group $`(\sigma _t^\phi )_t`$ to $`G`$. It follows from \[C1\] (see the proof of \[HW, Lemma 2.9\]) that $`𝒩_G`$ is a factor of type III<sub>0</sub>, which is also injective (by \[C2\]). Moreover $`T(𝒩_G)=G`$, where $`T`$ is Connes $`\pi `$-invariant. Hence $`GG^{}`$ implies, that $`𝒩_G`$ and $`𝒩_G^{}`$ are not von Neumann-algebra isomorphic. It is easy to check, that there are uncountably many dense countable subgroups of $``$. Put $`𝒫=𝒩_G\overline{}R_{\mathrm{}}`$. Since $`R_{\mathrm{}}\overline{}R_\lambda R_{\mathrm{}}`$ for $`0<\lambda <1`$, we have $`𝒫\overline{}R_\lambda 𝒫`$, $`0<\lambda <1`$, which by \[C1, Theorem 3.6.1\] implies that $`𝒫`$ is a factor of type III<sub>1</sub>. Since $`𝒫`$ is also injective we have
$$𝒩_G\overline{}R_{\mathrm{}}R_{\mathrm{}}=$$
as von Neumann algebras. As in the proof of (a), it now follows, that $`_{}`$ is $`\mathrm{cb}`$-isomorphic to a $`\mathrm{cb}`$-complemented subspace of $`(𝒩_G)_{}`$. Moreover, since $`_\alpha G`$ is a crossed product with respect to a discrete group, there is again an embedding $`i:_\alpha G`$ and a normal faithful conditional expectation $`\epsilon :_\alpha Gi()`$, and the rest of the proof of (b) follows now exactly as in the proof of (a). ∎
Let $`L(F_n)`$ denote the von Neumann algebra associated with the free group $`F_n`$ on $`n`$ generators. Then for $`2n\mathrm{}`$ $`L(F_n)`$ is a factor of type II<sub>1</sub>. It is a long standing open problem to decide whether these II<sub>1</sub>-factors are isomorphic as von Neumann algebras. Due to work of Voiculescu, Dykema and Radulescu, it is known that either these factors are all isomorphic or $`L(F_{n_1})\cong ̸L(F_{n_2})`$ whenever $`2n_1,n_2\mathrm{}`$ and $`n_1n_2`$ (cf. \[VDN\]). In \[Ar\] Arias proved that the von Neumann algebras $`L(F_n)`$, $`2n\mathrm{}`$ are isomorphic as operator spaces. We show below, that also their preduals are isomorphic as operator spaces. While Arias’ proof uses mainly group theoretical considerations, the proof of Theorem 7.3 below relies on one rather deep result of Voiculescu, that $`L(F_{\mathrm{}})M_k(L(F_{\mathrm{}}))`$ as von Neumann algebras for $`k=2,3,\mathrm{}`$ (cf. \[Vo\] or \[VDN\]).
###### Theorem 7.3.
$`L(F_n)_{}`$ is $`\mathrm{cb}`$-isomorphic to $`L(F_{\mathrm{}})_{}`$ for $`n=2,3,\mathrm{}`$.
###### Proof.
Let $`n`$, $`n2`$ and put $`𝒩=L(F_n)`$ and $`=L(F_{\mathrm{}})`$. Since $`F_n`$ is isomorphic to a subgroup of $`F_{\mathrm{}}`$ and vice versa, $`𝒩`$ is von Neumann-algebra isomorphic to a subfactor $`𝒩_1`$ of $``$ and $``$ is von Neumann-algebra isomorphic to a subfactor $`_1`$ of $`𝒩`$ (see \[Ar\] for details). Moreover, let $`\tau _{}`$ and $`\tau _𝒩`$ be the unique normal faithful tracial states on $``$ and $`𝒩`$ respectively. Then there is a unique normal faithful conditional expectation $`\epsilon :\stackrel{\text{onto}}{}𝒩_1`$ preserving the trace $`\tau _{}`$ (resp. a unique normal faithful conditional expectation $`\epsilon ^{}:𝒩\stackrel{\text{onto}}{}`$, preserving the trace $`\tau _𝒩`$). As in the proof of Theorem 7.2, this implies that $`X=_{}`$ and $`Y=𝒩_{}`$ satisfy condition (i) in Lemma 6.13. We next prove that (ii) in Lemma 6.13 is satisfied with $`q=1`$. Since $`=L(F_{\mathrm{}})`$ is a II<sub>1</sub>-factor, we can choose a sequence of orthogonal projections $`(p_i)_{i=1}^{\mathrm{}}`$ in $``$, such that $`\tau (p_i)=2^i`$ and $`_{i=1}^{\mathrm{}}p_i=1`$ (convergence in the strong operator topology). By Voiculescu’s result quoted above, $`L(F_{\mathrm{}})M_{2^i}(L(F_{\mathrm{}}))`$ for $`i=1,2,\mathrm{}`$ as von Neumann-algebras, which implies that $`p_ip_i`$ as von Neumann-algebras.
Indeed, Voiculescu’s result yields that there are orthogonal equivalent projections $`q_1,\mathrm{},q_{2^i}`$ in $``$ with $`_{j=1}^{2^i}q_j=\text{1}`$ so that $`q_1q_1`$. It follows (by uniqueness of $`\tau _{}`$) that $`\tau (q_j)=\tau (q_j^{})`$, for all $`j`$ and $`j^{}`$, and so $`\tau (q_1)=2^i`$. Since also $`\tau _{}(P_i)=2^i`$ and $``$ is a finite factor, $`q_1`$ and $`p_i`$ are equivalent, and hence $`p_ip_iq_1q_1`$ as desired.
Put
$$Q=(\mathrm{})_{\mathrm{}^{\mathrm{}}}=\overline{}\mathrm{}^{\mathrm{}}.$$
Then $`Q`$ is a von Neumann algebra isomorphic to $`Q_1=^{}p_ip_i`$, which is a von Neumann subalgebra of $``$. Moreover, there is a $`\tau _{}`$-preserving normal faithful conditional expectation $`\epsilon ^{\prime \prime }:\stackrel{\text{onto}}{}Q_1`$. Hence $`Q_{}`$ is $`\mathrm{cb}`$-isomorphic to a $`\mathrm{cb}`$-complemented subspace of $`_{}`$. Put as above $`X=_{}`$. Then $`Q_{}=(XX\mathrm{})_\mathrm{}^1`$ as operator spaces. Hence we have shown that $`(XX\mathrm{})_\mathrm{}^1`$ completely factors through $`X`$, so $`X`$ and $`(XX\mathrm{})_\mathrm{}^1`$ are completely isomorphic by Corollary 6.14. This proves (ii) iin Lemma 6.13 with $`q=1`$. Hence $`X=_{}`$ and $`Y=𝒩_{}`$ are completely isomorphic. ∎
In the rest of this section, we will show how Theorem 7.2 and Theorem 7.3 can be generalized to the non-commutative $`L^p`$-spaces associated with the von Neumann algebras in question. In \[Ko\], Kosaki proved, that the abstract $`L^p`$-spaces $`L^p()`$, $`1<p<\mathrm{}`$ associated with a $`\sigma `$-finite ($`=`$ countably decomposable) von Neumann algebra $``$, can be obtained by the complex interpolation method applied to the pair $`(,_{})`$ with the imbedding $`M_{}`$ given by the map $`xx\phi `$, $`x`$, for a fixed normal faithful state $`\phi `$ on $``$. Assume next that $`𝒩`$ is a von Neumann subalgebra of $``$ and $`\epsilon :𝒩`$ is a normal faithful conditional expectation of $``$ onto $`𝒩`$. By replacing $`\phi `$ by $`\phi \epsilon `$, we can assume, that the state $`\phi `$ used in Kosaki’s imbedding is $`\epsilon `$-invariant. Next, the adjoint of $`\epsilon `$ defines an imbedding of $`𝒩_{}`$ in $`_{}`$ and $`i^{}`$, the adjoint of the inclusion map $`i:𝒩`$ defines a $`\mathrm{cb}`$-contraction of $`_{}`$ onto $`𝒩_{}`$. Moreover, we have the following commuting diagram:
$$\begin{array}{ccccc}𝒩& \stackrel{𝑖}{}& & \stackrel{𝜀}{}& 𝒩\\ & & & & \\ 𝒩_{}& \stackrel{\epsilon ^{}}{}& _{}& \stackrel{i^{}}{}& 𝒩_{}\end{array}$$
where the vertical arrows are the Kosaki inclusions with respect to $`\phi _{1𝒩}`$, $`\phi `$ and $`\phi _{1𝒩}`$ respectively. By the complex interpolation method we now get contractions $`i_p:L^p(𝒩)L^p()`$ and $`\epsilon _p:L^p()L^p(𝒩)`$, such that the following diagram commutes:
$$\begin{array}{cccccc}𝒩& \stackrel{𝑖}{}& & \stackrel{𝜀}{}& 𝒩& \\ & & & & & \\ L^p(𝒩)& \stackrel{i_p}{}& L^p()& \stackrel{\epsilon _p}{}& L^p(𝒩)& \\ & & & & & \\ 𝒩_{}& \stackrel{\epsilon ^{}}{}& _{}& \stackrel{i^{}}{}& 𝒩& .\end{array}$$
Further, if we consider $`L^p(𝒩)`$ and $`L^p()`$ as operator spaces with the operator spaces structure introduce by Pisier in \[Pi1\], we get that $`i_p`$ and $`\epsilon _p`$ are complete contractions. Hence we have proved:
###### Lemma 7.4.
Let $``$ be a $`\sigma `$-finite von Neumann algebra, and $`𝒩`$ a sub von Neumann algebra, which is the range of a normal faithful conditional expectation $`\epsilon :𝒩`$. Then for every $`1<p<\mathrm{}`$, $`L^p(𝒩)`$ is $`\mathrm{cb}`$-isometrically isomorphic to a $`\mathrm{cb}`$-contractively complemented subspace of $`L^p()`$.
Lemma 7.4 implies that the proofs of Theorem 7.2 and Theorem 7.3 can be repeated almost word for word to cover the $`L^p`$-case. Note that the argument for $`𝒩_{}𝒩_{}𝒩_{}`$ and $`_{}_{}_{}`$ in the beginning of Theorem 7.2 also works for the $`L^p`$-spaces, when $`L^p(𝒩)`$ (resp. $`L^p()`$) are equipped with the natural left $``$-module structure (resp. left $`𝒩`$-module structure). Hence we get:
###### Theorem 7.5.
Let $`R_\lambda `$, $`0<\lambda <1`$ and $`R_{\mathrm{}}`$ be as in Theorem 7.2 and let $`1p<\mathrm{}`$. Then
* $`L^p(R_\lambda )_{\mathrm{cb}}L^p(R_{\mathrm{}})`$.
* There is an uncountable family of mutually non-isomorphic (in the von Neumann algebra sense) injective type III<sub>0</sub>-factors on a separable Hilbert space, for which $`L^p(N_i)_{\mathrm{cb}}L^p(R_{\mathrm{}})`$ for all $`iI`$.
* For every $`n`$, $`n2`$, $`L^p(L(F_n))_{\mathrm{cb}}L^p(L(F_{\mathrm{}}))`$.
|
warning/0005/cond-mat0005011.html
|
ar5iv
|
text
|
# Non-linear conduction in charge-ordered Pr0.63 Ca0.37 MnO3 : Effect of magnetic fields
\[
## Abstract
Non-linear conduction in a single crystal of charge-ordered Pr<sub>0.63</sub>Ca<sub>0.37</sub>MnO<sub>3</sub> has been investigated in an applied magnetic field. In zero field, the non-linear conduction, which starts at T$`<`$T<sub>CO</sub> can give rise to a region of negative differential resistance (NDR) which shows up below the Néel temperature. Application of a magnetic field inhibits the appearance of NDR and makes the non-linear conduction strongly hysteritic on cycling of the bias current. This is most severe in the temperature range where the charge ordered state melts in an applied magnetic field. Our experiment strongly suggests that application of a magnetic field in the CO regime causes a coexistance two phases.
\]
Rare earth manganites with general chemical formula Re<sub>1-x</sub>Ae<sub>x</sub>MnO<sub>3</sub> have attracted current interest because of rich variety of phenomena like colossal magnetoresistance (CMR) and charge ordering (CO) . For certain values of x, close to 0.5, these compounds undergo a first order charge-ordering (CO) transition where the Mn<sup>3+</sup> and Mn<sup>4+</sup> species arrange themselves alternately in the lattice. This transition leading to charge localization occurs on cooling below a temperature T<sub>CO</sub>, refrerred to as the charge ordering temperature. Orbital ordering also accompanies the charge ordering and a long range AFM order sets in at lower temperature (T<sub>N</sub> $`<`$ T<sub>CO</sub>).(For some systems, like Nd<sub>0.5</sub> Sr<sub>0.5</sub> MnO<sub>3</sub> T<sub>N</sub>$`=`$T<sub>CO</sub> .) The system studied by us, Pr<sub>1-x</sub>Ca<sub>x</sub>MnO<sub>3</sub> (x = 0.37), happens to be an unique charge ordered system in which the charge ordered state remains insulating for all values of x due to its low tolerance factor. It shows both charge and orbital order below T<sub>CO</sub> = 240K and the AFM order occurs at T<sub>N</sub> = 175K.
A fascinating aspect of the charge ordered/orbital ordered state is that the charge ordered insulator (COI) phase is unstable to a number of external perturbations (like magnetic field , electric field , optical radiation etc.) and can be melted by them. An application of the magnetic field leads to a collapse of the charge ordering gap( $`\mathrm{\Delta }_{CO}`$ ) and there is an insulator -metal transition (melting ) of the COI state to an FMM state . Optical radiation seems to create conducting filaments which at low temperatures leads to non-linear transport and even at much lower temperatures a region of negative differential resistance (with V $``$ I<sup>-n</sup> with 0 $`<`$ n $`<`$ 1) . Application of an electric field also gives rise to non-linear conduction, which seems to have a threshold field associated with it and a broad band noise of substantial magnitude with a power spectrum $``$ 1/frequency .
A schematic of the phase diagram of the material in the H-T plane studied is shown in the inset of figure 1. In a given field H, the COI state melts into the FMM phase at the melting temperature T<sub>MH</sub> (marked in the figure). We designate the region T<sub>MH</sub> $`<`$ T $`<`$ T<sub>CO</sub> with H $``$ 0 (shaded in inset)as the ”mixed charge ordered” (MCO)region to distinguish it from the COI (H = 0, T $`<`$ T<sub>CO</sub>) state as well as from the FMM state (H $``$ 0, T $`<`$ T<sub>MH</sub>). We call it MCO region because we will show below that thie region has co-existence of two phases. In this paper we address the specific question of the nature of non-linear electronic transport in the COI as well as MCO region for T $`<`$ T<sub>MH</sub>). We have performed experiments along a constant H line. In particular, we observe that in the COI state, transport is strongly non-linear with a negative differential resistance (NDR) developing for T$`<`$T<sub>N</sub> while in the MCO state, the NDR is strongly inhibited.In addition, in the MCO state there is a large hysteresis in the I-V characteristics seen on cycling of the bias current. We ascribe this to creation of a coexisting two phase region in the mixed state.(Hysteresis seen in the I-V curve on cycling the bias current is distinct from the hysteresis in resistivity seen on H field cycling in most past experiments)
Our sample was a (4 X 2 X 0.3 $`mm^3`$) single crystal of Pr<sub>0.63</sub>Ca<sub>0.37</sub>MnO<sub>3</sub> grown by the floating zone technique. Four linear contact pads of Ag-In alloy was soldered on to the sample in linear four-probe configuration with separation $``$ 0.25 mm. I-V data were taken with current biasing and with a temperature control better than 10mK. To avoid any memory effects, the sample was heated well above T<sub>CO</sub>, before recording the I-V data at each temperature and magnetic field. We have also measured electrical noise by digitizing the voltage across the sample voltage probes. The power spectrum was obtained by fourier transform of the time autocorrelation function of the voltage after subtracting out the mean. The magnetic field measurements were done using cryogen-free superconducting magnet capable of producing field upto 15T. In this communication, to limit the scope of data being presented, we report data only for a field of 8T, which in the temparature range of investigation encompasses all the phases that we need to study.
Fig.1 shows the resistivity ($`\rho `$) as a function of temperature for various magnetic fields. For H = 0, T$`{}_{CO}{}^{}`$ 240K. In a field of 8T, the T$`{}_{CO}{}^{}`$ 210K and the CO state completely melts to the FMM state at 80K. We identify this temperature as T<sub>MH</sub> for 8T field. A magnetic field of 12T arrests the formation of CO state and the sample remains metallic at all temperatures. The observed data are in agreement with previous studies done on the Pr<sub>1-x</sub> Ca<sub>x</sub> MnO<sub>3</sub> system . \[ Note: Magnetic susceptibility shows a transition near T $``$240K T<sub>N</sub> $``$175K.\]
Fig.2 shows the I-V characteristics at different temperatures in zero magnetic field. The I-V characteristics are linear for T $`>`$ 260K ($`>`$ T<sub>CO</sub>) and at all T $`<`$ T<sub>CO</sub>(i.e, in the COI state), non-linear conduction is observed. One can identify a threshold electric field or a current density beyond which the conduction is strongly non-linear as reported earlier in films of Nd<sub>0.5</sub>Ca<sub>0.5</sub>MnO<sub>3</sub> . In the inset of fig.2, we show the noise power measured at 6Hz as a function of biasing current I. At the onset of the non-linear conduction the noise increases rather rapidly. The noise has a 1/f character for 0.02 Hz $`<`$ f $`<`$ 20 Hz. Appearance of a strong noise component is a characteristics of the melting process.
At T $`<`$ 170K, a new component gets added to the non-linearity and a region of negative differential resistance ($`dV/dI`$$`<`$0) is observed when the bias current I exceeds a current threshold (I<sub>TH</sub>). I<sub>TH</sub> decreases as T decreases and at T = 85K , I<sub>TH</sub> $``$ 1.2mA corresponding to a current density $``$ 0.5 A/cm<sup>2</sup>.In the region of negative differential resistance(NDR), for I $`>`$ I<sub>TH</sub>, V $``$ I<sup>-n</sup>, where 0 $`<`$ n $`<`$ 1. The I-V characteristics are symmetric and exhibits no hysteresis on current cycling except for T close to T$`co`$. The observed data are highly reproducible. Interestingly, there are several similarities of NDR state created by current and that created by optical radiation using a laser . We find that the exponent $`n`$, which is a measure of the NDR ($`dlnV/dlnI`$ = -$`n`$), has a strong temperature dependence. In fig.3 we show $`n`$ as a function of temperature T scaled by T<sub>N</sub>. We find that $`n`$ $``$ 0 as T / T<sub>N</sub> $``$ 1 and for T / T<sub>N</sub> $``$ 1, $`n`$ $``$ constant( $``$ 0.5). A likely explanation for occurence of NDR can be that for I $`>`$ I<sub>TH</sub> metallic filaments open up. These filaments being of lower resistance will provide parallel paths of conduction. This extra current path will decrease the voltage drop across the sample for a given current . As the current increases presumably more such channels open up leading up to a further decrease in V. This will manifest itself in the I-V curve with a region of NDR. Non-linear conduction can also occur due to depinning of CO domains above a threshold applied field as has been reported earlier . But in the region of NDR, we assume conduction through metallic filaments to be the dominant mechanism for non-linear conduction. The formation of the conducting filament is not a breakdown as the I-V curve is reproducible on cycling. A similar mechanism is proposed for optically produced NDR . In an optically melted CO state in the NDR regime, I $``$ V<sup>-n</sup> with n $``$ 2/3 at T = 10K . Given the similarity, it is expected that both have the same origin.
We have measured the temperature rise of the sample with respect to the base ($`\mathrm{\Delta }T`$) by attaching directly a thermometer to the sample. $`\mathrm{\Delta }T`$ $``$20K at the lowest temperature ($``$ 80K) and at the highest power dissipation level (0.1W). At 150K, $`\mathrm{\Delta }<`$10K and its is negligible for T $``$ 180K. The power dissipation level where the NDR sets in at the lowest temperature leads to a $`\mathrm{\Delta }T`$ $``$ 5K. We also investigated whether the NDR is caused by this heating effect and rule out heating as the cause for the NDR. However, the heating can have some influence on the value of $`n`$ at the highest measuring current.
It appears that a strong correlation exists between the onset of the NDR regime and the magnetic order at T<sub>N</sub> . We illustrate this in figure 3 where we plot the temperature dependence of the intensity of the (0.5.0.5.0) line as obtained from the neutron diffraction in a sample of Pr<sub>0.65</sub> Ca<sub>0.35</sub> MnO<sub>3</sub>. The temperature dependence of the intensity of this line is a measure of the growth of the AFM order below T<sub>N</sub>. It is clearly seen from the figure that not only $`n`$ $``$ 0 at T = T<sub>N</sub>, it also follows a temperature dependence which closely matches that of the growth of the AFM order as observed through neutron diffraction. There may be two likely reasons for the appearance of NDR below T<sub>N</sub>. One reason can be that for T $`<`$ T<sub>N</sub> there may be incommensurate to commensurate transition of the CO or orbital order as has been reported in a closely related composition Pr<sub>0.5</sub> Ca<sub>0.5</sub> MnO<sub>3</sub> . The incommensuration which is due to disorder in orbital ordering can inhibit formatiom of such conducting filaments as are needed for the NDR. Alternatively, the AFM order in these materials being of pseudo-CE type , there is a FM coupling between the planes which contain zig-zag AFM chains. It may be that below T<sub>N</sub>, this interplane FM coupling enhances the formation of the metallic filaments which can be made up of FMM phases. If the later hypothesis is true we will not see formation of NDR region in CO systems with CE -type AFM order which has only AFM coupling between planes.
Figure 4 shows the I-V characteristics at different temperatures in a 8T magnetic field. In the MCO state, the non-linear conduction persists and is qualitatively similar to that found in COI. For both T $`>`$ T<sub>CO</sub> and T $`<`$ T<sub>MH</sub> the I-V behaviour is linear, but the two states have $`\rho `$ differing by two orders of magnitude. ( All the data were taken after field cooling (FC) in H = 8T from T $`=`$ 300K. At each temperature, the sample was freshly prepared as a FC sample after warming it upto room temperature to avoid any memory effect.) There is a distinct memory of the previously applied field when the field is changed in the COI state and on warming up beyond T<sub>CO</sub> one can erase the memory. Experiments were also carried out under zero-field cooled (ZFC) condition. The I-V characteristics for both FC and ZFC are qualitatively similar, and the ZFC curve consistently showing higher V (for a given I) than the corresponding FC curve at each T. Given the limited scope of this paper we donot discuss the details here.
The inset of figure 4 shows a comparison of the I-V curves of the COI and the MCO states at a representative temperature. The COI and the MCO state differ significantly in two aspects, there is no NDR at any temperature in the MCO state at 8T and there is strong hysteresis in the I-V curve in the MCO region.
An important observation of this investigation is that the melting of the CO state either by temperature ( T $``$ T<sub>CO</sub> ) or by a magnetic field is accompanied by strong hysteresis behaviour in the I-V curve (i.e, the I-V curves do not follow each other during current ramping up and down the cycle. This is different from the hysteresis seen on H cycling). This can be seen in figures 4 and 5.The hysteresis does not depend on the speed with which the current is ramped up and down. Typically in our experiment one cycle is taken over 40 minutes. The area under the hysteresis curve is defined as the hatched region in figure 5. We plot the area under the hysteresis curve as a function of T in the inset of fig 5. It is clearly seen that in the COI state, the hysteresis is observed only for $`\mathrm{\Delta }`$ T $``$ 10K below T<sub>CO</sub>. In contrast, in the MCO state, hysteresis persists over an extensive temperature range above T<sub>MH</sub>. In the inset (b) of fig. 5, we show an example of the hysteresis in the MCO state as T $``$ T<sub>MH</sub> from above.
The appearance of hysteresis in the I-V curves can be interpreted as due to the coexistance of two phases in the MCO region. There is evidence of such coexisting phases in the TEM data taken near T<sub>CO</sub> in (La,Ca)MnO<sub>3</sub> and (Pr,Ca)MnO<sub>3</sub> systems near the melting transition(T$``$T<sub>co</sub>. In this material under study, it has been seen that the ferromagnetic spin correlations persist below T<sub>CO</sub> . This is stabilized by the applied field and thus can be the nucleus of the FMM phase. In the MCO state these nuclei get stabilized on cooling and then grow as the second phase.Eventually as T$``$T<sub>MH</sub>, the MCO state collapses to the FMM phase. Existence of such coexisting phases will prevent the formation of the metallic filaments needed for the occurence of NDR. As a result the NDR will be strongly suppressed in this region.
To summarize, we have carried out a systematic investigation of non-linear transport in a CO system in a magnetic field. We find that a negative differential resistance region shows up below T<sub>N</sub> which is inhibited by application of a magnetic field. In the region T<sub>co</sub> $`<`$ T $`<`$ T<sub>MH</sub>, application of a magnetic field creates a region of coexisting phases leading to a strong hysteretic I-V curve which become substantial in the temperature regions close to T<sub>co</sub> and T<sub>MH</sub>.
A.Guha thanks the CSIR Center of Excellence in Chemistry, JNCASR, for financial support.
FIGURE CAPTIONS
(1) FIG. 1.Resistivity as a function of T in different magnetic fields for the sample Pr<sub>0.63</sub>Ca<sub>0.37</sub>MnO<sub>3</sub>. The inset shows the schematic phase diagram.
(2) FIG. 2. The I-V curves showing non-linear conduction and negative differential resistance. The inset shows the appearance of large noise at the onset of nonlinear conduction.
(3) FIG. 3. Temperature dependence of the exponent n. The line shows the temperature dependence of the intensity of the (1/2,0/1/2) line obtained from the neutron experiment .
(4) FIG. 4. The non-linear transport in a magnetic field of 8T. The inset shows a comparison of the I-V data at T = 90K in zero field and in H = 8T.
(5) FIG 5. Hysteresis observed in the I-V curves on cycling of the bias current. The inset (a) shows the temperature dependence of the area under hysteresis close to T$`co`$. The inset (b) shows the hysteresis in I-V curve in a magnetic field.
|
warning/0005/cond-mat0005078.html
|
ar5iv
|
text
|
# Current-spin-density-functional study of persistent currents in quantum rings
## I Introduction
Nanoscopic quantum rings small enough to be in the true quantum limit can nowadays be realized experimentally . Among the quantum effects manifested in such systems in the presence of an external magnetic field, is the Aharonov-Bohm (AB) effect, leading to periodic oscillations in the energy spectrum and thus persistent currents. This phenomenon, first predicted by Hund, was discussed in connection with superconducting rings and more recently predicted to occur also in one-dimensional metallic rings . In the ideal case of one electron in a clean, one-dimensional ring, the Aharonov-Bohm phase picked up by the electron modifies the periodic boundary conditions, leading to single particle energies given by
$`E_n={\displaystyle \frac{\mathrm{}^2}{2m}}\left({\displaystyle \frac{2\pi }{L}}\right)^2\left(n\alpha \right)^2`$ (1)
where $`\alpha =\varphi /\varphi _0`$ is the number of flux quanta penetrating the ring and $`L`$ is the length of the ring. The single-particle spectrum is periodic in $`\varphi `$ with periodicity $`\varphi _0`$. The persistent current associated with state $`n`$ is
$`J_n=c{\displaystyle \frac{E_n}{\varphi }}.`$ (2)
However, in realistic systems, interactions, lateral dimension, impurities and spin effects complicate the picture. In particular, interactions may shift different energy levels relative to each other, leading to complicated ground state patterns with transitions between states with different spin and/or angular momentum as the Aharonov-Bohm flux is increased. In this way, interactions may decrease the period of the oscillations in the ground state energy (“fractional Aharonov-Bohm effect”). The first systematic study of persistent currents in ideal, one-dimensional metallic rings, including temperature- and impurity effects but neglecting interactions, was reported in . Subsequent approaches include Hubbard model calculations , use of Hartree- and Hartree-Fock methods , exact diagonalization studies and very recently density-functional calculations. Experiments in the early nineties reported observations of persistent currents in an ensemble of $`10^7`$ Cu rings, in single gold rings and in a single loop in a GaAs heterojunction, all in the mesoscopic range. Very recently, Lorke et al. reported the first spectroscopic data on nanoscopic, self-assembled InGaAs quantum rings containing only one or two electrons.
Most of the theoretical approaches mentioned above, have the limitation that, e.g. , interactions, spin effects, impurities or lateral dimension had to be neglected to simplify the calculations. In this paper, we apply the so-called current spin density functional theory (CSDFT) including gauge fields in the energy functional. CSDFT was earlier applied to describe the electronic structure of quantum dots . This method, while being a mean field approach and thus not exact, has the advantage that one can take into account all the above effects, being more accurate than Hartree-Fock, and it should also be possible to take it to higher particle numbers than the exact diagonalization methods reported in the literature. Our aim is, first of all, to examine to what extent the CSDFT formalism captures the physics due to the Aharonov-Bohm effect, namely the periodic variations in the energy spectra as function of flux, and the corresponding persistent currents. Thus, after introducing the basics of our model in section II, we present, in section III, the spectra of impurity-free two- and four particle rings and discuss how they compare qualitatively and quantitatively to corresponding exact diagonalization- and experimental results. In section IV we study the effects of a symmetry breaking impurity potential on the density profile and persistent current of six-electron rings. It is also predicted that, for a fixed impurity strength, narrowing the confining potential will tend to localize the electrons and suppress the persistent current. Finally, section V is devoted to discussion and conclusions.
## II Model and numerical method
We consider $`N`$ electrons of effective mass $`m^{}`$, confined to a ring with radius $`R_0`$ by a potential
$`V(r)={\displaystyle \frac{1}{2}}m^{}\omega _0^2\left(rR_0\right)^2.`$ (3)
When we examine the system in the presence of an impurity, we will introduce an additional Gaussian potential centered at the bottom of the potential well $`(x=R_0,y=0)`$,
$`V_I(𝐫)=V_0\mathrm{exp}\left({\displaystyle \frac{(xR_0)^2}{a^2}}{\displaystyle \frac{y^2}{b^2}}\right).`$ (4)
The AB flux is provided by a flux tube of constant field $`B_0`$ and with radius $`r_i`$ in the center of the ring; the corresponding vector potential is chosen as
$`A_\phi `$ $`=`$ $`\{\begin{array}{c}B_0r/2,rr_i\\ \\ B_0r_i^2/(2r),r>r_i\end{array}`$ (8)
$`A_r`$ $`=`$ $`0,`$ (9)
with $`r_i`$ chosen small enough that the electrons themselves move essentially in a field-free region.
Instead of using the quantities $`N,\omega _0`$ and $`R_0`$ to describe the properties of the system, it is convenient to introduce the average $`1D`$ interparticle spacing $`r_{s,1D}`$ (related to the average density $`\rho _{1D}`$ by $`r_{s,1D}=1/(2\rho _{1D})=\pi R_0/N`$), and the “degree of one-dimensionality” $`C_F`$ . $`C_F`$ is a dimensionless number defined as the ratio between the “transverse” (oscillator) gap $`\mathrm{}\omega _0`$ and the Fermi energy; the latter is approximated by the Fermi energy of a free, one dimensional Fermi gas with the same density, $`ϵ_F=(\pi ^2\mathrm{}^2N^2)/(8m^{}L^2)`$, where $`L`$ is the length of the ring. Thus,
$`C_F=\omega _0{\displaystyle \frac{32m^{}r_{s,1D}^2}{\pi ^2\mathrm{}}}.`$ (10)
The higher the value of $`C_F`$, the narrower is the ring. We will be using values of $`C_F`$ for which the system is essentially described by a single channel, but with a smooth charge distribution in space.
For a given set of parameters, we compute the ground state charge- and current densities using CSDFT. In this formalism, originally introduced by Vignale and Rasolt , one solves the self-consistent Kohn-Sham type equations,
$`\left[{\displaystyle \frac{𝐩^2}{2m^{}}}+{\displaystyle \frac{e}{2m^{}}}\left(𝐩𝒜+𝒜𝐩\right)+𝒱_\delta \right]\mathrm{\Psi }_{i\delta }=\epsilon _{i\delta }\mathrm{\Psi }_{i\delta }`$ (11)
(We have dropped the arguments $`𝐫`$ for simplicity). The index $`i`$ labels the eigenstates with spin $`\delta =(,)`$, and $`𝒜:=𝐀+𝐀_{\mathrm{xc}}`$ and $`𝒱_\delta :=(e^2/2m^{})A^2+V_\delta +V_H+V_{\mathrm{xc}\delta }`$ are the effective vector and scalar potentials. Here, $`V_H`$ is the ordinary Hartree potential and $`V_\delta =V+()g^{}\mu _BB/2`$ is the external potential, including the Zeeman energy (which in our case is set to zero, as the electrons do not experience the magnetic field; $`\mu _B=e\mathrm{}/(2m_e)`$ is the Bohr magneton). The exchange-correlation vector and scalar potentials are
$`e𝐀_{\mathrm{xc}}={\displaystyle \frac{1}{\rho }}\{{\displaystyle \frac{}{y}}{\displaystyle \frac{[\rho e_{\mathrm{xc}}(\rho _\delta ,\gamma )]}{\gamma }},{\displaystyle \frac{}{x}}{\displaystyle \frac{[\rho e_{\mathrm{xc}}(\rho _\delta ,\gamma )]}{\gamma }}\}`$ (12)
and
$`V_{\mathrm{xc}\delta }={\displaystyle \frac{[\rho e_{\mathrm{xc}}(\rho _\delta ,\gamma )]}{\rho _\delta }}{\displaystyle \frac{e}{\rho }}𝐣_p𝐀_{\mathrm{xc}},`$ (13)
where $`\rho `$ is the particle density $`\rho =\rho _{}+\rho _{}`$ with $`\rho _\delta =_i|\mathrm{\Psi }_{i\delta }|^2`$ . The paramagnetic current density is given by $`𝐣_p=i\mathrm{}/(2m^{})_{i\delta }[\mathrm{\Psi }_{i\delta }^{}\mathrm{\Psi }_{i\delta }\mathrm{\Psi }_{i\delta }\mathrm{\Psi }_{i\delta }^{}]`$, and the real current density equals $`𝐣=𝐣_p+(e/m^{})𝐀\rho `$. The exchange-correlation energy $`e_{\mathrm{xc}}`$ depends on the so-called vorticity $`\gamma =\times (𝐣_p/\rho )|_z`$ of the wave function. For the details of the formalism, we refer to Ref.. The practical computational techniques that we found necessary to obtain convergent solutions of the CSDFT mean field equations, are given in Ref.. The parameters and results will be given in effective atomic units ($`a.u.^{}`$) with energy measured in $`Ha^{}=2Ry^{}=m^{}e^4/(4\pi \mathrm{}ϵϵ_0)^2`$ and length measured in $`a_B^{}=\mathrm{}^2(4\pi ϵϵ_0)/m^{}e^2`$, where $`ϵ`$ is the dielectric constant and $`m^{}`$ the effective mass. The results can then be scaled to the actual values for typical semiconductor materials.
## III Comparison to exact results
In order to check whether CSDFT provides a good description of the physics related to Aharonov-Bohm oscillations, we first apply it to some cases where exact diagonalization results are available, namely impurity-free rings containing two and four electrons, respectively. Exact diagonalization studies of these systems close to the ideal, one dimensional limit, were presented by Niemelä et al. , who calculated the energy spectra as function of the Aharonov-Bohm flux in the presence of interactions. Starting from the flux-free, non-interacting case, one can roughly understand the structure of the full (interacting) spectra from the following effects: As we have seen, the presence of an AB flux induces persistent currents and thus favors increasing total angular momentum $`L`$ (with the corresponding sign) of the electrons. E.g., for an even number of particles, as $`\varphi `$ is increased from zero, $`L<0`$ states come down in energy whereas $`L0`$ are pushed up in energy. The ground state of the non-interacting $`N`$-electron ring contains a series of crossings between different angular momentum states as the flux increases. As interactions are turned on, states with the highest possible symmetry in the spin part of the wave function are favored due to the gain in exchange energy. Thus, for example, triplet states come down in energy as compared to singlet states. This may lead to additional crossovers, not present in the non-interacting case, between states with different spin, in the ground state spectrum. Hence, as a first step, we proceed to check to what extent CSDFT can produce these features.
We start by considering a ring with two electrons, choosing a radius $`R_0=1.5`$ and $`\omega _0=1.0`$. This corresponds to the actual values estimated for the experimentally realized rings recently reported by Lorke et al. . Figure 1 shows the ground state energy per particle of the two-electron system vs. $`\varphi /\varphi _0`$ as computed from CSDFT. (Note that, for this choice of parameters, the electron density does not go entirely to zero in the center of the ring, so the electrons get partly exposed to the external field. This causes the energy spectrum to tilt upwards instead of being strictly periodic in $`\varphi `$, in contrast to, e.g. , Fig.3, where the ring radius is larger.)
For the lowest values of $`\varphi `$, the $`S=0`$, $`L=0`$ state is lowest in energy, with $`S`$ denoting the spin. Then a crossover takes place to the triplet state with $`L=1`$ and then back to $`S=0`$ with $`L=2`$. Exactly the same features were found in the exact diagonalization study , though for a more one-dimensional ring. The main difference between the two methods is that the cusp corresponding to the transition between $`L=0`$ and $`L=2`$ in the $`S=0`$ state (which is not the ground state here) near $`\varphi /\varphi _0=1/2`$ is rounded off in the mean field calculation, and the transition between the different $`L`$ states is gradual. The reason is the explicit breaking of the rotational symmetry in the internal structure of the wave function that is mapped out by the self-consistent mean-field solution. Due to this symmetry breaking, the angular momentum is no longer a “good” quantum number, and non-integer $`L`$-values are allowed. On the other hand, cusps at transitions between different spin states are not smeared out by the LSDA.
Furthermore, one can also check by direct comparison with the experimental results of Lorke et al. that even quantitatively, our approach provides reasonable predictions: The energy per particle in the two-electron ring at zero external flux can be roughly estimated from the experimental data to be around $`14`$meV. Converting our calculated result from effective to physical units, using the estimated values $`m^{}=0.07m`$ and $`ϵ=12.6`$, the energy per particle at zero flux is about $`9`$meV.
The good agreement between CSDFT and the exact results is encouraging and perhaps not too surprising: It has previously been shown by Ferconi and Vignale that CDFT provides an astonishingly accurate description of even very small quantum dots, and the same seems to be the case here, in the presence of spin. Similar conclusions were reached in a recent, independent, density-functional study by Emperador et al. of the same two-electron rings.
Next, we consider a four-electron ring with $`r_{s,1D}=2.5`$, $`C_F=10`$. Fig.2 shows the energy per particle for the lowest $`S=0`$ and $`S=1`$ states, respectively, as computed from CSDFT. We see that the singlet state is always the ground state, except for $`\varphi /\varphi _0`$ near $`0`$ and $`1`$, where the two states are nearly degenerate. For both $`S=0`$ and $`S=1`$, the total angular momentum changes continuously from $`L=0`$ at zero flux via $`L=2`$ at $`\varphi /\varphi _0=1/2`$ to $`L=4`$ when one flux quantum penetrates the ring. These features agree well with the exact diagonalization results in , which again were performed for a more one-dimensional ring.
We note in passing that changing the flux, keeping all other parameters fixed, can introduce a weak spin- or charge density wave in the ring.
As before, the main difference between the exact spectra and our mean field results is that in our case, the angular momentum changes continuously as a function of flux, thus smearing out the cusps at transitions between different $`L`$-states. This effect, which is due to the approximations made in our calculation, is similar to the effect an impurity has in the exact spectrum: An impurity potential, breaking the rotational symmetry, typically smears out the sharp cusps and opens up gaps at level crossings between different $`L`$-states. In the next section, we shall study systematically systems with such an explicitly symmetry-breaking potential of variable strength.
As an additional test of our method, we have compared, for given flux, the numerically computed current (integrated over a cross-section of the ring) to the derivative of the total energy with respect to flux, cfr. Eq.(2). We find very good agreement (to three decimal places, taking, for example $`N=4`$, $`r_{s,1D}=2.5`$, $`C_F=10`$, $`S=1`$ and $`\varphi =0.25\varphi _0`$) with the theoretical prediction (2). We thus conclude that CSDFT correctly captures the effect of Aharonov- Bohm phases picked up by the electrons in the ring.
For completeness, we also include the spectrum of a six-particle ring with $`r_{s,1D}=2`$, $`C_F=7`$ (Fig.3), which will be discussed in the next section in the presence of an impurity potential.
## IV Effects of impurity
In this section we will study the effects of a single Gaussian impurity, given by (4), on the persistent current and charge density of a six-electron ring. The main feature is that the impurity tends to localize the electrons in the ring, thus reducing the persistent current. This effect has previously been studied by Cheung et al. in the case of an ideal, one dimensional metal ring. We will see that our method produces results which are in good qualitative agreement with those in .
There is another mechanism which tends to deform the electron density: As we shall see, making the ring more and more narrow, i.e. increasing the parameter $`C_F`$, creates a charge-density-wave (CDW) in the ring. In the following, we shall study the dependence of the density and persistent current on impurity strength and on $`C_F`$.
We start with a direct study of the charge density in a ring containing six electrons. Fig.4 shows the spin down- and total densities of a ring with $`r_{s,1D}=2`$, $`C_F=7`$ and zero flux at various impurity strengths $`V_0`$ (with $`a=b=2`$, see Eq. (4)).
The spin, $`S`$, is zero in the ground state, see Fig.3. In this case, the density of both spin components is homogeneous and rotationally invariant for $`V_0=0`$ and $`0.1`$, whereas one can see that the electrons start to get localized at larger $`V_0`$. This localization is accompanied by the formation of a spin-density wave (SDW), with alternating spin up- and spin down electrons. We also see that the density at the impurity center decreases with increasing $`V_0`$. A more systematic analysis of this effect is presented in Fig.5, where we show the total electron density at ($`x=R_0`$, $`y=0`$) as function of impurity strength for the same ring at three different flux values, and also a more narrow ring with $`C_F=16`$, all normalized by the density at $`V_0=0`$. The curves roughly fall on top of each other, independent of $`\varphi `$ and $`C_F`$ and seem to fall off exponentially as function of the impurity strength.
Fig.6 shows an example of how the persistent current falls off with increasing impurity strength $`V_0`$ ($`a=b=2`$) in a six-electron ring ($`r_{s,1D}=2`$, $`C_F=7`$) in the $`S=1`$ state at $`\varphi =0.25\varphi _0`$. This choice of parameters, though not corresponding to the ground state (see Fig.3), is particularly convenient due to the large absolute value of the current, making numerical error less significant. We have checked that other (ground state) sets of parameters give qualitatively the same behavior. We note that the form of this curve, with a plateau at small $`V_0`$, followed by steeper fall-off, is very similar to the result by Cheung et al. , obtained by numerical diagonalization of an ideal one-dimensional ring with 20 electrons in the tight binding approximation.
Finally, we examine another effect which tends to localize the electrons in the ring: It turns out that making the ring more one-dimensional, i.e. increasing $`C_F`$ keeping everything else fixed, gradually localizes the electrons, creating a strong SDW along with a spatial modulation of the total charge density (CDW). This is illustrated in Fig.7 where we show density profiles of a six-electron ring with $`r_{s,1D}=2`$ at zero flux and without an impurity for three different values of $`C_F`$.
Such localization and antiferromagnetic ordering of the electron spin in a quasi one-dimensional ring confinement was recently confirmed by exact diagonalization studies: the many-body spectra of quantum rings with up to 6 electrons could be described by a spin model combined with a rigid center-of mass rotation . The stronger the ring confinement (i.e. the more narrow the quasi one-dimensional ring), or the lower the average particle density, the more pronounced are the effects of localization and formation of charge density waves in the internal structure of the many-body wave function.
One might expect that an impurity potential would “pin” a charge density wave in the ring, i.e. the persistent current of a ring with localized electrons should be reduced as compared to the non-localized case. We have examined the possibility of such a “pinning-depinning transition” by computing the persistent current at fixed impurity strength as a function of $`C_F`$. Fig.8 shows the result for several different values of the impurity strength $`V_0`$. We see that the current indeed decreases as the ring becomes more narrow; however, there is no “abrupt” transition since, as we have seen, localization happens gradually with increasing $`C_F`$.
Note the interesting scaling behavior suggested by Fig.8: The ratio between any two of the curves is just a constant, independent of $`C_F`$; in particular, the ratio $`I(V_0)/I(V_0=0)`$ is independent of $`C_F`$ for any $`V_0`$. Also note that the persistent current decreases with increasing $`C_F`$ even in the zero impurity case, thus making no qualitative distinction between a “clean” and a “dirty” ring. This may again be due to the explicit symmetry breaking by the LSDA which, as we have discussed previously, in a sense mimics disorder even in the case $`V_0=0`$.
## V Conclusion
The numerical analysis presented here demonstrates that the current-spin-density-functional formalism provides a suitable tool for describing the effects of Aharonov-Bohm phases and impurities in realistic quantum rings. In particular, we have shown that it reproduces the main qualitative features of the many-body spectra and persistent currents, taking into account the effects of interactions, spin and deviations from perfect one-dimensionality. Furthermore, we have shown that the persistent current in the ring may be suppressed by narrowing the confining potential at fixed impurity strength.
The main difference between this model and exact diagonalization results is that the LSDA introduces a breaking of the rotational symmetry in the ring even in the impurity-free case. One may hope that this is not a problem when describing experimentally realizable quantum rings, as a certain amount of disorder and non-perfect symmetry is expected to be present in any realistic system.
With the present experimental progress in fabricating and studying few-electron quantum rings, the methods described here may turn out useful for suggesting and describing future experiments.
ACKNOWLEDGEMENT: This work was financially supported by the Academy of Finland, the TMR programme of the European Community under contract ERBFMBICT972405, the “Bayerische Staatsministerium für Wissenschaft, Forschung und Kunst”, and the NORDITA Nordic project “Confined electronic systems”.
|
warning/0005/astro-ph0005261.html
|
ar5iv
|
text
|
# Using the acoustic peak to measure cosmological parameters
## Acknowledgements
I would like to thank David Spergel and Wayne Hu for patiently and expertly answering all my questions.
## appendix
We made two major approximations in arriving at equation (Using the acoustic peak to measure cosmological parameters). The first was to treat the sound speed as a constant, when a more accurate approximation would be to set $`c_s=1/\sqrt{3(1+\xi )}`$, where $`\xi =3a\mathrm{\Omega }_b/4\mathrm{\Omega }_\gamma `$ is the baryon-photon momentum density ratiohu . Keeping the next to leading term in $`\xi _{sls}`$, the size of the sound horizon is given by
$`\chi _{sh}={\displaystyle \frac{2}{H_0\sqrt{\mathrm{\Omega }_m}\sqrt{3z_{sls}}}}[(\sqrt{1+{\displaystyle \frac{z_{sls}}{z_{eq}}}}\sqrt{{\displaystyle \frac{z_{sls}}{z_{eq}}}})`$
$`{\displaystyle \frac{\xi _{sls}}{6}}(\sqrt{1+{\displaystyle \frac{z_{sls}}{z_{eq}}}}2{\displaystyle \frac{z_{sls}}{z_{eq}}}(\sqrt{1+{\displaystyle \frac{z_{sls}}{z_{eq}}}}\sqrt{{\displaystyle \frac{z_{sls}}{z_{eq}}}}))]`$
Assuming that there are three light neutrino species, the quantity $`\xi _{sls}`$ can be written as
$$\xi _{sls}=\frac{3}{4}\left(1+\frac{21}{8}\left(\frac{4}{11}\right)^{4/3}\right)\left(\frac{\mathrm{\Omega }_b}{\mathrm{\Omega }_m}\right)\left(\frac{z_{eq}}{z_{sls}}\right).$$
(27)
For reasonable values of the cosmological parameters, we find $`\xi _{sls}\mathrm{}<1`$, so we can negelect the second term in equation (appendix).
The second major approximation was to expand the expression for $`r_{sls}=R_c\mathrm{sinh}(\chi _{sls}/R_c)`$ in terms of $`\mathrm{\Omega }_c/\mathrm{\Omega }_m`$ and $`\mathrm{\Omega }_w/\mathrm{\Omega }_m`$:
$`r_{sls}={\displaystyle \frac{2}{H_0\sqrt{\mathrm{\Omega }_m}}}(1+{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mathrm{\Omega }_c}{\mathrm{\Omega }_m}}{\displaystyle \frac{1}{2(16w)}}{\displaystyle \frac{\mathrm{\Omega }_w}{\mathrm{\Omega }_m}}`$
$`{\displaystyle \frac{1}{8}}\left({\displaystyle \frac{\mathrm{\Omega }_c}{\mathrm{\Omega }_m}}\right)^2+{\displaystyle \frac{3}{8(112w)}}\left({\displaystyle \frac{\mathrm{\Omega }_w}{\mathrm{\Omega }_m}}\right)^2`$
$`{\displaystyle \frac{32w}{4(12w)(16w)}}\left({\displaystyle \frac{\mathrm{\Omega }_c}{\mathrm{\Omega }_m}}{\displaystyle \frac{\mathrm{\Omega }_w}{\mathrm{\Omega }_m}}\right)+\mathrm{}).`$ (28)
So long as $`\mathrm{\Omega }_m>\mathrm{\Omega }_c`$ and $`\mathrm{\Omega }_m>\mathrm{\Omega }_w`$, the higher order terms can safely be neglected. Figure 1. shows the percentage error in the first order truncation of $`r_{sls}`$ as compared to a full numerical evaluation. The fractional error is less than $`8\%`$ across a wide portion of parameter space, including the interesting region around $`(\mathrm{\Omega }_c,\mathrm{\Omega }_\mathrm{\Lambda })=(0,0.7)`$.
Our final task is to show that the period of the wave responsible for the first acoustic peak is large compared to the time taken for matter and radiation to decouple. If this were not the case, the anisotropy would not be frozen in and the acoustic peak would be washed out. The conformal period of the wave is given by $`T2\pi \chi _{sh}/c_s`$ and the conformal time interval taken to decouple is $`\mathrm{\Delta }\eta =\eta (z_{sls})\eta (z_{sls}+\mathrm{\Delta }z)`$, where $`\mathrm{\Delta }z300`$ is the redshift interval for decoupling. The ratio of $`T`$ to $`\mathrm{\Delta }\eta `$ is given by
$$\frac{T}{\mathrm{\Delta }\eta }4\pi \left(\frac{z_{sls}}{\mathrm{\Delta }z}\right)\left(1+\frac{z_{sls}}{z_{eq}}\sqrt{\frac{z_{sls}}{z_{eq}}\left(1+\frac{z_{sls}}{z_{eq}}\right)}\right).$$
(29)
For reasonable cosmological parameters, we find $`T/\mathrm{\Delta }\eta \mathrm{}>30`$, which tells us that the acoustic waves are effectively snap frozen when matter and radiation decouple.
|
warning/0005/hep-th0005126.html
|
ar5iv
|
text
|
# 1 Introduction
## 1 Introduction
In ref. I proposed the N=2 supersymmetric extension of the four-dimensional Born-Infeld (BI) action. I interpreted it as the Goldstone-Maxwell action associated with spontaneous (partial) breaking of (rigid) N=4 supersymmetry down to N=2, and the N=2 (abelian) vector supermultiplet of Goldstone fields. The basic idea behind this interpretation was the anticipated equivalence (modulo a non-linear field redefinition) between the N=2 super-BI action in four dimensions and the gauge-fixed world-volume action of a D3-brane propagating in six dimensions. This equivalence was verified in ref. , in the leading and subleading orders only (see ref. too), while no direct argument was presented. In this Letter I report on a progress in obtaining the transformation laws of the hidden non-linearly realized symmetries (including spontaneously broken translations and extra N=2 supersymmetry) which determine the form of the N=2 super-BI action and prove its Goldstone nature. The uniqueness of N=2 superextension of the BI action is also discussed. I give an N=2 superconformal extension of the BI theory, and speculate about its possible relation to noncommutative geometry.
## 2 Featuring the bosonic BI action
In this introductory section I recall some well-known facts about the bosonic BI action, in order to provide a basis for the subsequent discussion of the N=2 supersymmetric extension in sect. 3.
The bosonic BI action in flat four-dimensional spacetime with Minkowski metric $`\eta _{\mu \nu }`$, $`\mu ,\nu =0,1,2,3`$, reads <sup>3</sup><sup>3</sup>3The overall normalization of the BI action yields the Maxwell term, $`\frac{1}{4}F_{\mu \nu }F^{\mu \nu }`$, as the leading
contribution. The D3-brane action has, in addition, the inverse string coupling constant in front
of the action.
$$S_{\mathrm{BI}}=\frac{1}{b^2}d^4x\sqrt{det\left(\eta _{\mu \nu }+bF_{\mu \nu }\right)},$$
$`(1)`$
where $`F_{\mu \nu }=_\mu A_\nu _\nu A_\mu `$, and $`b>0`$ is the dimensionful parameter. For instance, in string theory one has $`b=2\pi \alpha ^{}`$, whereas in N=1 supersymmetric QED one has $`b=e^2/(2\sqrt{6}\pi m^2)`$. In what follows, I choose $`b=1`$ for simplicity.
The BI theory (1) can be thought of as the particular covariant deformation of Maxwell electrodynamics by higher order terms depending upon $`F`$ only. In fact, the BI theory also shares with the Maxwell theory some other physical properties, such as causal propagation, positive energy density and electric-magnetic duality (see, e.g., refs. and references therein). Unlike the Maxwell theory, its BI generalization gives rise to the celebrated taming of the Coulomb self-energy , i.e. it smears the singularity associated with a point-like charge in classical electrodynamics. Supersymmetry is known to be compatible with causality, positive energy and duality, so that one expects from supersymmetric BI actions the similar (properly generalized) properties. It is indeed the case for the N=1 BI action , and it should be the case for the N=2 BI action too.
As is also quite clear from its origin, either in open string theory or in N=1 scalar QED, the BI action is the effective action obtained by summing up certain quantum corrections (to all orders in $`b`$) that are independent upon spacetime derivatives $`(F)`$ of the Maxwell field strength $`F`$. The effective action is dictated by S-matrix, being defined modulo local field redefinitions. This does not, however, make the BI action to be ambiguous since it depends upon the vector gauge potential $`A`$ only via its field strength $`F`$, while any local reparametrization of $`A`$ merely results in the additive $`F`$-dependent terms which are to be disregarded by definition of the BI action,
$$\delta S=d^4x\delta (F)=d^4x\frac{}{F_{\mu \nu }}_\mu \delta A_\nu =d^4x_\mu \frac{}{F_{\mu \nu }}\delta A_\nu =O(F).$$
$`(2)`$
In other words, the BI action is the effective action of slowly varying (but not necessarily small) abelian gauge fields, which is dependent upon $`F`$, being independent upon $`F`$. In supersymmetric BI theories the rôle of $`F`$ is played by the gauge superfield strength $`W`$, so that the super-BI actions in superspace are defined modulo spacetime derivatives of $`W`$.
## 3 N=2 BI action and its symmetries
The N=1 BI action is well-known to be the Goldstone-Maxwell action associated with spontaneous partial supersymmetry breaking N=2$``$N=1 and the N=1 vector supermultiplet of Goldstone fields . Both N=1 and N=2 gauge field theories are most naturally formulated in superspace, with manifest off-shell N=1 or N=2 supersymmetry, respectively, which makes a study of partial breaking N=2$``$N=1 rather straightforward, by starting from a linear off-shell realization of N=2 supersymmetry and imposing a non-linear constraint. Partial breaking N=4$``$N=2 in N=2 superspace is more complicated since a natural (off-shell and N=4 supersymmetric) formulation of N=4 gauge theories does not exist.
The N=2 supersymmetric BI action can be formulated in the standard N=2 superspace parametrized by $`Z=(x^{\alpha \stackrel{_{{}_{}{}^{}}}{\alpha }},\theta _i^\alpha ,\overline{\theta }_\stackrel{_{{}_{}{}^{}}}{\alpha }^i)`$, where $`\alpha =1,2`$ and $`i=1,2`$. The N=2 flat superspace covariant derivatives $`(_{\alpha \stackrel{_{{}_{}{}^{}}}{\alpha }},D_\alpha ^i,\overline{D}_i^\stackrel{_{{}_{}{}^{}}}{\alpha })`$ satisfy the algebra
$$\{D_\alpha ^i,\overline{D}_{\stackrel{_{{}_{}{}^{}}}{\alpha }j}\}=2i\delta _j^i_{\alpha \stackrel{_{{}_{}{}^{}}}{\alpha }},\{D_\alpha ^i,D_\beta ^j\}=\{\overline{D}_i^\stackrel{_{{}_{}{}^{}}}{\alpha },\overline{D}_j^\stackrel{_{{}_{}{}^{}}}{\beta }\}=0.$$
$`(3)`$
The standard realization is given by
$$D_\alpha ^i=\frac{}{\theta _i^\alpha }+i\overline{\theta }^{\stackrel{_{{}_{}{}^{}}}{\alpha }i}_{\alpha \stackrel{_{{}_{}{}^{}}}{\alpha }},\overline{D}_{\stackrel{_{{}_{}{}^{}}}{\alpha }i}=\frac{}{\overline{\theta }^{\stackrel{_{{}_{}{}^{}}}{\alpha }i}}i\theta _i^\alpha _{\alpha \stackrel{_{{}_{}{}^{}}}{\alpha }}.$$
$`(4)`$
The $`SL(2,𝐂)`$ and $`SU(2)`$ indices are raised and lowered by the use of Levi-Civita ($`\epsilon `$) symbols, as usual. I use the notation
$$D^{ij}=\frac{1}{2}D^{\alpha i}D_\alpha ^j,D_{\alpha \beta }=\frac{1}{2}D_{\alpha i}D_\beta ^i,(D^3)_i^\alpha =\frac{}{D_\alpha ^i}D^4,D^4=\underset{\alpha ,i}{}D_\alpha ^i,$$
$`(5)`$
and similarly for $`\overline{D}`$’s and $`\theta `$’s ($`\overline{\theta }`$’s). The standard (Berezin) integration rules imply
$$d^4xd^8\theta d^4xd^4\theta d^4\overline{\theta }=d^4xd^4\theta \overline{D}^4=d^4xd^4\overline{\theta }D^4.$$
$`(6)`$
The abelian N=2 superfield strength is described by an N=2 restricted chiral superfield $`W`$ satisfying the off-shell N=2 superspace constraints
$$\overline{D}_i^\stackrel{_{{}_{}{}^{}}}{\alpha }W=0\mathrm{and}D^4W=\mathrm{}\overline{W}.$$
$`(7)`$
The second constraint (7) is just the N=2 Bianchi identity that implies $`\mathrm{}(D_{ij}W\overline{D}_{ij}\overline{W})=0`$ and, hence, $`D_{ij}W=\overline{D}_{ij}\overline{W}`$<sup>4</sup><sup>4</sup>4The fields are assumed to vanish at infinity. A solution to eq. (7) in components reads (in N=2 chiral superspace parametrized by $`y^\mu =x^\mu \frac{i}{2}\theta _i^\alpha \sigma _{\alpha \stackrel{_{{}_{}{}^{}}}{\alpha }}^\mu \overline{\theta }^{\stackrel{_{{}_{}{}^{}}}{\alpha }i}`$ and $`\theta _i^\alpha `$)
| $`W(y,\theta )`$ | $`=a(y)+\theta _i^\alpha \psi _\alpha ^i(y)\frac{1}{2}\theta _i^\alpha (\stackrel{}{\tau })^i{}_{j}{}^{}\theta _{\alpha }^{j}\stackrel{}{D}(y)`$ |
| --- | --- |
| | $`i(\theta ^3)^{i\alpha }_{\alpha \stackrel{_{{}_{}{}^{}}}{\beta }}\overline{\psi }_i^\stackrel{_{{}_{}{}^{}}}{\beta }(y)+\theta ^4\mathrm{}\overline{a}(y),`$ |
$`(8)`$
where I have introduced the complex (physical) scalar $`a`$, the chiral (physical) spinor isodoublet $`\psi _\alpha ^i`$, the real (auxiliary) isotriplet $`\stackrel{}{D}=\frac{1}{2}(\stackrel{}{\tau })_i{}_{}{}^{j}D_{}^{j}_i`$, and the Maxwell field strength $`F_{\mu \nu }`$ subject to the Bianchi identity
$$\epsilon ^{\mu \nu \lambda \rho }_\nu F_{\lambda \rho }=0,$$
$`(9)`$
whose solution is just $`F_{\mu \nu }=_\mu A_\nu _\nu A_\mu `$.
The N=2 supersymmetric extension of the BI action, proposed in ref. , reads
$$S=\frac{1}{2}d^4xd^4\theta W^2+\frac{1}{8}d^4xd^8\theta 𝒴(K,\overline{K})W^2\overline{W}^2,$$
$`(10)`$
where $`K=D^4W^2`$ and $`\overline{K}=\overline{D}^4\overline{W}^2`$, and
| $`𝒴(K,\overline{K})`$ | $`={\displaystyle \frac{1\frac{1}{4}(K+\overline{K})\sqrt{(1\frac{1}{4}K\frac{1}{4}\overline{K})^2\frac{1}{4}K\overline{K}}}{K\overline{K}}}`$ |
| --- | --- |
| | $`=1+\frac{1}{4}(K+\overline{K})+O(K^2).`$ |
$`(11)`$
The action (10) can be rewritten to the form
$$S=\frac{1}{4}d^4xd^4\theta X+\frac{1}{4}d^4xd^4\overline{\theta }\overline{X}+O(W),$$
$`(12)`$
where the N=2 chiral lagrangian $`X`$ is the iterative solution to the N=2 non-linear constraint
$$X=\frac{1}{4}X\overline{D}^4\overline{X}+W^2.$$
$`(13)`$
The uniqueness of the N=2 BI action was questioned in ref. by presenting a calculation of some terms in the iterative solution to eq. (13), which are absent in the perturbative expansion of the action (10), for example,
$$d^4xd^8\theta W^2\overline{W}^2\left[(D^4W^2)\overline{D}^4(\overline{W}^2D^4W^2)+(\overline{D}^4\overline{W}^2)D^4(W^2\overline{D}^4\overline{W}^2)\right]$$
$`(14a)`$
versus
$$d^4xd^8\theta W^2\overline{W}^2\left[(D^4W^2)^2\overline{D}^4\overline{W}^2+(\overline{D}^4\overline{W}^2)^2D^4W^2\right].$$
$`(14b)`$
However, it is not difficult to verify, by the use of eqs. (3) and (7), that the difference between eqs. (14a) and (14b) amounts to the $`W`$-dependent terms which do not belong to the N=2 BI action because they are ambiguous (cf. sect. 2). It was also explicitly demonstrated in ref. that the N=2 BI action (12) is self-dual with respect to an N=2 supersymmetric electric-magnetic duality (claimed in ref. too), by keeping all terms in the solution to eq. (13), including the $`W`$-dependent ones. This means that taking into account some $`W`$-dependent terms is apparently needed to demonstrate the N=2 supersymmetric electric-magnetic duality of the N=2 BI action. In general, however, it does not make sense to keep some $`W`$ (or $`F`$) dependent terms in the effective BI action originating either from a quantized open superstring theory or from a quantized supersymmetric gauge theory, while ignoring other possible $`W`$ (or $`F`$) dependent quantum corrections. Perhaps, the N=2 electric-magnetic self-duality may, nevertheless, be useful for a study of derivative corrections to the N=2 BI action in a more fundamental framework than just N=2 supersymmetry.
The Goldstone interpretation of the N=2 BI action implies that the complex scalar $`W|=a=P+iQ`$ is the Goldstone field associated with two spontaneously broken translations (in the directions orthogonal to a D3-brane world-volume in six dimensions). Hence, the action (10) or (12) should possess hidden invariance with respect to spontaneously broken (non-linearly realized) translations, $`\delta a=\lambda +\mathrm{},`$ where $`\lambda `$ is the complex (rigid) parameter. This symmetry is obvious from the viewpoint of a (1,0) supersymmetric BI action in six dimensions , which is related to the four-dimensdional N=2 BI action via dimensional reduction. Indeed, the six-dimensional action depends upon its gauge fields via their field strength only, while one can identify $`A_4+iA_5=a`$. Hence, the dimensionally reduced action actually depends upon the derivatives of $`a`$, and not upon $`a`$ itself, though it is not manifest in eq. (10). Similarly, the spinor components $`\psi _\alpha ^i`$ of $`W`$ in eq. (8) are supposed to be the Goldstone fermions associated with two spontaneously broken (non-linearly realized) supersymmetries in four dimensions, $`\delta \psi _\alpha ^i=\lambda _\alpha ^i+\mathrm{}`$, where $`\lambda _\alpha ^i`$ are the (rigid) spinor parameters.
Spontaneously broken symmetries determine the associated Goldstone action. Since, in our case, the N=2 BI action is fixed by the non-linear constraint (13), there should be a relation between the non-linear transformations in question and the constraint (8). It is now not difficult to find the relevant (without spacetime derivatives) terms in the N=2 superfield transformation laws,
$$\delta X=2\mathrm{\Lambda }W,\delta W=\mathrm{\Lambda }\left(1\frac{1}{4}\overline{D}^4\overline{X}\right)\frac{X}{W}\overline{D}^4\left(\overline{W}\overline{\mathrm{\Lambda }}\right)+\mathrm{},$$
$`(15)`$
where $`\mathrm{\Lambda }`$ is the spacetime-independent (rigid) N=2 superfield parameter,
$$\mathrm{\Lambda }=\lambda +\theta _i^\alpha \lambda _\alpha ^i+\frac{i}{8}\theta _i^\alpha (\sigma ^{\mu \nu })_\alpha {}_{}{}^{\beta }\theta _{\beta }^{i}\lambda _{\mu \nu },$$
$`(16)`$
$`X`$ is the iterative (to all orders in $`W`$ and $`\overline{W}`$, but modulo ambiguous $`W`$\- and $`\overline{W}`$-dependent terms) solution to the non-linear constraint (13). The dots in eq. (15) stand for the $`W`$-dependent terms needed for the consistency with the second equation (7). Since those terms are ambiguous in the N=2 BI action, I ignore them both in the action and in the transformation laws for simplicity. The invariance of the action (12) under the transformations (15) follows from the fact that
$$D^4W,(D^3)_i^\alpha W\mathrm{and}D_{\alpha \beta }W$$
$`(17)`$
are all total derivatives in $`x`$-space, because of eqs. (8) and (9). The second relation (15) now follows from the first one by varying the constraint (13). Comparing eqs. (8) and (15) shows that $`\lambda `$ is the rigid parameter of broken translations, whereas $`\lambda _\alpha ^i`$ are the rigid parameters of two broken supersymmetries. Surprisingly enough, there exists yet another non-linear symmetry with rigid (real and antisymmetric tensor) parameter $`\lambda _{\mu \nu }`$, which is apparently related to the Goldstone nature of the Maxwell field itself.
It is possible to rewrite the action (12) into the ‘free-field’ form (subject to the non-linear constraint)
$$S=\frac{1}{2}d^4xd^4\theta W^2+\frac{1}{8}d^4xd^8\theta \overline{X}X,$$
$`(18)`$
where I have merely substituted the constraint (13) into eq. (12) and used eq. (6). Equation (18) is the N=2 analogue to the known ‘free-field’ form of the N=1 BI action, given by the sum of free actions for an N=1 vector multiplet and an N=1 chiral multiplet, related by a non-linear constraint . There is, however, an obvious difference between the N=1 and N=2 ‘free-field’ actions, because the second term in eq. (18) gives rise to higher derivatives in components. The existence of the field redefinition eliminating the higher derivatives is guaranteed by the existence of the equivalent (gauge-fixed) D3-brane action without higher derivatives but with non-manifest (non-linearly realized or ‘deformed’) unbroken N=2 supersymmetry .
## 4 Outlook
The N=2 BI theory can be made (rigidly) N=2 superconformally invariant by modifying the constraint (13) as (cf. ref. )
$$X=\frac{1}{4}\frac{X}{\mathrm{\Phi }^2}\overline{D}^4\left(\frac{\overline{X}}{\overline{\mathrm{\Phi }}^2}\right)+W^2,$$
$`(19)`$
where the N=2 ‘superconformal compensator’ $`\mathrm{\Phi }`$ is an N=2 restricted chiral superfield obeying the constraints (7). The original N=2 BI action is obtained from eqs. (12) and (19) by ‘freezing’ $`\mathrm{\Phi }`$ to $`\mathrm{\Phi }=1/\sqrt{b}=(2\pi \alpha ^{})^{1/2}`$.
The bosonic BI lagrangian (sect. 2) interpolates between the Maxwell lagrangian, $`\frac{1}{4}F^{\mu \nu }F_{\mu \nu }\frac{1}{4}F^2`$, for small $`F`$ and the total derivative, $`\frac{i}{8}\epsilon ^{\mu \nu \lambda \rho }F_{\mu \nu }F_{\lambda \rho }\frac{i}{4}F\stackrel{~}{F}`$, for large $`F`$ . In the $`b0`$ limit the BI lagrangian reduces to
$$\frac{F^2}{F\stackrel{~}{F}},$$
$`(20)`$
whose N=2 supersymmetric extension
$$d^8\theta \frac{W^2\overline{W}^2}{K\overline{K}}\left(\frac{K+\overline{K}}{K\overline{K}}\right)$$
$`(21)`$
follows from eq. (11) in the $`b0`$ limit.
One may, therefore, think of $`\mathrm{\Phi }`$ as a constant non-covariant background containing a constant antisymmetric tensor $`B_{\mu \nu }`$ on the place of $`F_{\mu \nu }`$ in eq. (8). The $`b0`$ limit is then described by sending $`\mathrm{\Phi }`$ to infinity, at large $`B_{\mu \nu }`$ in particular. The N=2 BI action in this limit is believed to be equivalent to a rank-one (Maxwell) noncommutative N=2 supersymmetric gauge field theory via Seiberg-Witten map , with $`B_{\mu \nu }`$ being the measure of noncommutativity in $`x`$-space,
$$x^\mu ,x^\nu =i(B^1)^{\mu \nu }.$$
$`(22)`$
## Acknowledgement
It is my pleasure to thank Professor S. J. Gates Jr. for kind hospitality extended to me at the University of Maryland in College Park during preparation of this paper, and Professor E. Ivanov for discussions.
|
warning/0005/nucl-th0005049.html
|
ar5iv
|
text
|
# Nucleon Resonance Transition Couplings to Vector Mesons
## 1 Introduction
The observed enhancement of low mass dilepton production in relativistic nucleon-nucleon collisions has stimulated considerable theoretical activity. The concensus in the theoretical treatment is that the dileptons seen in the CERES experiments to date originate with densities close to the density $`\rho _0`$ of normal nuclear matter. At these densities hadronic variables are more suitable than quarks in the theoretical treatment. A rather convincing description is that based on the coupling of nucleon isobar-hole excitations was initiated by Peters et al. and later systematized by Rapp and Wambach . A relation between this approach and the overall scaling relations proposed by Brown and Rho , formulated in quark language and possibly applicable at higher densities has been given in ref.. More recently the empirical data have been used to fix the parameters in an effective field theory approach . In this relatively model independent treatment a good description of the dilepton excess seen by in the CERES experiments is given.
It has been argued that at higher densities, in the region between nuclear matter density, and the critical density of chiral symmetry restoration, hadronic variables should be replaced by quark variables . The example of the Nambu-Jona-Lasinio model, in which the constituent quark mass is the order parameter, suggests that hadron masses go to zero in the limit of bare current quark masses. In ref. this is brought about by replacing the $`\rho `$meson mass $`m_\rho `$, which sets the scales of the denominator by the Rapp-Wambach Lagrangian , by the effective $`\rho `$meson mass $`m_\rho ^{}`$. While this result may be obtained on the basis of self-consistency, the issue of the most appropriate treatment at densities above that of normal nuclear matter remains largely open. The chiral quark model, which implies relations between the vector meson transition couplings to the nucleon resonances and the vector meson couplings to nucleons, may impose useful constraints on the theoretical treatment.
The chiral quark model, in which constituent quarks couple directly to mesons, is known to describe the properties of the ground state octet and decuplet baryons fairly well . In particular it gives expressions for the $`N\mathrm{\Delta }(1232)`$ transition couplings, which are good at the level of $`25\%`$ accuracy, or better. Moreover, when augmented with a linear confining interaction, and two-pion and vector meson interactions between the quarks, it describes the whole empirical baryon spectrum satisfactorily .
We here use this model to calculate the $`\rho `$\- and $`\omega `$-meson transition couplings to the nucleon resonances up to 1700 MeV. These coupling constants cannot be determined directly from empirical decay widths, as they lie below the thresholds for vector meson decay. The quality of the model is tested by a calculation of the corresponding pion-resonance transition coupling constants in the impulse approximation, which are then compared to the values that are determined from the pionic decay widths. That the single quark coupling model for pion decays should provide a fair overall description is indicated by the fact that it implies that the ratios of the decay widths for $`\mathrm{\Xi }^{}\mathrm{\Xi }\pi `$, $`\mathrm{\Sigma }^{}\mathrm{\Sigma }\pi `$ and $`\mathrm{\Delta }N\pi `$ should be 1:4:9, which compares well to the empirical ratio 1:3.9:12.6 (that the number for the $`\mathrm{\Delta }`$ exceeds 9 is due to the larger phase space volume available). For the excited $`S`$ and the $`P`$shell resonances the quark model values, which are calculated here in the impulse approximation, fall within factors 1.5 – 2 of the empirically extracted transition couplings. The situation for the $`D`$shell resonances is less satisfactory. Improved agreement requires taking into account higher order contributions from two-quark operators.
The vector meson transition coupling constants to nucleon resonances are defined in terms of Lagrangian densities for the transition couplings, which involve generalized Rarita-Schwinger vector spinors. Comparison of the matrix elements of these Lagrangians to the corresponding matrix elements in the quark model makes it possible to express the transition coupling constants in terms of the corresponding vector meson coupling constants to the nucleons. The latter are determined - albeit within a liberal uncertainty range - by fits to nucleon-nucleon scattering data with phenomenological boson exchange interaction models. These expressions involve $`SU(2)`$ Clebsch-Gordan coefficients as well as orbital matrix elements of quark wave functions, which connect the $`P`$ and $`D`$shell and the excited $`S`$shell states to the ground states. The latter depend on the Hamiltonian model for the 3-quark system. We here employ a simple covariant harmonic oscillator model based on linear confining interaction with a flavor-spin dependent hyperfine interaction, which describes the empirical spectrum very well .
There is some freedom in the choice of the resonance transition coupling Lagrangians. This freedom is constrained by comparison with the corresponding quark model expressions, especially because of the orthogonality of the resonance and nucleon wave functions in the quark model. The matching of matrix elements of covariant Rarita-Schwinger type Lagrangians and quark model matrix elements will here be made for off-mass shell vector mesons with zero energy. This choice of kinematics is made with application of resonance propagation in nuclear matter in mind. The coupling of spin-isospin modes that propagate in matter to the $`P`$shell resonances has recently been shown to be both significant and intricate . A key aim of the present study is to determine the strength of this coupling.
Only rough correspondences can be made between the couplings in the the chiral quark model and in the hadronic model. Nevertheless, the coherence in the $`\omega `$meson coupling to the nucleon, where the factor 3 in the $`SU(3)`$ relation $`g_{\omega NN}3g_{\rho NN}`$ arises in the sum over the three quarks in the nucleon, seems to disappear as the coupling $`g_{\omega NN^{}}^{(1520)}`$ in the quark model is roughly equal to $`g_{\rho NN}`$ in the hadronic model, which lacks coherence in the quark sum. In the quark model we find that $`g_{\omega NN^{}}^{(1520)}1/2g_{\rho NN^{}}^{(1520)}`$. The same ratio of $`\omega `$ to $`\rho `$ transition couplings is found in hadronic resonance models as discussed below. A simple explanation for this is still wanting.
This paper is divided into 4 sections. In section 2 we derive the pion-resonance couplings and compare them to empirical data. The $`\rho `$-meson and $`\omega `$-meson resonance transition couplings are derived in section 3. A summarizing discussion is given in section 4.
## 2 Pion coupling constants
Before deriving the expressions for the vector meson transition couplings to nucleon resonances it is instructive to derive the corresponding pion transition coupling constants within the quark model. Under the assumption that the pion transitions are described by single quark operators one may also derive these coupling strengths directly from experiment. Overall the single quark operator approximation for the pion coupling to nucleon resonances underestimates the pion decay widths of the resonances, and therefore mainly has qualitative value . It does however determine the phases of the coupling constants, and, as shown below, if these are multiplied by about a factor 2, the decay widths are within a few ten percent of the empirical values.
The coupling of pions to constituent $`u`$ and $`d`$ quarks may be described by the pseudovector coupling
$$_{\pi qq}=i\frac{f_{\pi qq}}{m_\pi }\overline{\psi }_q\gamma _5\gamma _\mu _\mu \stackrel{}{\varphi }_\pi \stackrel{}{\tau }\psi _q.$$
$`(2.1)`$
Here $`\stackrel{}{\varphi }_\pi `$ is the isovector pion field operator, $`\psi _q`$ is the constituent quark field and $`m_\pi `$ is the pion mass. The pseudovector pion-quark coupling constant may be determined from the corresponding pion-nucleon coupling constant by comparison to the $`\pi NN`$ coupling:
$$_{\pi NN}=i\frac{f_{\pi NN}}{m_\pi }\overline{\psi }_N\gamma _5\gamma _\mu _\mu \stackrel{}{\varphi }_\pi \stackrel{}{\tau }\psi _N.$$
$`(2.2)`$
Comparison of the matrix elements of the Lagrangian (2.1) and (2.2) for the case of a proton with spin up, using the $`SU(6)`$ quark model wave functions in the case of the former, yields
$$<p,\frac{1}{2}|_{\pi NN}|p,\frac{1}{2}>=if_{\pi NN}\frac{k_3}{m_\pi }$$
$$<p,\frac{1}{2}|\underset{q=1}{\overset{3}{}}_{\pi qq}|p,\frac{1}{2}>=i\frac{5}{3}f_{\pi qq}\frac{k_3}{m_\pi },$$
$`(2.3)`$
where $`k_3`$ is the third component of the pion momentum. This gives the relation
$$f_{\pi qq}=\frac{3}{5}f_{\pi NN}.$$
$`(2.4)`$
As $`f_{\pi NN}1`$ it follows that $`f_{\pi qq}0.6`$. This result for the pion-quark coupling constant is close to that, which is obtained from the Goldberger-Treiman relation for quarks:
$$f_{\pi qq}=\frac{g_A^q}{2}\frac{m_\pi }{f_\pi }.$$
$`(2.5)`$
With $`g_A^q0.88`$ and the value $`f_\pi =93`$ MeV for the pion decay constant this relation gives $`f_{\pi qq}0.66`$. This coupling model does give a reasonably satisfactory account of the pion decay widths of the $`D`$meson resonances, for which the single quark current model should be adequate .
The chiral quark model may be employed to express the transition coupling constant of pions to nucleon resonance in terms of the pion quark coupling constant $`f_{\pi qq}`$. Since this is given in terms of the pion-nucleon coupling constant $`f_{\pi NN}`$ by (2.3), it thus becomes possible to express all pion-resonance transition coupling constants in terms of the pion-nucleon coupling constant.
These relations will depend on the orbital wave functions of the 3 constituent quarks that form the baryons. We shall here use the covariant harmonic oscillator model for the mass operator of the 3 quarks constructed in ref. to generate the 3 quark wave functions. That model is formed by an, in effect, linear confining interaction with a spin, flavor and orbital angular momentum dependent hyperfine interaction. It describes the baryon spectrum up to $``$ 1700 MeV to an accuracy of a few percent, which is quite adequate for the present application.
The 3-quark wave functions will be bilinear combinations of spin-flavor and orbital wave functions. The latter are products of harmonic oscillator functions of the two Jacobi coordinates of the 3-quark system:
$$\stackrel{}{r}=\frac{1}{\sqrt{2}}(\stackrel{}{r}_1\stackrel{}{r}_2),$$
$`(2.6a)`$
$$\stackrel{}{\rho }=\sqrt{\frac{2}{3}}(\stackrel{}{r}_3\frac{\stackrel{}{r}_1+\stackrel{}{r}_2}{2}).$$
$`(2.6b)`$
The 3-quark wave functions in the model of ref. are eigenfunctions of the mass operator
$$_0=\sqrt{3(\stackrel{}{\kappa }^2+\stackrel{}{k}^2+\omega ^4(\stackrel{}{\rho }^2+\stackrel{}{r}^2))},$$
$`(2.7)`$
where $`\omega `$ is a parameter that determines the strength of the confining interaction. Jacobi momenta $`\stackrel{}{\kappa }`$ and $`\stackrel{}{k}`$ are the canonical momentum operators that are conjugate to the Jacobi coordinates (2.5). The numerical value for the oscillator parameter $`\omega `$ is 311 MeV, as determined from the baryon spectrum. This value will be used here. In the impulse approximation the empirical value of the rms radius of the proton (0.8 fm) obtains with $`\omega =245`$ MeV .
The eigenfunctions of the mass operator (2.6) are products of harmonic oscillator functions of $`\stackrel{}{\rho }`$ and $`\stackrel{}{r}`$: $`\phi _{n_1l_1m_1}(\stackrel{}{\rho })\phi _{n_2l_2m_2}(\stackrel{}{r})`$. The appropriate symmetrized combinations of these with spin and isospin wave functions for the $`S`$, $`P`$, $`D`$ and lowest excited $`S`$state resonances are listed in Table 1. In the table $`|\frac{1}{2},s>_\pm `$ and $`|\frac{1}{2},t>_\pm `$ denote spin and isospin wave functions of mixed symmetry, which are symmetric $`(+)(\mathrm{"}(112)\mathrm{"})`$ or antisymmetric $`()(\mathrm{"}(121)\mathrm{"})`$ under exchange of the spins or isospins of the first two quarks. The states $`|\frac{3}{2},s>`$ and $`|\frac{3}{2},t>`$ denote spin and isospin states with total spin and isospin $`\frac{3}{2}`$, which are symmetric under exchange of any set of two coordinates.
The explicit expressions for the states of mixed symmetry with spin-$`z`$ projection $`s_z`$ are
$$|\frac{1}{2},s_z>_+=\underset{abcm}{}(\frac{1}{2},\frac{1}{2},a,b|1,m)(\frac{1}{2},1,c,m|\frac{1}{2},s_z)|a,b,c>,$$
$`(2.8a)`$
$$|\frac{1}{2},s_z>_{}=\underset{ab}{}(\frac{1}{2},\frac{1}{2},a,b|0,0)|a,b,s_z>.$$
$`(2.8b)`$
Here $`|a,b,c>`$ represent product states of three spins with $`z`$projections $`a`$,$`b`$ and $`c`$ respectively. The corresponding isospin states with isospin-$`z`$ projection $`t_z`$ are readily constructed by analogy . The rule for construction of symmetric combinations of product states of spin, isospin and spatial wave functions is based on the outer products of $`S_3`$.
Given the 3 quark wave functions for the nucleon resonances in Table 1, it becomes possible to calculate the matrix elements,
$$<p,\frac{1}{2}|\underset{q=1}{\overset{3}{}}_{\pi qq}|N^+,\frac{1}{2}>=i\frac{f_{\pi qq}}{m_\pi }<p,\frac{1}{2}|\underset{q=1}{\overset{3}{}}\stackrel{}{\sigma }^q(\stackrel{}{k}\omega _\pi \stackrel{}{v}_q)\tau _3^qe^{i\stackrel{}{k}\stackrel{}{r}_q}|N^+,\frac{1}{2}>,$$
$`(2.9)`$
of the pion-quark Lagrangian (2.1). Here $`N^{}`$ represents a nucleon or a $`\mathrm{\Delta }`$ resonance, and $`\stackrel{}{k}`$ is the momentum and $`\omega _\pi `$ the energy of the pion. The velocity operator for quark $`q`$ is denoted $`\stackrel{}{v}_q`$. These matrix elements, with the overall factor $`if_{\pi qq}/m_\pi `$ divided out, are listed in Table 2 for the resonances in Table 1.
In Table 2 the orbital matrix elements have been calculated with the eigenfunctions of the mass operator (2.7). These are harmonic oscillator functions, with the explicit forms
$$\phi _{000}(\stackrel{}{\rho })=(\frac{\omega ^2}{\pi })^{3/4}e^{\rho ^2\omega ^2/2},$$
$$\phi _{01m}(\stackrel{}{\rho })=\sqrt{2}\omega \stackrel{}{\rho }_m\phi _{000}(\stackrel{}{\rho }),$$
$$\phi _{200}(\stackrel{}{\rho })=\sqrt{\frac{2}{3}}\omega ^2(\rho ^2\frac{3}{2\omega ^2})\phi _{000}(\stackrel{}{\rho }),$$
$$\phi _{02m}(\stackrel{}{\rho })=\sqrt{\frac{4}{15}}\omega ^2\rho ^2\phi _{000}(\stackrel{}{\rho })\sqrt{4\pi }Y_{2m}(\widehat{\rho }).$$
$`(2.10)`$
In the Table 2 values of the following overlap integrals have been used:
$$(\phi _{000}(\stackrel{}{\rho }),e^{i\sqrt{\frac{2}{3}}\stackrel{}{\rho }\stackrel{}{k}}\phi _{000}(\stackrel{}{\rho }))=e^{k^2/6\omega ^2},$$
$`(2.11a)`$
$$(\phi _{000}(\stackrel{}{\rho }),e^{i\sqrt{\frac{2}{3}}\stackrel{}{\rho }\stackrel{}{k}}\phi _{200}(\stackrel{}{\rho }))=\frac{\sqrt{6}}{18}\frac{k^2}{\omega ^2}e^{k^2/6\omega ^2},$$
$`(2.11b)`$
$$(\phi _{000}(\stackrel{}{\rho }),e^{i\sqrt{\frac{2}{3}}\stackrel{}{\rho }\stackrel{}{k}}\phi _{01m}(\stackrel{}{\rho }))=i\frac{\sqrt{3}}{3}\frac{k_m}{\omega }e^{k^2/6\omega ^2},$$
$`(2.11c)`$
$$(\phi _{000}(\stackrel{}{\rho }),e^{i\sqrt{\frac{2}{3}}\stackrel{}{\rho }\stackrel{}{k}}\phi _{02m}(\stackrel{}{\rho }))=\frac{\sqrt{3}}{9}\frac{k^2}{\omega ^2}e^{k^2/6\omega ^2}\sqrt{\frac{4\pi }{5}}Y_{2m}(\widehat{k}).$$
$`(2.11d)`$
The goal here is not, however, these quark model relations per $`se`$, but expressions for the pion transitions couplings to the resonances, when these are described as (generalized) Rarita-Schwinger field operators . These couplings may be described by the following Lagrangians:
$$_{\pi N\mathrm{\Delta }}^{(1232)}=\frac{f_{\pi N\mathrm{\Delta }}^{(1232)}}{m_\pi }\overline{\psi }\chi ^{}_\mu \stackrel{}{\varphi }\stackrel{}{\chi }\psi _\mu +h.c.,$$
$`(2.12a)`$
$$_{\pi NN^{}}^{(1440)}=i\frac{f_{\pi NN^{}}^{(1440)}}{m_\pi }\overline{\psi }\chi ^{}\gamma _5\gamma _\mu _\mu \stackrel{}{\varphi }\stackrel{}{\tau }\chi \psi _N^{}+h.c.,$$
$`(2.12b)`$
$$_{\pi N\mathrm{\Delta }^{}}^{(1600)}=\frac{f_{\pi N\mathrm{\Delta }^{}}^{(1600)}}{m_\pi }\overline{\psi }\chi ^{}_\mu \stackrel{}{\varphi }\stackrel{}{\chi }\psi _\mu +h.c.,$$
$`(2.12c)`$
$$_{\pi NN^{}}^{(1535)}=i\frac{f_{\pi NN^{}}^{(1535)}}{m_\pi }\overline{\psi }\chi ^{}\gamma _\mu _\mu \stackrel{}{\varphi }\stackrel{}{\tau }\chi \psi _N^{}+h.c.,$$
$`(2.12d)`$
$$_{\pi NN^{}}^{(1520)}=\frac{f_{\pi NN^{}}^{(1520)}}{m_\pi }\overline{\psi }\chi ^{}\gamma _5_\mu \stackrel{}{\varphi }\stackrel{}{\tau }\chi \psi _\mu +h.c.,$$
$`(2.12e)`$
$$_{\pi N\mathrm{\Delta }^{}}^{(1620)}=i\frac{f_{\pi N\mathrm{\Delta }^{}}^{(1620)}}{m_\pi }\overline{\psi }\chi ^{}\gamma _\mu _\mu \stackrel{}{\varphi }\stackrel{}{\chi }\psi _\mathrm{\Delta }^{},$$
$`(2.12f)`$
$$_{\pi N\mathrm{\Delta }^{}}^{(1700)}=\frac{f_{\pi NN^{}}^{(1700)}}{m_\pi }\overline{\psi }\chi ^{}\gamma _5_\mu \stackrel{}{\varphi }\stackrel{}{\chi }\psi _\mu +h.c.,$$
$`(2.12g)`$
$$_{\pi NN^{}}^{(1650)}=i\frac{f_{\pi NN^{}}^{(1650)}}{m_\pi }\overline{\psi }\chi ^{}\gamma _\mu _\mu \stackrel{}{\varphi }\stackrel{}{\tau }\chi \psi _N^{}+h.c.,$$
$`(2.12h)`$
$$_{\pi NN^{}}^{(1700)}=\frac{f_{\pi NN^{}}^{(1700)}}{m_\pi }\overline{\psi }\chi ^{}\gamma _5_\mu \stackrel{}{\varphi }\stackrel{}{\tau }\chi \psi _\mu +h.c.,$$
$`(2.12i)`$
$$_{\pi NN^{}}^{(1675)}=\frac{f_{\pi NN^{}}^{(1675)}}{m_\pi ^2}\overline{\psi }\chi ^{}_\mu _\nu \stackrel{}{\varphi }\stackrel{}{\tau }\chi \psi _{\mu \nu }+h.c,$$
$`(2.12j)`$
$$_{\pi NN^{}}^{(1720)}=\frac{f_{\pi NN^{}}^{(1720)}}{m_\pi }\overline{\psi }\chi ^{}_\mu \stackrel{}{\varphi }\stackrel{}{\tau }\chi \psi _\mu +h.c.,$$
$`(2.12k)`$
$$_{\pi NN^{}}^{(1680)}=i\frac{f_{\pi NN^{}}^{(1680)}}{m_\pi ^2}\overline{\psi }\chi ^{}\gamma _5_\mu _\nu \stackrel{}{\varphi }\stackrel{}{\tau }\chi \psi _{\mu \nu }+h.c.$$
$`(2.12l)`$
Here $`\psi ,\psi _\mu `$ and $`\psi _{\mu \nu }`$ represent spin $`1/2,3/2`$ and spin $`5/2`$ Rarita-Schwinger spinor fields respectively, and $`\chi `$ and $`\stackrel{}{\chi }`$ represents isospin $`1/2`$ spinor and $`3/2`$ vector-spinors respectively.
The Rarita-Schwinger spinor field operators are defined as
$$\psi _\mu ^M=\mathrm{\Sigma }(1,\frac{1}{2},m,s|\frac{3}{2},M)ϵ_\mu (m)u_s,$$
$`(2.13a)`$
$$\psi _{\mu \nu }^M=\mathrm{\Sigma }(\frac{3}{2},1,r,n|\frac{5}{2}M)(1,\frac{1}{2},m,s|\frac{3}{2},r)ϵ_\mu (m)ϵ_\nu (n)u_s.$$
$`(2.13b)`$
The spin $`5/2`$ spinor $`\psi _{\mu \nu }`$ is symmetric in the two 4-vector indices. Above $`u_s`$ represents a Dirac spinor with $`s_z=s`$.
The resonance couplings (2.12) have been written in the chiral symmetry mandated form, which requires that the couplings vanish with pion 4-momentum. For baryons on their mass shell the couplings to the negative parity resonances may be simplified by use of the Dirac and corresponding Rarita-Schwinger equations.
The matrix elements of the transition couplings (2.12) between resonance states with charge $`+e`$ and spin-$`z`$ projection $`+\frac{1}{2}`$ and the proton with spin up are listed in Table 3. Direct comparison with the quark model couplings is possible only for those terms in (2.12) that have the corresponding dependence on pion momentum. Comparison of the terms, which depend on pion energy, would require the corresponding quark model couplings to be spelled out explicitly.
With this qualification comparison of the matrix elements in Tables 2 and 3 yield the sought for expressions for the resonance transition couplings to pions, $`f_{\pi NN^{}}`$, in terms of the pion-quark coupling constant, $`f_{\pi qq}`$, and then by (2.4) in terms of the $`\pi N`$ pseudovector coupling constant $`f_{\pi NN}`$. The resulting expressions are listed in Table 4. In order to have real coupling constants in the Rarita-Schwinger formalism the phase factors $`(i)^l`$ that appear in the quark model matrix elements have been dropped. These calculated values for the $`\pi NN^{}`$ coupling constants may be compared to the corresponding values determined from the empirically known widths for $`N\pi `$ decay of these resonances.
The empirical decay widths for $`N\pi `$ decay of the positive parity $`\mathrm{\Delta }`$ resonances $`\mathrm{\Delta }(1232)`$ and $`\mathrm{\Delta }(1600)`$ are obtained as
$$\mathrm{\Gamma }=\frac{1}{3}\frac{f_{\pi N\mathrm{\Delta }}^2}{4\pi }\frac{E^{}+m_N}{m_\mathrm{\Delta }}\frac{k^3}{m_\pi ^2}$$
$`(2.14)`$
Here $`E^{}`$ is the energy of the final nucleon. Insertion of the empirical decay widths and the corresponding kinematical factors then yields the value $`f_{\pi N\mathrm{\Delta }}^{(1232)}=2.2\pm 0.04`$ and $`f_{\pi N\mathrm{\Delta }}^{(1600)}=0.51\pm 0.07`$. These values exceed the quark model values (1.55 and 0.47) in Table 4 by factors 1.5 and 1.1 respectively. These underestimates are typical of the quark model for the pion decays in the single quark approximation.
For the $`N(1440)`$ the decay width for $`N(1440)N\pi `$ is obtained as
$$\mathrm{\Gamma }=3\frac{(f_{\pi NN^{}}^{(1440)})^2}{4\pi }\frac{E^{}m_N}{m^{}}\frac{k}{m_\pi ^2}(m^{}+m_N)^2$$
$`(2.15)`$
Here $`m^{}`$ is the resonance mass. From the empirical decay width for $`N\pi `$ decay $`227\pm 0.65`$ MeV of the $`N(1440)`$ one obtains $`f_{\pi NN^{}}^{(1440)}=0.39\pm 0.06`$. In this case the quark model value 0.26 again represent an underestimate of about a factor 1.5. This underestimate is a consequence of the fact that the $`\pi NN(1440)`$ coupling vanishes at $`k=0`$ in the quark model.
The $`N\pi `$ decay width for the spin $`\frac{1}{2}^{}`$ resonances are obtained as
$$\mathrm{\Gamma }=\frac{f_{\pi NN^{}}^2}{4\pi }\frac{E^{}+m_N}{m^{}}\frac{k}{m_\pi ^2}(m^{}m_N)^2.$$
$`(2.16)`$
Here the factor $`\alpha `$ is 3 for isospin $`1/2`$ resonances and 1 for isospin $`3/2`$ resonances with spin $`\frac{1}{2}^{}`$. For the $`N(1535)`$, $`N(1650)`$ and the $`\mathrm{\Delta }(1620)`$ this expression yields the values $`f_{\pi NN^{}}^{(1535)}=0.36\pm 0.05`$, $`f_{\pi NN^{}}^{(1650)}=0.31\pm 0.03`$ and $`f_{\pi N\mathrm{\Delta }}^{(1620)}=0.34\pm 0.06`$. The quark model results for these coupling constants in Table 4 are within 30% of these values.
The $`N\pi `$ decay widths for the spin $`\frac{3}{2}^{}`$ resonances are obtained as
$$\mathrm{\Gamma }=\frac{1}{3}\frac{f_{\pi NN^{}}^2}{4\pi }\frac{E^{}m_N}{m^{}}\frac{k^3}{m_\pi ^2}.$$
$`(2.17)`$
From this expressions and the empirical $`N\pi `$ decay widths we obtain the values $`f_{\pi NN^{}}^{(1520)}=1.56\pm 0.06`$, $`f_{\pi NN^{}}^{(1700)}=0.36`$ and $`f_{\pi N\mathrm{\Delta }}^{(1700)}=1.31`$. The quark model values in Table 4 are fairly close to the first and the last of these values, but falls below that for the $`N(1700)`$.
Finally the width for the decay $`N(1675)N\pi `$ is obtained as
$$\mathrm{\Gamma }=\frac{2}{5}\frac{(f_{\pi NN^{}}^{(1675)})^2}{4\pi }\frac{E^{}+m_N}{m^{}}\frac{k^5}{m_\pi ^4}.$$
$`(2.18)`$
The empirical decay width fraction $`67`$ MeV yields the value $`f_{\pi NN^{}}^{(1675)}=0.10`$. This is close to the corresponding quark model value (0.09) in Table 4. Note that the empirical decay width does not determine the phase of the pion resonance transition coupling constants.
The expressions for the pionic decay widths of the $`D`$shell resonances $`N(1720)`$ and $`N(1680)`$ are
$$\mathrm{\Gamma }(N(1720)N\pi )=\frac{1}{3}\frac{(f_{\pi NN^{}}^{(1720)})^2}{4\pi }\frac{E^{}+m_N}{m^{}}\frac{k^3}{m_\pi ^2},$$
$`(2.19)`$
$$\mathrm{\Gamma }(N(1680)N\pi )=\frac{2}{5}\frac{(f_{\pi NN^{}}^{(1680)})^2}{4\pi }\frac{E^{}m_N}{m^{}}\frac{k^5}{m_\pi ^4},$$
$`(2.20)`$
respectively. The fractional widths for $`N\pi `$ decay of these two resonances are $`22\pm 10`$ MeV and $`84\pm 10`$ MeV respectively . Given these widths we obtain the following coupling constant magnitudes: $`|(f_{\pi NN^{}}^{(1720)}|0.25\pm 0.06`$ and $`|(f_{\pi NN^{}}^{(1680)}|0.42\pm 0.04`$. Comparison of the calculated coupling constants in Table 4 shows that the quark model, in the present approximation, overestimates the former one of these coupling constants by about a factor 5 and underestimates the latter one by almost a factor 4. This problem with the pion decays of the $`D`$shell resonances has been noted before .
The overall situation that emerges is that with the one-quark transition operators, the quark model mostly underpredicts the resonance transition couplings by factors 1 - 1.5 with the present wave function model. The conclusion is that two-quark operators have to be significant for the description of pionic transitions of the baryon resonances . If neglected these have to be compensated for by multiplication of the $`\pi NN^{}`$ couplings by factors of the order 2-3.
## 3 Vector meson coupling constants
A universal $`SU(2)`$ symmetric model for the vector meson couplings to constituent quarks would be the following:
$$_{Vqq}=ig_{\rho qq}\overline{\psi }\gamma _\mu \stackrel{}{\tau }\stackrel{}{\rho }_\mu \psi +ig_{\omega qq}\overline{\psi }\gamma _\mu \omega _\mu \psi .$$
$`(3.1)`$
Here $`\stackrel{}{\rho }_\mu `$ and $`\omega _\mu `$ are the $`\rho `$-meson and $`\omega `$-meson field operators respectively.
The vector meson coupling constant $`g_{Vqq}`$ may be determined from either the $`\omega `$\- or $`\rho `$-nucleon coupling constants by writing the vector meson-nucleon coupling in the conventional form
$$_{VNN}=ig_{\omega NN}\overline{\psi }_N[\gamma _\mu +i\frac{\kappa _\omega }{2m}\sigma _{\mu \nu }_\nu ]\omega _\mu \psi _N$$
$$+ig_{\rho NN}\overline{\psi }_N[\gamma _\mu +i\frac{\kappa _\rho }{2m}\sigma _{\mu \nu }_\nu ]\stackrel{}{\rho }_\mu \stackrel{}{\tau }\psi _N.$$
$`(3.2)`$
Comparison of the matrix elements of the charge components of these Lagrangians for e.g. protons with spin up to the same matrix elements of the quark coupling operator (3.2) yields the relations
$$g_{\omega NN}=3g_{Vqq},$$
$`(3.3a)`$
$$g_{\rho NN}=g_{Vqq}.$$
$`(3.3b)`$
The tensor couplings $`\kappa _\omega `$ and $`\kappa _\rho `$ in (3.3) may be determined by comparing the matrix elements of the transverse part of the current couplings in (3.1) and (3.3). This yields the relations
$$\frac{g_{\omega qq}}{m_q}=\frac{1}{m_N}g_{\omega NN}(1+\kappa _\omega ),$$
$`(3.4a)`$
$$\frac{5}{3}\frac{g_{\rho qq}}{m_q}=\frac{1}{m_N}g_{\rho NN}(1+\kappa _\rho ).$$
$`(3.4b)`$
Boson exchange models for the nucleon-nucleon interaction indicate that $`\kappa _\omega `$ is small. From (3.4a) it then follows that $`m_q=m_N/3=313`$ MeV, in agreement with conventional quark model phenomenology. Equations (3.3b) and (3.3b) then imply that $`\kappa _\rho =4`$. This is close to the value $`\kappa _\rho =4.22`$ in a recent realistic boson exchange model for the nucleon-nucleon interaction , but somewhat smaller than the value 6.6 indicated in earlier interaction models . The values for the $`\omega NN`$ coupling constants differ between different potential models. In the recent Nijmegen model it is $`g_{\omega NN}=10.35`$, while in the Bonn model it is as big as $`g_{\omega NN}=15.85`$. The values for the $`\rho NN`$ vector coupling constant are more stable, ranging from $`g_{\rho NN}=2.97`$ to 3.19 . These uncertainties are commensurate with the expected uncertainties in the quark model.
These numbers are consistent with assuming equality between the $`\rho `$ and $`\omega `$ couplings to constituent quarks, and with taking
$$g_{\rho qq}=g_{\omega qq}3.$$
$`(3.5)`$
The anomalous tensor couplings of the constituent quarks are so small that they may be taken to be 0. We shall use these values here.
In Table 5 the explicit matrix elements of the charge and transverse current components of the $`\omega `$quark coupling (3.1) are given for all the nucleon and $`\mathrm{\Delta }`$ resonances in Table 1 are listed. Here only the non-vanishing terms of lowest order in the $`v/c`$ expansion have been included. The corresponding matrix elements of the $`\rho `$quark coupling are listed in Table 6.
As the aim here is to calculate the vector meson transition couplings to nucleon resonances by means of the quark model, the effective coupling Lagrangians will have to be expressed in a form, which has the same momentum dependence as the corresponding transition matrix elements in the quark model. The standard form of the generalized Rarita-Schwinger vector current couplings in ref. does not meet this criterion. The form of the transition couplings below have been chosen to have have the same momentum dependence as the quark model couplings. As a consequence the overall meson momentum factors drop out in the expressions for the transition couplings in terms of the corresponding vector meson coupling constants to nucleons.
The $`\omega `$-meson transition couplings to the nucleon resonances may be described by the following effective Lagrangians:
$$_{\omega NN^{}}^{(1440)}=i\frac{g_{\omega NN^{}}^{(1440)}}{m_\omega ^2}\overline{\psi }_N[\gamma _\mu \frac{m^{}m_N}{m_\omega ^2}_\mu ]^2\omega _\mu \psi _N^{}+h.c.,$$
$`(3.6a)`$
$$_{\omega NN^{}}^{(1535)}=i\frac{g_{\omega NN^{}}^{(1535)}}{m_\omega ^2}\overline{\psi }_N\gamma _5[\gamma _\mu ^2(m^{}+m_N)_\mu ]\omega _\mu \psi _N^{}+h.c.,$$
$`(3.6b)`$
$$_{\omega NN^{}}^{(1520)}=i\frac{g_{\omega NN^{}}^{(1520)}}{m_\omega ^2}\overline{\psi }_N\sigma _{\mu \nu }_\nu _\kappa \omega _\mu \psi _\kappa +h.c.,$$
$`(3.6c)`$
$$_{\omega NN^{}}^{(1650)}=i\frac{g_{\omega NN^{}}^{(1650)}}{m_\omega ^2}\overline{\psi }_N\gamma _5[\gamma _\mu ^2(m^{}+m_N)_\mu ]\omega _\mu \psi _N^{}+h.c.,$$
$`(3.6d)`$
$$_{\omega NN^{}}^{(1700)}=i\frac{g_{\omega NN^{}}^{(1700)}}{m_\omega ^2}\overline{\psi }_N\sigma _{\mu \nu }_\nu _\kappa \omega _\mu \psi _\kappa +h.c.,$$
$`(3.6e)`$
$$_{\omega NN^{}}^{(1675)}=i\frac{g_{\omega NN^{}}^{(1675)}}{m_\omega ^2}ϵ_{\mu \alpha \beta \delta }\overline{\psi }_N_\alpha _\nu \omega _\beta \gamma _\delta \psi _{\mu \nu }+h.c.,$$
$`(3.6f)`$
$$_{\omega NN^{}}^{(1720)}=\frac{g_{\omega NN^{}}^{(1720)}}{m_\omega ^2}\overline{\psi }_N\gamma _5[\delta _{\mu \nu }\frac{1}{m^{}+m_N}\gamma _\mu _\nu ]^2\omega _\mu \psi _\nu +h.c.,$$
$`(3.6g)`$
$$_{\omega NN^{}}^{(1680)}=i\frac{g_{\omega NN^{}}^{(1680)}}{m_\omega ^2}\overline{\psi }_N[\gamma _\mu \frac{m^{}m_N}{m_\omega ^2}_\mu ]_\alpha _\beta \omega _\mu \psi _{\alpha \beta }+h.c..$$
$`(3.6h)`$
The matrix element of the charge and transverse current coupling terms of these Lagrangians between resonances with charge state $`+e`$ and the proton with spin up are listed in Table 7. As these matrix elements relate to virtual vector meson production, we have here considered the vector mesons as having zero energy. For the comparison with the quark model operators, we have in addition dropped terms of order $`(m^{}m_N)/(m^{}+m_N)`$. For the heavier resonances this introduces a theoretical uncertainty range of almost 30%, which however is inherent in the mismatch between the non-relativistic quark model expressions and the covariant Rarita-Schwinger formalism. Numerical estimates for the $`\omega `$ meson transition couplings to the nucleon resonances are given in Table 8. In the calculation of the numerical values, the quark wave function factors exp$`\{\stackrel{}{k}^2/6\omega ^2\}`$ were set to unity.
The $`\rho `$-meson transition coupling to the nucleon and $`\mathrm{\Delta }`$-resonances are described by the following coupling Lagrangians in the generalized Rarita-Schwinger formalism:
$$_{\rho N\mathrm{\Delta }}^{(1232)}=g_{\rho N\mathrm{\Delta }}^{(1232)}\overline{\psi }\chi ^{}\gamma _5[\delta _{\mu \nu }\frac{1}{m_\mathrm{\Delta }+m_N}\gamma _\mu _\nu ]\stackrel{}{\rho }_\mu \stackrel{}{\chi }\psi _\nu +h.c.,$$
$`(3.7a)`$
$$_{\rho NN^{}}^{(1440)}=i\frac{g_{\rho NN^{}}^{(1440)}}{m_\rho ^2}\overline{\psi }\chi ^{}[\gamma _\mu \frac{m^{}m_N}{m_\rho ^2}_\mu +i\frac{\kappa _{\rho NN}^{(1440)}}{m^{}m_N}\sigma _{\mu \nu }_\nu ]^2\stackrel{}{\tau }\stackrel{}{\rho }_\mu \chi \psi _N^{}$$
$$+h.c.,$$
$`(3.7b)`$
$$_{\rho N\mathrm{\Delta }^{}}^{(1600)}=g_{\rho N\mathrm{\Delta }^{}}^{(1600)}\overline{\psi }\chi ^{}\gamma _5[\delta _{\mu \nu }\frac{1}{m_\mathrm{\Delta }^{}+m}\gamma _\mu _\nu ]^2\stackrel{}{\rho }_\mu \stackrel{}{\chi }\psi _\nu +h.c.,$$
$`(3.7c)`$
$$_{\rho NN^{}}^{(1535)}=i\frac{g_{\rho NN^{}}^{(1535)}}{m_\rho ^2}\overline{\psi }\chi ^{}\gamma _5[\gamma _\mu ^2(m^{}+m_N)_\mu ]\stackrel{}{\tau }\stackrel{}{\rho }_\mu \chi \psi _N^{}+h.c.,$$
$`(3.7d)`$
$$_{\rho NN^{}}^{(1520)}=i\frac{g_{\rho NN^{}}^{(1520)}}{m_\rho ^2}\overline{\psi }\chi ^{}\sigma _{\mu \nu }_\nu _\kappa \stackrel{}{\tau }\stackrel{}{\rho }_\mu \chi \psi _\kappa +h.c.,$$
$`(3.7e)`$
$$_{\rho N\mathrm{\Delta }^{}}^{(1620)}=i\frac{g_{\rho N\mathrm{\Delta }^{}}^{(1620)}}{m_\rho ^2}\overline{\psi }\chi ^{}\gamma _5[\gamma _\mu ^2(m_\mathrm{\Delta }^{}+m_N)_\mu ]\stackrel{}{\rho }_\mu \stackrel{}{\chi }\psi _\mathrm{\Delta }^{}+h.c.,$$
$`(3.7f)`$
$$_{\rho N\mathrm{\Delta }^{}}^{(1700)}=i\frac{g_{\rho N\mathrm{\Delta }^{}}^{(1700)}}{m_\rho ^2}\overline{\psi }\chi ^{}\sigma _{\mu \nu }_\nu _\kappa \stackrel{}{\rho }_\mu \stackrel{}{\chi }\psi _\kappa +h.c.,$$
$`(3.7g)`$
$$_{\rho NN^{}}^{(1650)}=i\frac{g_{\rho NN^{}}^{(1650)}}{m_\rho ^2}\overline{\psi }\chi ^{}\gamma _5[\gamma _\mu ^2(m^{}+m_N)_\mu ]\stackrel{}{\tau }\stackrel{}{\rho }_\mu \chi \psi _N+h.c.,$$
$`(3.7h)`$
$$_{\rho NN^{}}^{(1700)}=i\frac{g_{\rho NN^{}}^{(1700)}}{m_\rho ^2}\overline{\psi }\chi ^{}\sigma _{\mu \nu }_\nu _\kappa \stackrel{}{\tau }\stackrel{}{\rho }_\mu \chi \psi _\kappa +h.c.,$$
$`(3.7i)`$
$$_{\rho NN^{}}^{(1675)}=i\frac{g_{\rho NN^{}}^{(1675)}}{m_\rho ^2}ϵ_{\mu \alpha \beta \delta }\overline{\psi }_N\chi ^{}_\alpha _\nu \stackrel{}{\tau }\stackrel{}{\rho }_\beta \gamma _\delta \chi \psi _{\mu \nu }+h.c.,$$
$`(3.7j)`$
$$_{\rho NN^{}}^{(1720)}=\frac{g_{\rho NN^{}}^{(1720)}}{m_\rho ^2}\overline{\psi }_N\chi ^{}\gamma _5[\delta _{\mu \nu }\frac{1}{m^{}+m_N}\gamma _\mu _\nu ]^2\stackrel{}{\tau }\stackrel{}{\rho }_\mu \chi \psi _\nu +h.c.,$$
$`(3.7k)`$
$$_{\rho NN^{}}^{(1680)}=i\frac{g_{\rho NN^{}}^{(1680)}}{m_\rho ^2}\overline{\psi }_N\chi ^{}[\gamma _\mu \frac{m^{}m_N}{m_\rho ^2}_\mu ]_\alpha _\beta \stackrel{}{\tau }\stackrel{}{\rho }_\mu \chi \psi _{\alpha \beta }+h.c..$$
$`(3.7l)`$
The matrix elements of these transition couplings are listed in Table 9 for resonances with charge state $`e`$ and protons with spin $`1/2`$. The explicit expressions for the $`\rho `$meson transition couplings in terms of the corresponding $`\rho NN`$ coupling constant are listed in Table 10, along with numerical estimates. These expressions are obtained by comparison of the quark model transition matrix elements in Table 6 with the corresponding matrix elements of the transition couplings in Table 9.
Of the resonances considered only the $`N(1720)`$ lies above threshold for $`N\rho `$ decay, for which the decay branch is in fact large ($`80\%`$). This makes it possible to determine magnitude of the transition coupling constant $`g_{\rho NN^{}}^{(1720)}`$ directly from its decay width. For a real $`\rho `$ meson the coupling Lagrangian (3.7k) simplifies as the differential operator $`^2`$ may be replaced by $`m_\rho ^2`$. To lowest order in $`v/c`$ the Lagrangian (3.7c) then reduces to that of ref. , with the identification
$$f_{N^{}N\rho }=2\frac{\mu }{m_\rho }g_{\rho NN^{}}^{(1720)},$$
$`(3.8)`$
where $`f_{N^{}N\rho }`$ is the transition coupling constant defined in ref. . This was determined from the partial decay width for $`N(1720)N\rho `$ to be $`f_{N^{}N\rho }7.2`$ in ref.. This is somewhat smaller than the value 13.8 that is obtained from eqn.(3.8) with the value for $`g_{\rho NN^{}}^{(1720)}`$ obtained with the quark model in Table 9. The fact that the quark model leads to an overprediction for the $`\rho NN(1720)`$ transition strength in clearly related to the corresponding overprediction of the $`\pi NN(1720)`$ coupling constant noted above.
## 4 Discussion
The fact that the pion resonance transition couplings are underestimated by factors 1.5–2 by the single quark operator approximation suggests that the quark model, in the same approximation, may also lead to similar underestimates for the vector meson transition couplings to nucleon resonances. For the subthreshold resonance transition couplings considered here, and which are required in a dynamical treatment of nuclear matter, there is however no alternative to calculation based on a dynamical model.
The expressions for the vector meson transition coupling constants were derived here by comparing the matrix elements of the transverse components of the vector meson transition currents to the corresponding quark model matrix elements. The generalized Rarita-Schwinger coupling Lagrangians in eqs. (3.6) and (3.7) are, however, invariant and may be applied for virtual vector mesons of arbitrary momentum and energy in nuclear matter. As an example, consider the coupling (3.7e) of the $`\rho `$meson to the $`N(1520)`$ $`3/2^{}`$ resonance, which admits an interpretation as a $`\rho `$nucleon resonance. For non-zero $`\rho `$meson energy this coupling, to lowest order in the inverse baryon masses takes the form
$$_{\rho NN^{}}\frac{g_{\rho NN^{}}^{(1520)}}{m_\rho ^2}\psi ^{}\chi ^{}\{\frac{\stackrel{}{k}^2}{2\mu }\rho _0^a+i(1\frac{\omega _\rho }{2\mu })\stackrel{}{\sigma }\stackrel{}{k}\times \stackrel{}{\rho }^a\}\stackrel{}{\tau }^a\stackrel{}{k}\stackrel{}{\psi }\chi +h.c.$$
$`(4.1)`$
Here $`\mu `$ is defined as the baryon mass combination $`\mu =2m^{}m_N/(m^{}+m_N)`$.
This form of the $`\rho NN(1520)`$ coupling, which takes into account the $`L=1`$ aspect of the 3 quark description of the $`N(1520)`$ resonance differs from the form commonly employed for the same coupling :
$$_{\rho NN^{}}^{(1520)}=\frac{f_{\rho NN^{}}^{(1520)}}{m_\rho }\psi ^{}\chi ^{}(\omega \stackrel{}{\rho }^a\rho _0^a\stackrel{}{k})\tau ^a\stackrel{}{\psi }\chi +h.c.$$
$`(4.2)`$
If in (4.1) one sets $`\stackrel{}{k}^2=k^2`$, as appropriate for zero energy vector mesons, and then imposes the on-shell condition $`k^2=m_\rho ^2`$, along with the relation $`\stackrel{}{\psi }=i\stackrel{}{\sigma }\times \stackrel{}{\psi }`$, a formal equivalence between the expressions (4.1) and (4.2) obtains, provided that
$$f_{\rho NN^{}}^{(1520)}=\frac{m_\rho }{2\mu }g_{\rho NN^{}}^{(1520)}.$$
$`(4.3)`$
Comparing numbers, with the value 4.5 for $`g_{\rho NN^{}}^{(1520)}`$ given in Table 10, we obtain $`f_{\rho NN^{}}(1520)=1.5`$, which is about half of the value 3.2 obtained in ref.. Given the fact that the quark model underestimates the pion resonance transition couplings by factors 1.5–2 in the single quark operator approximation, this small quark model value is not unexpected.
This comparison may also be extended to the case of the $`\omega NN(1520)`$ transition coupling. By comparing the isospin independent versions of the transition Lagrangians (4.1) and (4.2), we obtain
$$f_{\omega NN^{}}^{(1520)}=\frac{m_\omega }{2\mu }g_{\omega NN^{}}^{(1520)}.$$
$`(4.4)`$
With the value $`g_{\omega NN^{}}^{(1520)}=7.7`$ obtained by means of the quark model in Table 7, one obtains $`f_{\omega NN^{}}^{(1520)}=2.6`$. This value is somewhat less than one half of that obtained in ref.. The smaller value may be a consequence of the fact that the chiral quark model result should apply to higher densities as discussed above. It roughly corresponds to about $`1/3`$ of $`g_{\omega NN}`$, which is the $`\omega `$nucleon coupling at zero density. The factor $`1/3`$ in the $`SU(3)`$ relation
$$g_{\omega NN}=3g_{\rho NN},$$
$`(4.5)`$
arises from coherence in the sum of the couplings of the three nonstrange quarks in the $`\omega `$. We would expect this coherence to disappear at higher density or temperature scales. We see no obvious simple reason for why for this resonance $`g_{\rho NN^{}}`$ is only about $`1/2`$ of $`g_{\omega NN^{}}`$, however.
This work should only be viewed as a first attempt at calculating the transition couplings for virtual vector mesons, and therefore the numerical values obtained should be viewed as suggestive, rather than as definite quantitative predictions. In an earlier study we found that two-pion exchange between constituent quarks furnished a significant contribution to the hyperfine interaction between constituent quarks, which, when combined with the one-pion exchange interaction, provides a dynamical basis for the effective spin-flavor structure that is required for a satisfactory description of the empirical spectra, when combined with a linear confining interaction.. A study of how such higher order corrections in the chiral quark model affects the problem at hand is now being undertaken. The goal is a better understanding of the change from meson to quark variables, given the general view that there is a region where the two descriptions overlap .
The present method for calculating the resonance transition couplings to the nucleon and $`\mathrm{\Delta }`$ resonances in the $`P`$ and $`SD`$ shells may be directly generalized to the higher lying resonances. The explicit quark model wave functions for all the resonances in the $`SD`$shell may be constructed by reference to the symmetry classification for the higher resonances in ref.. The construction of the corresponding transition couplings in the generalized Rarita-Schwinger formalism may be carried out with the methods outlined in ref. once care is taken to match the momentum dependence of the quark model matrix elements.
Acknowledgement
D. O. R. thanks the Physics Department of Brookhaven National Laboratory for its hospitality during the completion of this work. Research supported in part by the U.S. Department of Energy under grant DE-FG02-88ER40388 and the Academy of Finland by grants No. 43982 and 44903.
| Table 1. Explicit wave functions for the nucleon and $`\mathrm{\Delta }`$ resonances in the $`S`$, $`P`$ | |
| --- | --- |
| and $`D`$ shells including the lowest excited $`S`$-shell states. | |
| $`p,n,\frac{1}{2}^+`$ | $`\frac{1}{\sqrt{2}}\phi _{000}(\stackrel{}{\rho })\phi _{000}(\stackrel{}{r})\{|\frac{1}{2},t_3>_+|\frac{1}{2},s_3>_++|\frac{1}{2},t_3>_{}|\frac{1}{2},s_3>_{}\}`$ |
| $`\mathrm{\Delta }(1232),\frac{3}{2}^+`$ | $`\phi _{000}(\stackrel{}{\rho })\phi _{000}(\stackrel{}{r})|\frac{3}{2},t_3>|\frac{3}{2},s_3>`$ |
| $`N(1440),\frac{1}{2}^+`$ | $`\frac{1}{2}\{\phi _{200}(\stackrel{}{\rho })\phi _{000}(\stackrel{}{r})+\phi _{000}(\stackrel{}{\rho })\phi _{200}(\stackrel{}{r})\}`$ |
| | $`\{|\frac{1}{2},t_3>_+|\frac{1}{2},s_3>_++|\frac{1}{2},t_3>_{}|\frac{1}{2},s_3>_{}\}`$ |
| $`\mathrm{\Delta }(1600),\frac{3}{2}^+`$ | $`\frac{1}{\sqrt{2}}\{\phi _{200}(\stackrel{}{\rho })\phi _{000}(\stackrel{}{r})+\phi _{000}(\stackrel{}{\rho })\phi _{200}(\stackrel{}{r})\}|\frac{3}{2},t_3>|\frac{3}{2},s_3>`$ |
| $`N(1535),\frac{1}{2}^{}`$ | $`\frac{1}{2}_{ms}(1,\frac{1}{2},m,s|J,s_3)\{\phi _{01m}(\stackrel{}{\rho })\phi _{000}(\stackrel{}{r})`$ |
| $`N(1520),\frac{3}{2}^{}`$ | $`[|\frac{1}{2},t_3>_+|\frac{1}{2},s>_+|\frac{1}{2},t_3>_{}|\frac{1}{2},s>_{}]`$ |
| | $`\phi _{000}(\stackrel{}{\rho })\phi _{01m}(\stackrel{}{r})[|\frac{1}{2},t_3>_+|\frac{1}{2},s>_{}+|\frac{1}{2},t_3>_{}|\frac{1}{2}s>_+]\}`$ |
| $`\mathrm{\Delta }(1620),\frac{1}{2}^{}`$ | $`\frac{1}{\sqrt{2}}_{ms}(1,\frac{1}{2},m,s|J,s_3)\{\phi _{01m}(\stackrel{}{\rho })\phi _{000}(\stackrel{}{r})|\frac{3}{2},t_3>|\frac{1}{2},s>_+`$ |
| $`\mathrm{\Delta }(1700),\frac{3}{2}^{}`$ | $`+\phi _{000}(\stackrel{}{\rho })\phi _{01m}(\stackrel{}{r})|\frac{3}{2},t_3>|\frac{1}{2},s>_{}\}`$ |
| $`N(1650),\frac{1}{2}^{}`$ | $`\frac{1}{\sqrt{2}}_{ms}(1,\frac{3}{2},m,s|J,s_3)\{\phi _{01m}(\stackrel{}{\rho })\phi _{000}(\stackrel{}{r})|\frac{1}{2},t_3>_+`$ |
| $`N(1700),\frac{3}{2}^{}`$ | $`+\phi _{000}(\stackrel{}{\rho })\phi _{01m}(\stackrel{}{r})|\frac{1}{2},t_3>_{}\}|\frac{3}{2},s>`$ |
| $`N(1675),\frac{5}{2}^{}`$ | |
| $`N(1720),\frac{3}{2}^+`$ | $`\frac{1}{2}_{ms}(2,\frac{1}{2},m,s|J,s_3)\{\phi _{02m}(\stackrel{}{\rho })\phi _{000}(\stackrel{}{r})+\phi _{000}(\stackrel{}{\rho })\phi _{02m}(\stackrel{}{r})\}`$ |
| $`N(1680),\frac{5}{2}^+`$ | $`\{|\frac{1}{2},t_3>_+|\frac{1}{2},s_3>_++|\frac{1}{2},t_3>_{}|\frac{1}{2},s_3>_{}\}`$ |
| Table 2. Transition matrix elements of the quark operator $`O=_{q=1}^3\stackrel{}{\sigma }^q\stackrel{}{k}\tau _3^qe^{i\stackrel{}{k}\stackrel{}{r}_q}`$ | |
| --- | --- |
| between the nucleon resonances and the nucleon for charge states $`+1`$ with $`s_z=+\frac{1}{2}.`$ | |
| $`<p,\frac{1}{2}|O|p,\frac{1}{2}>`$ | $`\frac{5}{3}k_3`$ |
| $`<p,\frac{1}{2}|O|\mathrm{\Delta }(1232)^+,\frac{1}{2}>`$ | $`\frac{4\sqrt{2}}{3}k_3e^{k^2/6\omega ^2}`$ |
| $`<p,\frac{1}{2}|O|N(1440)^+,\frac{1}{2}>`$ | $`\frac{5\sqrt{3}}{54}k_3(\frac{k^2}{\omega ^2}+\frac{3}{2}\frac{\omega _\pi }{m_q})e^{k^2/6\omega ^2}`$ |
| $`<p,\frac{1}{2}|O|\mathrm{\Delta }(1600)^+,\frac{1}{2}>`$ | $`\frac{3\sqrt{6}}{27}k_3(\frac{k^2}{\omega ^2}+\frac{3}{2}\frac{\omega _\pi }{m_q})e^{k^2/6\omega ^2}`$ |
| $`<p,\frac{1}{2}|O|N(1535)^+,\frac{1}{2}>`$ | $`i\frac{2\sqrt{2}}{9}\omega (\frac{k^2}{\omega ^2}+\frac{9\omega _\pi }{2m_q})e^{k^2/6\omega ^2}`$ |
| $`<p,\frac{1}{2}|O|N(1520)^+,\frac{1}{2}>`$ | $`i\frac{2}{9}\frac{3k_3^2\stackrel{}{k}^2}{\omega }e^{k^2/6\omega ^2}`$ |
| $`<p,\frac{1}{2}|O|\mathrm{\Delta }(1620)^+,\frac{1}{2}>`$ | $`i\frac{\sqrt{2}}{9}\omega (\frac{k^2}{\omega ^2}+\frac{9\omega _\pi }{2m_q})e^{k^2/6\omega ^2}`$ |
| $`<p,\frac{1}{2}|O|\mathrm{\Delta }(1700)^+,\frac{1}{2}>`$ | $`\frac{i}{3}\frac{3k_3^2\stackrel{}{k}^2}{\omega }e^{k^2/6\omega ^2}`$ |
| $`<p,\frac{1}{2}|O|N(1650)^+,\frac{1}{2}>`$ | $`i\frac{\sqrt{2}}{9}\omega (\frac{k^2}{\omega ^2}+\frac{9\omega _\pi }{2m_q})e^{k^2/6\omega ^2}`$ |
| $`<p,\frac{1}{2}|O|N(1700)^+,\frac{1}{2}>`$ | $`i\frac{\sqrt{10}}{90}\frac{3k_3^2\stackrel{}{k}^2}{\omega }e^{k^2/6\omega ^2}`$ |
| $`<p,\frac{1}{2}|O|N(1675)^+,\frac{1}{2}>`$ | $`i\frac{\sqrt{10}}{15}\frac{3k_3^2\stackrel{}{k}^2}{\omega }e^{k^2/6\omega ^2}`$ |
| $`<p,\frac{1}{2}|O|N(1720)^+,\frac{1}{2}>`$ | $`\frac{\sqrt{15}}{27}k_3(\frac{k^2}{\omega ^2}+\frac{15\omega _\pi }{2m_q})e^{k^2/6\omega ^2}`$ |
| $`<p,\frac{1}{2}|O|N(1680)^+,\frac{1}{2}>`$ | $`\frac{\sqrt{45}}{108}\frac{k_3(5k_3^23k^2)}{\omega ^2}e^{k^2/6\omega ^2}`$ |
| Table 3. Transition matrix elements $`<p,\frac{1}{2}|_{\pi NN^{}}|N^+,\frac{1}{2}>`$ of the pion transition | |
| --- | --- |
| couplings in eqs. (2.12). Here $`m^{}`$ denotes the mass of the corresponding resonance | |
| The expressions after the vertical bars correspond to zero energy pions. | |
| $`<p,\frac{1}{2}|_{\pi N\mathrm{\Delta }}^{(1232)}|\mathrm{\Delta }(1232)^+,\frac{1}{2}>`$ | $`\frac{2}{3}i\frac{f_{\pi N\mathrm{\Delta }}}{m_\pi }k_3`$ |
| $`<p,\frac{1}{2}|_{\pi NN^{}}^{(1440)}|N(1440)^+,\frac{1}{2}>`$ | $`i\frac{f_{\pi NN^{}}^{(1440)}}{m_\pi }k_3`$ |
| $`<p,\frac{1}{2}|_{\pi N\mathrm{\Delta }^{}}^{(1600)}|\mathrm{\Delta }(1600)^+,\frac{1}{2}>`$ | $`i\frac{2}{3}\frac{f_{\pi N\mathrm{\Delta }^{}}^{(1600)}}{m_\pi }k_3`$ |
| $`<p,\frac{1}{2}|_{\pi NN^{}}^{(1535)}|N(1535)^+,\frac{1}{2}>`$ | $`if_{\pi NN^{}}^{(1535)}\frac{m^{}m}{m_\pi }|\{\frac{i}{4}\frac{f_{\pi NN^{}}^{(1535)}}{m_\pi }\frac{m^{}+m}{m^{}m}\stackrel{}{k}^2\}`$ |
| $`<p,\frac{1}{2}|_{\pi NN^{}}^{(1520)}|N(1520)^+,\frac{1}{2}>`$ | $`\frac{i}{4\sqrt{6}}\frac{f_{\pi NN^{}}^{(1520)}}{m_\pi }\frac{m^{}m}{m^{}m}(3k_3^2\stackrel{}{k}^2)`$ |
| $`<p,\frac{1}{2}|_{\pi N\mathrm{\Delta }^{}}^{(1620)}|\mathrm{\Delta }(1620)^+,\frac{1}{2}>`$ | $`i\sqrt{\frac{2}{3}}f_{\pi N\mathrm{\Delta }^{}}^{(1620)}\frac{m^{}+m}{m_\pi }|\{\frac{i}{4}\sqrt{\frac{2}{3}}\frac{f_{\pi N\mathrm{\Delta }^{}}^{(1620)}}{m_\pi }\frac{m^{}m}{m^{}m}\stackrel{}{k}^2\}`$ |
| $`<p,\frac{1}{2}|_{\pi N\mathrm{\Delta }^{}}^{(1700)}|\mathrm{\Delta }(1700)^+,\frac{1}{2}>`$ | $`\frac{i}{12}\frac{f_{\pi N\mathrm{\Delta }^{}}^{(1620)}}{m_\pi }\frac{m^{}m}{m^{}m}(3k_3^2\stackrel{}{k}^2)`$ |
| $`<p,\frac{1}{2}|_{\pi NN}^{(1650)}|N(1650)^+,\frac{1}{2}>`$ | $`if_{\pi NN^{}}^{(1650)}\frac{m^{}m}{m_\pi }|\{\frac{i}{4}\frac{f_{\pi NN^{}}^{(1650)}}{m_\pi }\frac{m^{}+m}{m^{}m}\stackrel{}{k}^2\}`$ |
| $`<p,\frac{1}{2}|_{\pi NN^{}}^{(1700)}|N(1700)^+,\frac{1}{2}>`$ | $`\frac{i}{4\sqrt{6}}\frac{f_{\pi NN^{}}^{(1700)}}{m_\pi }\frac{m^{}m}{m^{}m}(3k_3^2\stackrel{}{k}^2)`$ |
| $`<p,\frac{1}{2}|_{\pi NN^{}}^{(1675)}|N(1675)^+,\frac{1}{2}>`$ | $`\frac{i}{\sqrt{10}}\frac{f_{\pi NN^{}}^{(1675)}}{m_\pi }\frac{(3k_3^2\stackrel{}{k}^2)}{m_\pi }`$ |
| $`<p,\frac{1}{2}|_{\pi NN^{}}^{(1720)}|N(1720)^+,\frac{1}{2}>`$ | $`i\sqrt{\frac{2}{3}}\frac{f_{\pi NN^{}}^{(1720)}}{m_\pi }\frac{k_3}{m_\pi }`$ |
| $`<p,\frac{1}{2}|_{\pi NN^{}}^{(1680)}|N(1680)^+,\frac{1}{2}>`$ | $`i\frac{1}{4}\sqrt{\frac{2}{5}}\frac{f_{\pi NN^{}}^{(1680)}}{m_\pi ^2}\frac{m^{}+m}{m^{}m}k_3(5k_3^23k^2)`$ |
| Table 4. The resonance transition coupling constants in terms of the |
| --- |
| pion-nucleon pseudoscalar coupling constant $`f_{\pi NN}`$. |
| $`f_{\pi N\mathrm{\Delta }}^{(1232)}=\frac{6\sqrt{2}}{5}e^{k^2/6\omega ^2}f_{\pi NN}1.55f_{\pi NN}`$ |
| $`f_{\pi NN^{}}^{(1440)}=\frac{\sqrt{3}}{18}(\frac{k^2}{\omega ^2}+\frac{3}{2}\frac{\omega _\pi }{m_q})e^{k^2/6\omega ^2}f_{\pi NN}0.26f_{\pi NN}`$ |
| $`f_{\pi N\mathrm{\Delta }}^{(1600)}=\frac{\sqrt{6}}{10}(\frac{k^2}{\omega ^2}+\frac{3}{2}\frac{\omega _\pi }{m_q})e^{k^2/6\omega ^2}f_{\pi NN}0.47f_{\pi NN}`$ |
| $`f_{\pi NN^{}}^{(1535)}=\frac{2\sqrt{2}}{15}\frac{\omega }{(m^{}m_N)}(\frac{k^2}{\omega ^2}+\frac{9\omega }{2m_q})e^{k^2/6\omega ^2}f_{\pi NN}0.49f_{\pi NN}`$ |
| $`f_{\pi NN^{}}^{(1520)}=\frac{8\sqrt{6}}{15}\frac{m^{}m_N}{(m^{}m_N)\omega }e^{k^2/6\omega ^2}f_{\pi NN}1.71f_{\pi NN}`$ |
| $`f_{\pi N\mathrm{\Delta }}^{(1620)}=\frac{\sqrt{3}}{15}\frac{\omega }{(m^{}m_N)}(\frac{k^2}{\omega ^2}+\frac{9\omega }{2m_q})e^{k^2/6\omega ^2}f_{\pi NN}0.34f_{\pi NN}`$ |
| $`f_{\pi N\mathrm{\Delta }}^{(1700)}=\frac{12}{5}\frac{m^{}m_N}{(m^{}m_N)\omega }e^{k^2/6\omega ^2}f_{\pi NN}2.6f_{\pi NN}`$ |
| $`f_{\pi NN^{}}^{(1650)}=\frac{\sqrt{2}}{15}\frac{\omega }{(m^{}m_N)}(\frac{k^2}{\omega ^2}+\frac{9\omega }{2m_q})e^{k^2/6\omega ^2}f_{\pi NN}0.28f_{\pi NN}`$ |
| $`f_{\pi NN^{}}^{(1700)}=\frac{4\sqrt{15}}{75}\frac{m^{}m_N}{(m^{}m_N)\omega }e^{k^2/6\omega ^2}f_{\pi NN}0.22f_{\pi NN}`$ |
| $`f_{\pi NN^{}}^{(1675)}=\frac{10}{25}\frac{m_\pi }{\omega }e^{k^2/6\omega ^2}f_{\pi NN}0.09f_{\pi NN}`$ |
| $`f_{\pi NN^{}}^{(1720)}=\frac{\sqrt{10}}{30}(\frac{k^2}{\omega ^2}+\frac{15\omega _\pi }{2m_q})e^{k^2/6\omega ^2}f_{\pi NN}1.05`$ |
| $`f_{\pi NN^{}}^{(1680)}=\frac{\sqrt{2}}{6}\frac{m^{}}{(m^{}+m_N)}(\frac{m_N}{m_\pi })(\frac{m_\pi }{\omega })^2e^{k^2/6\omega ^2}f_{\pi NN}0.12`$ |
| Table 5. Transition matrix elements of the isoscalar quark charge $`_qe^{i\stackrel{}{k}\stackrel{}{r}_q}`$ and | | |
| --- | --- | --- |
| transverse spin current $`_q\stackrel{}{\sigma }^q(\stackrel{}{k}\times \stackrel{}{ϵ})e^{i\stackrel{}{k}\stackrel{}{r}_q}`$ operators between the proton and nucleon | | |
| resonances for charge states $`+1`$ with $`s_z=+1/2`$. | | |
| transition | charge | spin current |
| $`<p,\frac{1}{2}|O|p,\frac{1}{2}>`$ | 3 | $`(\stackrel{}{k}\times \stackrel{}{ϵ})_3`$ |
| $`<p,\frac{1}{2}|O|N(1440)^+,\frac{1}{2}>`$ | $`\frac{\sqrt{3}}{6}\frac{\stackrel{}{k}^2}{\omega ^2}e^{k^2/6\omega ^2}`$ | $`\frac{\sqrt{3}}{18}\frac{k^2}{\omega ^2}e^{k^2/6\omega ^2}(\stackrel{}{k}\times \stackrel{}{ϵ})_3`$ |
| $`<p,\frac{1}{2}|O|N(1535)^+,\frac{1}{2}>`$ | 0 | $`\frac{\sqrt{2}}{3}\frac{\stackrel{}{k}^2}{\omega }e^{k^2/6\omega ^2}ϵ_3`$ |
| $`<p,\frac{1}{2}|O|N(1520)^+,\frac{1}{2}>`$ | 0 | $`i\frac{2}{3}\frac{1}{\omega }e^{k^2/6\omega ^2}[\stackrel{}{k}_3(\stackrel{}{k}\times \stackrel{}{ϵ})_3+k_+(\stackrel{}{k}\times \stackrel{}{ϵ})_{}]`$ |
| $`<p,\frac{1}{2}|O|N(1650)^+,\frac{1}{2}>`$ | 0 | $`\frac{\sqrt{2}}{6}\frac{\stackrel{}{k}^2}{\omega }e^{k^2/6\omega ^2}ϵ_3`$ |
| $`<p,\frac{1}{2}|O|N(1700)^+,\frac{1}{2}>`$ | 0 | $`i\frac{\sqrt{10}}{15}\frac{1}{\omega }e^{k^2/6\omega ^2}[k_3(\stackrel{}{k}\times \stackrel{}{ϵ})_3+k_+(\stackrel{}{k}\times \stackrel{}{ϵ})_{}`$ |
| | | $`+3iϵ_3\stackrel{}{k}^2]`$ |
| $`<p,\frac{1}{2}|O|N(1675)^+,\frac{1}{2}>`$ | 0 | $`i\frac{3\sqrt{10}}{10}\frac{k_3}{\omega }e^{k^2/6\omega ^2}(\stackrel{}{k}\times \stackrel{}{ϵ})_3`$ |
| $`<p,\frac{1}{2}|O|N(1720)^+,\frac{1}{2}>`$ | $`0`$ | $`\frac{\sqrt{30}}{180}\frac{k^2}{\omega ^2}e^{k^2/6\omega ^2}(\stackrel{}{k}\times \stackrel{}{ϵ})_3`$ |
| $`<p,\frac{1}{2}|O|N(1680)^+,\frac{1}{2}>`$ | $`\frac{\sqrt{10}}{20}\frac{3k_3^2\stackrel{}{k}^2}{\omega ^2}e^{k^2/6\omega ^2}`$ | $`i\frac{\sqrt{30}}{45}\frac{k^3}{\omega ^2}e^{k^2/6\omega ^2}`$ |
| | | $`_q(3,1,q,q|3,0)\sqrt{\frac{4\pi }{7}}Y_{3q}(\widehat{k})\stackrel{}{ϵ}_q`$ |
| Table 6. Transition matrix elements of the isovector quark charge $`_q\tau _3^qe^{i\stackrel{}{k}\stackrel{}{r}_q}`$ | | |
| --- | --- | --- |
| and transverse current operator $`_q\tau _3^q\stackrel{}{\sigma }^q(\stackrel{}{k}\times \stackrel{}{ϵ})e^{i\stackrel{}{k}\stackrel{}{r}_q}`$ between the proton and | | |
| nucleon resonances with charge state $`+1`$ with $`s_z=+1/2`$. | | |
| transition | charge | spin current |
| $`<p,\frac{1}{2}|O|p,\frac{1}{2}>`$ | 1 | $`\frac{5}{3}(\stackrel{}{k}\times \stackrel{}{ϵ})_3`$ |
| $`<p,\frac{1}{2}|O|\mathrm{\Delta }(1232)^+,\frac{1}{2}>`$ | 0 | $`\frac{4\sqrt{2}}{3}e^{k^2/6\omega ^2}(\stackrel{}{k}\times \stackrel{}{ϵ})_3`$ |
| $`<p,\frac{1}{2}|O|N(1440)^+,\frac{1}{2}>`$ | $`\frac{\sqrt{3}}{18}\frac{k^2}{\omega ^2}e^{k^2/6\omega ^2}`$ | $`\frac{5\sqrt{3}}{54}\frac{\stackrel{}{k}^2}{\omega ^2}e^{k^2/6\omega ^2}(\stackrel{}{k}\times \stackrel{}{ϵ})_3`$ |
| $`<p,\frac{1}{2}|O|\mathrm{\Delta }(1600)^+,\frac{1}{2}>`$ | 0 | $`\frac{3\sqrt{6}}{27}\frac{k^2}{\omega ^2}e^{\stackrel{}{k}^2/6\omega ^2}(\stackrel{}{k}\times \stackrel{}{ϵ})_3`$ |
| $`<p,\frac{1}{2}|O|N(1535)^+,\frac{1}{2}>`$ | $`i\frac{\sqrt{2}}{3}\frac{k_3}{\omega }e^{k^2/6\omega ^2}`$ | $`\frac{2\sqrt{2}}{9}\frac{\stackrel{}{k}^2}{\omega }e^{k^2/6\omega ^2}ϵ_3`$ |
| $`<p,\frac{1}{2}|O|N(1520)^+,\frac{1}{2}>`$ | 0 | $`i\frac{4}{9}\frac{1}{\omega }e^{k^2/6\omega ^2}[k_3(\stackrel{}{k}\times \stackrel{}{ϵ})_3+k_+(\stackrel{}{k}\times \stackrel{}{ϵ})_{}]`$ |
| $`<p,\frac{1}{2}|O|\mathrm{\Delta }(1620)^+,\frac{1}{2}>`$ | $`i\frac{\sqrt{2}}{3}\frac{k_3}{\omega }e^{k^2/6\omega ^2}`$ | $`\frac{\sqrt{2}}{18}\frac{\stackrel{}{k}^2}{\omega }e^{k^2/6\omega ^2}ϵ_3`$ |
| $`<p,\frac{1}{2}|O|\mathrm{\Delta }(1700)^+,\frac{1}{2}>`$ | 0 | $`\frac{i}{9}\frac{1}{\omega }e^{k^2/6\omega ^2}[k_3(\stackrel{}{k}\times \stackrel{}{ϵ})_3+k_+(\stackrel{}{k}\times \stackrel{}{ϵ})_{}]`$ |
| $`<p,\frac{1}{2}|O|N(1650)^+,\frac{1}{2}>`$ | $`\frac{i}{6}\frac{k_3}{\omega }e^{k^2/6\omega ^2}`$ | $`\frac{\sqrt{2}}{18}\frac{\stackrel{}{k}^2}{\omega }e^{k^2/6\omega ^2}ϵ_3`$ |
| $`<p,\frac{1}{2}|O|N(1700)^+,\frac{1}{2}>`$ | 0 | $`i\frac{\sqrt{10}}{45}\frac{1}{\omega }e^{k^2/6\omega ^2}[k_3(\stackrel{}{k}\times \stackrel{}{ϵ})_3+k_+(\stackrel{}{k}\times \stackrel{}{ϵ})_{}`$ |
| | | $`+3iϵ_3\stackrel{}{k}^2]`$ |
| $`<p,\frac{1}{2}|O|N(1675)^+,\frac{1}{2}>`$ | 0 | $`i\frac{\sqrt{10}}{10}\frac{g_{\rho NN^{}}^{(1675)}}{m_\rho ^2}k_3(\stackrel{}{k}\times \stackrel{}{ϵ})_3`$ |
| $`<p,\frac{1}{2}|O|N(1720)^+,\frac{1}{2}>`$ | $`0`$ | $`\frac{\sqrt{30}}{108}\frac{k^2}{\omega ^2}e^{k^2/6\omega ^2}(\stackrel{}{k}\times \stackrel{}{ϵ})_3`$ |
| $`<p,\frac{1}{2}|O|N(1680)^+,\frac{1}{2}>`$ | $`\frac{\sqrt{10}}{60}\frac{3k_3^2\stackrel{}{k}^2}{\omega ^2}e^{k^2/6\omega ^2}`$ | $`i\frac{\sqrt{30}}{27}\frac{k^3}{\omega ^2}e^{k^2/6\omega ^2}`$ |
| | | $`_q(3,1,q,q|3,0)\sqrt{\frac{4\pi }{7}}Y_{3q}(\widehat{k})\stackrel{}{ϵ}_q`$ |
| Table 7. Transition matrix elements $`<p,\frac{1}{2}|_{\omega NN^{}}|N^+,\frac{1}{2}>`$ of the $`\omega `$-meson | | |
| --- | --- | --- |
| transition couplings in eqs. (3.6). Here $`\mu `$ is defined as | | |
| $`\mu =2m_Nm^{}/(m_N+m^{})`$. | | |
| transition | charge | spin current |
| $`<p,\frac{1}{2}|_{\omega NN^{}}^{(1440)}|N(1440)^+,\frac{1}{2}>`$ | $`\frac{\stackrel{}{k}^2}{m_\omega ^2}g_{\omega NN^{}}^{(1440)}`$ | $`i\frac{1}{2\mu }\frac{\stackrel{}{k}^2}{m_\omega ^2}(\stackrel{}{k}\times \stackrel{}{ϵ})_3g_{\omega NN^{}}^{(1440)}`$ |
| $`<p,\frac{1}{2}|_{\omega NN^{}}^{(1535)}|N(1535)^+,\frac{1}{2}>`$ | 0 | $`\frac{\stackrel{}{k}^2}{m_\omega ^2}ϵ_3g_{\omega NN^{}}^{(1535)}`$ |
| $`<p,\frac{1}{2}|_{\omega NN^{}}^{(1520)}|N(1520)^+,\frac{1}{2}>`$ | $`0`$ | $`i\frac{\sqrt{6}}{3m_\omega ^2}[k_3(\stackrel{}{k}\times \stackrel{}{ϵ})_3+k_+(\stackrel{}{k}\times \stackrel{}{ϵ})_{}]g_{\omega NN^{}}^{(1520)}`$ |
| $`<p,\frac{1}{2}|_{\omega NN^{}}^{(1650)}|N(1650)^+,\frac{1}{2}>`$ | 0 | $`\frac{\stackrel{}{k}^2}{m_\omega ^2}ϵ_3g_{\omega NN^{}}^{(1650)}`$ |
| $`<p,\frac{1}{2}|_{\omega NN^{}}^{(1700)}|N(1700)^+,\frac{1}{2}>`$ | $`0`$ | $`i\frac{\sqrt{6}}{3m_\omega ^2}[k_3(\stackrel{}{k}\times \stackrel{}{ϵ})_3+k_+(\stackrel{}{k}\times \stackrel{}{ϵ})_{}]g_{\omega NN^{}}^{(1700)}`$ |
| $`<p,\frac{1}{2}|_{\omega NN^{}}^{(1675)}|N(1675)^+,\frac{1}{2}>`$ | 0 | $`i\frac{3\sqrt{10}}{10}\frac{1}{m_\omega ^2}k_3(\stackrel{}{k}\times \stackrel{}{ϵ})_3g_{\omega NN^{}}^{(1675)}`$ |
| $`<p,\frac{1}{2}|_{\omega NN^{}}^{(1720)}|N(1720)^+,\frac{1}{2}>`$ | $`0`$ | $`i\frac{1}{2\mu }\sqrt{\frac{2}{3}}\frac{\stackrel{}{k}^2}{m_\omega ^2}(\stackrel{}{k}\times \stackrel{}{ϵ})_3g_{\omega NN^{}}^{(1720)}`$ |
| $`<p,\frac{1}{2}|_{\omega NN^{}}^{(1680)}|N(1680)^+,\frac{1}{2}>`$ | $`\frac{\sqrt{10}}{10}\frac{3k_3^2\stackrel{}{k}^2}{m_\omega ^2}g_{\omega NN^{}}^{(1680)}`$ | $`\frac{1}{\mu }\frac{\sqrt{30}}{15}\frac{\stackrel{}{k}^3}{m_\omega ^2}g_{\omega NN^{}}^{(1680)}`$ |
| | | $`(3,1,q,q|3,0)\sqrt{\frac{4\pi }{7}}Y_{3q}(\widehat{k})ϵ_q`$ |
| Table 8. The $`\omega `$-meson transition coupling constants in terms |
| --- |
| of the $`\omega NN`$ coupling constant $`g_{\omega NN}`$. The numerical values |
| correspond to $`|\stackrel{}{k}|=0`$. |
| $`g_{\omega NN^{}}^{(1440)}=\frac{\sqrt{3}}{18}\frac{m_\omega ^2}{\omega ^2}e^{\stackrel{}{k}^2/6\omega ^2}g_{\omega NN}5.5`$ |
| $`g_{\omega NN^{}}^{(1535)}=\frac{\sqrt{2}}{18}\frac{m_\omega ^2}{\omega m_q}e^{\stackrel{}{k}^2/6\omega ^2}g_{\omega NN}4.5`$ |
| $`g_{\omega NN^{}}^{(1520)}=\frac{\sqrt{6}}{18}\frac{m_\omega ^2}{\omega m_q}e^{\stackrel{}{k}^2/6\omega ^2}g_{\omega NN}7.7`$ |
| $`g_{\omega NN^{}}^{(1650)}=\frac{\sqrt{2}}{36}\frac{m_\omega ^2}{\omega m_q}e^{\stackrel{}{k}^2/6\omega ^2}g_{\omega NN}2.2`$ |
| $`g_{\omega NN^{}}^{(1700)}=\frac{\sqrt{15}}{90}\frac{m_\omega ^2}{\omega m_q}e^{\stackrel{}{k}^2/6\omega ^2}g_{\omega NN}2.4`$ |
| $`g_{\omega NN^{}}^{(1675)}=\frac{1}{6}\frac{m_\omega ^2}{\omega m_q}e^{\stackrel{}{k}^2/6\omega ^2}g_{\omega NN}9.4`$ |
| $`g_{\omega NN^{}}^{(1720)}=\frac{\sqrt{5}}{180}\frac{\mu }{m_q}\frac{m_\omega ^2}{\omega ^2}e^{\stackrel{}{k}^2/6\omega ^2}g_{\omega NN}2.7`$ |
| $`g_{\omega NN^{}}^{(1680)}=\frac{1}{18}\frac{\mu }{m_q}\frac{m_\omega ^2}{\omega ^2}e^{\stackrel{}{k}^2/6\omega ^2}g_{\omega NN}12.2`$ |
| Table 9. Transition matrix elements $`<p,\frac{1}{2}|_{\rho NN^{}}|N^+,\frac{1}{2}>`$ of the $`\rho `$-meson | | |
| --- | --- | --- |
| transition couplings in eqs. (3.7) Here $`\mu `$ is defined as | | |
| $`\mu =2m^{}m_N/(m^{}+m_N)`$. | | |
| transition | charge | spin current |
| $`<p,\frac{1}{2}|_{\rho N\mathrm{\Delta }}|\mathrm{\Delta }(1232)^+,\frac{1}{2}>`$ | 0 | $`i\frac{1}{3\mu }(\stackrel{}{k}\times \stackrel{}{ϵ})_3g_{\rho N\mathrm{\Delta }}`$ |
| $`<p,\frac{1}{2}|_{\rho NN^{}}^{(1440)}|N(1440)^+,\frac{1}{2}>`$ | $`\frac{\stackrel{}{k}^2}{m_\rho ^2}g_{\rho NN^{}}^{(1440)}`$ | $`i\frac{1}{2\mu }\frac{\stackrel{}{k}^2}{m_\rho ^2}(\stackrel{}{k}\times \stackrel{}{ϵ})_3g_{\rho NN^{}}^{(1440)}(1+\kappa _{\rho NN}^{(1440)})`$ |
| $`<p,\frac{1}{2}|_{\rho N\mathrm{\Delta }^{}}^{(1600)}|\mathrm{\Delta }(1600)^+,\frac{1}{2}>`$ | 0 | $`i\frac{1}{3\mu }\frac{\stackrel{}{k}^2}{m_\rho ^2}(\stackrel{}{k}\times \stackrel{}{ϵ})_3g_{\rho N\mathrm{\Delta }^{}}^{(1600)}`$ |
| $`<p,\frac{1}{2}|_{\rho NN^{}}^{(1535)}|N(1535)^+,\frac{1}{2}>`$ | $`0`$ | $`\frac{\stackrel{}{k}^2}{m_\rho ^2}ϵ_3g_{\rho NN^{}}^{(1535)}`$ |
| $`<p,\frac{1}{2}|_{\rho NN^{}}^{(1520)}|N(1520)^+,\frac{1}{2}>`$ | $`0`$ | $`i\frac{\sqrt{6}}{3m_\rho ^2}[k_3(\stackrel{}{k}\times \stackrel{}{ϵ})_3+k_+(\stackrel{}{k}\times \stackrel{}{ϵ})_{}]g_{\rho NN^{}}^{(1520)}`$ |
| $`<p,\frac{1}{2}|_{\rho N\mathrm{\Delta }^{}}^{(1620)}|\mathrm{\Delta }(1620)^+,\frac{1}{2}>`$ | $`0`$ | $`\sqrt{\frac{2}{3}}\frac{\stackrel{}{k}^2}{m_\rho ^2}g_{\rho N\mathrm{\Delta }^{}}^{(1620)}ϵ_3`$ |
| $`<p,\frac{1}{2}|_{\rho N\mathrm{\Delta }^{}}^{(1700)}|\mathrm{\Delta }(1700)^+,\frac{1}{2}>`$ | $`0`$ | $`i\frac{2}{3m_\rho ^2}[k_3(\stackrel{}{k}\times \stackrel{}{ϵ})_3+k_+(\stackrel{}{k}\times \stackrel{}{ϵ})_{}]g_{\rho N\mathrm{\Delta }^{}}^{(1700)}`$ |
| $`<p,\frac{1}{2}|_{\rho NN^{}}^{(1650)}|N(1650)^+,\frac{1}{2}>`$ | $`0`$ | $`\frac{\stackrel{}{k}^2}{m_\rho ^2}ϵ_3g_{\rho N\mathrm{\Delta }^{}}^{(1650)}`$ |
| $`<p,\frac{1}{2}|_{\rho NN^{}}^{(1700)}|N(1700)^+,\frac{1}{2}>`$ | $`0`$ | $`i\frac{\sqrt{6}}{3m_\rho ^2}[k\times \stackrel{}{ϵ})_3+k_+(\stackrel{}{k}\times \stackrel{}{ϵ})_{}]g_{\rho NN^{}}^{(1700)}`$ |
| $`<p,\frac{1}{2}|_{\rho NN^{}}^{(1675)}|N(1675)^+,\frac{1}{2}>`$ | 0 | $`i\frac{3\sqrt{10}}{10}\frac{1}{m_\rho ^2}k_3(\stackrel{}{k}\times ϵ)_3g_{\rho NN^{}}^{(1675)}`$ |
| $`<p,\frac{1}{2}|_{\rho NN^{}}^{(1720)}|N(1720)^+,\frac{1}{2}>`$ | $`0`$ | $`i\frac{1}{2\mu }\sqrt{\frac{2}{3}}\frac{\stackrel{}{k}^2}{m_\rho ^2}(\stackrel{}{k}\times \stackrel{}{ϵ})_3g_{\rho NN^{}}^{(1720)}`$ |
| $`<p,\frac{1}{2}|_{\rho NN^{}}^{(1680)}|N(1680)^+,\frac{1}{2}>`$ | $`\frac{\sqrt{10}}{10}\frac{3k_3^2\stackrel{}{k}^2}{m_\rho ^2}g_{\rho NN^{}}^{(1680)}`$ | $`\frac{1}{\mu }\frac{\sqrt{30}}{15}\frac{\stackrel{}{k}^3}{m_\rho ^2}g_{\rho NN^{}}^{(1680)}`$ |
| | | $`(3,1,q,q|3,0)\sqrt{\frac{4\pi }{7}}Y_{3q}(\widehat{k})ϵ_q`$ |
| Table 10. The $`\rho `$-meson transition coupling constant in terms of the |
| --- |
| $`\rho NN`$ coupling constants $`g_{\rho NN}`$. The numerical values correspond |
| to $`|\stackrel{}{k}|=0`$. |
| $`g_{\rho N\mathrm{\Delta }}=\frac{6\sqrt{2}}{5}\frac{\mu }{m_q}e^{\stackrel{}{k}^2/6\omega ^2}g_{\rho NN}8.7`$ |
| $`g_{\rho NN^{}}^{(1440)}=\frac{\sqrt{3}}{18}\frac{m_\rho ^2}{\omega ^2}e^{\stackrel{}{k}^2/6\omega ^2}g_{\rho NN}1.76,\kappa _{\rho NN^{}}^{1440}=4`$ |
| $`g_{\rho N\mathrm{\Delta }^{}}^{(1600)}=\frac{\sqrt{6}}{10}\frac{\mu }{m_q}\frac{m_\rho ^2}{\omega ^2}e^{\stackrel{}{k}^2/6\omega ^2}g_{\rho NN}17.0`$ |
| $`g_{\rho NN^{}}^{(1535)}=\frac{\sqrt{2}}{9}\frac{m_\rho ^2}{\omega m_q}e^{\stackrel{}{k}^2/6\omega ^2}g_{\rho NN}2.9`$ |
| $`g_{\rho NN^{}}^{(1520)}=\frac{\sqrt{6}}{9}\frac{m_\rho ^2}{\omega m_q}e^{\stackrel{}{k}^2/6\omega ^2}g_{\rho NN^{}}4.5`$ |
| $`g_{\rho N\mathrm{\Delta }^{}}^{(1620)}=\frac{\sqrt{3}}{36}\frac{m_\rho ^2}{\omega m_q}e^{\stackrel{}{k}^2/6\omega ^2}g_{\rho NN}0.88`$ |
| $`g_{\rho N\mathrm{\Delta }^{}}^{(1700)}=\frac{1}{12}\frac{m_\rho ^2}{\omega m_q}e^{\stackrel{}{k}^2/6\omega ^2}g_{\rho NN}1.5`$ |
| $`g_{\rho NN^{}}^{(1650)}=\frac{\sqrt{2}}{36}\frac{m_\rho ^2}{\omega m_q}e^{\stackrel{}{k}^2/6\omega ^2}g_{\rho NN}0.72`$ |
| $`g_{\rho NN^{}}^{(1700)}=\frac{\sqrt{5}}{90}\frac{m_\rho ^2}{\omega m_q}e^{\stackrel{}{k}^2/6\omega ^2}g_{\rho NN}0.45`$ |
| $`g_{\rho NN^{}}^{(1675)}=\frac{1}{6}\frac{m_\rho ^2}{\omega m_q}e^{\stackrel{}{k}^2/6\omega ^2}g_{\rho NN}3.0`$ |
| $`g_{\rho NN^{}}^{(1720)}=\frac{\sqrt{5}}{36}\frac{\mu }{m_q}\frac{m_\rho ^2}{\omega ^2}e^{\stackrel{}{k}^2/6\omega ^2}g_{\rho NN}4.4`$ |
| $`g_{\rho NN^{}}^{(1680)}=\frac{5}{18}\frac{\mu }{m_q}\frac{m_\rho ^2}{\omega ^2}e^{\stackrel{}{k}^2/6\omega ^2}g_{\rho NN}19.6`$ |
|
warning/0005/astro-ph0005426.html
|
ar5iv
|
text
|
# Spectroscopic confirmation of a white dwarf companion to the B star 16 Dra
## 1 Introduction
Unresolved Sirius-type binary systems consisting of a white dwarf and a main sequence star (spectral type B$``$K) are difficult to identify optically, since the bright main sequence companion completely swamps the degenerate star’s flux. However, through the ROSAT Wide Field Camera (WFC, Pounds et al. 1993) and Extreme Ultraviolet Explorer (EUVE, Bowyer et al. 1994) surveys, EUV radiation with the spectral signature of a hot white dwarf has been detected originating from apparently inactive main sequence stars, giving a clue to the existence of a previously unidentified population of Sirius-type binaries. Over 20 new systems have now been identified (e.g. Barstow et al. 1994, Burleigh et al. 1997, Vennes et al. 1998). For companions of spectral type $``$A5 or later, far-ultraviolet spectra obtained with the International Ultraviolet Explorer (IUE) have been used to confirm the identifications, since the white dwarf is actually the brighter component at these wavelengths. Unfortunately, stars of spectral types O, B and early A will still dominate any emission from a white dwarf in the far-UV regime, and IUE or HST cannot be used to identify any putative degenerate companions to these objects.
Two bright B stars, $`\theta `$ Hya (HR3665) and y Pup (HR2875), were unexpectedly detected in the ROSAT and EUVE surveys. Since their soft X-ray and EUV colours were similar to many known hot white dwarfs, it was suspected that they too were hiding hot white dwarf companions. Fortunately, both EUV sources were bright enough to be observed by EUVE’s spectrometers. y Pup was observed by EUVE in 1996, $`\theta `$ Hya in 1998, and the formal discovery of these Sirius-type systems was subsequently reported by Vennes et al. (1997, y Pup), Burleigh & Barstow (1998, y Pup), and Burleigh & Barstow (1999, $`\theta `$ Hya).
White dwarf companions to B stars are of significant importance since they must have evolved from massive progenitors, perhaps close to the maximum mass for white dwarf progenitor stars. They are also likely to be significantly more massive than the mean for white dwarf stars in general ($`0.57M_{}`$, Bergeron et al. 1992). The value of the maximum mass for a white dwarf progenitor star, and hence the minimum mass for producing a Type II supernova through core collapse in a single star, is a long-standing astrophysical problem. Weidemann (1987) gives the limit as $`8M_{}`$ in his semi-empirical initial-final mass relation. Observationally, this limit is best set by the white dwarf companion to y Pup (HR2875). Echelle spectroscopy of this object by Vennes (2000) has recently revealed that this system comprises two main sequence B stars (B3.5V$`+`$B6V) in an eccentric $`15`$ day orbit, with the white dwarf forming a third, wider component. The white dwarf must then have evolved from a star more massive than B3.5V, $`5.5M_{}`$ (Vennes 2000). We also note that Berghöfer et al. (2000) have recently suggested that the spectroscopic companion to the B1.5IV star $`\lambda `$ Sco might be a hot ultramassive white dwarf ($`1.25M_{}<M_{\mathrm{WD}}<1.4M_{}`$), based on an excess of EUV and soft X-ray radiation detected in ROSAT and EUVE photometric observations. If the existence of a white dwarf in this system was confirmed it would obviously set the lower limit on the maximum mass for white dwarf progenitors at a value near Weidemann’s semi-empirical $`8M_{}`$ limit, although unfortunately $`\lambda `$ Sco is a close binary ($`P=5.959`$ days) and mass transfer may have taken place at some stage.
White dwarfs in Sirius-type binaries can also be used to investigate the relationship between the mass of a main sequence star and its white dwarf progeny, the initial-final mass relation. In particular, white dwarf $`+`$ B star binaries can be used to investigate the upper end of this relation. Likewise, if the two components can eventually be resolved and an astrometric mass determined for the white dwarf, these systems can potentially be used to investigate the high mass end of the theoretical white dwarf mass-radius relation, for which few observational data points currently exist (Vauclair et al. 1997, Provencal et al. 1998).
In addition to y Pup and $`\theta `$ Hya, another B star was also detected in the ROSAT WFC and EUVE surveys, 16 Dra (HD150100, B9.5V). In this paper we present an EUVE spectrum of 16 Dra, which proves that it too has a hot white dwarf companion.
## 2 The 16 Dra system
16 Dra (HR6184, HD150100, ADS10129C) is a V=5.51 B9.5V star in a visual triple system with two other early-type stars, 17 Dra A (HR6185, HD150117, ADS10129A, B9V, V=5.08) and 17 Dra B (HR6186, HD150118, ADS10129B, A1V, V=6.58). The Hipparcos Catalogue (Perryman et al. 1997) gives the separation of 17 Dra A&B as just 3.208 arcsec; 17 Dra A and 16 Dra are then separated by 90.17 arcsec (Fig.1). The Hipparcos parallax for 16 Dra is 8.16$`\pm `$0.55 mas, corresponding to a distance of 122.5 pc (114.8$``$131.4 pc), and the parallax for 17 Dra A is similar, 8.22$`\pm `$0.60 mas, corresponding to a distance of 121.6 pc (113.4$``$131.2 pc). No solution is given for 17 Dra B. The proper motions of 16 Dra and 17 Dra A are also very similar: for 16 Dra $`\mu `$<sub>α</sub>$`=`$$``$12.9$`\pm `$0.6 mas/yr and $`\mu `$<sub>δ</sub>$`=`$28.7$`\pm `$0.6 mas/yr; for 17 Dra A $`\mu `$<sub>α</sub>$`=`$$``$12.3$`\pm `$0.6 mas/yr and
$`\mu `$<sub>δ</sub>$`=`$27.4$`\pm `$0.6 mas/yr. All three stars are therefore almost certainly related. In that case, any white dwarf in the system must have descended from a progenitor more massive than the earliest-type star extant, B9V (17 Dra A).
## 3 Detection of EUV radiation from 16 Dra in the ROSAT WFC and EUVE surveys
The ROSAT EUV and X-ray all-sky surveys were conducted between July 1990 and January 1991; the mission and instruments are described elsewhere (e.g. Trümper 1992, Sims et al. 1990). 16 Dra is associated with the WFC source RE J1636$`+`$525, and was later also detected in the EUVE all-sky survey (conducted between July 1992 and January 1993). This source is also coincident with a ROSAT Position Sensitive Proportional Counter (PSPC) soft X-ray detection. The count rates from all three instruments and associated filters are given in Table 1. The WFC count rates are taken from the revised 2RE Catalogue (Pye et al. 1995), which was constructed using improved methods for source detection and background screening. The EUVE count rates are taken from the First EUVE Source Catalog (Bowyer et al. 1994). The source is not included in the revised Second EUVE Source Catalog (Bowyer et al. 1996, see discussion below). The PSPC count rate was obtained from the on-line ROSAT All Sky Survey Bright Source Catalogue source browser, maintained by the Max Planck Institute in Germany (Voges et al. 1999)<sup>1</sup><sup>1</sup>1http://www.rosat.mpe-garching.mpg.de/cats/src-browser/.
Fig.1 shows an optical image of the 16 Dra field from the Digitized Sky Survey<sup>2</sup><sup>2</sup>2http://ledas-www.star.le.ac.uk/DSSimage/, including the nearby pair 17 Dra A/B. Also shown are the ROSAT WFC, EUVE and PSPC source error boxes. Clearly, no obvious optical counterpart is visible or resolved within the intersection of these three boxes, other than 16 Dra itself.
The EUV and soft X-ray colours are similar to known hot white dwarfs, and the EUV radiation is too strong for it to be the result of UV leakage through the WFC filters, although in the EUVE source catalogs (e.g. Bowyer et al. 1994) it is flagged as such (and, indeed, omitted from the second EUVE source catalog as a result, Bowyer et al. 1996). Far-UV leakage is a known problem for EUVE, but the effect is almost negligible in the WFC, especially for a late B star like 16 Dra. Assuming a temperature of 10,000K, we estimate the far-UV leakage contribution to the WFC S2 flux at just 3$`\times `$10<sup>-5</sup> counts/sec, compared with the $``$0.05 counts/sec detected. Add the fact that the EUV detection is also coincident with a PSPC soft X-ray detection, and it can be safely assumed that it is real.
16 Dra is only detected in the soft 0.1$``$0.4 keV PSPC band; only one (rather unusual) white dwarf has ever been detected at higher energies (KPD0005$`+`$5105, Fleming et al. 1993). Most active stars are also hard X-ray sources, and indeed Berghöfer et al. (1996) found only three of the B stars detected by the ROSAT PSPC were not hard X-ray sources. Interestingly, these are y Pup, $`\theta `$ Hya (both of which have confirmed white dwarf companions) and 16 Dra. Therefore, Burleigh & Barstow (1999) confidently predicted that 16 Dra might also be hiding a hot white dwarf companion, and using the ROSAT count rates demonstrated that it probably had a surface temperature between 25,000$``$37,000K.
## 4 EUVE pointed observation and data reduction
16 Dra was observed twice by EUVE in dither mode, firstly for $``$220,000 sec between 1999 February 28th and 1999 March 7th, and then for $``$230,000 sec between 1999 June 27th and 1999 July 6th, giving a total exposure time of 453,985.5 sec. We have extracted the spectra from the detector images using standard IRAF procedures. Our general reduction techniques are described in earlier work (e.g. Barstow et al. 1997).
The target was only detected, weakly, in the short (70$``$190Å) wavelength spectrometer. To improve the signal/noise ratio, we combined the two observations before extracting this spectrum, and then binned the data by a factor 32. The resultant spectrum is shown in Fig.2. The flux distribution is similar in shape to the familiar EUV continuum expected from hot white dwarfs in this spectral region (Fig.2, inset). No emission lines are visible in the raw data (Fig.2, inset), despite the apparent excess of flux at $``$120Å in the binned spectrum.
The only stars other than white dwarfs whose photospheric EUV radiation has been detected by the ROSAT WFC and EUVE are the bright early B giants $`\beta `$ CMa (B1II$``$III, Cassinelli et al. 1996) and $`ϵ`$ CMa (B2II, Cohen et al. 1996). The photospheric continuum of $`ϵ`$ CMa is visible down to $``$300Å, although no continuum flux from $`\beta `$ CMa is visible below the HeI edge at 504Å. Both stars also have strong EUV and X-ray emitting winds, and in $`ϵ`$ CMa emission lines are seen by EUVE in the short and medium wavelength spectrometers from e.g. high ionisation features of iron. Similarly, strong narrow emission features of e.g. oxygen, nickel and calcium are commonly seen in EUV spectra of active stars and RS CVn systems. Since no such features are visible in this spectrum of 16 Dra, we can categorically rule out a hot wind or unresolved active late-type companion to 16 Dra as an alternative source for the EUV radiation.
## 5 Analysis of the hot white dwarf’s EUV spectrum
We have matched the EUV spectrum of 16 Dra with a grid of hot white dwarf $`+`$ ISM model atmospheres, in order to constrain the atmospheric parameters (temperature and surface gravity) of the degenerate star and the interstellar column densities of HI, HeI and HeII. Unfortunately there are no spectral features in this wavelength region to give us an unambiguous determination of $`T_{\mathrm{eff}}`$ and log $`g`$. However, by making a range of assumptions to reduce the number of free parameters in our models, we can place constraints on some of the white dwarf’s physical parameters. Our method is similar to that used in the analysis of y Pup (Burleigh & Barstow 1998) and $`\theta `$ Hya (Burleigh & Barstow 1999).
Firstly, we assume that the white dwarf has a pure-hydrogen atmosphere. This is a reasonable assumption to make, since Barstow et al. (1993) first showed that for T$`{}_{\mathrm{eff}}{}^{}<40,000`$K hot DA white dwarfs have an essentially pure-H atmospheric composition. We can then fit a range of models, each fixed at a discrete value of the surface gravity log $`g`$. However, before we can do this we need to know the normalisation parameter of each model, which is equivalent to $`(R_{\mathrm{WD}}/D)^2`$. For this, we can use the Hipparcos parallax to give us the distance, and the Hamada & Salpeter (1961) zero-temperature mass-radius relation to give us the radius of the white dwarf corresponding to each value of the surface gravity (see Table 2).
We can also reduce the number of unknown free parameters in the ISM model. From EUVE spectroscopy, Barstow et al. (1997) measured the line-of-sight interstellar column densities of HI, HeI and HeII to a number of hot white dwarfs. They found that the mean H ionisation fraction in the local ISM was 0.35$`\pm `$0.1, and the mean He fraction was 0.27$`\pm `$0.04. From these estimates, and assuming a cosmic H/He abundance, we calculate the ratio $`N_{\mathrm{HI}}/N_{\mathrm{HeI}}`$ in the local ISM$`=`$8.9 and $`N_{\mathrm{HeI}}/N_{\mathrm{HeII}}=2.7`$. We can then fix these column density ratios in our model, leaving us with just two free parameters - temperature and the HI column density.
The model fits at a range of surface gravities from log $`g=7.59.0`$ are summarized in Table 3. Note that our range of fitted temperatures is in agreement with those of Burleigh & Barstow (1999), who modelled the ROSAT EUV and soft X-ray photometric data for 16 Dra on the assumption that the source was indeed an unresolved white dwarf.
## 6 Discussion
We have analysed the weak EUVE spectrum of the B9.5V star 16 Dra, and confirm that there is an unresolved hot white dwarf in the field.
Fig.1 clearly shows that the white dwarf is not resolved from 16 Dra in the Digitized Sky Survey image, and their angular separation can be no more than $``$30 arcsec. However, if the white dwarf lies at the same distance as 16 Dra and its proper motion companions 17 Dra A & B, and is related to them, then it must have evolved from a progenitor more massive than the earliest extant star in the system, 17 Dra A (B9V). Thus this degenerate has the second most massive progenitor among known white dwarfs. Table 4 lists the earliest type stars known to have white dwarf companions, including all three B star $`+`$ white dwarf binaries and Sirius.
EUVE spectra provide us with little information with which to constrain a white dwarf’s surface gravity, and hence its mass, but we can use a theoretical initial-final mass relation between main sequence stars and white dwarfs, e.g. that of Wood (1992), to estimate the mass of the white dwarf if the progenitor was slightly more massive than a B9V star: $`M_{\mathrm{WD}}=Aexp(B\times M_{\mathrm{MS}})`$, where $`A=\mathrm{\hspace{0.17em}0.49}M_{}`$ and $`B=\mathrm{\hspace{0.17em}0.094}M_{}^1`$.
For $`M_{\mathrm{MS}}=\mathrm{\hspace{0.17em}3.7}M_{}`$, we find $`M_{\mathrm{WD}}=\mathrm{\hspace{0.17em}0.69}M_{}`$. This would suggest the surface gravity of the white dwarf log $`g>8.0`$ and, therefore, its surface temperature most likely lies between $`29,000`$K and $`35,000`$K.
Finally, we note that if this white dwarf can be resolved from 16 Dra, then an optical spectrum may potentially be obtained (e.g. with HST/STIS) from which its temperature and gravity can be tightly constrained. The mass can then be estimated, and this binary could be used to investigate the initial-final mass relation and to test the high mass end of the mass-radius relation.
###### Acknowledgements.
Matt Burleigh is the UK ROSAT Support Scientist and acknowledges the support of PPARC, UK. We thank Detlev Koester (Kiel) for the use of his model atmosphere grids, and Jurek Madej (Warsaw University) and Victor Bychkov (Special Astrophysical Observatory, Russian Academy of Sciences) for their help in obtaining and reducing the optical spectrum of 16 Dra. This research has made us of the SIMBAD database operated by CDS, Strasbourg, France, and the Leicester Database and Archive Service (LEDAS).
## 7 Appendix: HD93847 - another B star $`+`$ white dwarf binary in the ROSAT WFC catalogue?
A fourth B star, HD93847 (B9, V$`=`$7.46) was detected in the WFC survey, with a count rate of 9$`\pm `$4 counts/ksec in the S1 filter and 18$`\pm `$5 counts/ksec in the S2 filter. Since the S2/S1 count rate ratio is similar to known hot white dwarfs, it would be reasonable to suggest that perhaps this B star is also hiding a degenerate companion. The detected flux is highly unlikely to be due to far-UV leakage: we estimate this contribution as only 6$`\times `$10<sup>-6</sup> counts/sec. Unfortunately, this weak WFC source was not detected by the ROSAT PSPC or by EUVE, and thus it is not clear whether the detection is real. With a combined S1$`+`$S2 detection significance of only 5.8$`\sigma `$ it is in the regime where a few spurious detections are expected (see Pye et al. 1995). Alternatively, the B star may be hiding an unresolved active late-type companion that flared in the EUV waveband. We are therefore unwilling to claim that this source is due to another hidden white dwarf.
|
warning/0005/cond-mat0005419.html
|
ar5iv
|
text
|
# Strings in charge-transfer Mott insulators: effects of lattice vibrations and the Coulomb interaction
## Abstract
Applying the canonical transformation with the $`1/\lambda `$ perturbation expansion in the nonadiabatic and intermediate regime and the discrete generalisation of Pekar’s continuous nonlinear equation in the extreme adiabatic regime we show that there are no strings in narrow-band ionic insulators due to the Fröhlich electron-phonon interaction alone. The multi-polaron system is a homogeneous state in a wide range of physically interesting parameters, no matter how strong correlations are. At the same time the Fröhlich interaction allows the antiferromagnetic interactions and/or a short-range electron-phonon interactions to form short strings in doped antiferromagnetic insulators if the static dielectric constant is large enough.
The electron-phonon interaction is strong in ionic cuprates and manganites as established both experimentally and theoretically . The carriers, doped into the Mott insulator, are coupled with the antiferromagnetic background as well. The antiferromagnetic interactions are thought to give rise to spin and charge segregation (stripes) . There is growing experimental evidence that stripes occur in slightly doped insulators. Their theoretical studies were restricted so far to the repulsive strongly correlated models , or to an extreme adiabatic limit of the electron-phonon interaction in narrow and wide band polar semiconductors and polymers. On the other hand there is strong evidence that the nonadiabatic electron-phonon interaction and small polarons are involed in the physics of stripes . Also the role of the long-range Coulomb and Fröhlich interactions remains to be properly addressed.
In this letter we prove that the Fröhlich electron-phonon interaction combined with the direct Coulomb repulsion does not lead to charge segregation like strings in doped narrow-band insulators, both in the nonadiabatic and adiabatic regimes. However, this interaction significantly reduces the Coulomb repulsion, which might allow much weaker antiferromagnetic and/or short-range electron-phonon interactions to segregate charges in the doped insulators, as suggested by previous studies .
To begin with, we consider a generic Hamiltonian, including, respectively, the kinetic energy of carriers, the Fröhlich electron-phonon interaction, phonon energy, and the Coulomb repulsion as
$`H`$ $`=`$ $`{\displaystyle \underset{ij}{}}t(𝐦𝐧)\delta _{s,s^{}}c_i^{}c_j+{\displaystyle \underset{𝐪,i}{}}\omega _𝐪n_i[u_i(𝐪)d_𝐪+H.c.]`$ (1)
$`+`$ $`{\displaystyle \underset{𝐪}{}}\omega _𝐪(d_𝐪^{}d_𝐪+1/2)+{\displaystyle \frac{1}{2}}{\displaystyle \underset{ij}{}}V(𝐦𝐧)n_in_j`$ (2)
with bare hopping integral $`t(𝐦)`$, and matrix element of the electron-phonon interaction
$$u_i(𝐪)=\frac{1}{\sqrt{2N}}\gamma (𝐪)e^{i𝐪𝐦}.$$
(3)
Here $`i=(𝐦,s)`$, $`j=(𝐧,s^{})`$ include site $`𝐦,𝐧`$ and spin $`s,s^{}`$ quantum numbers, $`n_i=c_i^{}c_i`$, $`c_i,d_𝐪`$ are the electron (hole) and phonon operators, respectively, and $`N`$ is the number of sites. At large distances ( or small $`q`$) one finds
$$\gamma (𝐪)^2\omega _𝐪=\frac{4\pi e^2}{\kappa q^2},$$
(4)
and
$$V(𝐦𝐧)=\frac{e^2}{ϵ_{\mathrm{}}|𝐦𝐧|}.$$
(5)
The phonon frequency $`\omega _𝐪`$ and the static and high-frequency dielectric constants in $`\kappa ^1=ϵ_{\mathrm{}}^1ϵ_0^1`$ are those of the host insulator ($`\mathrm{}=c=1`$).
One can apply the canonical transformation and the $`1/\lambda `$ multi-polaron perturbation theory to integrate out phonons,
$$S=\underset{𝐪,i}{}n_i[u_i(𝐪)d_𝐪H.c.].$$
(6)
The result is
$`\stackrel{~}{H}`$ $`=`$ $`e^SHe^S={\displaystyle \underset{ij}{}}\widehat{\sigma }_{ij}c_i^{}c_jE_p{\displaystyle \underset{i}{}}n_i`$ (7)
$`+`$ $`{\displaystyle \underset{𝐪}{}}\omega _𝐪(d_𝐪^{}d_𝐪+1/2)+{\displaystyle \frac{1}{2}}{\displaystyle \underset{ij}{}}v_{ij}n_in_j,`$ (8)
where
$$\widehat{\sigma }_{ij}=t(𝐦𝐧)\delta _{s,s^{}}\mathrm{exp}(\underset{𝐪}{}[u_i(𝐪)u_j(𝐪)]d_𝐪H.c.)$$
(9)
is the renormalised hopping integral depending on the phonon variables, $`E_p=zt\lambda `$ is the polaron level shift and
$$v_{ij}=V(𝐦𝐧)\frac{1}{N}\underset{𝐪}{}\gamma (𝐪)^2\omega _𝐪\mathrm{cos}[𝐪(𝐦𝐧)]$$
(10)
is the net interaction of polarons comprising the long-range Coulomb repulsion and the long-range attraction due to ionic lattice deformations. Here $`\lambda =_𝐪\gamma (𝐪)^2\omega _𝐪/2Nzt`$ is the dimensionless coupling constant, $`t`$ is the nearest neighbour hopping integral and $`z`$ is the coordination lattice number.
The extention of the deformation surrounding (Fröhlich) polarons is large, so their deformation fields overlap at finite density. However, taking into account both the long-range attraction of polarons due to the lattice deformations $`and`$ the direct Coulomb repulsion, the net long-range interaction is repulsive . At distances larger than the lattice constant, $`|𝐦𝐧|a1`$, this interaction is significantly reduced to
$$v_{ij}=\frac{e^2}{ϵ_0|𝐦𝐧|}.$$
(11)
Optical phonons nearly nullify the bare Coulomb repulsion in ionic solids if $`ϵ_0>>1`$, which is normally the case in oxides. The kinetic energy term in the exact Hamiltonian, Eq.(6) involves multiphonon events generating a residual $`polaron`$-phonon interaction . Below we show that in the two opposite limits, the nonadiabatic ($`\omega _𝐪t`$) and in the extreme adiabatic ( $`\omega _𝐪0`$) regimes, there is no charge segregation or any other instability of the polaronic liquid due to the Fröhlich interaction in doped insulators, but only Wigner crystallization at very low densities.
First we consider the nonadiabtic and intermediate regime. The properties of a single small polaron with the Fröhlich electron-phonon interaction were discussed a long time ago . Exact Quantum Monte-Carlo simulations showed that the first order $`1/\lambda `$ perturbation theory is numerically accurate for $`any`$ coupling if the phonon frequency is sufficiently large, $`\omega _𝐪>t/2`$. The characteristic frequency of phonons strongly coupled with carriers is about $`\omega _𝐪=75`$ meV in cuprates, so cuprates are in this regime. Hence, one can replace the hopping operator in Eq.(6) for its phonon average, reducing the problem to narrow-band fermions with the weak repulsive interaction, Eq.(9). Next order corrections in $`1/\lambda `$ increase the polaron binding energy with little effect on the bandwidth . Because the net long-range repulsion is relatively weak, the relevant dimensionless parameter $`r_s=m^{}e^2/ϵ_0(4\pi n/3)^{1/3}`$ is not very large in doped cuprates. Wigner crystallization appears around $`r_s100`$ or larger, which corresponds to the atomic density of polarons $`n10^6`$ with $`ϵ_0=30`$ and the polaronic mass $`m^{}=5m_e`$ typical for cuprates and manganites. This estimate shows that small polarons in cuprates and manganites are in the homogeneous state at physically interesting densities.
In the opposite adiabatic limit one can apply a discrete version of the continuos nonlinear equation proposed in Ref. for the Holstein (molecular) model of the electron-phonon interaction and extended to the case of the deformation and Fröhlich interactions in Ref. . Applying the Hartree approximation for the Coulomb repulsion, the single-particle wave-function, $`\psi _𝐧`$ (the amplitude of the Wannier state $`|𝐧`$) obeys the following equation
$$\underset{𝐦0}{}t(𝐦)[\psi _𝐧\psi _{𝐧+𝐦}]e\varphi _𝐧\psi _𝐧=E\psi _𝐧.$$
(12)
The potential $`\varphi _{𝐧,k}`$ acting on a fermion $`k`$ at the site $`𝐧`$ is created by the polarization of the lattice $`\varphi _{𝐧,k}^l`$ and by the Coulomb repulsion with the other $`M1`$ fermions, $`\varphi _{𝐧,k}^c`$,
$$\varphi _{𝐧,k}=\varphi _{𝐧,k}^l+\varphi _{𝐧,k}^c.$$
(13)
Both potentials satisfy the descrete Poisson equation as
$$\kappa \mathrm{\Delta }\varphi _{𝐧,k}^l=4\pi e\underset{p=1}{\overset{M}{}}|\psi _{𝐧,p}|^2,$$
(14)
and
$$ϵ_{\mathrm{}}\mathrm{\Delta }\varphi _{𝐧,k}^c=4\pi e\underset{p=1,pk}{\overset{M}{}}|\psi _{𝐧,p}|^2,$$
(15)
with $`\mathrm{\Delta }\varphi _𝐧=_𝐦(\varphi _𝐧\varphi _{𝐧+𝐦})`$. Differently from Ref. we include the Coulomb interaction in Pekar’s functional $`J`$ , describing the total energy, in a selfconsistent manner using the Hartree approximation, so that
$`J`$ $`=`$ $`{\displaystyle \underset{𝐧,p,𝐦0}{}}\psi _{𝐧,p}^{}t(𝐦)[\psi _{𝐧,p}\psi _{𝐧+𝐦,p}]`$ (16)
$``$ $`{\displaystyle \frac{2\pi e^2}{\kappa }}{\displaystyle \underset{𝐧,p,𝐦,q}{}}|\psi _{𝐧,p}|^2\mathrm{\Delta }^1|\psi _{𝐦,q}|^2`$ (17)
$`+`$ $`{\displaystyle \frac{2\pi e^2}{ϵ_{\mathrm{}}}}{\displaystyle \underset{𝐧,p,𝐦,qp}{}}|\psi _{𝐧,p}|^2\mathrm{\Delta }^1|\psi _{𝐦,q}|^2.`$ (18)
If we assume, following Ref. that the single-particle function of a fermion trapped in a string of the length $`N`$ is a simple exponent, $`\psi _n=N^{1/2}\mathrm{exp}(ikn)`$ with the periodic boundary conditions, then the functional $`J`$ is expressed as $`J=T+U`$, where $`T=2t(N1)\mathrm{sin}(\pi M/N)/[N\mathrm{sin}(\pi /N)]`$ is the kinetic energy (for an $`odd`$ number $`M`$ of spinless fermions) , proportional to $`t`$, and
$$U=\frac{e^2}{\kappa }M^2I_N+\frac{e^2}{ϵ_{\mathrm{}}}M(M1)I_N,$$
(19)
corresponds to the polarisation and the Coulomb energies. Here the integral $`I_N`$ is given by
$`I_N`$ $`=`$ $`{\displaystyle \frac{\pi }{(2\pi )^3}}{\displaystyle _\pi ^\pi }𝑑x{\displaystyle _\pi ^\pi }𝑑y{\displaystyle _\pi ^\pi }𝑑z{\displaystyle \frac{\mathrm{sin}(Nx/2)^2}{N^2\mathrm{sin}(x/2)^2}}`$ (20)
$`\times `$ $`(3\mathrm{cos}x\mathrm{cos}y\mathrm{cos}z)^1.`$ (21)
It has the following asymptotics, Fig.1,
$$I_N=\frac{1.31+\mathrm{ln}N}{N},$$
(22)
which is also derived analyticially at large $`N`$ by the use of the fact that $`\mathrm{sin}(Nx/2)^2/(2\pi N\mathrm{sin}(x/2)^2)`$ can be replaced by a $`\delta `$\- function. If we split the first (attractive) term in Eq.(15) into two parts by replacing $`M^2`$ for $`M+M(M1)`$, then it becomes clear that the net interaction between polarons remains repulsive in the adiabatic regime as well because $`\kappa >ϵ_{\mathrm{}}`$. Hence, there are no strings within the Hartree approximation for the Coulomb interaction. Strong correlations do not change this conclusion. Indeed, if we take the Coulomb energy of spinless one-dimensional fermions comprising both Hartree and exchange terms as
$$E_C=\frac{e^2M(M1)}{Nϵ_{\mathrm{}}}[0.916+\mathrm{ln}M],$$
(23)
the polarisation and Coulomb energy per particle becomes (for large $`M>>1`$)
$$U/M=\frac{e^2M}{Nϵ_{\mathrm{}}}[0.916+\mathrm{ln}M\alpha (1.31+\mathrm{ln}N)],$$
(24)
where $`\alpha =1ϵ_{\mathrm{}}/ϵ_0<1`$. Minimising this energy with respect to the length of the string $`N`$ we find
$$N=M^{1/\alpha }\mathrm{exp}(0.31+0.916/\alpha ),$$
(25)
and
$$(U/M)_{min}=\frac{e^2}{\kappa }M^{11/\alpha }\mathrm{exp}(0.310.916/\alpha ).$$
(26)
Hence, the potential energy per particle increases with the number of particles so that the energy of $`M`$ well separated polarons is lower than the energy of polarons trapped in a string no matter correlated or not. The opposite conclusion of Ref. originates in an incorrect approximation of the integral $`I_NN^{0.15}/N`$. The correct asymptotic result is $`I_N=\mathrm{ln}(N)/N`$.
One can argue that a finite kinetic energy ($`t`$) can stabilise a string of a finite length. Unfortunately, this is not correct either. We performed exact (numerical) calculations of the total energy $`E(M,N)`$ of $`M`$ spinless fermions in a string of the length $`N`$ including both kinetic and potential energy with the typical values of $`ϵ_{\mathrm{}}=5`$ and $`ϵ_0=30`$. The local energy minima (per particle) in the string of the length $`1N69`$ containing $`MN/2`$ particles are presented in the Table. Strings with the even fermion numbers carry a finite current and hence the local minima are found for odd $`M`$. In the extreme wide band regime with $`t`$ as large as 1 eV the global string energy minimum is found at $`M=3,N=25`$ ($`E=2.1167`$ eV), and at $`M=3,N=13`$ for $`t=0.5`$ eV ($`E=1.2138`$ eV). However, this is $`not`$ the ground state energy in both cases. The energy of well separated $`d2`$-dimensional polarons is well below, less than $`2dt`$ per particle (i.e. $`6`$ eV in the first case and $`3`$ eV in the second one in the three dimensional cubic lattice, and $`4`$ eV and $`2`$ eV, respectively, in the two-dimensional square lattice). This argument is applied for any values of $`ϵ_0,ϵ_{\mathrm{}}`$ and $`t`$. As a result we have proved that strings are impossible with the Fröhlich interaction alone contrary to the erroneous Ref. .
The Fröhlich interaction is, of course, not the only electron-phonon interaction in ionic solids. As discussed in Ref. , any short range electron-phonon interaction, like, for example, the Jahn-Teller (JT) distortion can overcome the residual weak repulsion of Fröhlich polarons to form small bipolarons. At large distances small nonadibatic bipolarons weakly repel each other due to the long-range Coulomb interaction, four times of that of polarons, Eq.(9). Hence, they form a liquid state , or bipolaronic-polaronic crystal-like structures depending on their effective mass and density. The fact, that the Fröhlich interaction almost nullifies the Coulomb repulsion in oxides justifies the use of the Holstein-Hubbard model . The ground state of the 1D Holstein-Hubbard model is a liquid of intersite bipolarons with a significantly reduced mass (compared with the on-site bipolaron) as shown recently . The bound states of three or more polarons are not stable in this model, thus ruling out phase separation. However, the situation might be different if the antiferromagnetic and JT interaction or any short (but finite)-range electron-phonon interaction are strong enough. Due to long-range nature of the Coulomb repulsion the length of a string should be finite (see, also Ref.). One can readely estimate its length by the use of Eq.(8) for any type of the short-range electron-phonon interaction. If, for example, we take dispersive phonons, $`\omega _𝐪=\omega _0+\delta \omega (\mathrm{cos}q_x+\mathrm{cos}q_y+\mathrm{cos}q_z)`$ with a $`q`$-independent matrix element $`\gamma (𝐪)=\gamma `$, we obtain a short-range polaron-polaron attraction as
$$v_{att}(𝐧𝐦)=E_{att}(\delta \omega /\omega )\delta _{|𝐧𝐦|,1},$$
(27)
where $`E_{att}=\gamma ^2\omega _0/2`$. Taking into account the long-range repulsion as well, Eq.(9), the potential energy of the string with $`M=N`$ polarons becomes
$$U=\frac{e^2}{ϵ_0}N^2I_N\frac{NE_{att}\delta \omega }{\omega }.$$
(28)
Minimization of this energy yields the length of the string as
$$N=\mathrm{exp}\left(\frac{ϵ_0E_{att}\delta \omega }{e^2\omega }2.31\right).$$
(29)
Actually, this expression provides a fair estimate of the string length for any kind of attraction (not only generated by phonon dispersion), but also for the antiferromagnetic exchange and/or Jahn-Teller type of interactions . Due to the numerical coefficient in the exponent in Eq.(24) one can expect only short strings (if any) with the realistic values of $`E_{att}`$ (about a few hundreds millivolts), and the static dielectric constant $`ϵ_0100`$.
We conclude that there are no strings in ionic doped insulators with the Fröhlich interaction alone. Depending on their density and mass polarons remain in a liquid state or Wigner crystal. On the other hand the short-range electron-phonon and/or antiferromagnetic interactions might provide a liquid bipolaronic state and/or charge segregation (strings of a finite length) because the long-range Fröhlich interaction significantly reduces the Coulomb repulsion in highly polarizable ionic insulators.
We greatly appreciate enlightening discussions with Antonio Bianconi, Janez Bonca, Alex Bratkovsky, David Eagles, Pavel Kornilovitch, Fedor Kusmartsev, Dragan Mihailovic, and Jan Zaanen. One of us (V.V.K.) acknowledges support of the work by INTAS grant No.97-963.
Figure captions
Fig.1. The polarisation energy of small Fröhlich polarons trapped in a string depends on its length as $`\mathrm{ln}(N)/N`$.
|
warning/0005/nlin0005061.html
|
ar5iv
|
text
|
# Spectra of Random Contractions and Scattering Theory for Discrete-Time Systems
## Abstract
Random contractions (sub-unitary random matrices) appear naturally when considering quantized chaotic maps within a general theory of open linear stationary systems with discrete time. We analyze statistical properties of complex eigenvalues of generic $`N\times N`$ random matrices $`\widehat{A}`$ of such a type, corresponding to systems with broken time-reversal invariance. Deviations from unitarity are characterized by rank $`MN`$ and a set of eigenvalues $`0<T_i1,i=1,\mathrm{},M`$ of the matrix $`\widehat{T}=\widehat{\mathrm{𝟏}}\widehat{A}^{}\widehat{A}`$. We solve the problem completely by deriving the joint probability density of $`N`$ complex eigenvalues and calculating all $`n`$ point correlation functions. In the limit $`N>>M,n`$ the correlation functions acquire the universal form found earlier for weakly non-Hermitian random matrices.
The theory of wave scattering can be looked at as an integral part of the general theory of linear dynamic open systems in terms of the input-output approach. These ideas and relations were developed in system theory and engineering mathematics many years ago, see papers and references therein. Unfortunately, that development went almost unnoticed by the majority of physicists working in the theory of chaotic quantum scattering and related phenomena, see and references therein. For this reason we feel it could be useful to recall some basic facts of the input-output approach in such a context.
An Open Linear System is characterized by three Hilbert spaces: the space $`E_0`$ of internal states $`\mathrm{\Psi }E_0`$ and two spaces $`E_\pm `$ of incoming (-) and outgoing (+) signals or waves also called input and output spaces, made of vectors $`\varphi _\pm E_\pm `$. Acting in these three spaces are four operators, or matrices: a) the so-called fundamental operator $`\widehat{A}`$ which maps any vector from internal space $`E_0`$ onto some vector from the same space $`E_0`$ b) two operators $`\widehat{W}_{1,2}`$, with $`\widehat{W}_1`$ mapping incoming states onto an internal state and $`\widehat{W}_2`$ mapping internal states onto outgoing states and c) an operator $`\widehat{S}_0`$ acting from $`E_{}`$ to $`E_+`$.
We will be interested in describing the dynamics $`\mathrm{\Psi }(t)`$ of an internal state with time $`t`$ provided we know the state at initial instant $`t=0`$ and the system is subject to a given input signal $`\varphi _{}(t)`$. In what follows we consider only the case of the so-called stationary (or time-invariant) systems when the operators are assumed to be time-independent. Let us begin with the case of continuous-time description. The requirements of linearity, causality and stationarity lead to a system of two dynamical equations:
$$\begin{array}{c}i\frac{d}{dt}\mathrm{\Psi }=\widehat{A}\mathrm{\Psi }(t)+\widehat{W}_1\varphi _{}(t)\\ \varphi _+(t)=\widehat{S}_0\varphi _{}(t)+i\widehat{W}_2\mathrm{\Psi }(t)\end{array}$$
(1)
Interpretation of these equations depends on the nature of the state vector $`\mathrm{\Psi }`$ as well as of the vectors $`\varphi _\pm `$ and is different in different applications. In the context of quantum mechanics one relates the scalar product $`\mathrm{\Psi }^{}\mathrm{\Psi }`$ with the probability to find a particle inside the ”inner” region at time $`t`$, whereas $`\varphi _\pm ^{}\varphi _\pm `$ stays for probability currents flowing in and out of the region of internal states (the number of particles coming or leaving the inner domain per unit time). The condition of particle conservation then reads as:
$$\frac{d}{dt}\mathrm{\Psi }^{}\mathrm{\Psi }=\varphi _{}^{}\varphi _{}\varphi _+^{}\varphi _+$$
(2)
It is easy to verify that Eq.(2) is compatible with the dynamics Eq.(1) only provided the operators satisfy the following relations:
$`\widehat{A}^{}\widehat{A}=i\widehat{W}\widehat{W}^{},\widehat{S}_0^{}\widehat{S}_0=\widehat{\mathrm{𝟏}}\text{and}\widehat{W}^{}\widehat{W}_2=\widehat{S}_0\widehat{W}_1^{}`$
which shows, in particular, that $`\widehat{A}`$ can be written as $`A=\widehat{H}\frac{i}{2}\widehat{W}\widehat{W}^{}`$, with a Hermitian $`\widehat{H}=\widehat{H}^{}`$.
The meaning of $`\widehat{H}`$ is transparent: it governs the evolution $`i\frac{d}{dt}\mathrm{\Psi }=\widehat{H}\mathrm{\Psi }(t)`$ of an inner state $`\mathrm{\Psi }`$ when the coupling $`\widehat{W}`$ between the inner space and input/output spaces is absent. As such, it is just the Hamiltonian describing the closed inner region. The fundamental operator $`\widehat{A}`$ then has a natural interpretation of the effective non-selfadjoint Hamiltionian describing the decay of the probability from the inner region at zero input signal: $`\varphi _{}(t)=0`$ for any $`t0`$. If, however, the input signal is given in the Fourier-domain by $`\varphi _{}(\omega )`$, the output signal is related to it by:
$$\varphi _+(\omega )=\left[\widehat{S}(\omega )\widehat{S}_0\right]\varphi _{}(\omega ),\widehat{S}(\omega )=\widehat{\mathrm{𝟏}}i\widehat{W}^{}\frac{1}{\omega \widehat{\mathrm{𝟏}}\widehat{A}}\widehat{W}$$
(3)
where we assumed $`\mathrm{\Psi }(t=0)=0`$. The unitary matrix $`\widehat{S}(\omega )`$ is known in the mathematical literature as the characteristic matrix-function of the non-Hermitian operator $`\widehat{A}`$. In the present context it is just the scattering matrix whose unitarity is guaranteed by the conservation law Eq(2).
The contact with the theory of chaotic scattering is now apparent: the expression Eq.(3) was frequently used in the physical literature as a starting point for extracting universal properties of the scattering matrix for a quantum chaotic system within the so-called random matrix approach. The main idea underlying such an approach is to replace the actual Hamiltonian $`\widehat{H}`$ by a large random matrix and to calculate the ensuing statistics of the scattering matrix. The physical arguments in favor of such a replacement can be found in the cited literature.
In particular, most recently the statistical properties of complex eigenvalues of the operator $`\widehat{A}`$ as well as related quantities were studied in much detail. Those eigenvalues are poles of the scattering matrix and have the physical interpretation of resonances \- long-lived intermediate states to which discrete energy levels of the closed system are transformed due to coupling to continua.
In the theory presented above the time $`t`$ was a continuous parameter. On the other hand, a very useful instrument in the analysis of classical Hamiltonian systems with chaotic dynamics are the so-called area-preserving chaotic maps. They appear naturally either as a mapping of the Poincaré section onto itself, or as result of a stroboscopic description of Hamiltonians which are periodic functions of time. Their quantum mechanical analogues are unitary operators which act on Hilbert spaces of finite large dimension $`N`$. They are often referred to as evolution, scattering or Floquet operators, depending on the physical context where they are used. Their eigenvalues consist of $`N`$ points on the unit circle (eigenphases). Numerical studies of various classically chaotic systems suggest that the eigenphases conform statistically quite accurately the results obtained for unitary random matrices of a particular symmetry (Dyson circular ensembles).
Let us now imagine that a system represented by a chaotic map (”inner world”) is embedded in a larger physical system (”outer world”) in such a way that it describes particles which can come inside the region of chaotic motion and leave it after some time. Models of such type appeared, for example, in where a kind of scattering theory for ”open quantum maps” was developed based on a variant of Lipmann-Schwinger equation.
On the other hand, in the general system theory dynamical systems with discrete time are considered as frequently as those with continuous time. For linear systems a ”stroboscopic” dynamics is just a linear map $`(\varphi _{}(n);\mathrm{\Psi }(n))(\varphi _+(n);\mathrm{\Psi }(n+1))`$ which can be generally written as:
$$\left(\begin{array}{c}\mathrm{\Psi }(n+1)\\ \varphi _+(n)\end{array}\right)=\widehat{V}\left(\begin{array}{c}\mathrm{\Psi }(n)\\ \varphi _{}(n)\end{array}\right),\widehat{V}=\left(\begin{array}{cc}\widehat{A}& \widehat{W}_1\\ \widehat{W}_2& \widehat{S}_0\end{array}\right)$$
(4)
Again, we would like to consider a conservative system, and the discrete-time analogue of Eq.(2) is:
$`\mathrm{\Psi }^{}(n+1)\mathrm{\Psi }(n+1)\mathrm{\Psi }^{}(n)\mathrm{\Psi }(n)=\varphi _{}^{}(n)\varphi _{}(n)\varphi _+^{}(n)\varphi _+(n)`$
which amounts to unitarity of the matrix $`\widehat{V}`$ in Eq.(4). In view of such a unitarity $`\widehat{V}`$ of the type entering Eq.(4) can always be parametrized as (cf.):
$$\widehat{V}=\left(\begin{array}{cc}\widehat{u}_1& 0\\ 0& \widehat{v}_1\end{array}\right)\left(\begin{array}{cc}\sqrt{1\widehat{\tau }\widehat{\tau }^{}}& \widehat{\tau }\\ \widehat{\tau }^{}& \sqrt{1\widehat{\tau }^{}\widehat{\tau }}\end{array}\right)\left(\begin{array}{cc}\widehat{u}_2& 0\\ 0& \widehat{v}_2\end{array}\right)$$
(5)
where the matrices $`u_{1,2}`$ and $`v_{1,2}`$ are unitary and $`\widehat{\tau }`$ is a rectangular $`N\times M`$ diagonal matrix with the entries $`\tau _{ij}=\delta _{ij}\tau _j,\mathrm{\hspace{0.17em}1}iN,\mathrm{\hspace{0.17em}1}jM0\tau _j1`$.
Actually, it is frequently convenient to redefine input, output and internal state as: $`\varphi _{}(n)\widehat{v}_2^1\varphi _{}(n),\varphi _+\widehat{v}_1\varphi _+(n)`$ and $`\mathrm{\Psi }(n)\widehat{u}_2\mathrm{\Psi }(n)`$ which amounts just to choosing an appropriate basis in the corresponding spaces. The transformations bring the matrix $`\widehat{V}`$ to a somewhat simplier form:
$$\widehat{V}=\left(\begin{array}{cc}\widehat{u}\sqrt{1\widehat{\tau }\widehat{\tau }^{}}& u\widehat{\tau }\\ \widehat{\tau }^{}& \sqrt{1\widehat{\tau }^{}\widehat{\tau }}\end{array}\right)$$
(6)
where $`\widehat{u}=\widehat{u}_2^{}\widehat{u}_1`$. Such a form suggests a clear interpretation of the constituents of the model. Indeed, for $`\widehat{\tau }=0`$ the dynamics of the system amounts to: $`\mathrm{\Psi }(n+1)=\widehat{u}\mathrm{\Psi }(n)`$. We therefore identify $`\widehat{u}`$ as a unitary evolution operator of the ”closed” inner state domain decoupled both from input and output spaces. Correspondingly, $`\widehat{\tau }0`$ just provides a coupling that makes the system open and converts the fundamental operator $`\widehat{A}=\widehat{u}\sqrt{1\widehat{\tau }\widehat{\tau }^{}}`$ to a contraction: $`1\widehat{A}^{}\widehat{A}=\tau \tau ^{}0`$. As a result, the equation $`\mathrm{\Psi }(n+1)=\widehat{A}\mathrm{\Psi }(n)`$ describes an irreversible decay of any initial state $`\mathrm{\Psi }(0)`$ for zero input $`\varphi _{}(n)=0`$, whereas for a nonzero input and $`\mathrm{\Psi }(0)=0`$ the Fourier-transforms $`\varphi _\pm (\omega )=_{n=0}^{\mathrm{}}e^{in\omega }\varphi _\pm (n)`$ are related by a unitary scattering matrix $`\widehat{S}(\omega )`$ given by:
$$\widehat{S}(\omega )=\sqrt{1\widehat{\tau }^{}\widehat{\tau }}\widehat{\tau }^{}\frac{1}{e^{i\omega }\widehat{A}}\widehat{u}\widehat{\tau }$$
(7)
Assuming further that the motion outside the inner region is regular, we should be able to describe generic features of open quantized chaotic maps choosing the matrix $`\widehat{u}`$ to be a member of a Dyson circular ensemble. Then one finds: $`\widehat{\tau }^{}\widehat{\tau }=1\left|\overline{\widehat{S}(\omega )}\right|^2`$ , with the bar standing for the averaging of $`\widehat{S}(\omega )`$ in Eq.(7) over $`\widehat{u}`$. Comparing this result with we see that the $`M`$ eigenvalues $`0T_i1`$ of the $`M\times M`$ matrix $`\widehat{\tau }^{}\widehat{\tau }`$ play the role of the so-called transmission coefficients and describe a particular way the chaotic region is coupled to the outer world.
In fact, this line of reasoning is motivated by recent papers . The authors of considered the Floquet description of a Bloch particle in a constant force and periodic driving. After some approximations the evolution of the system is described by a mapping: $`𝐜_{n+1}=\mathrm{𝐅𝐜}_n`$, where the unitary Floquet operator $`𝐅=\widehat{S}\widehat{U}`$ is the product of a unitary ”M-shift” $`\widehat{S}:S_{kl}=\delta _{l,kM},l,k=\mathrm{},\mathrm{},\mathrm{}`$ and a unitary matrix $`\widehat{U}`$. The latter is effectively of the form $`\widehat{U}=\text{diag}(\widehat{d_1},\widehat{u},d_2)`$, where $`\widehat{d}_{1,2}`$ are (semi)infinite diagonal matrices and $`\widehat{u}`$ can be taken from the ensemble of random $`N\times N`$ unitary matrices.
One can check that such a dynamics can be easily brought to the standard Eqs.(4,5) with the fundamental operator being an $`N\times N`$ random matrix of the form $`\widehat{A}=\sqrt{\mathrm{𝟏}\widehat{\tau }^{}\widehat{\tau }}\widehat{u}`$, and all $`M`$ diagonal elements of the $`N\times M`$ matrix $`\tau `$ are equal to unity. Actually, the original paper employed a slightly different but equivalent construction dealing with an ”enlarged” internal space of the dimension $`N+M`$. We prefer to follow the general scheme because of its conceptual clarity.
Direct inspection immediately shows that the non-vanishing eigenvalues of the fundamental operator $`\widehat{A}`$ as above coincide with those of a $`(NM)\times (NM)`$ subblock of the random unitary matrix $`u`$. Complex eigenvalues of such ”truncations” of random unitary matrices were studied in much detail by the authors of a recent insightful paper . They managed to study eigenvalue correlations analytically for arbitrary $`N,M`$. In particular, they found that in the limit $`N\mathrm{}`$ for fixed $`M`$ these correlation functions practically coincide with those obtained earlier for operators of the form $`\widehat{A}=\widehat{H}\frac{i}{2}\widehat{W}\widehat{W}^{}`$ occuring in the theory of open systems with continuous-time dynamics.
Such a remarkable universality, though not completely unexpected, deserves to be studied in more detail. In fact, truncated unitary matrices represent only a particular case of random contractions $`\widehat{A}`$. Actually, some statistical properties of general subunitary matrices were under investigation recently as a model of scattering matrix for systems with absorption, see .
The particular case of rank-one deviations from unitarity is the simplest one to investigate and was considered in a recent preprint. However, generalization to arbitrary $`M`$ along the lines of seemed to be problematic. The main goal of the present paper is to suggest a regular way of studying the specta of random contractions for a given deviation from unitarity.
The ensemble of general $`N\times N`$ random contractions $`\widehat{A}=\widehat{u}\sqrt{1\widehat{\tau }\widehat{\tau }^{}}`$ describing the chaotic map with broken time-reversal symmetry can be described by the following probability measure in the matrix space:
$$𝒫(\widehat{A})d\widehat{A}\delta (\widehat{A}^{}\widehat{A}\widehat{G})d\widehat{A},\widehat{G}1\widehat{\tau }\widehat{\tau }^{}$$
(8)
where $`d\widehat{A}=_{ij}d\widehat{A}_{ij}dA_{ij}^{}`$ and we assumed that the unitary matrix $`\widehat{u}`$ is taken from the Dyson circular unitary ensemble. The $`N\times N`$ matrix $`\widehat{\tau }\widehat{\tau }^{}=\mathrm{𝟏}\widehat{G}0`$ is natural to call the deviation matrix and we denote it $`\widehat{T}`$. It has $`M`$ nonzero eigenvalues coinciding with the transmission coefficients $`T_a`$ introduced above. The particular choice $`T_{iM}=1,T_{i>M}=0`$ corresponds to the case considered in . In what follows we assume all $`T_i<1`$, but the resulting expressions turn out to be valid in the limiting case $`T_i=1`$ as well.
Our first step is, following , introduce the Schur decomposition $`\widehat{A}=\widehat{U}(\widehat{Z}+\widehat{R})\widehat{U}^{}`$ of the matrix $`A`$ in terms of a unitary $`\widehat{U}`$, diagonal matrix of the eigenvalues $`\widehat{Z}`$ and a lower triangular $`\widehat{R}`$. One can satisfy oneself, that the eigenvalues $`z_1,\mathrm{},z_N`$ are generically not degenerate, provided all $`T_i<1`$. Then, the measure written in terms of new variables is given by $`d\widehat{A}=|\mathrm{\Delta }(\{z\})|^2d\widehat{R}d\widehat{Z}d\mu (U)`$, where the first factor is just the Vandermonde determinant of eigenvalues $`z_i`$ and $`d\mu (U)`$ is the invariant measure on the unitary group. The joint probability density of complex eigenvalues is then given by:
$`𝒫(\{z\})`$ $``$ $`|\mathrm{\Delta }(\{z\})|^2{\displaystyle 𝑑\mu (U)𝑑\widehat{R}\delta \left((\widehat{Z}+\widehat{R})(\widehat{Z}+\widehat{R})^{}\widehat{U}^{}\widehat{G}\widehat{U}\right)}`$ (9)
The integration over $`\widehat{R}`$ can be performed with some manipulations using its triangularity (some useful hints can be found in ). As the result, we arrive at:
$`𝒫(\{z\})|\mathrm{\Delta }(\{z\})|^2{\displaystyle 𝑑\mu (U)\underset{l=1}{\overset{N}{}}\delta \left(|z_1|^2\mathrm{}|z_l|^2det\left[1\widehat{T}\widehat{U}\widehat{P}_l\widehat{U}^{}\right]\right)}`$ (10)
where $`\widehat{P}=\text{diag}(1,\mathrm{}1,0,\mathrm{},0)`$, with first $`l`$ entries being equal to unity, is a projector.
To perform the remaining integration over the unitary group was mentioned as the main technical problem in , and proved to be more difficult than the corresponding procedure for non-Hermitian matrices, see . Here we outline the main steps and present some intermediate expressions relegating details of the calculation to a more extended publication. First, one considers the columns of $`N\times N`$ unitary matrices as $`N`$ component of (mutually orthogonal) vectors $`𝐚_l,l=1,\mathrm{},N`$ and introduces $`N\times N`$ matrices of rank $`l`$: $`\widehat{Q}_l=_{i=1}^l𝐚_i𝐚_i^{}`$ , so that $`\widehat{U}\widehat{P}_l\widehat{U}^{}=\widehat{Q}_l`$. Writing down the corresponding constraints in a form of $`\delta `$ functions, one can represent the expression Eq.(9) in the form:
$`𝒫(\{z\})|\mathrm{\Delta }(\{z\})|^2{\displaystyle }\left({\displaystyle \underset{l=1}{\overset{N1}{}}}d\widehat{Q}_l\right){\displaystyle \underset{l=1}{\overset{N}{}}}\delta (|z_1|^2\mathrm{}|z_l|^2det(1\widehat{T}\widehat{Q}_l){\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle }d𝐚d𝐚^{}\delta (\widehat{Q}_i\widehat{Q}_{i1}𝐚𝐚^{})`$ (11)
where the matrices $`\widehat{Q}_l`$ are considered to be unconstrained $`N\times N`$ Hermitian. We also used the orthonormality condition $`\widehat{Q}_N=1`$ as well as the convention $`\widehat{Q}_0=0`$.
Due to the fact that only $`M`$ out of $`N`$ eigenvalues of the matrix $`\widehat{T}`$ are non-zero, both the matrices $`\widehat{Q}_l`$ and the vectors $`𝐚`$ can be effectively taken to be of the size $`M`$ ( it amounts to changing the unspecified normalisation constant in Eq.(11)) and redefine the matrix $`\widehat{T}`$ as $`\widehat{T}=\text{diag}(T_1,\mathrm{},T_M)=\widehat{\tau }^{}\widehat{\tau }`$. Then it is convenient to change: $`\widehat{T}^{1/2}\widehat{Q}_l\widehat{T}^{1/2}\widehat{Q}_l`$ and separate integration over eigenvalues and eigenvectors of matrices $`\widehat{Q}_l`$. The latter can be performed in a recursive way $`ll+1`$, with the multiple use of the Itzykson-Zuber-Harish-Chandra formula. After quite an elaborate manipulation, one finally arrives at the following representation:
$`𝒫(\{z\})`$ $``$ $`{\displaystyle \frac{det^{MN}(\widehat{T})}{det(1\widehat{T})_{c_1<c_2}\left(T_{c_1}T_{c_2}\right)}}{\displaystyle \underset{c_1<c_2}{}}\left({\displaystyle \frac{}{\tau _{c_1}}}{\displaystyle \frac{}{\tau _{c_2}}}\right){\displaystyle 𝑑\widehat{\lambda }e^{i\text{Tr}\widehat{\tau }\widehat{\lambda }}|\mathrm{\Delta }(\{z\})|^2\underset{k=1}{\overset{N}{}}f(\mathrm{ln}|z_k|^2,\widehat{\lambda })},`$ (12)
where we defined the diagonal matrices of size $`M`$ as: $`\widehat{\tau }=\text{diag}(\tau _1,\mathrm{},\tau _M),\widehat{\lambda }=\text{diag}(\lambda _1,\mathrm{},\lambda _M)`$ and used the notations: $`\tau _c=\mathrm{ln}(1T_c)`$ and
$$f(a,\widehat{\lambda })=i^{M1}\underset{q=1}{\overset{M}{}}\frac{e^{i\lambda _qa}}{_{s(q)}(\lambda _q\lambda _s)}.$$
(13)
The distribution Eq.(12) is written for $`|z_k|^21`$ for any $`k=1,\mathrm{},N`$ and vanishes otherwise. The remarkable feature of such a distribution is that it allows for calculation of all $`n`$point correlation functions for arbitary $`N,n,M`$ with help of the method of orthogonal polynomials. Again, the particular case $`M=1`$ is quite instructive and can be recommended to follow first for understanding of the general formulae outlined below.
To this end, we write
$`|\mathrm{\Delta }(\{z\})|^2{\displaystyle \underset{k=1}{\overset{N}{}}}f(\mathrm{ln}|z_k|^2,\widehat{\lambda })={\displaystyle \underset{k=1}{\overset{N}{}}}N_k(\widehat{\lambda })det\left[{\displaystyle \underset{n=1}{\overset{N}{}}}{\displaystyle \frac{(z_iz_j^{})^{n1}}{N_n(\widehat{\lambda })}}f(\mathrm{ln}|z_j|^2,\widehat{\lambda })\right]_{i,j=1,\mathrm{}N}.`$ (14)
where the constants $`N_n(\widehat{\lambda })`$ are provided by the orthonormality condition:
$$_{|z|^21}d^2zz^{m1}(z^{})^{n1}f(\mathrm{ln}|z|^2,\widehat{\lambda })=\delta _{m,n}N_n(\widehat{\lambda }),$$
(15)
which yields after a simple calculation $`N_n(\widehat{\lambda })=\pi _{c=1}^M\frac{1}{(n+i\lambda _c)}`$.
Now, by applying the standard machinery of orthogonal polynomials one can find the correlation function:
$`R_n(z_1,\mathrm{},z_n)={\displaystyle \frac{N!}{(Nn)!}}{\displaystyle d^2z_{n+1}\mathrm{}d^2z_N𝒫(\{z\})}`$ (16)
as given by:
$`R_n(z_1,\mathrm{},z_n)\widehat{𝒟}{\displaystyle 𝑑\widehat{\lambda }e^{i\text{Tr}\widehat{\tau }\widehat{\lambda }}\underset{k=1}{\overset{N}{}}N_k(\widehat{\lambda })det\left[K(z_i,z_j;\widehat{\lambda })\right]_{(i,j)=1,\mathrm{},n}},`$ (17)
where the kernel $`K`$ is defined as:
$$K(z_1,z_2;\widehat{\lambda })=\frac{1}{\pi }\underset{n=1}{\overset{N}{}}det(i\widehat{\lambda }+n)(z_1z_2^{})^{n1}f(\mathrm{ln}|z_2|^2,\widehat{\lambda })$$
(18)
and the differential operator $`\widehat{𝒟}`$ is just the expression in front of the $`\lambda `$ integral in Eq.(12).
In principle, all $`\lambda `$ integrations in the equation Eq.(17) can be performed explicitly and the resulting formulae provide the desired general solution of the problem. However, for arbitary $`N,M,n`$ the results obtained in that way are still quite cumbersome. We present below as an example the lowest correlation function $`R_1(z)`$, which is just the mean eigenvalue density inside the unit circle $`|z|<1`$. It can be calculated from the following recursive relation connecting the density for $`M`$ and $`M1`$ open channels:
$$R_1^{(M)}(z)=R_1^{(M1)}(z)+\frac{1}{\pi }\frac{}{|z|^2}_1^{(M)}\{T_c;|z|^2\}_2^{(M1)}\{T_c;|z|^2\},$$
(19)
where
$`_1^{(M)}\{T_c;|z|^2\}`$ $`=`$ $`{\displaystyle \underset{c=1}{\overset{M}{}}}\left[1\left({\displaystyle \frac{1}{|z|^2}}1\right)\left({\displaystyle \frac{1}{T_c}}1\right)\right]^{N1}{\displaystyle \frac{\theta (|z|^21+T_c)}{_{sc}\left(\frac{1}{T_s}\frac{1}{T_c}\right)}},`$ (20)
$`_2^{(M1)}\{T_c;|z|^2\}`$ $`=`$ $`{\displaystyle \frac{|z|^{2N}}{(N1)!}}{\displaystyle _0^{\mathrm{}}}𝑑te^{t|z|^2}t^{N1}{\displaystyle \underset{c=1}{\overset{M1}{}}}\left({\displaystyle \frac{1}{T_c}}1+{\displaystyle \frac{1}{t}}{\displaystyle \frac{}{|z|^2}}\right){\displaystyle \frac{1|z|^{2N}}{1|z|^2}}.`$ (21)
For the case of all equivalent channels, i.e. when all the transmission coefficients $`T_c`$ are equal: $`T_c=T`$, such a recursive relation can be represented in a more compact form:
$$R_1^{(M)}(z)=R_1^{(M1)}(z)+\frac{1}{(M1)!}\left(\frac{}{t}\right)^{M1}(t,z)|_{t=0},$$
(22)
where the generating function $`(t,z)`$ is given by:
$$(t,z)=\frac{1}{\pi }\frac{}{|z|^2}\frac{\xi ^N\eta ^N}{\xi \eta }$$
(23)
and
$$\xi =1+(t1)(\frac{1}{|z|^2}1)(\frac{1}{T}1),\eta =1+(t1)(1|z|^2)\frac{1}{T}.$$
(24)
All the equations above are valid for arbitrary $`NM,n`$. In the theory of quantum chaotic scattering we, however, expect a kind of universality in the semiclassical limit. Translated to the random matrix language such a limit corresponds to $`N\mathrm{}`$ at fixed $`n,M`$. Still, extracting the asymptotic behaviour of the correlation function $`R_n(z_1,\mathrm{},z_n)`$ from Eq.(17) in such a limit is not a completely straightforward task. A useful trick is to notice that Eq.(17) can be rewritten as:
$`R_k(z_1,\mathrm{},z_n)`$ $``$ $`{\displaystyle \underset{q_1=1,\mathrm{},q_n=1}{\overset{M}{}}}F_{q_1,\mathrm{},q_n}\left(\{T_c;z\}\right),`$ (25)
$`F_{q_1,\mathrm{},q_n}\left(\{T_c;z\}\right)`$ $`=`$ $`\{T_c\}det\left[{\displaystyle \underset{k=1}{\overset{N}{}}}\left({\displaystyle \frac{}{\tau _1}}+k\right)\mathrm{}\left({\displaystyle \frac{}{\tau _M}}+k\right)(z_iz_j^{})^{k1}\right]_{i,j=1,\mathrm{},n}`$ (26)
$`\times `$ $`{\displaystyle \underset{c_1<c_2}{}}\left({\displaystyle \frac{}{\tau _{c_1}}}{\displaystyle \frac{}{\tau _{c_2}}}\right){\displaystyle \underset{c}{}\left(d\lambda _c\frac{\mathrm{exp}i\lambda _c\tau _c}{_{l=1}^N(l+i\lambda _c)}\right)\frac{e^{i_{j=1}^n\lambda _{q_j}\mathrm{ln}|z_j|^2}}{_{j=1}^n_{sq_j}^M(\lambda _{q_j}\lambda _s)}},`$ (27)
where we introduced the notation:
$`\{T_c\}={\displaystyle \frac{1}{_{c_1<c_2}\left(T_{c_1}T_{c_2}\right)}}{\displaystyle \underset{c=1}{\overset{M}{}}}{\displaystyle \frac{T_c^{MN}}{(1T_c)}}.`$
Introducing now the auxiliary differential operator $`\widehat{𝒟}_{q_1,\mathrm{},q_n}=_{j=1}^n_{sq_j}^M\left(\frac{}{\tau _{q_j}}\frac{}{\tau _s}\right)`$ and considering its action upon the ratio $`F_{q_1,\mathrm{},q_n}/\{T_c\}`$ one can satisfy oneself that in the limit $`NM,n`$ the leading contribution to $`F_{q_1,\mathrm{},q_n}`$ is given by:
$`F_{q_1,\mathrm{},q_n}`$ $``$ $`{\displaystyle \underset{c=1}{\overset{M}{}}}\theta (1\stackrel{~}{T}_c){\displaystyle \frac{(1\stackrel{~}{T}_c)}{(1T_c)}}\left({\displaystyle \frac{\stackrel{~}{T}_c}{T_c}}\right)^{NM}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \underset{sq_j}{\overset{M}{}}}\left({\displaystyle \frac{1}{T_{q_j}}}{\displaystyle \frac{1}{T_s}}\right)^1det\left[K(z_i,z_j;\{\stackrel{~}{T_c}\})\right]_{(i,j)=1,\mathrm{},n},`$ (28)
where the kernel is given by
$`K(z_i,z_j;\{\stackrel{~}{T_c}\})`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{N}{}}}{\displaystyle \underset{c=1}{\overset{M}{}}}\left[(NM){\displaystyle \frac{1\stackrel{~}{T}_c}{\stackrel{~}{T}_c}}+k1\right](z_iz_j^{})^{k1}`$ (29)
$`=`$ $`{\displaystyle \underset{c=1}{\overset{M}{}}}\left[(NM){\displaystyle \frac{1\stackrel{~}{T}_c}{\stackrel{~}{T}_c}}+x{\displaystyle \frac{d}{dx}}\right]{\displaystyle \frac{1x^N}{1x}}|_{x=z_1z_2^{}}`$ (30)
where we used the notation: $`\stackrel{~}{T}_c=1\mathrm{exp}\left(\tau _c_{j=1}^n\delta _{q_j,c}\mathrm{ln}|z_j|^2\right)`$.
Further simplifications occur after taking into account that eigenvalues $`z_i`$ are, in fact, concentrated typically at distances of order of $`1/N`$ from the unit circle. Then it is natural to introduce new variables $`y_i,\varphi _i`$ according to $`z_i=(1y_i/N)e^{i\varphi _i}`$ and consider $`y_i`$ to be of the order of unity when $`N\mathrm{}`$. First of all, one immediately finds that:
$$\underset{N\mathrm{}}{lim}\underset{c=1}{\overset{M}{}}\left(\frac{\stackrel{~}{T}_c}{T_c}\right)^{NM}=\mathrm{exp}\left[2\underset{j=1}{\overset{n}{}}y_j\frac{1T_{q_j}}{T_{q_j}}\right]$$
(31)
As to the phases $`\varphi _i`$, we expect their typical separation scaling as: $`\varphi _i\varphi _j=O(1/N)`$. Now it is straightforward to perform explicitly the limit $`N\mathrm{}`$ in Eq.(29). Combining all factors together, one brings the correlation function Eq.(25) to the final form:
$`R_n(z_1,\mathrm{},z_n)`$ $``$ $`{\displaystyle \underset{k=1}{\overset{n}{}}}{\displaystyle \underset{q=1}{\overset{M}{}}}{\displaystyle \frac{e^{g_qy_k}}{_{sq}(g_qg_s)}}det\left[{\displaystyle _1^1}𝑑\lambda {\displaystyle \underset{c=1}{\overset{M}{}}}(\lambda +g_c)e^{\frac{i}{2}\lambda \delta _{ij}}\right]_{i,j=1,n}`$ (32)
with $`g_c=2/T_c1`$ and $`\delta _{ij}=N(\varphi _i\varphi _j)i(y_i+y_j)`$. The expression above coincides in every detail with that obtained in for random GUE matrices deformed by a finite rank anti-Hermitian perturbation.<sup>*</sup><sup>*</sup>*One should remember that the mean density of phases $`\varphi _i`$ along the unit circle is $`\nu =1/(2\pi )`$ and take into account that the constants $`g_c`$ defined in are, in fact, $`\pi \nu g_c=g_c/2`$ in the notations of the present paper.. This completes the proof of universality for finite-rank deviations.
YVF is obliged to B. Khoruzhenko for discussions and suggestions on the earlier stage of the research. The financial support by SFB 237 ”Unordnung und grosse Fluktuationen” as well as of the grant No. INTAS 97-1342 is acknowledged with thanks.
|
warning/0005/cond-mat0005387.html
|
ar5iv
|
text
|
# Circuit theory of multiple Andreev reflections in diffusive SNS junctions: the incoherent case
## I Introduction
The concept of multiple Andreev reflections (MAR) was first introduced by Klapwijk, Blonder, and Tinkham in order to explain the subharmonic gap structure (SGS) on current-voltage characteristics of superconducting junctions. The theory was originally formulated for perfect SNS junctions and then extended to include the effect of resistance of the SN interface (OTBK theory). Within this approach, the subgap current transport is described in terms of ballistic propagation of quasiclassical electrons through the normal metal region, accompanied by Andreev and normal reflections from specular NS boundaries. During every passage across the junction, the electrons and the retro-reflected holes gain energy equal to $`eV`$, which allows them eventually to escape from the SNS well. This energy gain results in strong quasiparticle nonequilibrium within the subgap energy region $`|E|<\mathrm{\Delta }`$.
OTBK theory gives a qualitatively adequate description of dc current transport in voltage biased SNS junctions; however, its quantitative results have a rather limited range of applicability. In short ballistic junctions with length $`d`$ comparable with or smaller than the coherence length (e.g., atomic-size junctions), the quantum coherence of subsequent Andreev reflections plays a crucial role leading to the ac Josephson effect. It has been shown that such a coherence also strongly modifies the dc current and SGS (coherent MAR regime). In fact, even in long ballistic SNS junctions (e.g. 2DEG-based devices), the coherence effects are important and give rise to resonant structures in the current due to Andreev quantization. In this respect, the quasiclassical OTBK theory, which does not include any coherence effects, may be qualified as a model for the incoherent MAR regime.
One might expect that impurities could provide the conditions for incoherent MAR by washing out the Andreev spectrum. However, this is not the case for a short diffusive junction, where appreciable Josephson coupling gives rise to coherent MAR. The electron-hole coherence in the normal metal holds over a distance of the coherence length $`\xi _E=\sqrt{\mathrm{}𝒟/2E}`$ from the superconductor ($`𝒟`$ is the diffusion constant). The overlap of coherent proximity regions induced by both SN interfaces creates an energy gap in the electron spectrum of the normal metal, which plays the role of the level spacing in the ballistic case. In short junctions with a wide proximity gap of the order of the energy gap $`\mathrm{\Delta }`$ in the superconducting electrodes, the phase coherence covers the entire normal region.
An incoherent MAR regime will occur in long diffusive SNS junctions with a small proximity gap of the order of Thouless energy $`E_{\text{Th}}=\mathrm{}𝒟/d^2\mathrm{\Delta }`$. If the applied voltage is large on a scale of the Thouless energy, $`eVE_{\text{Th}}`$, then the coherence length $`\xi _E`$ is much smaller than the junction length at all relevant energies $`E\mathrm{m}in(eV,\mathrm{\Delta })`$. In this case, the proximity regions near the SN interfaces become virtually decoupled and the Josephson oscillations are strongly suppressed. At the same time, as soon as the inelastic mean free path exceeds the junction length, the subgap electrons must undergo many incoherent Andreev reflections before they enter the reservoir. We emphasize that such incoherency is provided by the small coherence length at large enough voltages, while the intrinsic dephasing length can be arbitrarily large. In order to describe such an incoherent MAR regime, one has to operate with the electron and hole diffusion flows across the junction rather than with ballistic quasiparticle trajectories, and to consider the Andreev reflections as the relationships between these diffusive flows.
The first step in extending the OTBK approach to diffusive SNS structures was taken by Volkov and Klapwijk, who derived recurrence relations between the boundary values of the distribution functions. In that study, only a weak nonequilibrium was considered, which implies suppression of MAR by inelastic relaxation. In the present paper, we focus on the opposite case of strong nonequilibrium in the developed MAR regime, which results in the appearance of SGS on $`I`$-$`V`$ characteristics of the diffusive SNS junctions. Following the interpretation of MAR as a transport problem in energy space, we analyze it by formulating an equivalent network in the spirit of Nazarov’s circuit theory. Within this approach, the energy-dependent tunnel and Andreev resistances of an equivalent circuit play roles similar to the normal and Andreev reflection probabilities in OTBK theory, and the effective voltage source is represented by Fermi reservoirs.
The paper is organized as follows. In Section II, we derive the equations for incoherent MAR from the general Keldysh equations. In Sections III and IV, the circuit representation is formulated; some applications are considered in Section V. The SGS in junctions with resistive interfaces is calculated in Section VI. The complete solution of the problem suitable for numerical calculation of the $`I`$-$`V`$ characteristics is obtained in Section VII by using a chain-fraction technique. In Section VIII, we discuss limitations on the MAR regime imposed by inelastic processes.
## II Microscopic background
The system under consideration consists of a normal channel ($`0<x<d`$) confined between two voltage biased superconducting electrodes, with the elastic mean free path $`l`$ much shorter than any characteristic size of the problem. In this limit, the microscopic analysis of current transport can be performed within the framework of the diffusive equations of nonequilibrium superconductivity for the $`4\times 4`$ supermatrix Keldysh-Green function $`\stackrel{ˇ}{G}(t_1t_2,x)`$:
$$[\stackrel{ˇ}{H},\stackrel{ˇ}{G}]=i\mathrm{}𝒟_x\stackrel{ˇ}{J},\stackrel{ˇ}{J}=\stackrel{ˇ}{G}_x\stackrel{ˇ}{G},\stackrel{ˇ}{G}^2=\stackrel{ˇ}{1},$$
(1)
$$\stackrel{ˇ}{H}=\stackrel{ˇ}{1}[i\mathrm{}\sigma _z_te\varphi (t)+\widehat{\mathrm{\Delta }}(t)],\widehat{\mathrm{\Delta }}=\mathrm{\Delta }e^{i\sigma _z\chi }i\sigma _y,$$
(2)
where $`\mathrm{\Delta }`$ is the modulus and $`\chi `$ is the phase of the order parameter, and $`\varphi `$ is the electric potential. The Pauli matrices $`\sigma _i`$ operate in the Nambu space of $`2\times 2`$ matrices denoted by “hats”, and the products of two-time functions are interpreted as their time convolutions. The junction length $`d`$ is assumed to be smaller than the inelastic and phase-breaking lengths, which allows us to exclude the inelastic collisions from our consideration at this stage; their role will be discussed later. The electric current $`I`$ per unit area is expressed through the Keldysh component $`\widehat{J}^K`$ of the supermatrix current $`\stackrel{ˇ}{J}`$:
$$I(t)=(\pi \mathrm{}\sigma _N/4e)Tr\sigma _z\widehat{J}^K(tt,x),$$
(3)
where $`\sigma _N`$ is the conductivity of the normal metal.
At the SN interface, the supermatrix $`\stackrel{ˇ}{G}`$ satisfies the boundary condition
$$(\sigma _N\stackrel{ˇ}{J})_{\pm 0}=(2R_{SN})^1[\stackrel{ˇ}{G}_0,\stackrel{ˇ}{G}_{+0}],$$
(4)
where the indices $`\pm 0`$ denote the right and left sides of the interface and $`R_{SN}`$ is the interface resistance per unit area in the normal state, which relates to, e.g., a Schottky barrier or mismatch between the Fermi velocities. Within the model of infinitely narrow potential of the interface barrier, $`U(x)=H\delta (x)`$, the interface resistance is related to the barrier strength $`Z=H(\mathrm{}v_F)^1`$ as $`R_{SN}=2lZ^2/3\sigma _N`$. It has been shown in Ref. that Eq. (4) is valid either for a completely transparent interface ($`R_{SN}0`$, $`\stackrel{ˇ}{G}_{+0}=\stackrel{ˇ}{G}_0`$) or for an opaque barrier whose resistance is much greater than the resistance $`R(l)=l/\sigma _N`$ of a metal layer with the thickness formally equal to $`l`$.
According to the definition of the supermatrix $`\stackrel{ˇ}{G}`$,
$$\stackrel{ˇ}{G}=\left(\begin{array}{ccc}\widehat{g}^R& \widehat{G}^K& \\ 0& \widehat{g}^A& \end{array}\right),\widehat{G}^K=\widehat{g}^R\widehat{f}\widehat{f}\widehat{g}^A,$$
(5)
Eqs. (1) and (4) represent a compact form of separate equations for the retarded and advanced Green’s functions $`\widehat{g}^{R,A}`$ and the distribution function $`\widehat{f}=f_++\sigma _zf_{}`$. Their time evolution is imposed by the Josephson relation $`\chi (t)=2eVt`$ for the phase of the order parameter in the right electrode (we assume $`\chi =0`$ in the left terminal). This implies that the function $`\stackrel{ˇ}{G}(t_1t_2,x)`$ consists of a set of harmonics $`\stackrel{ˇ}{G}(E_n,E_m,x)`$, $`E_n=E+neV`$, which interfere in time and produce the ac Josephson current. However, when the junction length $`d`$ is much larger than the coherence length $`\xi _E`$ at all relevant energies $`EeV`$, we may consider coherent quasiparticle states separately at both sides of the junction, neglecting their mutual interference and the ac Josephson effect. Thus, the Green’s function in the vicinity of left SN interface can be approximated by the solution $`\widehat{g}=\sigma _z\mathrm{cosh}\theta +i\sigma _y\mathrm{sinh}\theta `$ of the static Usadel equations for a semi-infinite SN structure, with the spectral angle $`\theta (E,x)`$ satisfying the equation
$$\mathrm{tanh}[\theta (E,x)/4]=\mathrm{tanh}[\theta _N(E)/4]\mathrm{exp}(x/\xi _E\sqrt{i}),$$
(6)
with the boundary condition
$$W\sqrt{i\mathrm{\Delta }/E}\mathrm{sinh}(\theta _N\theta _S)+2\mathrm{sinh}(\theta _N/2)=0.$$
(7)
The indices $`S`$, $`N`$ in these equations refer to the superconducting and the normal side of the interface, respectively.
The dimensionless parameter $`W`$ in Eq. (7),
$$W=\frac{R(\xi _\mathrm{\Delta })}{R_{SN}}=\frac{\xi _\mathrm{\Delta }}{dr},r=\frac{R_{SN}}{R_N},$$
(8)
where $`R_N=R(d)=d/\sigma _N`$ is the resistance of the normal channel per unit area, has the meaning of an effective barrier transmissivity for the spectral functions. Note that even at large barrier strength $`Z1`$ ensuring the validity of the boundary conditions Eq. (4), the effective transmissivity $`W(\xi _\mathrm{\Delta }/l)Z^2`$ of the barrier in a “dirty” system, $`l\xi _\mathrm{\Delta }`$, could be large. In this case, the spectral functions are virtually insensitive to the presence of a barrier and, therefore, the boundary conditions Eqs. (4) can be applied to an arbitrary interface if we approximately consider high-transmissive interfaces with $`W\xi _\mathrm{\Delta }/l1`$ as completely transparent, $`W=\mathrm{}`$. For low transmissivity, $`W1`$, Eq. (7) can be analyzed within a perturbative approach (see the Appendix). At arbitrary $`W`$, Eq. (7) should be solved numerically.
The distribution functions $`f_\pm (E,x)`$ are to be considered as global quantities within the whole normal channel determined by the diffusive kinetic equations
$$_x[D_\pm (E,x)_xf_\pm (E,x)]=0,$$
(9)
with dimensionless diffusion coefficients
$$D_+=(1/4)Tr(1\widehat{g}^R\widehat{g}^A)=\mathrm{cos}^2Im\theta ,$$
(11)
$$D_{}=(1/4)Tr(1\sigma _z\widehat{g}^R\sigma _z\widehat{g}^A)=\mathrm{cosh}^2Re\theta .$$
(12)
Assuming the normal conductance of electrodes to be much greater than the junction conductance, we consider them as equilibrium reservoirs with unperturbed spectral characteristics, $`\theta _S=Arctanh(\mathrm{\Delta }/E)`$, and equilibrium quasiparticle distribution, $`\widehat{f}_S(E)=f_0(E)\mathrm{tanh}(E/2T)`$. Within this approximation, the boundary conditions for the distribution functions in Eq. (9) at $`x=0`$ read
$$\sigma _ND_+_xf_+(E,0)=G_+(E)[f_+(E,0)f_0(E)],$$
(13)
$$\sigma _ND_{}_xf_{}(E,0)=G_{}(E)f_{}(E,0),$$
(14)
where
$$G_\pm (E)=R_{SN}^1(N_SN_NM_S^\pm M_N^\pm ),$$
(15)
$$N(E)=Re(\mathrm{cosh}\theta ),M^+(E)+iM^{}(E)=\mathrm{sinh}\theta .$$
(16)
At large energies, $`|E|\mathrm{\Delta }`$, when the normalized density of states $`N(E)`$ approaches unity and the condensate spectral functions $`M^\pm (E)`$ turn to zero at both sides of the interface, the conductances $`G_\pm (E)`$ coincide with the normal barrier conductance; within the subgap region $`|E|<\mathrm{\Delta }`$, $`G_+(E)=0`$.
Similar considerations are valid for the right NS interface, if we eliminate the explicit time dependence of the order parameter in Eq. (1), along with the potential of right superconducting electrode, by means of a gauge transformation
$$\stackrel{ˇ}{G}(t_1t_2,x)=\mathrm{exp}(i\sigma _zeVt_1)\stackrel{~}{\stackrel{ˇ}{G}}(t_1t_2,x)\mathrm{exp}(i\sigma _zeVt_2).$$
(17)
As a result, we arrive at the same static equations and boundary conditions, Eqs. (6)-(16), with $`xdx`$, for the gauge-transformed functions $`\stackrel{~}{\widehat{g}}(E,x)`$ and $`\stackrel{~}{\widehat{f}}(E,x)`$. Thus, to obtain a complete solution for the distribution function $`f_{}`$, which determines the dissipative current
$$I=\frac{\sigma _N}{2e}_{\mathrm{}}^{\mathrm{}}𝑑ED_{}_xf_{},$$
(18)
we must solve the boundary problem for $`\widehat{f}(E,x)`$ at the left SN interface, and a similar boundary problem for $`\stackrel{~}{\widehat{f}}(E,x)`$ at the right interface, and then match the distribution function asymptotics deep inside the normal region by making use of the relationship following from Eqs. (5), (17):
$$\widehat{f}(E,x)=\stackrel{~}{\widehat{f}}(E+\sigma _zeV,x).$$
(19)
## III Circuit representation of boundary conditions
In order for this kinetic scheme to conform to the conventional physical interpretation of Andreev reflection in terms of electrons and holes, we introduce the following parameterization of the matrix distribution function,
$$\widehat{f}(E,x)=1\left(\begin{array}{ccc}n^e(E,x)& 0& \\ 0& n^h(E,x)& \end{array}\right),$$
(20)
where $`n^e`$ and $`n^h`$ will be considered as the electron and hole population numbers. Deep inside the normal metal region, they acquire rigorous meaning of distribution functions of electrons and holes. In equilibrium, the functions $`n^{e,h}`$ approach the Fermi distribution. In this representation, Eqs. (9) take the form
$$D_\pm (E,x)_xn_\pm (E,x)=\text{const}I_\pm (E)/\sigma _N,$$
(21)
where $`n_\pm =n^e\pm n^h`$, and they may be interpreted as conservation equations for the (specifically normalized) net probability current $`I_+`$ of electrons and holes, and for the electron-hole imbalance current $`I_{}`$. Furthermore, the probability currents of electrons and holes, defined as $`I^{e,h}=(1/2)(I_+\pm I_{})`$, separately obey the conservation equations. The probability currents $`I^{e,h}`$ are naturally related to the electron and hole diffusion flows, $`I^{e,h}=\sigma _N_xn^{e,h}`$, at large distances $`x\xi _E`$ from the SN boundary. Within the proximity region, $`x\xi _E`$, each current $`I^{e,h}`$ generally consists of a combination of both the electron and hole diffusion flows,
$$I^{e,h}=(\sigma _N/2)\left[(D_+\pm D_{})_xn^e+(D_+D_{})_xn^h\right],$$
(22)
which reflects coherent mixing of normal electron and hole states in this region.
In terms of electrons and holes, the boundary conditions in Eqs. (13), (14) read
$$I^{e,h}=G_T(n_Fn^{e,h})G_A(n^en^h),$$
(23)
where
$$G_T=G_+,G_A=(G_{}G_+)/2.$$
(24)
Each of the equations Eq. (23) may be clearly interpreted as a Kirchhoff rule for the electron or hole probability current flowing through the effective circuit (tripole) shown in Fig. 1(a). Within such an interpretation, the nonequilibrium populations of electrons and holes $`n^{e,h}`$ at the interface correspond to “potentials” of nodes attached to the “voltage source” – the Fermi distribution $`n_F(E)`$ in the superconducting reservoir – by “tunnel resistors” $`R_T(E)=G_T^1(E)`$. The “Andreev resistor” $`R_A(E)=G_A^1(E)`$ between the nodes provides electron-hole conversion (Andreev reflection) at the SN interface.
The circuit representation of the diffusive SN interface is analogous to the scattering description of ballistic SN interfaces: the tunnel and Andreev resistances in the diffusive case play the same role as the normal and Andreev reflection coefficients in the ballistic case. For instance, for $`|E|>\mathrm{\Delta }`$ \[Fig. 1(a)\], the probability current $`I^e`$ is contributed by equilibrium electrons incoming from the superconductor through the tunnel resistor $`R_T`$, and also by the current flowing through the Andreev resistor $`R_A`$ as the result of hole-electron conversion. Within the subgap region, $`|E|<\mathrm{\Delta }`$, \[Fig. 1(b)\], the quasiparticles cannot penetrate into the superconductor, $`R_T=\mathrm{}`$, and the voltage source is disconnected, which results in detailed balance between the electron and hole probability currents, $`I^e=I^h`$ (complete reflection). For the perfect interface, $`R_A`$ turns to zero, and the electron and hole population numbers become equal, $`n^e=n^h`$ (complete Andreev reflection). The nonzero value of the Andreev resistance for $`R_{SN}0`$ accounts for suppression of Andreev reflection due to the normal reflection by the interface.
Detailed information about the boundary resistances can be obtained from asymptotic expressions for the bare interface resistances $`R_\pm (E)G_\pm ^1(E)`$ (see the Appendix) and numerical plots of $`R_\pm (E)`$ in Fig. 2. In particular, $`R_\pm (E)`$ turns to zero at the gap edges due to the singularity in the density of states which enhances the tunneling probability. Furthermore, the imbalance resistance $`R_{}(E)`$ approaches the normal value $`R_{SN}`$ at $`E0`$ due to the enhancement of the Andreev reflection at small energies, which results from multiple coherent backscattering of quasiparticles by the impurities within the proximity region. This property is the reason for the re-entrant behavior of the conductance of high-resistive SIN systems at low voltages. In the MAR regime, one cannot expect any reentrance since quasiparticles at all subgap energies participate in the charge transport even at small applied voltage.
The diffusion coefficients $`D_\pm `$ in Eq. (9) turn to unity far from the SN boundary, and therefore the population numbers $`n^{e,h}`$ become linear functions of $`x`$,
$$n^{e,h}(E,x)\overline{n}^{e,h}(E,0)R_NI^{e,h}(E)x/d.$$
(25)
This equation defines the renormalized population numbers $`\overline{n}^{e,h}(E,0)`$ at the NS interface, which differ from $`n^{e,h}(E,0)`$ due to the proximity effect, as shown in Fig. 3. These quantities have the meaning of the true electron/hole populations which would appear at the NS interface if the proximity effect had been switched off. It is possible to formulate the boundary conditions in Eq. (23) in terms of these population numbers by including
the proximity effect into renormalization of the tunnel and Andreev resistances. To this end, we will associate the node potentials with renormalized boundary values $`\overline{n}^{e,h}(E,0)=(1/2)[\overline{n}_+(E,0)\pm \overline{n}_{}(E,0)]`$ of the population numbers, where $`\overline{n}_\pm (E,0)`$ are found from the exact solutions of Eqs. (21),
$$\overline{n}_\pm (E,0)=n_\pm (E,0)m_\pm (E)I_\pm (E).$$
(26)
Here $`m_\pm (E)`$ are the proximity corrections to the normal metal resistance at given energy for the probability and imbalance currents, respectively,
$$m_\pm (E)=\pm R_N(\xi _\mathrm{\Delta }/d)\mu _\pm (E),$$
(28)
$$\mu _\pm (E)=_0^{\mathrm{}}\frac{dx}{\xi _\mathrm{\Delta }}\left|D_\pm ^1(E,x)1\right|>0,$$
(29)
see the insets in Fig. 2(a,b). It follows from Eq. (26) that the same Kirchhoff rules as in Eqs. (23), (24) hold for $`\overline{n}^{e,h}(E,0)`$ and $`I^{e,h}(E)`$, if the bare resistances $`R_\pm `$ are substituted by the renormalized ones,
$$R_\pm (E)\overline{R}_\pm (E)=R_\pm (E)+m_\pm (E).$$
(30)
The energy dependence of the renormalized boundary resistances $`\overline{R}_T(E)`$ and $`\overline{R}_A(E)`$ is illustrated in Fig. 4. In some cases, there is an essential difference between the bare and renormalized resistances, which leads to qualitatively different properties of the SN interface for normal electrons and holes compared to the properties of the bare boundary. Let us first discuss a perfect SN interface with $`R_{SN}0`$. Within the subgap region $`|E|<\mathrm{\Delta }`$, the bare tunnel resistance $`R_T`$ is infinite whereas the bare Andreev resistance $`R_A`$ turns to zero; this corresponds to complete Andreev reflection, as already explained. However, the Andreev resistance for normal electrons and holes, $`\overline{R}_A(E)=2m_{}(E)`$, is finite and negative, which leads to enhancement of the normal metal conductivity within the proximity region. At $`|E|>\mathrm{\Delta }`$, the bare tunnel resistance $`R_T`$ is zero, while the renormalized tunnel resistance $`\overline{R}_T(E)=m_+(E)`$ is finite (though rapidly decreasing at large energies). This leads to suppression of the probability currents of normal electrons and holes within the proximity region, which is to be attributed to the appearance of Andreev reflection. Such a suppression is a global property of the proximity region in the presence of sharp spatial variation of the order parameter, and it is similar to the over-the-barrier Andreev reflection in the ballistic systems. In the presence of normal scattering at the SN interface, the overall picture depends on the interplay between the bare interface resistances $`R_\pm `$ and the proximity corrections $`m_\pm `$; for example, the renormalized tunnel resistance $`\overline{R}_T(E)`$ diverges at $`|E|\mathrm{\Delta }`$, along with the proximity correction $`m_+(E)`$, in contrast to the bare tunnel resistance $`R_T(E)`$. This indicates complete Andreev reflection at the gap edge independently of the transparency of the barrier, which is similar to the situation in the ballistic systems where the probability of Andreev reflection at $`|E|=\mathrm{\Delta }`$ is always equal to unity.
## IV Assembling MAR networks
To complete the definition of an equivalent MAR network, we have to construct a similar tripole for the right NS interface and to connect boundary values of population numbers (node potentials) using the matching condition Eq. (19) expressed in terms of electrons and holes:
$$n^{e,h}(E,x)=\stackrel{~}{n}^{e,h}(E\pm eV,x).$$
(31)
Since the gauge-transformed distribution functions $`\stackrel{~}{f}_\pm `$ obey the same equations Eq. (9)-(16), the results of the previous Section can be applied to the functions $`\stackrel{~}{n}^{e,h}(E)`$ and $`\stackrel{~}{I}^{e,h}(E)`$ (the minus sign implies that $`\stackrel{~}{I}`$ is associated with the current incoming to the right-boundary tripole). In particular, the asymptotics of the gauge-transformed population numbers far from the right interface are given by the equation
$$\stackrel{~}{n}^{e,h}(E,x)\stackrel{~}{\overline{n}}^{e,h}(E,d)+R_N\stackrel{~}{I}^{e,h}(E)\left(1x/d\right).$$
(32)
After matching the asymptotics in Eqs. (25) and (32) by means of Eq. (31), we find the following relations:
$$I^{e,h}(E)=\stackrel{~}{I}^{e,h}(E\pm eV),$$
(33)
$$\overline{n}^{e,h}(E,0)\stackrel{~}{\overline{n}}^{e,h}(E\pm eV,d)=R_NI^{e,h}(E).$$
(34)
From the viewpoint of the circuit theory, Eq. (34) may be interpreted as Ohm’s law for the resistors $`R_N`$ which connect energy-shifted boundary tripoles, separately for the electrons and holes, as shown in Fig. 1(c).
The final step which essentially simplifies the analysis of the MAR network, is based on the following observation. The spectral probability currents $`I^{e,h}`$ yield opposite contributions to the electric current in Eq. (18),
$$I=\frac{1}{2e}_{\mathrm{}}^{\mathrm{}}𝑑E[I^e(E)I^h(E)],$$
(35)
due to the opposite charge of electrons and holes. At the same time, these currents, referred to the energy axis, transfer the charge in the same direction, viz., from bottom to top of Fig. 1(c), according to our choice of positive $`eV`$. Thus, by introducing the notation $`I_n(E)`$ for an electric current entering the node $`n`$ with the energy $`E_n=E+neV`$, as shown by arrows in Fig. 1(c),
$$I_n(E)=\{\begin{array}{ccc}I^e(E_{n1}),& n=2k+1,& \\ I^h(E_n),& n=2k,& \end{array}$$
(36)
we arrive at an “electrical engineering” problem of current distribution in an equivalent network in energy space plotted in Fig. 5, where the difference between electrons and holes becomes unessential. The bold curve in Fig. 5 represents a distributed voltage source – the Fermi distribution $`n_F(E)`$ connected periodically with the network nodes. Within the gap, $`|E_n|<\mathrm{\Delta }`$, the nodes are disconnected from the Fermi reservoir and therefore all partial currents associated with the subgap nodes are equal.
Since all resistances and potentials of this network depend on $`E_n=E+neV`$, the partial currents obey the relationship $`I_n(E)=I_k[E+(nk)eV]`$ which allows us to express the physical electric current, Eq. (35), through the sum of all partial currents $`I_n`$ flowing through the normal resistors $`R_N`$, integrated over an elementary energy interval $`0<E<eV`$:
$$I=\frac{1}{2e}_{\mathrm{}}^{\mathrm{}}𝑑E[I_1(E)+I_0(E)]=\frac{1}{e}_0^{eV}𝑑EJ(E),$$
(37)
$$J(E)=\underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}I_n(E).$$
(38)
The spectral density $`J(E)`$ is periodic in $`E`$ with the period $`eV`$ and symmetric, $`J(E)=J(E)`$, which follows from the symmetry of all resistances with respect to $`E`$.
As soon as the partial currents are found, the population numbers can be recovered by virtue of Eqs. (21), (23), (25), and (36):
$$n^{e,h}(E,x)=\overline{n}^{e,h}(E,0)R_NI_{1,0}(E)x/d,$$
(39)
$$\overline{n}^{e,h}(E,0)=n_F\frac{1}{2}\left[\overline{R}_+(I_1I_0)\pm \overline{R}_{}(I_1+I_0)\right]$$
(40)
at $`|E|>\mathrm{\Delta }`$. Within the subgap region, Eq. (40) is inapplicable due to the undeterminacy of product $`\overline{R}_+(I_1I_0)`$. In this case, one may consider the subgap part of the network as a voltage divider between the nodes nearest to the gap edges, having the numbers $`N_{}`$, $`N_+`$, respectively, where $`N_\pm =\text{Int}[(\mathrm{\Delta }E)/eV]+1`$, Int($`x`$) denoting the integer part of $`x`$. Then the boundary populations at $`|E|<\mathrm{\Delta }`$ become
$`\overline{n}^{e,h}(E,0)`$ $`=`$ $`n^{L,R}(E_{\pm N_\pm })`$ (41)
$`\pm `$ $`I_0\left[N_\pm R_N+{\displaystyle \underset{k=1}{\overset{N_\pm 1}{}}}R_A(E_{\pm k})\right],`$ (42)
where $`R,L`$ indicate the right (left) node of the tripole, irrespectively of whether it relates to the left (even $`n`$) or right (odd $`n`$) interface. The physical meaning of $`n^{R,L}(E_n)`$, however, depends on the parity of $`n`$:
$$n^{R,L}(E_n)=\{\begin{array}{ccc}\overline{n}^{e,h}(E_n,0),& n=2k,& \\ \stackrel{~}{\overline{n}}^{h,e}(E_n,d),& n=2k+1.& \end{array}$$
(43)
The values $`n^{R,L}`$ in Eq. (41) can be found from Eq. (40) which is generalized for any tripole of the network in Fig. 5 outside the gap as
$`n^{R,L}(E_n)=n_F(E_n)(1/2)[\overline{R}_+(E_n)(I_{n+1}I_n)`$ (44)
$`\pm \overline{R}_{}(E_n)(I_{n+1}+I_n)].`$ (45)
The circuit formalism can be simply generalized to the case of different transparencies of NS interfaces, as well as to different values of $`\mathrm{\Delta }`$ in the electrodes. In this case, the network resistances become dependent not only on $`E_n`$ but also on the parity of $`n`$. As a result, the periodicity of the current spectral density doubles: $`J(E)=J(E+2eV)`$, and, therefore, $`J(E)`$ is to be integrated in Eq. (37) over the period $`2eV`$, with an additional factor $`1/2`$.
## V Simple applications
A helpful example of an asymmetric junction which allows us immediately to obtain an analytical solution is given by the SNN structure. The problem of calculation of its conductivity is inherently static and therefore may be completely solved for any junction length. If the latter is much larger than the coherence length, the circuit approach of the previous Section is applicable without restrictions. Due to the absence of superconducting correlations at the right NN interface, odd Andreev resistors are eliminated and, therefore, the whole network may be split into separate finite circuits located around even (superconducting) nodes, as shown in Fig. 6(a); moreover, odd tunnel resistances are to be considered as normal ones. After some simple algebra, we obtain the sum of partial currents in each subcircuit,
$$I_{2k}+I_{2k+1}=\frac{n_F(E_{2k1})n_F(E_{2k+1})}{R_N+R_{SN}+\overline{R}_{}(E_{2k})},$$
(46)
which leads to a well known formula for the $`I`$-$`V`$ characteristics of a long SNN junction:
$`I={\displaystyle \frac{1}{2e}}{\displaystyle _0^{2eV}}𝑑E{\displaystyle \underset{k=\mathrm{}}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{n_F(E_{2k1})n_F(E_{2k+1})}{R_N+R_{SN}+\overline{R}_{}(E_{2k})}}`$ (47)
$`={\displaystyle \frac{1}{e}}{\displaystyle _0^{\mathrm{}}}𝑑E{\displaystyle \frac{n_F(EeV)n_F(E+eV)}{R_N+R_{SN}+\overline{R}_{}(E)}}.`$ (48)
If the junction is short enough ($`d`$ and $`\xi _E`$ are comparable), one might naively expect some kind of proximity-induced Andreev scattering at the right NN interface, followed by MAR and SGS anomalies in the $`I`$-$`V`$ characteristic. However, the circuit theory rejects this assumption at once: since the condensate spectral functions $`M^\pm `$, Eq. (16), disappear in the normal terminal, the conductivities $`G_\pm `$ become equal, and the Andreev channel becomes closed ($`G_A=0`$) at the NN interface at any length of the junction. Thus, the circuit model of charge transport in Fig. 6(a) remains valid, with a few modifications: (i) the spectral angle $`\theta (E,x)`$ has to be found from the Usadel equation for the finite interval $`0<x<d`$, (ii) the proximity corrections $`m_\pm `$ at the left SN interface are to be expressed through the integrals
$$m_\pm (E)=\pm R_N_0^d\frac{dx}{d}\left|D_\pm ^1(E,x)1\right|,$$
(49)
instead of Eq. (26), and (iii) odd tunnel resistances $`R_{SN}`$ should be replaced by the energy-dependent bare resistances $`R_+(E)`$ \[or, equivalently, $`R_{}(E)`$\]. From this point of view, the entire channel represents a global “Andreev reflector” for normal electrons and holes incoming from the right reservoir.
The next simple application of this circuitry is given by calculation of the excess (deficit) current in SNS junctions, i.e., the difference $`I_{ex}=I(V)V/R`$ between the currents in superconducting and normal state at large voltages $`V\mathrm{\Delta }/e`$. Assuming the integration in Eq. (37) to be reduced to the interval $`0<E<eV/2`$ by making use of the symmetry $`J(E)=J(E)=J(eVE)`$, we note that the Andreev conductances are negligibly small for all nodes with $`n0`$ ($`E_n\mathrm{\Delta }`$). Thus, the circuit in Fig. 5 may be split, as in the case of the SNN junction discussed above, into the chain of separate circuits shown in Fig. 6(b). The contribution of the central circuit is described by Eq. (46) with $`k=0`$, whereas the other parts are to be considered as normal circuits and represent contribution of thermally excited quasiparticles:
$$\underset{n0,1}{}I_n=[1+n_F(E_1)n_F(E_1)]R^1,$$
(50)
where $`R=R_N+2R_{SN}`$ is the net normal resistance of the junction. In summary, we obtain another well-known result,
$`I_{ex}={\displaystyle \frac{2}{eR}}{\displaystyle _0^{eV/2}}dE{\displaystyle \frac{n_F(EeV)n_F(E+eV)}{R_N+R_{SN}+\overline{R}_{}(E)}}[R_{SN}`$ (51)
$`\overline{R}_{}(E)]{\displaystyle \frac{2}{eR}}{\displaystyle }_0^{\mathrm{}}dE{\displaystyle \frac{R_{SN}\overline{R}_{}(E)}{R_N+R_{SN}+\overline{R}_{}(E)}}.`$ (52)
It is of interest to note that the net resistance $`R_T=\overline{R}_+`$ for the probability current never enters final results in these examples and, therefore, the superconducting modification involves only the imbalance resistance $`\overline{R}_{}`$. In other words, only the evolution of the imbalance $`n_{}`$ between the electron and hole populations is relevant for the charge transport in such cases.
## VI MAR transport
Proceeding with the analysis of current transport through the SNS junction at arbitrary voltages, we first discuss the case of low-transparent barriers, $`W1`$. We note that in practice this case is relevant for a wide range of junctions both with high interface resistance, $`R_{SN}R_N`$, and comparatively low interface resistance $`R_{SN}R_N`$. Indeed, according to Eq. (8), the ratio $`R_N/R_{SN}=Wd/\xi _\mathrm{\Delta }`$, being proportional to $`W`$, contains also the large parameter $`d/\xi _\mathrm{\Delta }`$. Therefore, the limit $`W1`$ covers most of the practically interesting situations, $`0<R_N/R_{SN}d/\xi _\mathrm{\Delta }`$, and only the case of very small interface resistances, $`R_N/R_{SN}>d/\xi _\mathrm{\Delta }1`$, requires special consideration.
At $`W1`$, the proximity effect is essentially suppressed and the calculations can be performed on the basis of a simplified model of the equivalent network, which nevertheless provides a quantitative description of $`I(V)`$. Due to the sharp increase in the Andreev resistance at $`|E|>\mathrm{\Delta }`$ \[see Fig. 4(b)\], all Andreev resistors outside the gap can be excluded, and we arrive at the sequence of subcircuits shown in Fig. 6(c). The central circuit between the nodes $`N_{}`$ and $`N_+`$ includes $`N_++N_{}`$ normal and $`N_++N_{}1`$ Andreev resistors within the gap, as well as two tunnel resistors at the circuit edges. The total resistance $`R_\mathrm{\Delta }`$ of this circuit is
$`R_\mathrm{\Delta }(E)=(N_++N_{})R_N+{\displaystyle \underset{N_{}<k<N_+}{}}\overline{R}_A(E_k)`$ (53)
$`+\overline{R}_T(E_{N_+})+\overline{R}_T(E_N_{}),`$ (54)
and the current $`I_0`$ through this circuit is given by Ohm’s law:
$$I_0(E)=[n_F(E_N_{})n_F(E_{N_+})]/R_\mathrm{\Delta }(E).$$
(55)
Thus, the contribution of this circuit to the current spectral density, Eq. (38), is $`(N_++N_{})I_0`$.
The current of thermal excitations is carried by the side circuits ($`n>N_+`$, $`nN_{}`$):
$$I_n=\frac{n_F(E_{n1})n_F(E_n)}{R_N+\overline{R}_T(E_n)+\overline{R}_T(E_{n1})}.$$
(56)
From Eq. (55) we obtain a simple formula at $`T\mathrm{\Delta }`$:
$$I=_0^{eV}\frac{dE}{eR_{\text{MAR}}(E)},R_{\text{MAR}}(E)=\frac{R_\mathrm{\Delta }(E)}{N_++N_{}},$$
(57)
where $`R_{\text{MAR}}(E)`$ has the meaning of the effective resistance of the subgap region for the physical electric current.
In Fig. 7 we present the $`I`$-$`V`$ characteristics and the differential resistance vs inverse voltage at zero temperature, calculated numerically by means of Eq. (57). The parameter $`W`$ was chosen to be equal to $`0.1`$ and $`0.2`$ at $`d/\xi _\mathrm{\Delta }=5`$, which corresponds to the resistance ratio $`r=R_{SN}/R_N`$ equal to $`2`$ and $`1`$, respectively. In our calculation of $`\overline{R}_{T,A}(E)`$ in Eq. (53), we used the asymptotic Eq. (Circuit theory of multiple Andreev reflections in diffusive SNS junctions: the incoherent case) for the bare resistances $`R_{T,A}(E)`$ at $`W1`$, neglecting small proximity corrections $`m_\pm (E)R_N(\xi _\mathrm{\Delta }/d)W^2`$, Eq. (93). The results are in good agreement with those obtained on the basis of exact calculations \[see further Eq. (83)\]. The smeared steps in the $`I`$-$`V`$ characteristic indicate steplike increase in the number of subgap Andreev resistors in Eq. (53). The sharp peaks (dips) in $`dV/dI`$ arise from the rapidly varying contribution of the nodes which simultaneously cross the gap edges, and therefore both the edge Andreev resistances undergo strong suppression. A certain contribution also comes from the edge tunnel resistances which also show singular behavior at $`|E|\mathrm{\Delta }`$. The peaks are more pronounced for even subharmonics, when the middle Andreev resistor crosses the Fermi level, $`E=0`$, and its value is suppressed simultaneously with the edge Andreev resistors. Careful analysis shows that at the gap subharmonics, $`eV_n=2\mathrm{\Delta }/n`$, the second derivative $`d^2V/dI^2`$ has sharp maxima.
The magnitude of $`dV/dI`$ strongly depends on the number of large Andreev resistors which contribute to $`R_{\text{MAR}}`$. At $`eV<\mathrm{\Delta }`$, at least one Andreev resistor appears far from the “resonant” points $`E=0,\pm \mathrm{\Delta }`$ where $`R_A`$ sharply decreases. Thus, the net subgap resistance $`R_{\text{MAR}}(E)`$ remains large ($`R_{SN}/W`$) at any energy, which results in large differential resistance $`dV/dIR_{SN}/W`$ at these voltages. In the vicinity of the second subharmonic, $`eV\mathrm{\Delta }`$, the current transport involves a high-transmissive circuit with three Andreev resistors located near the resonant points, which yields a much smaller value of $`dV/dIR_{SN}`$. The same effect occurs at $`eV2\mathrm{\Delta }`$ when the resonant circuit contains two suppressed Andreev resistances at the gap edges. At $`eV>2\mathrm{\Delta }`$ the differential resistance is basically determined by quasiparticles which overcome the energy gap without Andreev reflections, and it turns to the normal value $`R`$ at large voltages.
At low voltage, the amplitude of the oscillations of the differential resistance decreases and the asymptotic value of $`dV/dI`$ at $`V\mathrm{\Delta }/e`$ can be found analytically from Eqs. (53), (57), by replacing the sum in Eq. (53) with an energy-independent integral. As a result we get
$$dV/dI(0)R_N+\frac{1}{2\mathrm{\Delta }}_\mathrm{\Delta }^\mathrm{\Delta }𝑑ER_A(E)=R_N+\frac{16\sqrt{2}}{21W}R_{SN}.$$
(58)
Since each circuit in Fig. 6(c) represents a separate voltage divider, we easily obtain the boundary values of the population numbers. If the node $`n=0`$ belongs to the central circuit ($`\mathrm{\Delta }eV<E<\mathrm{\Delta }`$ for electrons and $`\mathrm{\Delta }<E<\mathrm{\Delta }+eV`$ for holes), we have
$`\overline{n}^{e,h}(E,0)=n_F(E_{\pm N_\pm })\pm {\displaystyle \frac{n_F(E_N_{})n_F(E_{N_+})}{R_\mathrm{\Delta }(E)}}`$ (59)
$`\times \left[R_T(E_{\pm N_\pm })+N_\pm R_N+{\displaystyle \underset{k=1}{\overset{N_\pm 1}{}}}R_A(E_{\pm k})\right],`$ (60)
and, in the opposite case,
$$\overline{n}^{e,h}(E,0)=n_F(E)\frac{R_T(E)[n_F(E)n_F(E\pm eV)]}{R_N+R_T(E)+R_T(E\pm eV)}.$$
(61)
As follows from Eqs. (59), (61), at low temperatures, the energy distribution of quasiparticles within the region $`\mathrm{\Delta }eV<E<\mathrm{\Delta }+eV`$ has a steplike form (Fig. 8), which is qualitatively similar to, but quantatively different from, that found in OTBK theory. The number of steps increases at low voltage, and the shape of the distribution function becomes resemblant to a “hot electron” distribution with the effective temperature of the order of $`\mathrm{\Delta }`$. This distribution is modulated due to the discrete nature of the heating mechanism of MAR, which transfers the energy from an external voltage source to the quasiparticles by energy quanta $`eV`$. Since the subgap probability current $`I_0`$, Eq. (55), is strongly suppressed by large subgap resistance $`R_\mathrm{\Delta }R_{SN}(N_++N_{})/WR_N`$, the spatial variations of the population numbers, Eqs. (25), (39), are negligibly small: $`n^{e,h}(E,x)n^{e,h}(E,0)R_N/R_\mathrm{\Delta }1`$.
## VII Exact solution
In the case of high transmittance of the NS interface, the partial currents outside the gap noticeably contribute to the net electric current even at $`T=0`$. In this case, the $`I`$-$`V`$ curves should be calculated on the basis of the exact solution of the recurrence relations between partial currents,
$$I_{n+1}r_n+I_{n1}r_{n1}I_n(\rho _n+r_{n1}+r_n)=U_nU_{n1},$$
(62)
following from the Kirchhoff rules for an infinite network in Fig. 5. Here $`U_n(E)=n_F(E_n)`$ and
$$\rho _n(E)=R_N+\overline{R}_{}(E_{n1})+\overline{R}_{}(E_n),$$
(64)
$$r_n(E)=[\overline{R}_+(E_n)\overline{R}_{}(E_n)]/2.$$
(65)
By analogy with differential equations, we introduce the following ansatz:
$$I_n(E)=A_n^+(E)I_n^+(E)+A_n^{}(E)I_n^{}(E),$$
(66)
where $`I_n^\pm `$ are the fundamental solutions of the corresponding uniform equation, decreasing at $`n\pm \mathrm{}`$, respectively,
$$I_n^+(E)=\{\begin{array}{ccc}S_{n1,0},& n>0,& \\ S_{1,n}^1,& n<0,& \end{array}$$
(68)
$$I_n^{}(E)=\{\begin{array}{ccc}P_{n1,0}^1,& n>0,& \\ P_{1,n},& n<0,& \end{array}$$
(69)
$`I_0^+=I_0^{}=1`$. The quantities
$$S_{mn}(E)=\underset{j=n}{\overset{m}{}}s_j(E),P_{mn}(E)=\underset{j=n}{\overset{m}{}}p_j(E)$$
(70)
are expressed through the products of chain fractions, $`s_n=I_{n+1}^+/I_n^+`$, $`p_n=I_n^{}/I_{n+1}^{}`$, defined by the recurrences
$$s_n=\frac{r_n}{\rho _{n+1}+r_n+a_{n+1}},p_n=\frac{r_n}{\rho _n+r_n+b_{n1}},$$
(72)
$$a_n=r_n(1s_n),b_n=r_n(1p_n),$$
(73)
with the boundary conditions $`s_{n+\mathrm{}}0`$, $`p_n\mathrm{}0`$. Within the gap, $`|E_n|<\mathrm{\Delta }`$, where $`r_n\mathrm{}`$, the values $`s_n`$, $`p_n`$ are equal to $`1`$, in accordance with the conservation of the subgap currents mentioned above.
The coefficients $`A_n^\pm `$ in Eq. (66) satisfy an equation following from Eqs. (62), (66), (66),
$`[(\rho _{n+1}r_ns_n^1)\delta A_{n+1}^+r_ns_n^1\delta A_n^+]S_{n,0}+[(\rho _{n+1}`$ (74)
$`r_np_n)\delta A_{n+1}^{}r_np_n\delta A_n^{}]P_{n,0}^1=U_{n+1}U_n,`$ (75)
for $`n>0`$ (and similar for $`n<0`$), where $`\delta A_n=A_{n+1}A_n`$. The requirement of cancellation of terms with $`\delta A_{n+1}^\pm `$ in Eq. (74) allows us to completely determine $`A_n^\pm `$. This yields first-order recurrences for $`A_n^\pm `$ which lead to the formula for partial currents at $`n>0`$,
$$I_n=A_0^+I_n^++A_0^{}I_n^{}+\underset{k=1}{\overset{n}{}}j_k\left(S_{n1,k}P_{n1,k}^1\right),$$
(76)
$$j_n(E)=(U_{n1}U_n)/(\rho _n+a_n+b_{n1}).$$
(77)
The undeterminacy of $`a_n`$, $`b_{n1}`$ in Eq. (77) within the subgap region \[where $`s_n=p_n=1`$ and $`r_n\mathrm{}`$, see Eq. (73)\], is resolved by the recurrences $`a_n=\rho _{n+1}+a_{n+1}`$, $`b_n=\rho _n+b_{n1}`$ following from Eqs. (70) for $`|E_n|<\mathrm{\Delta }`$. These recurrences are to be continued until the nodes $`N_+`$ and $`N_{}`$, respectively, are reached. As a result, we obtain a convenient representation for $`j_n`$ at $`(N_{}+1)nN_+`$:
$$j_n(E)=\frac{U_{n1}U_n}{R_\mathrm{\Delta }r_{N_+}(1+s_{N_+})r_N_{}(1+p_N_{})}.$$
(78)
The effective subgap resistance determined by the denominator in Eq. (78) differs from $`R_\mathrm{\Delta }`$ in Eq. (53) by extra terms describing leakage of the subgap current through the Andreev resistors outside the gap.
The coefficients $`A_0^\pm `$ have to provide finite values of $`I_n`$ in Eq. (76) at all $`n`$; for instance, for $`n+\mathrm{}`$, the divergent products of $`p_j^1`$ in $`I_n^{}`$ and $`P_{n1,k}^1`$ should compensate each other: $`A_0^{}=_{k=0}^+\mathrm{}j_kP_{k1,0}`$. A similar procedure for negative $`n`$ determines the value of $`A_0^+`$ and results in the final formula for the partial current with arbitrary index $`n`$,
$`I_n(E)=j_n(E)+{\displaystyle \underset{k=n+1}{\overset{+\mathrm{}}{}}}j_k(E)P_{k1,n}(E)`$ (79)
$`+{\displaystyle \underset{k=\mathrm{}}{\overset{n1}{}}}j_k(E)S_{n1,k}(E).`$ (80)
By making use of the relation $`p_n(E)=s_n(E)`$ following from Eq. (70) and taking into account the symmetry of all resistances with respect to $`E`$, we obtain the net current spectral density,
$`J(E)={\displaystyle \underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}}\{j_n(E)+{\displaystyle \underset{k=n}{\overset{+\mathrm{}}{}}}[j_n(E)S_{k,n}(E)`$ (81)
$`+j_n(E)S_{k+1,n+1}(E)]\}.`$ (82)
At low temperatures $`T\mathrm{\Delta }`$, only the term with $`n=0`$ contributes to the sum in Eq. (81),
$$J(E)=j_0(E)\left\{1+\underset{k=1}{\overset{+\mathrm{}}{}}\left[S_{k1,0}(E)+S_{k,1}(E)\right]\right\}.$$
(83)
Figure 9 shows the results of a numerical calculation of $`I(V)`$ and $`dV/dI`$ for an SNS junction with high-transmissive interfaces, $`W=1`$ and $`W\mathrm{}`$, at zero temperature. In practice, due to rapid decrease of the coefficients $`r_n(E)`$ in Eq. (62) at large energies, it is enough to calculate recurrences in Eq. (70) starting from the maximum number $`n_{\mathrm{max}}=N_{}+2`$ and assuming the chain fractions to be truncated, $`s_n=0`$, at $`n>n_{\mathrm{max}}`$. To avoid formal singularity in $`m_{}(E)`$ at $`E0`$, we introduce a small dephasing factor $`\mathrm{\Gamma }`$ which provides a cutoff of the coherence length $`\xi _E\sqrt{\mathrm{}𝒟/2(E+i\mathrm{\Gamma })}`$. The corresponding dephasing length $`l_\varphi =\sqrt{\mathrm{}𝒟/2\mathrm{\Gamma }}`$ was chosen equal to the junction length; the variation of $`l_\varphi `$ is not critical for the fine structure of $`dV/dI`$.
Similar to the case of low barrier transmittance, the current transport through an SNS junction with nearly perfect interfaces can be qualitatively explained within a simplified model of MAR, where the over-the-barrier ($`|E|>\mathrm{\Delta }`$) Andreev reflection is ignored. Indeed, as follows from Fig. 4, at $`R_{SN}R_N`$, the tunnel resistances $`\overline{R}_T(E)`$ outside the gap are much smaller than the Andreev and/or normal resistances, except for a narrow region around the gap edges, where $`\overline{R}_T(E)`$ diverges due to complete Andreev reflection. This allows us to assume all the normal resistors at $`|E|>\mathrm{\Delta }`$ to be connected directly to the “voltage source” $`n_F(E)`$ and therefore to exclude all Andreev resistors in Fig. 5 outside the gap. As a result, we arrive at the sequence of subcircuits shown in Fig. 6(c), with $`\overline{R}_T=0`$ for side circuits. Thus, at $`T\mathrm{\Delta }`$, the subgap current may be approximately described by Eqs. (53), (57), with the tunnel and Andreev resistances renormalized by the proximity effect.
Within this model, the SGS oscillations in the differential resistance in Fig. 9 can be explained in the following way. As the voltage decreases, the subgap current, which approximately follows Ohm’s law, undergoes an additional suppression in the vicinity of the gap subharmonics, due to the presence of a high-resistive circuit with two large tunnel resistors located just at the gap edges. These current steps, being almost invisible in the $`I`$-$`V`$ characteristic, manifest themselves as sharp dips in $`dV/dI`$. At even subharmonics, this effect is partially compensated by the middle negative Andreev resistor, which rapidly reduces the effective normal metal resistance due to the increase in the size of the proximity region at small energies. As a result, the even dips become less pronounced and, as long as the interface transparency increases, turn into small peaks. At low voltages, the differential resistance approaches a constant value, which can be estimated for perfect NS interface by the following expression similar to Eq. (58):
$`dV/dI(0)R_N\left[12{\displaystyle \frac{\xi _\mathrm{\Delta }}{d}}{\displaystyle _\mathrm{\Delta }^\mathrm{\Delta }}{\displaystyle \frac{dE}{2\mathrm{\Delta }}}\mu _{}(E)\right]`$ (84)
$`=R_N\left(12.64\xi _\mathrm{\Delta }/d\right),R_{SN}=0.`$ (85)
Unlike the ballistic SNS junction, but similar to short diffusive constriction, the SGS survives at zero temperature even for perfect NS interfaces. In this case, the SGS occurs due to coherent impurity scattering of quasiparticles inside the proximity region, with the amplitude approximately proportional to the characteristic length $`\xi _\mathrm{\Delta }`$ of this region. If we neglect the proximity corrections ($`m_\pm 0`$), the SGS disappears, along with the excess current, and the $`I`$-$`V`$ characteristic shows perfect Ohmic behavior. Thus, we conclude that in the both cases of resistive and transparent NS interfaces, SGS appears due to the normal electron backscattering in the proximity region. This formally corresponds to the finite value of the renormalized Andreev resistance of the interface.
## VIII Role of inelastic scattering
Since the relative SGS amplitude increases along with the NS barrier strength (though the current itself decreases), one might expect systems with high-resistive interfaces to be more favorable for the experimental observation of SGS. However, there exists an intrinsic restriction for the effect: to provide strong nonequilibrium of the subgap quasiparticles inherent to MAR, the time $`\tau _d`$ of their diffusion through the whole MAR staircase, from $`\mathrm{\Delta }`$ to $`\mathrm{\Delta }`$, must be smaller than the inelastic relaxation time $`\tau _\epsilon `$. The value of $`\tau _d`$ can be estimated as the time for diffusion over an effective length $`(2\mathrm{\Delta }/eV)d`$. At low interface resistance, $`R_{SN}R_N`$, the diffusion rate is basically determined by the impurity scattering: $`\tau _d(eV)(2\mathrm{\Delta }/eV)^2d^2/𝒟`$. For high interface resistance, $`R_{SN}R_N`$, the large Andreev resistance $`R_AR_{SN}/W`$ present bottleneck which renormalizes the diffusion coefficient in $`\tau _d`$ by a small factor $`WR_N/R_{SN}`$. In this way, the level of nonequilibrium of the subgap quasiparticles is controlled by the parameter
$$W_\epsilon =\frac{\tau _\epsilon }{\tau _d(2\mathrm{\Delta })}=\frac{R_N}{R_N+R_{SN}/W}E_{\mathrm{Th}}\tau _\epsilon ,$$
(86)
which must be large enough to allow observation of at least a few subharmonics in $`I(V)`$, i.e., the condition $`W_\epsilon >1`$ determines the lower boundary for the barrier transparency. An analogous estimate for the inelastic parameter, with the barrier resistance $`R_TR_N`$ substituting for the Andreev resistance $`R_{SN}/W`$ in Eq. (86), was obtained in Ref. for an SNINS structure with perfect NS interfaces. In this case, the tunnel barrier does not affect the Andreev reflection but produces renormalization of the diffusion coefficient, $`𝒟𝒟R_N/R_T`$.
At $`eV/2\mathrm{\Delta }W_\epsilon ^{1/2}`$, when $`\tau _\epsilon \tau _d(eV)`$, the normal channel may be considered as a reservoir with a certain effective temperature (depending on the details of inelastic scattering which are beyond our model approach), and the $`I`$-$`V`$ curve becomes structureless. At small $`W_\epsilon `$, the SNS junction behaves at two SN junctions connected in series through the equilibrium normal reservoir.
## IX Summary
We have developed a consistent theory of incoherent MAR in long diffusive SNS junctions with arbitrary transparency of the SN interfaces. We formulated a circuit representation for the incoherent MAR, which may be considered as an extension of Nazarov’s circuit theory to a system of voltage biased superconducting terminals connected by normal wires in the absence of supercurrent. We constructed an equivalent MAR network which includes a new resistive element, ”Andreev resistor”, which provides electron-hole conversion at the SN interfaces. Separate Kirchhoff rules are formulated for electron (hole) population numbers and diffusive flows. Within this approach, the electron and hole population numbers are considered as potential nodes connected through the tunnel and Andreev resistors with a distributed voltage source – the equilibrium Fermi reservoirs in the superconducting terminals.
The theory was applied to calculation of the $`I`$-$`V`$ characteristics. The subgap current decreases step-wise with decreasing applied voltage in junctions with resistive interfaces, while in junctions with transparent interfaces an appreciable SGS appears only on the differential resistance. In all cases, $`dV/dI`$ exhibits sharp structures whose maximum slopes correspond to the gap subharmonics, $`eV_n=2\mathrm{\Delta }/n`$. The amplitude and the shape of SGS oscillations strongly depend on the interface/normal-metal resistance ratio $`r=R_{SN}/R_N`$ and reveal a noticeable parity effect: difference of the amplitudes of the even and odd structures. This effect is specific for diffusive junctions: it comes from the strongly enhanced probability of Andreev reflection at zero energy. Inelastic scattering results in smearing of the SGS, which disappears at small applied voltage.
Our theory of incoherent MAR is valid as soon as the applied voltage $`eV`$ is much larger than the Thouless energy: in this case, one may neglect the overlap of the proximity regions near the NS interfaces. In the opposite case, $`eVE_{\mathrm{Th}}\mathrm{\Delta }`$, the problem of the low-energy part of the effective circuit should be considered more carefully, by taking into account the interference between the proximity regions. Aspects of this ac Josephson regime have been considered in Ref. in terms of adiabatic oscillations of the quasiparticle spectrum and distribution functions, which produce nonequilibrium ac supercurrent. Within our approach, this effect will introduce an effective boundary condition for the probability currents at small energy, which must be included into the circuit representation of MAR in energy space.
It is useful to discuss the effect of dephasing on MAR which has not been sufficiently investigated yet. In the present case of diffusive junction, the dephasing is modeled by an imaginary addition $`i\mathrm{\Gamma }`$ to the energy, $`EE+i\mathrm{\Gamma }`$. It is easy to see that this model leads to more pronounced SGS. Indeed, as it follows from Eq. (7), the effect of large dephasing rate $`\mathrm{\Gamma }`$ is similar to the effect of an opaque interface barrier (the transmissivity parameter $`W`$ becomes effectively small). This is easily understood since the dephasing suppresses the anomalous Green’s function within the normal metal similar to the effect of the interface barrier. Thus, even for highly transmissive interfaces, the $`I`$-$`V`$ curves for large $`\mathrm{\Gamma }E_{\text{Th}}`$ become similar to the ones in the case of low-transparent interfaces (see Fig. 7), with deficit current and pronounced SGS. This conclusion is also supported by direct numerical calculation. It is of interest that the parity effect in SGS almost disappears due to the cutoff of the long-range proximity effect at small energies. In such case, the incoherent MAR regime persists at arbitrarily small voltages and our theory is valid until the quasiparticle relaxation will affect MAR regime as described in Sec. VIII.
###### Acknowledgements.
Support from NFR, KVA and NUTEK (Sweden), from NEDO (Japan), and from Fundamental Research Foundation of Ukraine is gratefully acknowledged.
##
The analytical expressions for bare conductivities $`G_\pm `$ and proximity corrections $`m_\pm `$ can be obtained in the case of low-transparent NS interface, $`W1`$, by making use of a perturbative solution of Eq. (7):
$$\theta _N(E)=W\sqrt{i\mathrm{\Delta }/E}\mathrm{sinh}\theta _S(E),$$
(87)
$`R_{SN}G_\pm ={\displaystyle \frac{E\mathrm{\Theta }(E\mathrm{\Delta })}{\sqrt{E^2\mathrm{\Delta }^2}}}{\displaystyle \frac{W\mathrm{\Delta }^2\mathrm{\Theta }[\pm (E\mathrm{\Delta })]}{E^2\mathrm{\Delta }^2}}`$ (89)
$`\times \left[\sqrt{{\displaystyle \frac{\mathrm{\Delta }}{2E}}}{\displaystyle \frac{W\mathrm{\Delta }}{\sqrt{\mathrm{\Delta }^2E^2}}}\left(\begin{array}{ccc}0& & \\ 1& & \end{array}\right)\right],`$ (92)
$`{\displaystyle \frac{m_\pm }{R_N}}={\displaystyle \frac{\xi _\mathrm{\Delta }W^2\mathrm{\Delta }^2}{2d|E^2\mathrm{\Delta }^2|}}\left({\displaystyle \frac{\mathrm{\Delta }}{2E}}\right)^{3/2}[\pm 2`$ (93)
$`+\mathrm{\Theta }(E\mathrm{\Delta })\mathrm{\Theta }(\mathrm{\Delta }E)],`$ (94)
where $`\mathrm{\Theta }(x)`$ is the Heaviside function and $`E`$ is assumed for brevity to be positive. From Eq. (89), we obtain approximations for the tunnel and Andreev resistances:
$$R_T(E)R_{SN}\left(1\mathrm{\Delta }^2/E^2\right),$$
(96)
$`R_A(E){\displaystyle \frac{2R_{SN}|E^2\mathrm{\Delta }^2|}{W\mathrm{\Delta }}}\sqrt{{\displaystyle \frac{2E}{\mathrm{\Delta }}}}`$ (97)
$`\times \left[1+W\mathrm{\Theta }(\mathrm{\Delta }E)\left|2E\mathrm{\Delta }/(E^2\mathrm{\Delta }^2)\right|^{1/2}\right].`$ (98)
In the vicinity of the gap edges, $`|E\mathrm{\Delta }|W^2`$, and at small energies, $`EW^2`$, where $`\theta _N(E)`$ in Eq. (87) diverges, the following approximate solutions of Eq. (7) have to be used instead of Eq. (87):
$$\mathrm{exp}(\theta _N)=u^2(t),t=2|E\mathrm{\Delta }|/\mathrm{\Delta }W^2$$
(99)
at $`|E\mathrm{\Delta }|\mathrm{\Delta }`$, where $`u(t)`$ is the solution of a cubic equation $`u^3u=\sqrt{i/t}`$, and
$$\mathrm{sinh}\frac{\theta _N}{2}=\frac{i}{\sqrt{2}}\mathrm{exp}\left(Arcsinh\sqrt{iE/2W^2\mathrm{\Delta }}\right)$$
(100)
at $`|E|\mathrm{\Delta }`$. The asymptotics of $`m_\pm =\pm R_N(\xi _\mathrm{\Delta }/d)\mu _\pm `$ and $`G_\pm `$ near these “dangerous” points, are given by
$$R_{SN}G_+=\left(\sqrt{3}t^{1/6}/2W\right)\mathrm{\Theta }(E\mathrm{\Delta }),$$
(102)
$$R_{SN}G_{}=(t^{5/6}/4W)\left[\sqrt{3}\mathrm{\Theta }(E\mathrm{\Delta })+\mathrm{\Theta }(\mathrm{\Delta }E)\right],$$
(103)
$$\mu _+=\sqrt{2/3t^{1/3}}\left(\sqrt{2+\sqrt{3}}+\sqrt{2\sqrt{3}}\right),\mu _{}=\sqrt{2}$$
(104)
at $`t1`$. At $`EW^2`$,
$$R_{SN}G_{}=1,\mu _{}=\left(\sqrt{2}1\right)\sqrt{\mathrm{\Delta }/E},$$
(105)
For perfect SN interface, $`G_\pm ^1=0`$, the asymptotics of $`\mu _{}(E)`$ at $`E\mathrm{\Delta }\mathrm{\Delta }`$ and $`E\mathrm{\Delta }`$ are given by Eqs. (104), (105) respectively, whereas $`m_+(E)`$ diverges at the gap edge as $`[\mathrm{\Delta }/2(E\mathrm{\Delta })]^{1/4}`$. Several examples of these dependencies calculated numerically are presented in Fig. 2 (see discussion in Sec. IV).
|
warning/0005/quant-ph0005018.html
|
ar5iv
|
text
|
# Quantum Kolmogorov Complexity
## 1 Introduction
In classical computations, the Kolmogorov-Solomonoff-Chaitin (Kolmogorov, for short) complexity of a finite string is a measure of its randomness. The Kolmogorov complexity of $`x`$ is the length of the shortest program which produces $`x`$ as its output. It can be seen as a lower bound on the optimal compression that $`x`$ can undergo, and it is closely related to Shannon information theory.
Kolmogorov complexity has been shown to have a windfall of applications in fields as diverse as learning theory, complexity theory, combinatorics and graph theory, analysis of algorithms, to name just a few.
With the advent of quantum computation, it is natural to ask what is a good definition for the Kolmogorov complexity of quantum strings. Our goal in this paper is to argue that our definition is a natural and robust measure the amount of quantum information contained in a quantum string, which has several appealing properties.
Recently, Paul Vitányi has also proposed a definition for quantum algorithmic complexity. Our definition differs significantly from Vitányi’s: the definition he proposes is a measure of the amount of *classical* information necessary to approximate the quantum state.
The paper will be organized as follows: In Section 3, we give basic notation, definitions, prior work and some theorems that will be used in proofs in the paper. In Section 4 we give our definition of quantum Kolmogorov complexity. In Section 5 we prove the invariance theorem. Section 6 compares the properties of quantum and classical Kolmogorov complexity, including incompressibility, subadditivity, and the complexity of copies. Section 7 discusses the relationship with quantum information theory. We conclude with a discussion of possible extensions and future work.
## 2 What is a Good Definition?
A good definition of quantum Kolmogorov complexity should meet the following fundamental criteria. These are intended to insure that it gives an accurate representation of the information content of a quantum string.
* It should be robust, that is, invariant under the choice of the underlying quantum Turing machine.
* It should bear a strong relationship with quantum information theory.
* It should be closely related to classical complexity on classical strings.
However, quantum Kolmogorov complexity should not be expected to always behave the way classical Kolmogorov complexity does. The reader may want to bear in mind quantum phenomena such as the no-cloning theorem, whose consequences we will discuss later in the paper.
### 2.1 Critical issues
A first attempt at defining quantum Kolmogorov complexity of a qubit string $`X`$ is to consider the length of the shortest quantum program that produces $`X`$ as its output. There are many questions that arise from this ‘definition’.
Bits or qubits? The first question to consider is whether we want to measure the amount of algorithmic information of a string in bits, or in qubits. Note that bit strings (programs) are countable, whereas qubit strings are uncountable, so any definition that measures in bits would have to overcome this apparent contradiction. Paul Vitányi considers classical descriptions of qubit strings, whereas we consider qubit descriptions.
Exact or inexact? What does ‘produce’ mean? Is a minimal program required to produce the string $`X`$ exactly, or only up to some fidelity? In the latter case, is the fidelity a constant? Otherwise, how is it parameterized? (For exact simulation, we can only hope to simulate a subclass of the Turing machines, say by restricting the set of possible amplitudes. What would be a reasonable choice?) We will use an approximation scheme.
What model of computation? Size of quantum circuits is not an appropriate measure since large circuits may be very simple to describe. The Turing machine model is the appropriate one to consider.
What is meant by ‘quantum program?’ A program for a quantum Turing machine is its input, and if we want to count program length in qubits, we must allow for ‘programs’ to be arbitrary qubit strings. (These can be viewed as programs whose code may include some auxiliary ‘hard-coded’ qubit strings.)
One-time description or multiple generation? In the classical setting, the program that prints the string $`x`$ can be run as many times as desired. Because of the no-cloning theorem of quantum physics however, we cannot assume that the shortest program can be run several times to produce several copies of the same string. This may be due to the fact that it is not possible to recover the program without losing its output. There is also a second reason not to choose the multiple generation option. The complex-valued parameters $`\alpha `$ and $`\beta `$ of a qubit $`|q=\alpha |0+\beta |1`$ can contain an unbounded amount of information. If we would be able to reproduce $`q`$ over and over again, then we would have to conclude that the single qubit $`q`$ contains an unlimited amount of information. This contradicts the fact that the quantum mechanical system of $`q`$ can only contain one bit of information. For the above two reason, we will not require a ‘reusability’ condition.
## 3 Preliminaries
We start with some notation, definitions, and results that will be used to prove the results in this paper.
### 3.1 Notation
We use $`x`$,$`y`$,…to denote finite, classical Boolean strings. When we write $`|x`$, we mean the quantum state vector in the standard basis that corresponds to the classical string $`x`$. In general we use $`\varphi ,\psi ,\mathrm{}`$ to denote quantum pure states. Mixed states are represented by the letters $`\rho ,\sigma `$ etc. We also use uppercase letters $`X,Y,\mathrm{}`$ for (mixed) quantum states that are strings of qubits. The terms quantum state, qubit string, and quantum register are used interchangeably (sometimes to emphasize the purpose of the quantum state at hand.) Lower-case letters $`i,j,k,l,m,n`$ denote integer indices or string lengths.
For classical strings over the alphabet $`\mathrm{\Sigma }=\{0,1\}`$, $`\mathrm{}(x)`$ denotes the length of the string. For finite sets $`A,|A|`$ denotes the cardinality of the set. Concatenation of $`x,y`$ is written as the juxtaposition $`xy`$, and the $`n`$-fold concatenation of $`x`$ is written $`x^n`$.
For Hilbert spaces, we write $`_d`$ for the $`d`$-dimensional Hilbert space and $`^m`$ for the $`m`$-fold tensor product space $`\mathrm{}`$. A pure quantum state $`\varphi `$ represented as a vector in such a Hilbert space is denoted by the ket $`|\varphi `$. The *fidelity* between two pure states $`\varphi `$ and $`\psi `$ is the absolute value of the inner product of the two vectors: $`|\varphi |\psi |`$ (although some authors use the square of this value).
We slightly abuse notation by sometimes letting the state symbols $`\varphi ,\rho ,\mathrm{}`$ also stand for the corresponding density matrices. Hence, a pure state $`\varphi `$ as a Hilbert space vector is denoted by $`|\varphi `$, whereas its density matrix $`|\varphi \varphi |`$ can also be denoted by $`\varphi `$.
A density matrix can always be decomposed as a mixture of pure, orthogonal states: $`\rho =_ip_i|\varphi _i\varphi _i|`$, with $`p_1,p_2,\mathrm{}`$ a probability distribution over the mutually orthogonal states $`\varphi _1,\varphi _2,\mathrm{}`$. The matrix $`\rho `$ represents a pure state if and only if $`\rho ^2=\rho `$, in which case we can also say $`\sqrt{\rho }=\rho `$. The square root of a general mixed state is described by
$$\sqrt{\rho }=\sqrt{\underset{i}{}p_i|\varphi _i\varphi _i|}=\underset{i}{}\sqrt{p_i}|\varphi _i\varphi _i|.$$
We use the above rule for the generalization of the fidelity to mixed states. The fidelity between two density matrices $`\rho `$ and $`\sigma `$ is defined by
$`\mathrm{Fidelity}(\rho ,\sigma )`$ $`=`$ $`\mathrm{tr}\left(\sqrt{\sqrt{\rho }\sigma \sqrt{\rho }}\right).`$ (1)
For pure states $`\varphi `$ and $`\psi `$, the above definition coincides again with the familiar $`|\varphi |\psi |`$. If $`\mathrm{Fidelity}(\rho ,\sigma )=1`$, then $`\rho =\sigma `$, and vice versa.
An ensemble $``$ is specific distribution $`p_1,p_2,\mathrm{}`$ over a set of (mixed) states $`\rho _1,\rho _2,\mathrm{}`$. We denote this by $`=\{(\rho _i,p_i)\}`$. The average state of such an ensemble $``$ is $`\rho =_ip_i\rho _i`$. An average state corresponds to several different ensembles. When an ensemble is used to produce a sequence of states $`\rho _i`$ according to the probabilities $`p_i`$, we speak of a *source* $``$.
The length of a quantum state is denoted by $`\mathrm{}(X)`$, by which we mean the smallest $`l`$ for which $`X`$ sits in the $`2^l`$-dimensional Hilbert space (in the standard basis).
A transformation $`\$`$ on the space of density matrices is allowed by the laws of quantum mechanics if and if only it is a completely positive, trace preserving mapping.
### 3.2 Classical Kolmogorov complexity
The Kolmogorov complexity of a string, in the classical setting, is the length of the shortest program which prints this string on an empty input.
Formally, this is stated first relative to a partial computable function, which as we know can be computed by a Turing machine.
###### Definition 1
Fix a Turing machine $`T`$ that computes the partial computable function $`\mathrm{\Phi }`$. For any pair of strings $`x,y\{0,1\}^{}`$, the Kolmogorov complexity of $`x`$ relative to $`y`$ (with respect to $`\mathrm{\Phi }`$) is defined as
$`C_\mathrm{\Phi }(x\text{ }y)`$ $`=`$ $`\mathrm{Min}\{\mathrm{}(p):\mathrm{\Phi }(p,y)=x\}.`$
When $`y`$ is the empty string, we simply write $`C_\mathrm{\Phi }(x)`$. Also the notation $`C_T(x\text{ }y)`$ is used.
The key theorem on which rests the robustness of Kolmogorov complexity is the *invariance theorem*. This theorem states that the length of shortest programs does not depend by more than an additive constant on the underlying Turing machine. In the classical case, this theorem is proven with the existence of a universal Turing machine. This machine has two inputs: a finite description of the original Turing machine, and the program that this Turing machine executes to output the string.
More formally, the invariance theorem in the classical case can be stated as follows.
###### Theorem 1
There is a universal partial computable function $`\mathrm{\Phi }_0`$ such that for any partial computable $`\mathrm{\Phi }`$ and pair of strings $`x,y`$,
$`C_{\mathrm{\Phi }_0}(x\text{ }y)`$ $``$ $`C_\mathrm{\Phi }(x\text{ }y)+c,`$
where $`c`$ is a constant depending only on $`\mathrm{\Phi }`$.
Giving an invariance theorem will be key to showing that quantum Kolmogorov complexity is robust.
Since for any string $`x`$ of length $`n`$, $`C(x)n+O(1)`$, a string which has complexity at least $`n`$ is called *incompressible*. The existence of incompressible strings is a crucial fact of Kolmogorov complexity.
###### Proposition 1
For every string length $`n`$, there is a string $`x`$ of length $`n`$ such that $`C(x)n`$.
The proof that there exists incompressible strings is a simple application of the pigeonhole principle. By comparing the number of strings of length $`n`$ ($`2^n`$) and the number of programs of length smaller than $`n`$ ($`2^n1`$ in total), one must conclude that there is at least one string of length $`n`$ which is not the output of any of the program of length $`<n`$.
### 3.3 Entropy of classical sources
The Shannon entropy of a random source that emits symbols from an alphabet is a measure of the amount of randomness in the source.
###### Definition 2
Let $`A`$ be a random source that emits letter $`x_i`$ (independently) with probability $`p_i`$. The Shannon entropy $`H`$ of $`A`$ is $`H(A)=_ip_i\mathrm{log}p_i`$.
In the classical setting, Kolmogorov complexity and Shannon entropy are closely related, as we describe now. This is an important property of Kolmogorov complexity, and one would expect a similarly strong relationship to hold between quantum Kolmogorov complexity and quantum entropy.
Shannon’s noiseless coding theorem states that the entropy corresponds to the average number of bits required to encode sequences of character emitted by a random source.
###### Proposition 2
Shannon’s noiseless coding : Consider a classical channel $`\mathrm{A}`$ that is used to transmit letters taken from an ensemble $`\{(\mathrm{x}_\mathrm{i},\mathrm{p}_\mathrm{i})\}`$, where the $`\mathrm{x}_\mathrm{i}`$ are the letters and $`\mathrm{p}_\mathrm{i}`$ their corresponding probabilities. Then
1. for any $`ϵ,\delta `$, there is an $`n`$ such that there is an encoding that on $`n`$ letters encodes on average the letters with $`H(A)+\delta `$ bits for which the probability of successfully decoding $`P_{\mathrm{success}}1ϵ`$;
2. for any $`ϵ,\delta `$, there is an $`n`$ such that for any $`\delta ^{}`$, there is an $`ϵ^{}`$ such that if the channel encodes $`n`$ letters, each letter with less than $`H(A)\delta ^{}`$ bits per letter, then the probability of success $`P_{\mathrm{success}}2^{n(\delta ^{}\delta )}+ϵ^{}`$.
In the classical case, the Kolmogorov complexity of a string is bounded by the entropy of a source ‘likely to have emitted this string’. A brief summary of the argument is included here. (Details can be found in \[12, page 180\].)
Let $`x`$ be a (long) binary string. It can be broken down into $`m`$ blocks of length $`k`$, where each block is thought of as a character in an alphabet of size $`2^k`$. Define the frequency $`f_i`$ of a character $`c_i`$ to be the number of times it appears as a block in $`x`$, and let $`A`$ represent the source $`\{c_i,f_i/m\}`$. To reconstruct $`x`$, it suffices to provide the frequency of each character ($`_i\mathrm{log}f_i`$ bits) and then specify $`x`$ among the strings that share this frequency pattern. With some manipulations, it can be shown that
###### Proposition 3
$`C(x)`$ $`<`$ $`m(H(A)+\gamma ),`$
where $`A`$ is the source defined in the discussion above, and $`\gamma `$ vanishes as $`m`$ goes to infinity.
### 3.4 Quantum information theory
We have seen that in the classical setting, Kolmogorov complexity is very closely related to Shannon entropy. In this section we describe the quantum, or Von Neumann, entropy, related measures, and important properties which will be used in the proofs of our results.
###### Definition 3
Von Neumann entropy: The Von Neumann entropy of a mixed state $`\mathrm{\rho }`$ is defined as $`\mathrm{S}(\mathrm{\rho })=\mathrm{tr}(\mathrm{\rho }\mathrm{log}\mathrm{\rho })`$. If we decompose $`\mathrm{\rho }`$ into its mutually orthogonal eigenstates $`\mathrm{\varphi }_\mathrm{i}`$, we see that
$`S(\rho )=S\left({\displaystyle \underset{i}{}}p_i|\varphi _i\varphi _i|\right)=H(p),`$
where $`H(p)`$ is the Shannon entropy of the probability distribution $`p_1,p_2,\mathrm{}`$
A source $`=\{(\rho _i,p_i)\}`$ has an associated Von Neumann entropy $`S(\rho )`$ of the average state $`\rho =_ip_i\rho _i`$. Schumacher’s noiseless coding theorem shows how to obtain an encoding with average letter-length $`S(\rho )`$ for a source of pure states, where the fidelity of the encoding goes to $`1`$ as the number of letters emitted by the source goes to infinity. (A survey can be found in Preskill’s lecture notes \[15, page 190\] or in Nielsen’s thesis \[14, Chapter 7\].)
We will use a slightly stronger result, which gives a universal compression scheme. That is, one that does not depend on the source itself, but only on its entropy. This result is due to Jozsa et al. , building upon the work of Jozsa and Schumacher .
###### Theorem 2
Universal quantum compression (see ): Consider pure state sources $`=\{(\mathrm{\varphi }_\mathrm{i},\mathrm{p}_\mathrm{i})\}`$. For any $`\mathrm{ϵ},\mathrm{\delta }`$, there is an $`\mathrm{n}=\mathrm{n}(\mathrm{ϵ},\mathrm{\delta })`$ such that for any entropy bound $`\mathrm{S}`$, there is an encoding scheme that works for any source of Von Neumann entropy at most $`\mathrm{S}`$ that has the following properties. Let $`\mathrm{\rho }=_\mathrm{i}\mathrm{p}_\mathrm{i}|\mathrm{\varphi }_\mathrm{i}\mathrm{\varphi }_\mathrm{i}|`$ be the average state, with all $`|\mathrm{\varphi }_\mathrm{i}_\mathrm{d}`$, and $`\mathrm{\rho }`$ has entropy $`\mathrm{S}(\mathrm{\rho })\mathrm{S}`$, then
1. Each $`|\varphi _i`$ can be encoded by a code word $`\sigma _i`$, which has length $`S+\delta +\frac{1}{n}(d^2\mathrm{log}(n+1))`$.
2. For each $`i`$, $`\mathrm{Fidelity}(\varphi _i,\sigma _i)1ϵ`$.
We continue the section by defining the ‘$`\chi `$ quantity’ for ensembles.
###### Definition 4
Holevo’s chi quantity : For an ensemble $`=\{(\mathrm{\rho }_\mathrm{i},\mathrm{p}_\mathrm{i})\}`$, with $`\mathrm{\rho }=_\mathrm{i}\mathrm{p}_\mathrm{i}\mathrm{\rho }_\mathrm{i}`$, Holevo’s chi quantity equals
$`\chi ()`$ $`=`$ $`S(\rho ){\displaystyle \underset{i}{}}p_iS(\rho _i).`$
Note that the $`\chi `$ quantity depends not only on $`\rho `$, but also on the specific pairs $`(p_i,\rho _i)`$.
The following monotonicity property of Lindblad and Uhlmann will be very useful later in the paper.
###### Theorem 3
Lindblad-Uhlmann monotonicity : Let $`=\{(\mathrm{\rho }_\mathrm{i},\mathrm{p}_\mathrm{i})\}`$ be an ensemble, and $`\$`$ a completely positive, trace preserving mapping. For every such $``$ and $`\$`$, it holds that: $`\mathrm{\chi }(\$())\mathrm{\chi }()`$, where $`\$()`$ is the transformed ensemble $`\{(\$(\mathrm{\rho }_\mathrm{i}),\mathrm{p}_\mathrm{i})\}`$.
The entropy of finite systems is robust against small changes. This continuity of $`S`$ over the space of finite dimensional density matrices $`\rho `$ is also called *insensitivity*, and is expressed by the following lemma.
###### Lemma 1
Insensitivity of Von Neumann entropy (see Section II.A in ): If a sequence $`\mathrm{\rho }_1,\mathrm{\rho }_2,\mathrm{}`$, has $`lim_\mathrm{k}\mathrm{}\mathrm{\rho }_\mathrm{k}=\mathrm{\rho }`$, then also $`lim_\mathrm{k}\mathrm{}\mathrm{S}(\mathrm{\rho }_\mathrm{k})=\mathrm{S}(\mathrm{\rho })`$.
Proof: The convergence of $`\rho _1,\rho _2,\mathrm{}`$ to $`\rho `$ is understood to use some kind of norm for the density matrices that is continuous in the matrix entries $`i|\rho |j`$. (The operator norm $`|\rho |=\mathrm{tr}(\rho \rho ^{})`$, for example.) The entropy $`S(\rho )`$ is a continuous function of the finite set of eigenvalues of $`\rho `$. These eigenvalues are also continuous in the entries of $`\rho `$. $``$$``$
Further background on these measures of quantum information and their properties can be found in \[15, Chapter 5\]. Another good source is Nielsen’s thesis .
### 3.5 Symmetric spaces
We use the symmetric subspace of the Hilbert space to show some of our results on copies of quantum states. Let $`_D`$ be a Hilbert space of dimension $`D`$ with the basis states labeled $`|1,\mathrm{},|D`$. The symmetric subspace $`\mathrm{Sym}(_D^m)`$ of the $`m`$-fold tensor product space $`_D^m`$ is a subspace spanned by as many basis vectors as there are multisets of size $`m`$ of $`\{1,\mathrm{},D\}`$. Let $`A=\{i_1,\mathrm{},i_m\}`$ be such a multiset of $`\{1,\mathrm{},D\}`$. Then, $`|s_A`$ is the normalized superposition of all the different permutations of $`i_1,\mathrm{},i_m`$. The set of the different vectors $`|s_A`$ (ranging over the multisets $`A`$) is an orthogonal basis of the symmetric subspace $`\mathrm{Sym}(_D^m)`$. Hence the dimension of the symmetric subspace is $`\left(\genfrac{}{}{0pt}{}{m+D1}{D1}\right)`$. (This is because choosing a multiset is the same thing as splitting $`m`$ consecutive elements into $`D`$ (possibly empty) intervals, where the size of $`i`$th interval represents the number of times the $`i`$th element appears in the multiset. The number of ways of splitting an interval of size $`m`$ into $`D`$ intervals is $`\left(\genfrac{}{}{0pt}{}{m+D1}{D1}\right)`$.)
An equivalent definition of the symmetric subspace is that it is the smallest subspace that contains all the states of the form $`|\varphi ^m`$, for all $`|\varphi _D`$. (For more on the symmetric subspace and its properties, see the paper by Barenco et al. .)
### 3.6 Accumulation of errors
The following lemma is used to bound the error introduced when composing two inexact quantum procedures.
###### Lemma 2
Fidelity of composition: If $`\mathrm{Fidelity}(\mathrm{\rho },\mathrm{\rho }^{})1\mathrm{\delta }_1`$ and $`\mathrm{Fidelity}(\mathrm{\rho }^{},\mathrm{\rho }^{\prime \prime })1\mathrm{\delta }_2`$, then $`\mathrm{Fidelity}(\mathrm{\rho },\mathrm{\rho }^{\prime \prime })12\mathrm{\delta }_12\mathrm{\delta }_2`$.
Proof: This follows from the fact that the fidelity between two mixed states $`\rho `$ and $`\sigma `$ equals the maximum ‘pure state fidelity’ $`|\varphi |\psi |`$, where $`\varphi `$ and $`\psi `$ are ‘purifications’ of $`\rho `$ and $`\sigma `$. (See for more details on this.) $``$$``$
In order to give bounds on the complexity of several copies of a state, as we do in Section 6.3, we need the following bound on the total error in the $`n`$-fold tensor product of the approximation of a given state.
###### Lemma 3
Let $`\rho ^n`$ and $`\sigma ^n`$ be the $`n`$-fold copies of the mixed states $`\rho `$ and $`\sigma `$, then $`\mathrm{Fidelity}(\rho ^n,\sigma ^n)=(\mathrm{Fidelity}(\rho ,\sigma ))^n`$.
Proof: This follows directly from the definition $`\mathrm{Fidelity}(\rho ,\sigma )=\mathrm{tr}\left(\sqrt{\sqrt{\rho }\sigma \sqrt{\rho }}\right)`$. $``$$``$
## 4 Quantum Kolmogorov Complexity
We define the *quantum Kolmogorov complexity* $`\mathrm{𝑄𝐶}`$ of a string of qubits, relative to a quantum Turing machine $`M`$, as the length of the shortest qubit string which when given as input to $`M`$, produces on its output register the qubit string. (Note that we only allow $`M`$ that have computable transition amplitudes. See the articles , and particularly Definition 3.2.2 in , for a further description of this computational model.)
### 4.1 Input/Output Conventions
We give some precisions about what is meant by ‘input’ and ‘output’.
We consider quantum Turing machines with two heads on two one-way infinite tapes. We allow the input tape to be changed. This is required: for example, the contents of the input may have to be moved to the output tape.
For a QTM $`M`$ with a single input, when we say $`M`$ starts with input $`Y`$, we mean that M starts with the quantum state $`|Y\$00\mathrm{}`$ on its input tape, and $`|00\mathrm{}`$ on the output tape. The $`\$`$ symbol is a special endmarker (or blank) symbol.
Note that testing for the end of the input can be done without disturbing the input, since we assume that the ‘$’ state is orthogonal to the ‘0’ and ‘1’ states. (This is analogous to the classical case, where where Turing machine inputs are encoded in a three-letter alphabet; nevertheless we consider the actual input to be encoded only over the characters 0 and 1.)
A string is a proper input if the endmarker symbol appears only once and is not in superposition with any other position of the tape. We dismiss any non-proper inputs.
For a QTM with multiple inputs, we also assume that there is a convention for encoding the multiple inputs so that they can be individually recovered. For example, when we write $`M(P,Y)`$, we may assume that the input tape is initialized to $`|1^{\mathrm{}(P)}PY\$00\mathrm{}`$. We only count the length of $`X`$ and $`Y`$ for the length of the input. Likewise, for multiple outputs, if we write $`M(P,Y)=(X_1,X_2)`$, we mean that $`X_1`$ and $`X_2`$ must be encoded according to a prearranged convention so that $`X_1`$ and $`X_2`$ can be recovered individually from the output tape.
(Note that we do not define prefix-free complexity in this paper. The programs themselves need not be prefix-free.)
We let $`M^T(X)`$ denote the contents of the output tape after $`T`$ steps of computation. We consider only QTMs which do not modify their output tape after they have halted. (Because of reversibility, they may modify the input tape after reaching the halting state.) The output $`M(X)`$ is the content of the output tape at any time after $`M`$ has stopped changing its output tape.
### 4.2 Definitions
For some fidelity function $`f:[0,1]`$ we will now define the corresponding quantum Kolmogorov complexity.
###### Definition 5
Quantum Kolmogorov complexity with fidelity f: For any quantum Turing machine $`\mathrm{M}`$ and qubit string $`\mathrm{X}`$, the $`\mathrm{f}`$-approximation quantum Kolmogorov complexity, denoted $`\mathrm{QC}_\mathrm{M}^\mathrm{f}(\mathrm{X})`$, is the length of the smallest qubit string $`\mathrm{P}`$ such that for any fidelity parameter $`\mathrm{k}`$ we have $`\mathrm{Fidelity}(\mathrm{X},\mathrm{M}(\mathrm{P},1^\mathrm{k}))\mathrm{f}(\mathrm{k})`$.
Note that we require that the same string $`P`$ be used for all approximation parameters $`k`$.
We will say that program $`P`$ $`M`$-computes $`X`$ with fidelity $`f(k)`$ if $`k,\mathrm{Fidelity}(M(P,1^k),X)f(k)`$.
If $`f`$ is the constant function $`1`$, we have the following definition.
###### Definition 6
Quantum Kolmogorov complexity with perfect fidelity: The perfect fidelity quantum Kolmogorov complexity is $`\mathrm{QC}_\mathrm{M}^1(\mathrm{X})`$.
The problem with this definition is that it is not known whether an invariance theorem can be given for the ideal Kolmogorov complexity. This is because the invariance theorems that are known for quantum computers deal with *approximating* procedures. We therefore prove an invariance theorem for a weaker, limiting version, where the output of $`M`$ must have high fidelity with respect to the target string $`X`$: $`\mathrm{Fidelity}(X,M(P))1`$.
###### Definition 7
Quantum Kolmogorov complexity with bounded fidelity: For any constant $`\mathrm{ϵ}<1`$, $`\mathrm{QC}_\mathrm{M}^\mathrm{ϵ}(\mathrm{X})`$ is the constant-fidelity quantum Kolmogorov complexity.
There are two problems with this definition. First, it may be the case that some strings are very easy to describe up to a given constant, but inherently very hard to describe for a smaller error. Second, it may be the case that some strings are easier to describe up to a given constant on one machine, but not on another machine. For these two reasons, this definition does not appear to be robust.
A stronger notion of approximability is the existence of an approximation *scheme.* (See, for example, the book by Garey and Johnson \[7, Chapter 6\] for more on approximation algorithms and approximation schemes.)
For constant-approximability, different algorithms (with different sizes) can exist for different constants. In an approximation scheme, a single program takes as auxiliary input an approximation parameter $`k`$, and produces an output that approximates the value we want within the approximation parameter. This is the model we wish to adopt for quantum Kolmogorov complexity.
###### Definition 8
Quantum Kolmogorov complexity with fidelity converging to 1: The complexity $`\mathrm{QC}_\mathrm{M}^1(\mathrm{X})`$ is equal to $`\mathrm{QC}_\mathrm{M}^\mathrm{f}(\mathrm{X})`$, where $`\mathrm{f}(\mathrm{k})=1\frac{1}{\mathrm{k}}`$.
We choose to encode the fidelity parameter in unary, and the convergence function to be $`f(k)=1\frac{1}{k}`$ so that the model remains robust when polynomial time bounds are added. We discuss this further in Section 5.
We may also define $`\mathrm{𝑄𝐶}_M^1(X\text{ }Y)`$, the complexity of producing $`X`$ when $`Y`$ is given as an auxiliary input, in the usual way.
## 5 Invariance
To show that our definition is robust we must show that the complexity of a qubit string does not depend on the underlying quantum Turing machine.
We use the following result, proved in the paper of Bernstein and Vazirani . To be precise, we use the notation $`\overline{M}`$ to denote the classical description of the quantum Turing machine $`M`$. (Recall that we only consider quantum Turing machines whose amplitudes can be computed to arbitrary precision with a finite classical description.)
###### Theorem 4
Universal quantum Turing machine (see ): There exists a universal quantum Turing machine $`\mathrm{U}`$ that has a finite classical description such that the following holds. For any quantum Turing machine $`\mathrm{M}`$ (which has a finite classical description), for any pure state $`\mathrm{X}`$, for any approximation parameter $`\mathrm{k}`$, and any number of time steps $`\mathrm{T}`$, $`\mathrm{Fidelity}(\mathrm{U}(\overline{\mathrm{M}},\mathrm{X},1^\mathrm{k},\mathrm{T}),\mathrm{M}^\mathrm{T}(\mathrm{X}))1\frac{1}{\mathrm{k}}`$. Recall that $`\mathrm{M}^\mathrm{T}`$ is the contents of the output tape of $`\mathrm{M}`$ after $`\mathrm{T}`$ time steps.
###### Theorem 5
There is a universal quantum Turing machine $`U`$ such that for any quantum Turing machine $`M`$ and qubit strings $`X`$,
$`\mathrm{𝑄𝐶}_U^1(X)`$ $``$ $`\mathrm{𝑄𝐶}_M^1(X)+c_M,`$
where $`c_M`$ is a constant depending only on $`M`$.
Proof: The proof follows from the existence of a universal quantum Turing machine, as proven by Bernstein and Vazirani . Let $`U`$ be this UTM as mentioned above. The constant $`c_M`$ represents the size of the finite description that $`U`$ requires to calculate the transition amplitudes of the machine $`M`$. Let $`P`$ be the state that witness that $`\mathrm{𝑄𝐶}_M^1(X)=\mathrm{}(P)`$, and hence $`\mathrm{Fidelity}(X,M(P,1^k))1\frac{1}{k}`$ for every $`k`$.
With the description corresponding to $`c_M`$, $`U`$ can simulate with arbitrary accuracy the behavior of $`M`$. Specifically, $`U`$ can simulate machine $`M`$ on input $`(P,1^{4k})`$ with a fidelity of $`1\frac{1}{4k}`$. Therefore, by Lemma 2, $`\mathrm{Fidelity}(X,U(M,P,1^{4k}))1\frac{1}{k}`$. $``$$``$
The same holds true for the conditional complexity, that is, $`UM,X,Y`$, $`\mathrm{𝑄𝐶}_U^1(X\text{ }Y)\mathrm{𝑄𝐶}_M^1(X\text{ }Y)+c_M`$.
Henceforth, we will fix a universal quantum Turing machine $`U`$ and simply write $`\mathrm{𝑄𝐶}(X)`$ instead of $`\mathrm{𝑄𝐶}_U^1(X)`$. Likewise we write $`\mathrm{𝑄𝐶}(X|Y)`$ instead of $`\mathrm{𝑄𝐶}_U^1(X|Y)`$. We also abuse notation and write $`M`$ instead of $`\overline{M}`$ to represent the code of the quantum Turing machine $`M`$ used as an input to the universal Turing machine.
We may also define time-bounded $`\mathrm{𝑄𝐶}`$ is the usual way, that is, fix $`T:`$ a fully-time-computable function. Then $`\mathrm{𝑄𝐶}^T(X|Y)`$ is the length of the shortest program which on input $`Y,1^k`$, produces $`X`$ on its output tape after $`T(\mathrm{}(X)+\mathrm{}(Y))`$ computation steps. The Bernstein and Vazirani simulation entails a polynomial time blowup (polynomial in the length of the input and the length of the fidelity parameter encoded in unary), so there is a polynomial time blowup in the corresponding invariance theorem.
The simplest application of the invariance theorem is the following proposition.
###### Proposition 4
For any qubit string $`X`$, $`\mathrm{𝑄𝐶}(X)\mathrm{}(X)+c`$, where $`c`$ is a constant depending only on our choice of the underlying universal Turing machine.
Proof: Consider the quantum Turing machine $`M`$ that moves its input to the output tape, yielding $`\mathrm{𝑄𝐶}_M(X)=\mathrm{}(X)`$. The proposition follows by invariance. $``$$``$
## 6 Properties of Quantum Kolmogorov Complexity
In this section we compare classical and quantum Kolmogorov complexity by examining several properties of both. We find that many of the properties of the classical complexity, or natural analogues thereof, also hold for the quantum complexity. A notable exception is the complexity of $`m`$-fold copies of arbitrary qubit strings.
### 6.1 Correspondence for classical strings
We would like to show that for classical states, classical and quantum Kolmogorov complexity coincide, up to a constant additive term.
###### Proposition 5
For any finite, classical string $`x`$, $`\mathrm{𝑄𝐶}(x)C(x)+O(1)`$.
(The constant hidden by the big-$`O`$ notation depends only on the underlying universal Turing machine.)
Proof: This is clear: the universal quantum computer can also simulate any classical Turing machine. $``$$``$
We leave as a tantalizing open question whether the converse is also true, that is:
###### Open Problem 1
Is there a constant $`c`$ such that for every finite, classical string $`x`$, $`C(x)\mathrm{𝑄𝐶}(x)+c`$?
### 6.2 Quantum incompressibility
In this section, we show that there exist quantum-incompressible strings.
Our main theorem is a very general form of the incompressibility theorem. We state some useful special cases as corollaries.
Assume we want to consider the minimal-length programs that describe a set of quantum states. In general, these may be pure or mixed states. We will use the following notation throughout the proof. The mixed states $`\rho _1,\mathrm{},\rho _M`$ be are the target strings (those we want to produce as output). Their minimal-length programs will be $`\sigma _1,\mathrm{},\sigma _M`$, respectively. The central idea is that if the states $`\rho _i`$ are sufficiently different, then the programs $`\sigma _i`$ must be different as well. We turn this into a quantitative statement with the use of the insensitive chi quantity in combination with the monotonicity of quantum mechanics.
###### Theorem 6
For any set of strings $`\rho _1,\mathrm{},\rho _M`$ such that $`i,QC(\rho _i)l`$, this $`l`$ is bounded from below by
$`l`$ $``$ $`S(\rho )\frac{1}{M}{\displaystyle \underset{i}{}}S(\rho _i),`$
where $`\rho `$ is the ‘average’ density matrix $`\rho =\frac{1}{M}_i\rho _i`$.
(Stated slightly differently, this says that there is an $`i`$ such that $`\mathrm{𝑄𝐶}(\rho _i)S(\rho )\frac{1}{M}_iS(\rho _i)`$.)
Proof: Take $`\rho _1,\mathrm{},\rho _M`$ and their minimal programs $`\sigma _1,\mathrm{},\sigma _M`$ (and hence $`QC(\rho _i)=\mathrm{}(\sigma _i)`$). Let $`\$^k`$ be the completely positive, trace preserving map corresponding to the universal QTM $`U`$ with fidelity parameter $`k`$. With this, we define the following three uniform ensembles:
* the ensemble $`=\{(\rho _i,\frac{1}{M}\}`$ of the original strings,
* $`_\sigma `$ the ensemble of programs $`\{(\sigma _i,\frac{1}{M})\}`$, and
* the ensemble of the $`k`$-approximations $`\stackrel{~}{}^k=\$^k(_\sigma )=\{(\stackrel{~}{\rho }_i^k,\frac{1}{M})\}`$, with $`\stackrel{~}{\rho }_i^k=\$^k(\sigma _i)`$.
By the monotonicity of Theorem 3 we know that for every $`k`$, $`\chi (\stackrel{~}{}^k)\chi (_\sigma )`$. The chi factor of the ensemble $`_\sigma `$ is upper bounded by the maximum size of its strings: $`\chi (_\sigma )\mathrm{max}_i\{\mathrm{}(\sigma _i)\}l`$. Thus the only thing that remains to be proven is that $`\chi (\stackrel{~}{}^k)`$, for sufficiently big $`k`$, is ‘close’ to $`\chi ()`$. This will be done by using the insensitivity of the Von Neumann entropy.
By definition, for all $`i`$, $`lim_k\mathrm{}\mathrm{Fidelity}(\rho _i,\stackrel{~}{\rho }_i^k)=1`$, and hence $`lim_k\mathrm{}\stackrel{~}{\rho }_i^k=\rho _i`$. Because the ensembles $``$ and $`\stackrel{~}{}^k`$ have only a finite number ($`M`$) of states, we can use Lemma 1, to obtain $`lim_k\mathrm{}\chi (\stackrel{~}{}^k)=\chi ()`$. This shows that for any $`\delta >0`$, there exists a $`k`$ such that $`\chi ()\delta \chi (\stackrel{~}{}^k)`$. With the above inequalities we can therefore conclude that $`\chi ()\delta l`$ holds for arbitrary small $`\delta >0`$, and hence that $`l\chi ()`$. $``$$``$
The following four corollaries are straightforward with the above theorem.
###### Corollary 1
For every length $`n`$, there is an incompressible classical string of length $`n`$.
Proof: Apply Theorem 6 to the set of classical strings of $`n`$ bits: $`\rho _x=|xx|`$ for all $`x\{0,1\}^n`$. All $`\rho _x`$ are pure states with zero Von Neumann entropy, hence the lower bound on $`l`$ reads $`lS(\rho )`$. The average state $`\rho =2^n_x|xx|`$ is the total mixture $`2^nI`$ with entropy $`S(\rho )=n`$, hence indeed $`ln`$. $``$$``$
###### Corollary 2
For any set of orthogonal pure states $`|\varphi _1,`$ $`\mathrm{},`$ $`|\varphi _M`$, the smallest $`l`$ such that for all $`i`$, $`\mathrm{𝑄𝐶}(\varphi _i)l`$ is at least $`\mathrm{log}M`$. (Stated differently, there is an $`i`$ such that $`\mathrm{𝑄𝐶}(\varphi _i)\mathrm{log}M`$.)
Proof: All the pure states have zero entropy $`S(\varphi _i)=0`$, hence by Theorem 6: $`lS(\rho )`$. Because all $`\varphi _i`$s are mutually orthogonal, this Von Neumann entropy $`S(\rho )`$ of the average state $`\rho =\frac{1}{M}_i|\varphi _i\varphi _i|`$ equals $`\mathrm{log}M`$. $``$$``$
###### Corollary 3
For every length $`n`$, at least $`2^n2^{nc}+1`$ qubit strings of length $`n`$ have complexity at least $`nc`$.
###### Corollary 4
For any set of pure states $`|\varphi _1,\mathrm{},|\varphi _M`$, the smallest $`l`$ such that for all $`i`$, $`\mathrm{𝑄𝐶}(\varphi _i)l`$ is at least $`S(\rho )`$, where $`\rho =\frac{1}{M}_i|\varphi _i\varphi _i|`$.
### 6.3 The complexity of copies
A case where quantum Kolmogorov complexity behaves differently from classical Kolmogorov complexity is that, in general, the relation $`C(x^m)C(x)+O(\mathrm{log}m)`$ does not hold, as we show below. We give an upper and a lower bound for the Kolmogorov complexity of $`X^m`$.
###### Theorem 7
$`\mathrm{𝑄𝐶}(X^m)\mathrm{log}\left(\genfrac{}{}{0pt}{}{m+2^{\mathrm{𝑄𝐶}(X)}1}{2^{\mathrm{𝑄𝐶}(X)}1}\right)+O(\mathrm{log}m)+O(\mathrm{log}\mathrm{𝑄𝐶}(X)).`$
Proof: First we sketch the proof, omitting the effect of the approximation. Consider any qubit string $`X`$ whose minimal-length program is $`P_X`$. To produce $`m`$ copies of $`X`$, it suffices to produce $`m`$ copies of $`P_X`$ and make $`m`$ runs of $`P_X`$.
Let $`l`$ be the length of $`P_X`$; we call $``$ the $`2^l`$-dimensional Hilbert space. Consider $`^m=\mathrm{}`$, the $`m`$-fold tensor product of $``$. The symmetric subspace $`\mathrm{Sym}(^m)`$ is $`d`$-dimensional, where $`d=\left(\genfrac{}{}{0pt}{}{m+2^l1}{2^l1}\right)`$. The state $`P_X^m`$ sits in this symmetric subspace, and can therefore be encoded exactly using $`\mathrm{log}d+O(\mathrm{log}m)+O(\mathrm{log}l)`$ qubits, where the $`O(\mathrm{log}m)`$ and $`O(\mathrm{log}l)`$ terms are used to describe the rotation onto $`\mathrm{Sym}(^m)`$. Hence, the quantum Kolmogorov complexity of $`X^m`$ is bounded from above by $`\mathrm{log}d+O(\mathrm{log}m)+O(\mathrm{log}l)`$ qubits.
For the full proof, we will need to take into account the effect of the imperfect fidelities of the different computations.
To achieve a fidelity of $`1\frac{1}{k}`$, we will compute $`m`$ copies of the minimal program $`P_X`$ to a fidelity of $`1\frac{1}{4km}`$. On each copy, we simulate the program with fidelity of $`1\frac{1}{4km}`$, and thus obtain the strings $`\stackrel{~}{X}_i`$ ($`1im`$), each of which has (according to Lemma 2) fidelity $`1\frac{1}{km}`$ with the target string $`X`$. By Lemma 3 we get a total fidelity of at least $`1\frac{1}{k}`$.
We now proceed to the details of the proof. First we introduce some notation.
Assume that for some QTM $`M`$, $`QC_M(X)\mathrm{}(P_X)=l`$, where $`P_X`$ $`M`$-computes $`X`$ (with fidelity $`1\frac{1}{k}`$ for any $`k`$.)
Let $`R`$ be the rotation that takes qubit strings $`X^m\mathrm{Sym}(^m)`$ to qubit strings of length $`\mathrm{log}(\mathrm{dim}(\mathrm{Sym}(^m)))`$. More precisely, $`R`$ is the rotation that takes the $`i`$th basis state of $`\mathrm{Sym}(^m)`$ to the $`i`$th classical basis state of the Hilbert space of dimension $`2^{\mathrm{log}(\mathrm{dim}(\mathrm{Sym}(^m)))}`$.
For any fidelity parameter $`\delta `$, $`R^1`$ can be computed efficiently and to arbitrary precision. By that we mean that for any $`\delta `$, there is a transformation $`R_\delta ^1`$ for which the following holds: Let $`Z=R(X^m)`$ for some $`X`$. If $`\stackrel{~}{X^m}=R_\delta ^1(Z)`$, then for each $`i`$, the mixed state $`\stackrel{~}{X}_i`$ obtained from $`X`$ by tracing out all components that do not correspond to the $`i`$th copy of $`X`$, is such that $`\mathrm{Fidelity}(X,\stackrel{~}{X}_i)1\delta `$.
We now define the program that witnesses the upper bound on $`\mathrm{𝑄𝐶}(X^m)`$ claimed in the theorem.
Let $`M^{}`$ be the quantum Turing machine that does the following on input $`(Z,l,m,1^k)`$.
1. Computes $`Z^{}=R_{1/4km}^1(Z)`$. (When $`Z`$ is an $`m`$-proper input, which we specify below, then $`Z^{}Y^m`$ for some $`Y`$.)
2. On each ‘copy’ $`\stackrel{~}{Y}_i`$ of $`Y`$, runs the QTM $`M(\stackrel{~}{Y}_i,1^{4km})`$. (That is, $`\stackrel{~}{Y}_i`$ is the result of tracing out all but the positions of $`Z^{}`$ that correspond to the $`i`$th block of $`l`$ qubits.)
The input $`Z`$ is an ‘$`m`$-proper input’ if for some $`Y`$, $`Z=R(Y^m)`$. (Note that $`Z`$ is exactly $`R(Y^m)`$, not an approximation up to some fidelity.)
If we run the above QTM $`M^{}`$ on input $`(R(P_X^m),l,m,1^k)`$ then the output of this $`M^{}`$ is $`M^{}(R(P_x^m),l,m,1^k)=\stackrel{~}{X^m}=\stackrel{~}{X}_1\mathrm{}\stackrel{~}{X}_m`$. (Recall that $`l`$ is the length of $`P_X`$.)
It remains to show the following claims.
###### Claim 1
$`\mathrm{Fidelity}(\stackrel{~}{X^m},X^m)1\frac{1}{k}`$.
###### Claim 2
The length of the program above for $`M^{}`$ is $`\mathrm{log}d_{l,m}+O(\mathrm{log}l)+O(\mathrm{log}m)`$, where $`d_{l,m}=\left(\genfrac{}{}{0pt}{}{m+2^l1}{2^l1}\right)`$.
Claim 2 follows immediately from the fact that the total length of the inputs $`R(P_X^m),l,m`$ is $`\mathrm{log}d+O(\mathrm{log}l)+O(\mathrm{log}m)`$.
We prove Claim 1. Since we chose a precision $`\delta =\frac{1}{4km}`$ in step 1, $`i`$, $`\mathrm{Fidelity}(P_X,\stackrel{~}{Y}_i)1\frac{1}{4km}`$. Furthermore, since the computation at step 2 introduces at most an error of $`\frac{1}{4km}`$, $`i`$, $`\mathrm{Fidelity}(X,\stackrel{~}{X}_i)1\frac{1}{km}`$ (by Lemma 2.) Therefore by Lemma 3, $`\mathrm{Fidelity}(\stackrel{~}{X^m},X^m)(1\frac{1}{km})^m1\frac{1}{k}`$. This completes the proof of Claim 1.
Claim 1 and Claim 2 together give us that $`QC_M^{}(X^m)\mathrm{log}d_{l,m}+O(\mathrm{log}l)+O(\mathrm{log}m)\mathrm{log}d_{n,m}+O(\mathrm{log}n)+O(\mathrm{log}m)`$, where $`n`$ is the length of $`X`$ and an upper bound on its complexity. By invariance, we can conclude that $`QC(X^m)\mathrm{log}d_{n,m}+O(\mathrm{log}n)+O(\mathrm{log}m)+O(1)`$, which proves the theorem. $``$$``$
This upper bound is also very close to being tight for some $`X`$, as we show in the next theorem.
###### Theorem 8
For every $`m`$ and $`n`$, there is an $`n`$-qubit state $`X`$ such that $`\mathrm{𝑄𝐶}(X^m)\mathrm{log}\left(\genfrac{}{}{0pt}{}{m+2^n1}{2^n1}\right)`$.
Proof: Fix $`m`$ and $`n`$ and let $``$ be the $`2^n`$-dimensional Hilbert space. Consider the (continuous) ensemble of all $`m`$-fold tensor product states $`X^m`$: $`=\{(X^m,\mu )\}`$, where $`\mu ^1=_X𝑑X`$ is the appropriate normalization factor. The corresponding average state is calculated by the integral $`\rho =\mu _XX^m𝑑X`$. This mixture is the totally mixed state in the symmetric subspace $`\mathrm{Sym}(^m)`$ (see Section 3 in ), and hence has entropy $`S(\rho )=\mathrm{log}\left(\genfrac{}{}{0pt}{}{m+2^n1}{2^n1}\right)`$. Because all $`X^m`$ are pure states, we can use Corollary 4 to prove the existence of a $`X`$ for which $`QC(X^m)\mathrm{log}\left(\genfrac{}{}{0pt}{}{m+2^n1}{2^n1}\right)`$. $``$$``$
### 6.4 Subadditivity
Consider the following subadditivity property of classical Kolmogorov complexity.
###### Proposition 6
For any $`x`$ and $`y`$, $`C(x,y)C(x)+C(y\text{ }x)+O(1)`$.
In the classical case, we can produce $`x`$, and then produce $`y`$ from $`x`$, and print out the combination of $`x`$ and $`y`$. In the quantum case, producing $`Y`$ from $`X`$ may destroy $`X`$. In particular, with $`X=Y`$, the immediate quantum analogue of Proposition 6 would contradict Theorem 8 (for $`m=2`$).
A natural quantum extension of this result is as follows.
###### Proposition 7
For any $`X,Y`$, $`\mathrm{𝑄𝐶}(X,Y)\mathrm{𝑄𝐶}(X,X)+\mathrm{𝑄𝐶}(Y\text{ }X)+O(1)`$.
## 7 Quantum Information Theory
In this section we establish a relationship between quantum compression theory and the bounded-fidelity version of quantum Kolmogorov complexity.
One would like to give a direct analogue of Proposition 3. However, we believe that such a statement does not hold for quantum Kolmogorov complexity. The argument can be summarized as follows. In the classical case, given a string $`x`$, we can define a source $`A`$ such that $`x`$ is in the so-called ‘typical subspace’ of $`A`$. This allows us to give a short, exact description of $`x`$.
In the quantum case, we may also define a quantum source likely to have emitted a given qubit string $`X`$ (in an appropriate tensor space). However, we do not get that $`X`$ is *in* the typical subspace of this source, only that it is *close* to the typical subspace. How close it can be guaranteed to be depends on the length of $`X`$. Therefore, for a fixed string length $`n`$, we may not be able to get an encoding of arbitrary high fidelity.
We now prove a slightly weaker statement, for bounded-fidelity complexity.
###### Theorem 9
Let $`U`$ be the universal quantum Turing machine from . Then for any $`ϵ,\delta `$ there is an $`n`$ such that for any $`d`$-dimensional $``$, and any qubit string $`X=|\varphi _1\mathrm{}|\varphi _n^n,`$
$`\mathrm{𝑄𝐶}_U^ϵ(X)`$ $``$ $`n(S(\rho )+\delta +\frac{1}{n}(d^2\mathrm{log}(n+1))),`$
where $`\rho =\frac{1}{n}_i|\varphi _i\varphi _i|`$.
Proof: Fix $`ϵ,\delta `$. Apply Theorem 2 with $`ϵ^{}=\frac{ϵ}{4},\delta ^{}=\delta ,`$ and let $`n=n(ϵ^{},\delta ^{})`$ be the value from the theorem. Let $`|\varphi _1\mathrm{}|\varphi _n^n`$ be the string for whose quantum Kolmogorov complexity we want to give an upper bound. By Theorem 2, item 1, we get that the length of the encoding is what was given in the statement of the theorem. By simulating the decoding algorithm to a precision of $`\frac{ϵ}{4}`$, together with Theorem 2, item 2, and Lemma 2, we have that the fidelity of the encoding is at least $`1ϵ`$. That completes the proof. $``$$``$
## 8 Extensions and Future Work
We have argued that the $`\mathrm{𝑄𝐶}`$ of Definition 8 is a robust notion of Kolmogorov complexity for the quantum setting. Nevertheless, it would be interesting to see if an invariance theorem can be shown for the ideal quantum Kolmogorov complexity of Definition 6.
The number of applications of classical Kolmogorov complexity is countless, and it is our hope that this definition will lead to a similar wide variety of applications in quantum complexity theory.
## 9 Acknowledgements
We would like to thank several people for interesting discussions on this work: Paul Vitányi, Harry Buhrman, Richard Cleve, David Deutsch, Ronald de Wolf, John Watrous, Miklos Santha, Frédéric Magniez, and Jérémy Barbay.
This work has been supported by Wolfson College Oxford, Hewlett-Packard, European TMR Research Network ERP-4061PL95-1412, the Institute for Logic, Language and Computation in Amsterdam, an NSERC postdoctorate fellowship, and the EU fifth framework project QAIP IST-1999-11234.
|
warning/0005/astro-ph0005012.html
|
ar5iv
|
text
|
# FUSE Observations of O VI in High Velocity Clouds
## 1 Introduction
Since their discovery nearly 40 years ago, interstellar high velocity clouds (HVCs) have remained enigmatic. Their origins, fundamental physical properties, and distances are unknown in most cases (see review by Wakker & van Woerden 1997). Although HVCs are often described in terms of their H I 21 cm emission and their peculiar velocities<sup>1</sup><sup>1</sup>1Usually, clouds moving in excess of 100 km s<sup>-1</sup> with respect to the Local Standard of Rest fall into the category of high velocity cloud. with respect to the motions of the general interstellar medium, the traditional view of HVCs as completely neutral entities has given way to one that recognizes the importance of an ionized component to many of the clouds. H$`\alpha `$ emission has been detected toward several large HVC complexes (Tufte, Reynolds, & Haffner 1998; Bland-Hawthorn et al. 1998; Weiner & Williams 1996), and recent absorption line observations with the Hubble Space Telescope have shown that some HVCs are almost fully ionized (Sembach et al. 1995, 1999).
Understanding the ionization of HVCs is a key step toward a more complete description of the high velocity gas. The large peculiar velocities of HVCs make them excellent candidates for studying hot gas within the clouds and their surrounding environments. A prime diagnostic of gas at temperatures T $`10^510^6`$ K is the O VI $`\lambda \lambda 1031.93,1037.62`$ doublet, which can be observed in the spectra of distant QSOs and AGNs with the newly commissioned the Far Ultraviolet Spectroscopic Explorer (FUSE).
In this Letter we report the first detections of O VI absorption in several HVCs located in different regions of the sky. The directions considered are listed in Table 1, where we summarize sight line characteristics and basic measurements for the HVCs. The reader may find information about the overall extent and distribution of O VI in the Galactic halo in a companion paper (Savage et al. 2000). A description of the FUSE mission can be found in Moos et al. (2000).
## 2 Observations and Data Processing
The FUSE data for this investigation were obtained during the commissioning stage of the mission in late 1999. Each observation was obtained with the source centered in the $`30\mathrm{}\times 30\mathrm{}`$ aperture of the LiF1 spectrograph channel. Exposure times ranged from 13 ksec (Ton S210) to 55 ksec (Mrk 509) (see Savage et al. 2000). The time-tagged photon lists were processed through the standard FUSE calibration pipeline available at the Johns Hopkins University as of November 1999. The lists were screened for valid data with constraints imposed for earth limb angle avoidance and passage through the South Atlantic Anomaly. Corrections for detector backgrounds, Doppler shifts caused by spacecraft orbital motions, and geometrical distortions were applied (see Sahnow et al. 2000). No corrections were made for optical astigmatism aberrations or small spectral shifts introduced by thermal effects since the data were obtained prior to completion of in-orbit focusing activities. The primary effect of omitting these processing steps is to degrade the spectral resolution slightly.
The processed data have a nominal spectral resolution of 25 km s<sup>-1</sup> (FWHM), with a relative wavelength dispersion solution accuracy of $`6`$ km s<sup>-1</sup> ($`1\sigma `$). The zero point of the wavelength scale for each observation was determined by registering the H<sub>2</sub> (6–0) P(3) line at 1031.19 Å to the velocity of the peak H I 21 cm emission for the sight line. In cases where no H<sub>2</sub> absorption is detected, we compared the velocities of the strong Si II $`\lambda 1020.70`$ and O I $`\lambda 1039.23`$ lines to those observed for lines of the same species at longer wavelengths.
Fully reduced FUSE LiF1A data covering the 1020–1045 Å spectral region are shown for PKS 2155-304 and Mrk 509 in Figure 1. Note the presence of negative high velocity O VI 1031.93 Å absorption along both sight lines. O VI $`\lambda `$1037.62 absorption is present in both spectra; however, the HVC components are blended with C II $`\lambda `$1037.02. There is also weak, narrow C II $`\lambda `$1036.34 absorption toward both objects near the velocities of the O VI HVCs. Other figures showing the O VI HVCs considered herein can be found in Savage et al. (2000), Oegerle et al. (2000), and Murphy et al. (2000a).
## 3 High Velocity O VI Measurements
Continuum normalized intensity profiles for the O VI HVC sight lines are presented in Figure 2. Equivalent widths and column densities were computed by direct integration of the O VI $`\lambda 1031.93`$ intensity and apparent column density profiles (Sembach & Savage 1992). N<sub>a</sub>(O VI) \[cm<sup>-2</sup>\] = 2.748$`\times `$10<sup>12</sup> $`\tau _a`$(v) dv, where $`\tau _a`$(v) is the measured optical depth of the 1031.93 Å line at velocity v (in km s<sup>-1</sup>) (Savage & Sembach 1991). We have assumed an oscillator strength f = 0.133 (Morton 1991). N<sub>a</sub>(O VI) is a reasonable approximation to N(O VI) since the HVC lines are both weak and broad. The integration ranges lie beyond the velocities expected for O VI absorption arising in the general environment of the thick disk and low halo ($`|`$z$`|`$ $`<`$ 3 kpc) (see Savage et al. 2000). In Table 1 we list the fraction of total O VI along each sight line in HVCs, $`f_{HVC}`$. On average, $`f_{HVC}30\%`$.
The H<sub>2</sub> (6–0) P(3) 1031.19 Å and R(4) 1032.35 Å lines occur at velocities of –214 km s<sup>-1</sup> and +123 km s<sup>-1</sup>, respectively, relative to the O VI $`\lambda `$1031.93 line. We have modeled the impact of these lines on the observed O VI profiles by examining other H<sub>2</sub> (J = 3 or 4) lines in the (3–0) to (8–0) vibrational bands covered by the FUSE LiF1A spectra. Our estimates of the strengths and shapes of the P(3) 1031.19 Å and R(4) 1032.35 Å lines are shown as heavy solid lines in Figure 2. The values of W<sub>λ</sub> and log N(O VI) in Table 1 have had the illustrated H<sub>2</sub> contributions removed.
## 4 HVC Identifications
High velocity (V$`{}_{LSR}{}^{}<100`$ km s<sup>-1</sup>) O VI is seen along 7 of the 11 extended sight lines toward AGNs and QSOs observed by FUSE. The HVCs have negative velocities, with the exception of a weak HVC detected toward Ton S180 (see Table 1). In many cases, the O VI HVCs can be related to H I HVCs located at large distances ($`>`$3 kpc) from the Galactic plane.
### 4.1 The C IV HVCs: Local Group Clouds
High velocity O VI absorption toward Mrk 509 and PKS 2155-304 occurs at velocities similar to those of some of the C IV HVCs studied by Sembach et al. (1999). The Mrk 509 C IV HVCs ($``$V$`{}_{LSR}{}^{}`$ = –287, –228 km s<sup>-1</sup>) are believed to be located in the Local Group outside the Milky Way based upon their ionization properties and the very low thermal pressures (p/k $``$ 2 cm<sup>-3</sup> K) inferred if the clouds are in photoionization equilibrium. They display strong C IV absorption, with little lower ionization absorption and no detectable H I 21 cm emission. Furthermore, they show no evidence of H$`\alpha `$ emission (Sembach, Bland-Hawthorn, & Savage 2000). This ionization pattern is characteristic of gas clouds irradiated by extragalactic background radiation.
The O VI HVC absorption toward Mrk 509 is distributed in a broad component centered on $``$V$`{}_{LSR}{}^{}230`$ km s<sup>-1</sup>, with FWHM $``$ 120 km s<sup>-1</sup>. It exhibits a steep decline in strength at velocities where the C IV is strongest (i.e., in the –287 km s<sup>-1</sup> component). Most of the O VI appears to be associated with the lower velocity C IV HVC at –228 km s<sup>-1</sup>. The ratio of C IV to O VI is $`<`$1 in this cloud and the HVCs observed toward PKS 2155-304. The observed amount of O VI is more than an order of magnitude higher than predicted by the standard photoionization model<sup>2</sup><sup>2</sup>2Their photoionization model uses an AGN/QSO spectral energy distribution with a mean intensity at the Lyman limit J<sub>0</sub> = 1$`\times `$10<sup>-23</sup> erg cm<sup>-2</sup> s<sup>-1</sup> Hz<sup>-1</sup> sr<sup>-1</sup> (Haardt & Madau 1996). described by Sembach et al. (1999), which leads us to conclude that there are multiple ionization processes in these HVCs.
### 4.2 Complex C
Mrk 876 lies behind HVC Complex C, which is located more than 3.5 kpc from the Galactic plane (van Woerden et al. 1999). The low metallicity of Complex C, \[S/H\] $``$$`0.5`$ to $`1.0`$, indicates that it may be material falling onto the Milky Way rather than material ejected from the Galactic disk (Wakker et al. 1999; Gibson et al. 2000). In the direction of Mrk 876, Complex C has two distinct H I components centered on V<sub>LSR</sub> = $`175`$ and $`132`$ km s<sup>-1</sup> (Murphy et al. 2000a). The broad O VI absorption spans these neutral components and forms a smooth absorption trough (see Figure 2).
The presence of O VI at velocities similar to those of the neutral tracers of Complex C suggests that Complex C contains a substantial amount of ionized gas. For a gas in collisonal ionization equilibrium at T = 3$`\times `$10<sup>5</sup> K with Z $``$ 0.1–0.3 Z, N(O VI) = 1.5$`\times 10^{14}`$ cm<sup>-2</sup> translates into N(H<sup>+</sup>) $`(39)\times 10^{18}`$ cm<sup>-2</sup>, or roughly 10–30% of the observed H I column density of $`2.9\times 10^{19}`$ cm<sup>-2</sup>.
### 4.3 Magellanic Stream
Three of the 7 sight lines exhibiting high velocity O VI lie in the general direction of the Magellanic Stream. Ton S210 and Ton S180 lie about 10 off the Stream near Complex MSII. No high velocity 21 cm emission from the Stream is detected toward either object; however 21 cm emission is detected toward Ton S210 at $`196`$ km s<sup>-1</sup> with a width of 26 km s<sup>-1</sup> (Murphy, Sembach, & Lockman 2000b). The H I and O VI HVCs toward Ton S210 are distinct from the nearby Stream material, which has velocities V<sub>LSR</sub> $``$ –100 km s<sup>-1</sup> in this direction. A sensitive H I map of this region indicates that the 21 cm emission is isolated (M. Putman, private communication) and resembles the compact H I HVCs believed to be located in the Local Group (Braun & Burton 1999; but see also Charlton et al. 2000 and Zwaan & Briggs 2000).
In the direction of NGC 7469, the high velocity O VI absorption is related to Magellanic Stream material seen in H I 21 cm emission near $`350`$ km s<sup>-1</sup> (Murphy et al. 2000b). Ionized Stream gas has been detected previously through its H$`\alpha `$ emission (Weiner & Williams 1996). Most of that emission could be explained by photoionization by starlight (Bland-Hawthorn & Maloney 1999). However, the presence of O VI in the Stream toward NGC 7469 indicates that hot gas must be present since the low gas densities required to produce the observed amounts of O VI solely by photoionization requires very large path lengths ($`l100200`$ kpc $`>`$ d<sub>MS</sub>).
### 4.4 Outer Galaxy
Most of the O VI absorption toward H 1821+643 can be attributed to the thick disk and outer Galactic warp (Oegerle et al. 2000; Savage et al. 2000). The broad, shallow absorption between $`300`$ and $`175`$ km s<sup>-1</sup> is probably tracing gas in the most distant portions of the outer Galaxy. The limiting velocity of co-rotating interstellar gas in this direction is roughly –200 km s<sup>-1</sup>. Savage, Sembach, & Lu (1995) noted a very weak C IV feature with N(C IV) = (1.2$`\pm `$0.3)$`\times 10^{13}`$ cm<sup>-2</sup> near –213 km s<sup>-1</sup>. The velocity of the LSR with respect to the velocity centroid of the Local Group (l $``$ 105, b $``$ –8) is $`300`$ km s<sup>-1</sup> (see Mihalas & Binney 1982), so the projection of this relative velocity onto the H 1821+643 sight line is $``$ –220 km s<sup>-1</sup>, similar to that of the high velocity O VI.
## 5 Discussion
O VI is difficult to produce by photoionization (114 eV photons are required). The gas density must be so low that path lengths exceeding the distance to the Magellanic Stream are necessary to account for the observed quantities of O VI in the HVCs. See Sembach et al. (1999) for a discussion of the photoionization models.
Many of the O VI HVCs contain cooler material (e.g., Mrk 876, NGC 7469), or are in close proximity to H I HVCs at similar velocities (e.g., Mrk 509, PKS 2155-304, Ton S210). To produce O VI by shocks requires a shock speed of $``$ 170 km s<sup>-1</sup> (Hartigan, Raymond, & Hartmann 1987). The observed velocity gradients of the neutral gas in Complex C (van Woerden, Schwarz, & Hulsbosch 1985) and the Magellanic Stream (Mathewson & Ford 1984; Putman & Gibson 1999) are much smaller than this, so it seems unlikely that cloud-cloud collisions are responsible for the O VI.
One possibility for the production of the O VI observed in the HVCs is that the clouds are moving through a pervasive, hot (T $``$ 10<sup>6</sup> K), low density (n$`{}_{H}{}^{}10^410^5`$ cm<sup>-3</sup>) Galactic halo or Local Group medium. The existence of such gas has been considered in various contexts (see Fabian 1991; Weiner & Williams 1996; Blitz & Robishaw 2000; Murali 2000). The O VI HVCs have velocities comparable to the adiabatic sound speed ($`150`$ km s<sup>-1</sup>) for a hot, low density medium. Thus, strong shocks are probably not responsible for the O VI production. Rather, in such a scenario, the O VI would occur in the conductive interfaces or turbulently mixed regions of ionized gas between the hot medium and the Magellanic Stream or other cooler gas detected in 21 cm emission or ultraviolet absorption.
Alternatively, some of the O VI may be produced within cooling regions of hot gas structures associated with the assembly of the Milky Way. Much of the baryonic content of the low redshift universe is expected to be at temperatures of $`10^510^7`$ K in the vicinity of galaxies and groups of galaxies (Cen & Ostriker 1999). As the hot gas flows onto galaxies, portions of it should cool as the density increases. Complex C may represent a relatively advanced stage of such an accretion, while the HVCs toward Mrk 509 and PKS 2155-304 would represent an earlier evolutionary stage. FUSE data for a large number of AGN/QSO sight lines should help to distinguish between these various possibilities for the production of the O VI HVCs.
This work is based on data obtained for the Guaranteed Time Team by the NASA-CNES-CSA FUSE mission operated by the Johns Hopkins University. Financial support has been provided by NASA contract NAS5-32985.
|
warning/0005/astro-ph0005305.html
|
ar5iv
|
text
|
# The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope
## 1 Introduction
Individual stars are the visible building blocks of galaxies and direct tracers of the galaxy formation process. Massive elliptical galaxies are believed to contain the majority of the oldest stars in the Universe (Tinsley and Gunn, 1976). Two main scenarios have been proposed for the formation of these galaxies. In the traditional scenario of single “monolithic” collapse, the most massive early-type (E and S0) galaxies form at early times ($`z2`$) and evolve with little or no star formation thereafter (Tinsley, 1980; Bruzual and Kron, 1980; Koo, 1981; Shanks et al., 1984; King and Ellis, 1985; Yoshii and Takahara, 1988; Guiderdoni and Rocca-Volmerange, 1990). In hierarchical models of galaxy formation, massive galaxies assemble later ($`z2`$) from mergers of smaller subunits (e.g., White and Rees, 1978; White and Frenk, 1991; Lacey et al., 1993; Kauffmann, White and Guiderdoni, 1993; Cole et al., 1994; Somerville, 1997). At least some elliptical galaxies show signatures of intermediate-age stars in addition to an old stellar population (e.g., Worthey, Faber and Gonzalez, 1992). Dynamical studies also suggest that the pressure support of stellar populations in elliptical galaxies could result from mergers (Toomre and Toomre, 1972; Hibbard et al., 1994). Moreover, detailed observations of some elliptical galaxies reveal morphological and kinematic signs of a past merging event. These range from the observations of proto-elliptical merger remnants like NGC 7252 to evidence of counter-rotating or otherwise decoupled cores for nearby ellipticals (star-star or star-gas, Bertola 1997). The dominance of old stars and evidence of merging in elliptical galaxies can be understood simultaneously if the youngest stars contribute only a small fraction of the observed integrated light. In fact, Silva and Bothun (1998) show that nearby elliptical galaxies with morphological signatures of recent merger activity have near-IR colors similar to those of galaxies not showing signs of mergers. They conclude that intermediate-age ($`13`$ Gyr) stars contribute at most 10%$``$15% of the total stellar mass in galaxies with recent merger activity in their sample.
Resolving individual stars in elliptical galaxies has become feasible only recently with the advent of the Hubble Space Telescope (HST), since no suitable examples are near enough to be observed from the ground (this will change with the application of adaptive optics systems to large ground-based telescopes). Such detailed information about the stellar content of elliptical galaxies can help us reconstruct their star formation history, and hence, constrain their process of formation. With HST, it is possible to resolve the population of the massive elliptical galaxy NGC 5128 (Centaurus A) due to its proximity to us. NGC 5128 is part of a group of 25 galaxies composed mainly of dwarf galaxies extending over about 25 degree on the sky (see review by Israel 1998). Its distance estimate relies upon different measurement methods. A distance modulus of $`(mM)_0=27.53\pm 0.25`$ (Tonry and Schechter, 1990) is derived from the globular cluster luminosity function (LF) of Harris (1986). The planetary nebula luminosity function yields $`(mM)_0=27.73\pm 0.14`$ (Hui et al., 1993). The distance modulus of $`(mM)_0=27.48\pm 0.06`$ measured by Tonry and Schechter (1990) from surface brightness fluctuations is revised to $`(mM)_0=27.71\pm 0.10`$ by Israel (1998) using the results of Tonry (1991) (the more recent results of Tonry et al. (1997) yield a higher revised value of $`(mM)_0=28.18\pm 0.07`$). An estimate of $`(mM)_0=27.72\pm 0.20`$ is derived from the magnitude of the tip of the red giant branch (TRGB) for stars in the halo observed with HST WFPC2 by Soria et al. (1996; hereafter SMW96). A more recent measurement by Harris et al. (1999; hereafter HHP99) comes from the magnitude of the TRGB for stars in another halo field observed with HST WFPC2 (these data reach $`3`$ magnitudes down the RGB). They find a distance modulus of $`(mM)_0=27.98\pm 0.15`$, or $`D=3.9\pm 0.3`$ Mpc. The weighted average of distance moduli of $`(mM)_0=27.75\pm 0.06`$ is adopted throughout this paper and corresponds to a distance of $`D=3.5\pm 0.1`$ Mpc. At this distance, 1 arcmin corresponds to 1018 pc.
NGC 5128 is a clear case of an elliptical galaxy showing signs of past merger activity. It is one of the largest known radio galaxies (500 $`\times `$ 250 kpc wide) and a massive disk of gas, young stars, and dust is embedded in its center. Within a radius of about 18$``$ from the nucleus, the galaxy shows optical shell structures made up of old disk stars and associated H i shells (Schiminovich et al., 1994). This suggests that NGC 5128 might have experienced more than just one merger in its past (Weil and Hernquist, 1996). The H i shells detected in the outer part of the galaxy seem to show signs of recent star formation (Graham, 1998). The halo globular clusters in NGC 5128 (region in radius $`R<`$ 24′) have a mean metallicity of \[Fe/H\]=$`0.8\pm 0.2`$ (0.5 dex higher than their Milky Way counterparts; Harris et al. 1992) and show a bimodal distribution in color, with peaks at \[Fe/H\]$`1.1`$ and $`0.3`$ ($`R>`$ 4′; Harris et al. 1992), an effect commonly associated with a merging event. Jablonka et al. (1996) find no object with a metallicity higher than solar in a similar sample. The globular clusters in the inner 3 kpc are more metal rich with $`0.6`$\[Fe/H\]$`+0.1`$ than in the outer regions and show signs of a metallicity gradient (Jablonka et al., 1996; Minniti et al., 1996; Alonso and Minniti, 1997). Recent HST WFPC2 observations yield a mean value of \[Fe/H\]$`>0.9`$ (SMW96) and \[Fe/H\]=$`0.41`$ (HHP99) for red giant branch stars in the halo of NGC 5128.
The age of the current burst of star formation in the disk of NGC 5128 is typically a few times $`10^7`$ years (van den Bergh, 1976; Dufour et al., 1979). In the halo, the presence of $`200`$ stars found to be brighter than the TRGB prompted SMW96 to suggest the presence of an intermediate-age population of $`5`$ Gyr, making up at most 10% in number of the total halo population. Dynamical estimates based on the model of a merger of a late-type galaxy of mass a few times $`10^{10}M_{}`$ with NGC 5128 suggest a more recent merging timescale of a few hundred million years (Tubbs, 1980; Malin, Quinn and Graham, 1983). The total mass of the galaxy estimated from velocity dispersion measurements of the planetary nebula system is $`M=4\pm 1\times 10^{11}M_{}`$, with half of it estimated to be due to dark matter (Mathieu, Dejonghe and Hui, 1996).
We present in this paper the first IR color-magnitude diagram (CMD) for the halo of a giant elliptical galaxy. Our NICMOS observations of the halo of NGC 5128 probe a range of $`4`$ magnitudes down the luminosity function to our 50% completeness limit. Section 2 presents the details of the NICMOS observations we obtained in August 1998. The data analysis using the stellar photometry package DAOPHOT is described in Section 3. The importance of doing artificial-stars experiments is emphasized in Section 4, where the completeness functions are presented. The luminosity functions and discontinuities are derived in Section 5. Section 6 presents the first IR CMD of the halo of NGC 5128. This section is divided into four sub-sections discussing the total uncertainties, deriving an estimate of the metallicity spread of the halo stars, and discussing the nature of the bright stars above the TRGB of an old population. A summary of our results and conclusions appears in Section 7. The more technical details of the magnitude system transformations are given in Appendix A.
## 2 Observations
Images of a field in the halo of NGC 5128 were taken on 1998 August 31 with the NIC1 and NIC2 camera of the Near-Infrared and Multiobject Spectrometer (NICMOS) (Thompson et al., 1998) on board HST. The spacecraft pointing was chosen to image the existing WFPC2 CHIP-3 field of SMW96, at a distance of 08$``$ 50$``$ $``$ from the nucleus, with the NICMOS cameras (see Figure The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope). We chose our positions of NIC1 and NIC2 so they would lie inside the WFPC2 CHIP-3 field of view of SMW96 for any unrestricted value of position angle. Based on our adopted distance of 3.5 Mpc, this corresponds to a distance of $`9`$ kpc south of the center of the galaxy. The geometrical center of NIC2 was at position $`\alpha =+13`$<sup>h</sup>25<sup>m</sup>24$`^\mathrm{s}.`$323, $`\delta =43`$09$``$ 58$``$ $``$ .53 (J2000) and the NIC1 observations were obtained in parallel. The NIC2 detector has a field of view of 19$``$ $``$ .2 $`\times `$ 19$``$ $``$ .2 with 256 pixels on a side and 0$``$ $``$ .075 per pixel. A total of 8192 s (3 orbits) of integration time were obtained, 5376 s in the F160W filter and 2816 s in the F110W filter. The observations were taken in a four dither positions mode with offsets of (0,0), (0,15.5), (15.5,15.5), and (15.5,0) pixels, corresponding to a maximum shift of 1$``$ $``$ .17. This enables the replacement of a bad pixel in one frame with the average of good pixels from dithered frames. In addition, short exposures were taken at the beginning of each orbit in the MULTIACCUM mode to reduce the cosmic rays persistence effects (Najita, Dickinson and Holfeltz, 1998). All the observations were taken in the MULTIACCUM mode with the SPARS64 sequence.
A basic image reduction was done by STScI using the standard NICMOS pipeline procedure called CALNICA which performs bias subtraction, dark-count correction, and flat-fielding (MacKenty et al., 1997). Our own subsequent data reduction with IRAF<sup>1</sup><sup>1</sup>1IRAF is distributed by the National Optical Astronomy Observatories, which are operated by the Association of Universities for Research in Astronomy, Inc., under cooperative agreement with the National Science Foundation. consisted of masking bad pixels, correcting for a constant level offset between the four quadrants of the NIC2 camera, correcting for background variations by subtracting a median sky image and adding back a constant sky level, and averaging the images. Each resulting NIC2 mosaic covers a field of view of 20$``$ $``$ .4 $`\times `$ 20$``$ $``$ .4. The final combined NIC2-F110W and -F160W images of our field are shown in Figure The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope.
## 3 Stellar Photometry
Photometry was obtained for our NICMOS images by using the automated star-detection algorithm DAOPHOT (Stetson, 1987, 1992). The data were initially processed using the subroutine DAOFIND with a conservative detection threshold of $`5\sigma `$ above the local background level. The PSFs were derived by using many well-exposed, isolated stars sampling the frame uniformly. The full width at half-maximum (FWHM) of the F160W and F110W PSF profiles were measured to be equal to 1.7 and 1.2 pixel, respectively. The stellar photometry was accomplished by processing our images using the subroutine ALLSTAR once, and then another time on the first residual image. The fitting was done on pixels within a fitting radius of the centroid of a star equal to the FWHM of the PSF. The 3-parameter least-square fit to the star determines the position of the center (x,y), the brightness and its standard deviation. We measured magnitudes for 971 stars in the F110W image and 1321 stars in the F160W image.
The residual images generated by DAOPHOT show that stars near the edges of the mosaics suffer from bad PSF fits (their distance from the edge is of the order of their FWHM). Our final sample contains stars selected so that they are not too close to the edge of the mosaic, i.e., no less than 3 pixels away from the edge. We also applied a stellar fit $`\chi ^2`$ cut to our sample of stars. The cut was set to a $`\chi ^2`$ of 2.6, given that the probability of exceeding that value of $`\chi ^2`$ is 5%. The $`\chi ^2`$ test was done only on the F160W image for which the PSF is well sampled. In the case of the under-sampled F110W image for which the PSF is almost a Dirac delta function, $`\chi ^2`$ provides a poor discriminant because it is possible to fit every pixel that contains a signal. The \[F160W\] $`\chi ^2`$ cut eliminated 215 stars of which 3 were on the right hand side of the red envelope of the normal giants in our CMD, leaving a single extremely red star at \[F110W\]$``$\[F160W\]=2.5$`\pm 0.1`$ (see Section 6 and Figure The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope).
The instrumental magnitudes were transformed into the Vega-based system. Since the photometric keywords needed to transform countrates into magnitudes refer to a nominal infinite aperture, given in the HST Data Handbook as 1.15 times the flux in a 0.5 arcsec radius aperture, the measured countrates have to be corrected accordingly. This is done by first correcting the measured PSF-fitting photometry to the 0$``$ $``$ .13 radius aperture photometry (a 0$``$ $``$ .13 or 1.7 pixel radius aperture, of the order of the FWHM of the F160W PSF, was used to calibrate the PSF flux in both filters). Secondly, in order to estimate the aperture correction, we used both our observed PSF and the artificial PSF generated by the Space Telescope package TINYTIM (Krist, 1993). We found for our choice of aperture radius of 0$``$ $``$ .13 that the fraction of PSF-fitting photometry to 0$``$ $``$ .5 radius aperture was 60% for the F110W filter and 56% for the F160W filter. Putting in the final 1.15 correction factor, we calculated a correction to brighter magnitudes of 0.71 and 0.77, respectively. The photometric calibration of PHOTFNU = $`2.190\times 10^6`$ Jy sec/DN and ZP(Vega) = 1083 Jy (HST Data Handbook; updated values from Rieke 1999) produced the NIC2-F160W zeropoint in the Vega-based system of 21.74. After including the aperture correction, the Vega-based magnitudes for stars in our F160W image were therefore computed by using mag (Vega-based) = $`20.972.5`$log(counts(e-/s)) (gain is 5.0 e-/DN). The F110W zeropoint was calculated to be 21.64, given PHOTFNU = $`2.031\times 10^6`$ Jy sec/DN, ZP(Vega) = 1775 Jy and the aperture correction given above. The systematic errors in magnitude and color measurements are estimated to be less than 0.05 mag (Calzetti et al. 1999; see Colina and Rieke 1997 for more detail).
Colors were calculated for the stars in the NICMOS images in the following way. Each color image was analyzed separately. The identification of a star in both images was done by requesting that (1) the two positions agreed to a specific tolerance (matching) radius, and (2) the star assigned was the closest uniquely assigned star in either color image. A total of 666 stars (out of the sample of 1087 stars in F160W and 941 stars in F110W) were identified in both filters using a maximum matching radius $`r_m=1`$ pixel. This maximum matching radius was chosen because the number of stars matched show a clear cut at that radius, as seen in Figure The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope where we allowed the matching radius to be as large as $`r_m=5`$ pixels.
To justify our choice of the maximum matching radius and illustrate that the peak of the distribution lies between $`0.10.2`$ pixel, we modeled the number of stars with a match between the F110W and F160W images. The model is made up of a Gaussian component that reflects the uncertainty in the centroid of each image, due to undersampling in the F110W image, photon statistics, and the effects of crowding. To this we add a constant background probability of making a false match with an un-associated star. Hence the number of stars matching, $`N`$, between radii $`r`$ and $`r+dr`$ should follow:
$$N(r)dr=N_{true}\left(\frac{r}{\sigma ^2}\right)e^{0.5(\frac{r}{\sigma })^2}dr+\mathrm{\hspace{0.33em}2}\pi rN_{}dr,$$
where $`N_{true}`$ is the number of true matches, $`N_{}`$ is the number of un-associated stars per unit area, and $`\sigma `$ represents the two-dimensional relative error in star positions between the two frames. The best fit curve is shown over-plotted on the histogram of matching radius in Figure The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope. This implies that the net two-dimensional root-mean-square (rms) registration error between the F110W and F160W frames is 0.19 pixel, and that a natural choice of cut-off is at a radius of 1 pixel which is equivalent to $`5\sigma `$. The excess matches above the fit between radii of $`0.60.9`$ pixels is due to genuine matches of faint stars with poorly determined centroids. The false match rate implies that within our chosen cutoff of 1 pixel, 3% of our stars are matched with the wrong counterpart.
## 4 Artificial-Star Tests
In the absence of other systematic errors, the photometric errors depend only on photon statistics (from source and background) and detector noise. Simulations we ran prior to our observations showed that in order to be limited by these errors and not be affected by crowding, we needed to observe a region with a surface brightness no higher than 21.0 mag/arcsec<sup>2</sup> in the F160W filter (assuming a Baade’s window LF). This is consistent with analytic estimates of the effects of crowding (Renzini, 1998). In the simulations, we were able to recover the input luminosity function down to \[F160W\]$`23.5`$ with photometric accuracy to the 10% level. To assess the errors associated with doing photometry in our field, we simulated artificial stars in our NICMOS images.
| TABLE 1 |
| --- |
| COMPLETENESS |
| Magnitude | Completeness (%) | Completeness (%) |
| --- | --- | --- |
| | \[F110W\] | \[F160W\] |
| 20.0 | 100 | 100 |
| 20.5 | 100 | 99 |
| 21.0 | 100 | 98 |
| 21.5 | 98 | 98 |
| 22.0 | 99 | 98 |
| 22.5 | 98 | 98 |
| 23.0 | 98 | 94 |
| 23.5 | 95 | 78 |
| 24.0 | 87 | 33 |
| 24.5 | 52 | 5 |
| 25.0 | 17 | 3 |
| 25.5 | 4 | 2 |
| 26.0 | 2 | 2 |
| 26.5 | 2 | 2 |
Completeness tests were performed by adding artificial stars to each individual color image. We simulated stars with F110W and F160W magnitudes between 20.0 and 26.5, at a 0.1 mag interval. Each simulation consisted of adding 132 stars to the real image, corresponding to 10% in number of the stars recovered from the F160W image. The number of artificial stars was chosen to be large enough to compile accurate statistics on incompleteness and photometric errors, and small enough to increase the crowding negligibly. The positions of the added stars on the images were randomly chosen but identical for both the F110W and F160W images so that colors could be measured. The frames were then processed in a manner identical to the original data. The F110W and F160W completeness functions measured from these artificial-star tests are shown in Figure The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope and listed in Table 1. The 50% completeness level occurs at \[F110W\]=24.5 and \[F160W\]=23.8, respectively.
## 5 The Luminosity Functions
Our NICMOS observations probe $`4`$ magnitudes below the tip of the luminosity function down to our 50% completeness limit. The NICMOS data are less crowded by bright stars than the WFPC2 data since the IR luminosity functions tend to be steeper and so there are fewer bright stars. Figure The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope shows the luminosity functions for stars detected in the F110W and F160W image respectively, taking into account the edge and $`\chi ^2`$ cuts. Also shown in Figure The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope and listed in Table 2 are the luminosity functions for each respective color image obtained after applying the matching criterion ($`r_m=1`$ pixel). The red F160W stars that were not matched/detected in the F110W image contribute to the F160W LF only at faint magnitudes (\[F160W\]$`22.0`$).
The matched F160W LF displayed in Figure The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope clearly shows a discontinuity in the number of stars at \[F160W\]$``$20.0. Although no other clear discontinuity is visible at fainter magnitudes, the slope of the counts changes slightly at \[F160W\]$``$21.0. As will be discussed in detail in Section 6, these discontinuities are thought to be associated with the TRGB of an intermediate-age and old population in the halo of NGC 5128. In the IR, the presence of an intermediate-age population and an old population will be blended and make the TRGB of the old population difficult to detect while the younger population will be clearly brighter and more visible. Note that the stars brighter than the \[F160W\]$``$20.0 discontinuity are most probably due to Galactic contamination (see Section 6.3 for discussion).
Given that no IR surface brightness measurements for the halo of NGC 5128 can be found in the literature, it is only possible to estimate this based on the V-band measurements of van den Bergh (1976). The V-band surface brightness at the radial distance of our field is measured to be 23.2 mag/arcsec<sup>2</sup>. This corresponds to a \[F160W\] surface brightness value of 20.2$`\pm 0.1`$ mag/arcsec<sup>2</sup>, assuming the color transformation $`VH=3.1\pm 0.1`$ (Persson, Frogel and Aaronson, 1979) and using the magnitude system transformations given in Appendix A. This mean color and standard deviation are calculated based on the distribution in $`VH`$ of the Persson, Frogel and Aaronson (1979) field ellipticals. The total light in our resolved population that passes the edge and $`\chi ^2`$ cuts averages to 21.3 mag/arcsec<sup>2</sup>. Hence, $`40`$% of the light is resolved into the stars appearing in the \[F160W\] LF (see Figure The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope; top panel). This result is consistent with a Baade’s window luminosity function (used in our simulations) and confirms that we are not confused by crowding for \[F160W\]$`23.5`$.
## 6 IR Color Dispersion
The NICMOS CMD for NGC 5128 is shown in Figure The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope. An electronic version of the photometry table may be obtained on request from the first author. The foreground reddening in the direction of our NICMOS field is $`E(BV)=0.11\pm 0.02`$ (Frogel, 1984; Harris et al., 1992). For $`R_V=3.1`$, this corresponds to $`A_J=0.10`$, $`A_H=0.06`$ and $`E(JH)=0.04`$ (Rieke and Lebofsky, 1985). Also shown in this figure is the mean color and standard deviation measured at each magnitude bin. The color histograms and Gaussian fits from which these means and standard deviations were computed are shown in Figure The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope. The reader is referred to Table 2 for the list of statistical measurements computed for the real data as well as for the artificial-stars tests. The weighted average of the mean giant branch color above our 50% completeness limit is \[F110W\]$``$\[F160W\]=1.22$`\pm 0.08`$ ($`(JH)_{CIT}=0.78`$) with a dispersion of 0.19 mag.
### 6.1 Discussion of Uncertainties
In order to estimate the metallicity spread in our IR-selected sample of halo stars, we first need to check whether the spread in color we detect in our data is real. We can then use the Bruzual and Charlot (2000; hereafter BC00) isochrones, described in detail in Section 6.2, to estimate the metallicity spread associated with the observed real color spread.
We begin by comparing the observed spread with the total uncertainties (statistical+systematic errors) calculated from the artificial-star tests. The error in color cannot simply be measured by adding in quadrature the uncertainty in \[F110W\] and \[F160W\] because the covariance of these errors is not zero; i.e., the errors are correlated. This occurs because the photometric errors are not random but associated with systematic errors that occur in the photometry of crowded fields; faint stars next to brighter ones will have correlated errors in their F160W and F110W photometry. The photometric errors for the \[F160W\] magnitudes do not appear in Figure The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope as they are smaller than the symbols used to plot the points. The mean of the DAOPHOT \[F160W\] photometric errors for our artificial-star tests are added up in quadrature and listed in Table 2. The artificial-star tests total uncertainty measurements become larger than the DAOPHOT photometric errors only in the faintest magnitude bin where the systematic errors (mostly due to crowding of faint stars) become dominant.
The standard deviations in mean color for the seven magnitude bins, as listed in Table 2 and shown in Figure The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope, can be compared with the total uncertainties measured from the artificial-stars tests. The observed spread is larger than the total uncertainty in all bins except the faintest one near our 50% completeness limit. For the faintest magnitude bin, we are only $`0.30.8`$ magnitude above our 50% completeness limit. For the three brightest bins, the observed spread is larger than the total uncertainty by at least a factor of 1.5–3.5. Assuming that the total uncertainty and intrinsic color spread add in quadrature, the observed rms for the magnitude bin \[F160W\]=22.0–22.5 of 0.20$`\pm 0.02`$ mag implies a real color spread of 0.11$`\pm 0.02`$ mag. We find for the 12 Gyr isochrones of BC00 (see Section 6.2) that at $`M_{[F160W]}=5.5`$:
$$\frac{d[Fe/H]}{d([F160W][F110W])}=2.03\mathrm{dex}/\mathrm{mag}.$$
This means that a 0.1 mag error in color corresponds to a 0.2 dex error in metallicity. Hence, the real color spread of 0.11$`\pm 0.02`$ mag corresponds to a rms for \[Fe/H\] of 0.22$`\pm 0.04`$ dex. The FWHM of the metallicity distribution is then 0.5$`\pm 0.1`$ dex.
| TABLE 2 |
| --- |
| LUMINOSITY FUNCTIONS AND COMPARISON BETWEEN |
| OBSERVED \[F110W\]$``$\[F160W\] COLOR SPREAD AND TOTAL UNCERTAINTY |
| \[F160W\] | N<sub>F160W</sub> | Observed | Observed | Artificial | Artificial | DAOPHOT | N<sub>F110W</sub> |
| --- | --- | --- | --- | --- | --- | --- | --- |
| mag bin | ($`\chi ^22.6)`$ | mean color | rms color | mean color | rms color | rms color error | |
| 17.0$``$17.5 | 1 | $`\mathrm{}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | | 0 |
| 17.5$``$18.0 | 0 | $`\mathrm{}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | | 1 |
| 18.0$``$18.5 | 0 | $`\mathrm{}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | | 0 |
| 18.5$``$19.0 | 1 | $`\mathrm{}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | | 0 |
| 19.0$``$19.5 | 2 | $`\mathrm{}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | | 0 |
| 19.5$``$20.0 | 1 | $`\mathrm{}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | | 1 |
| 20.0$``$20.5 | 6 | $`\mathrm{}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | | 2 |
| 20.5$``$21.0 | 24 | 1.25 | 0.17$`\pm 0.02`$ | 1.3 | 0.05 | 0.10 | 1 |
| 21.0$``$21.5 | 39 | 1.24 | 0.14$`\pm 0.02`$ | 1.3 | 0.06 | 0.11 | 1 |
| 21.5$``$22.0 | 95 | 1.27 | 0.18$`\pm 0.02`$ | 1.3 | 0.11 | 0.11 | 20 |
| 22.0$``$22.5 | 163 | 1.24 | 0.20$`\pm 0.02`$ | 1.2 | 0.17 | 0.13 | 37 |
| 22.5$``$23.0 | 165 | 1.18 | 0.23$`\pm 0.01`$ | 1.2 | 0.17 | 0.15 | 59 |
| 23.0$``$23.5 | 129 | 1.05 | 0.24$`\pm 0.02`$ | 1.0 | 0.24 | 0.19 | 150 |
| 23.5$``$24.0 | 36 | 0.77 | 0.20$`\pm 0.08`$ | 0.8 | 0.29 | 0.22 | 212 |
| 24.0$``$24.5 | 4 | $`\mathrm{}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | | 201 |
| 24.5$``$25.0 | 0 | $`\mathrm{}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | | 52 |
| 25.0$``$25.5 | 0 | $`\mathrm{}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | | 2 |
### 6.2 Estimate of the Metallicity Spread
We now compare our IR CMD with the observed cluster giant branches and theoretical isochrones from Bertelli et al. (1994) and BC00 to estimate the upper and lower bound on the metallicity distribution of the stars in the halo of NGC 5128.
We compare our data with the red giant branches of the clusters M92, 47 Tuc, M67, and NGC 6553 in Figure The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope. The references for the cluster data, their distance moduli, and their metallicities are listed in Table 3. The star cluster magnitudes were transformed to \[F160W\] and \[F110W\] magnitudes according to the formulae given in Appendix A. The NGC 6553 giant branch of Davidge and Simons (1994) was corrected for reddening ($`E(JH)=0.231`$) and extinction ($`A_H=0.38`$) (Rieke and Lebofsky, 1985; Guarnieri et al., 1998). The metallicity bounds on the data are chosen to bracket $`\pm 1\sigma `$ of the real color spread. The most metal poor stars in this halo field have metallicities matching those of the globular cluster M92 with \[Fe/H\]=$`2.03`$. This is in good agreement with the lower metallicity bound derived by SMW96 and HHP99. At the other end of the metallicity range NGC 6553, with \[Fe/H\]=$`0.29`$, reproduces well the upper bound of our IR data as for the optical data of HHP99. The cluster M67 is consistent with being more metal rich than 47 Tuc but its low magnitude giant branch restricts us from using it as a comparison to our data. From this comparison we conclude that the metallicity spread based on the color spread in the IR halo data is $`2.0`$\[Fe/H\]$`0.3`$ (these limits cover slightly more than $`\pm 1\sigma `$ of the real color spread; see Figure The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope). This comparison confirms the presence of a metallicity spread as estimated in Section 6.1.
Assuming an old population, the IR data support the findings of SMW96 and HHP99 in the optical suggesting that the halo of NGC 5128 is composed of stars ranging from metal-poor to near-solar metallicities. To complicate things, the presence of an intermediate-age population, as will be discussed in Section 6.4, can be another contributor to the large spread in metallicities at magnitudes below the TRGB where the populations become intertwined. In fact, it is worth noting that previous work, including the analysis presented here, has not solved for the age of the population but assumed it to be old in order to solve for the metallicity spread. If we assume that the spread in metallicity is indeed associated with an old population, then this spread suggests that metal enrichment occurred during the primordial collapse of the galaxy, or alternative, that a low-metallicity component was accreted from an external source. Based on the hierarchical picture of galaxy formation, HHP99 proposed that the old metal-poor stars in the halo formed during a first burst of star formation occurring in the in-falling clumps of gas. A first wave of supernovae explosions then ejected and enriched the gas that eventually fell back into the potential well ($`12`$ Gyr later) to form the more metal-rich stars in the halo.
The comparison with star cluster data is the best evidence of a spread in metallicity in the old population of NGC 5128. Because of the model uncertainties, comparison with theoretical isochrones can only be used as an example of how to tie the spread in colors to a spread in absolute metallicity. We compare our data with the 12 Gyr old stellar population isochrones of Bertelli et al. (1994) in Figure The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope. For a meaningful comparison, the magnitudes were properly transformed using the mathematical expressions derived in Appendix A. The isochrones shown in Figure The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope are for the metallicities \[Fe/H\]=$`1.7`$, $`0.7`$, $`0.4`$, $`+0.0`$, and $`+0.4`$. The upper parts of these IR isochrones with \[Fe/H\]$`1.7`$ (e.g., \[F160W\] $`<`$ 22.19 or $`M_{H_{BB}}<5.59`$ for \[Fe/H\]=$`0.7`$; $`BB`$ refers to the Bessell and Brett (1988) system; see Appendix A and Figure The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope) are not yet satisfactory according to Bertelli (1999), due to the problems in atmosphere models and scales of effective temperature for M giants. The comparison shows that our data are clearly more metal rich than the \[Fe/H\]=$`1.7`$ isochrone and at least as metal rich as \[Fe/H\]=$`0.7`$. Unfortunately, we cannot put a stronger upper limit on the spread in metallicity because of the problems described above.
We also make use of the most up-to-date stellar evolution models of BC00 to interpret the color spread in our CMD. The new isochrones of BC00 are generated from the new library of metallicity-dependent spectra calibrated by Lejeune, Cuisinier and Buser (1997, 1998) and an improved color-effective temperature relation for cool stars. In addition, the BC00 isochrones include the carbon AGB stars (see BC00 and Liu, Charlot and Graham 2000 for details). The prescriptions for those are semi-empirical, and essentially based on models and observations of stars in the Small Magellanic Cloud (SMC), Large Magellanic Cloud (LMC), and Milky Way. The BC00 isochrones indicate more clearly a mean metallicity of \[Fe/H\]$`0.4`$ to $`0.7`$, depending on magnitude (see Figure The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope), and a range of $`1.7`$\[Fe/H\]$`+0.0`$, where the lower limit is clearly too low for the $`\pm 1\sigma `$ color spread. The theoretical isochrones indicate a slightly more metal-rich spread in metallicities than the estimate from cluster giant branches. This is also observed in the optical data of HHP99. The red upper limit indicates that the most metal rich stars in the halo field we observed with NICMOS are roughly solar in metallicity.
We use the 12 Gyr isochrones of BC00 to estimate the metallicity for individual stars in our NICMOS field. The metallicity histogram is therefore only derived for stars with \[Fe/H\]=$`1.7`$ to $`+0.4`$. The histogram shown in Figure The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope for each 0.5 magnitude bin in the range \[F160W\]=$`21.024.0`$ was generated for stars which fell on or in between isochrones and for which we could estimate the metallicity by using a simple linear interpolation method. This is a very crude and model dependent estimate of the metallicity, as it does not include stars above the 12 Gyr isochrone TRGBs (e.g., bright intermediate-age AGB stars) and ignores the stars redder (e.g., post-AGB stars in their superwind phase) and bluer than the predicted colors (see Figure The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope). The metallicity distribution for a total of 353 stars in our CMD, shown in Figure The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope, peaks at \[Fe/H\]=$`0.76`$ with a dispersion of $`\sigma =0.44`$. This peak and dispersion agree well with the overall metallicity distribution of the globular clusters and of the halo stars of HHP99 (see Figure The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope; Harris et al. 1992; HHP99). The \[Fe/H\] distribution of stars in our NICMOS field does not seem to be resolved into multiple peaks, as is the case for the halo stars of HHP99. We find no obvious sign of a “sub-peak” which would match the largest sub-peak of the halo metallicity distribution of HHP99 at \[Fe/H\]=$`0.32`$ with a dispersion of $`\sigma =0.22`$, corresponding nicely with the second largest sub-peak in the globular cluster metallicity distribution. It is not clear at this point if the difference in the shape and peak of the metallicity distribution between the IR and the optical data of HHP99 is real or due to the large uncertainty associated with the color-metallicity transformation. The general agreement of the metallicity spread is promising; the combination of the NICMOS observations and the F555W and F814W observations of SMW96 of the same halo field should help resolve the issue (Marleau et al., 2000).
| TABLE 3 |
| --- |
| ADOPTED PARAMETERS FOR M92, 47 TUC, M67, AND NGC 6553 |
| Star Cluster | Cluster Type | CMD Data | \[Fe/H\] | $`(mM)_0`$ | Magnitude System | References |
| --- | --- | --- | --- | --- | --- | --- |
| M92 | Globular | 4 | $`2.03`$ | 14.6 | CIT | 1 |
| 47 Tuc | Globular | 7 | $`0.65`$ | 13.4 | CIT | 2 |
| M67 | Old Open | 3 | $`0.09`$ | 9.38 | CIT | 3,4 |
| NGC 6553 | Globular | 5 | $`0.29`$ | 13.6 | CIT | 5,6 |
| (1) Stetson and Harris (1988); (2) Hesser et al. (1987); (3) Houdashelt, Frogel and Cohen (1992); |
| --- |
| (4) Cohen, Frogel and Persson (1978); (5) Davidge and Simons (1994); (6) Guarnieri, Renzini and Ortolani (1997); |
| (7) Frogel, Persson and Cohen (1981) |
### 6.3 Bright Stars
As can be seen in Figure The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope and Table 2, we detect a population of bright stars above the TRGB of globular clusters and old stellar populations. These bright stars in our CMD can be due to (1) blended images, (2) Galactic contamination, or (3) a young or intermediate-age halo population.
We searched for blended doubles by first examining the SHARPNESS parameter calculated by DAOPHOT. The SHARPNESS is defined so that is zero for stars, larger than zero for galaxies or unrecognized blended doubles, and less than zero for cosmic rays or single pixel defects. We plotted the SHARPNESS and $`\chi ^2`$ as a function of magnitude for our data and found the SHARPNESS parameter to be close to zero for those bright stars above the TRGB, with good values of $`\chi ^2`$ (as expected from our initial $`\chi ^2`$ cut). Therefore, we found no reason to reject them.
The five brightest, blue stars in the upper left region of Figure The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope with \[F160W\] $`<`$ 20.0 and 0.40 $`<`$ \[F110W\]$``$\[F160W\] $`<`$ 0.85, have colors and magnitudes that are consistent with foreground Galactic stars. The theoretical isochrone of Bertelli et al. (1994) for a 5 Gyr disk population with solar metallicity is shown in Figure The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope assuming a distance modulus for the Galactic stars of 15.0. Main sequence K$``$M dwarfs have $`0.4JH0.7`$ (Tokunaga, 1995; Davidge, 1998) and therefore the five bright stars have colors and magnitudes consistent with K$``$M dwarfs in our Galaxy. To assess the number of foreground Galactic stars expected to contaminate our data, we used the Galaxy model of Cohen (1993) for IR star counts for the direction $`(l,b)=(309.5^{},19.4^{})`$, adopting a solar displacement of 15 pc, and a halo scale factor of 0.5 (Cohen, 1995). Table 4 lists the predicted number of stars for 1 magnitude wide bins for the \[F160W\] filter. The expected number of foreground stars contaminating each magnitude bin in the F160W image based on Cohen’s model, scaled by a factor of 3.9, is shown in Figure The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope. No correction was made between H and \[F160W\].
| TABLE 4 |
| --- |
| PREDICTED NUMBER OF GALACTIC STARS |
| \[F160W\] | N<sub>p</sub> | N$`{}_{p}{}^{}\times `$ 3.9 | N<sub>obs</sub> |
| --- | --- | --- | --- |
| 17.0$``$18.0 | 0.30 | 1.17 | 1 |
| 18.0$``$19.0 | 0.42 | 1.64 | 1 |
| 19.0$``$20.0 | 0.55 | 2.14 | 3 |
| 20.0$``$21.0 | 0.68 | 2.65 | 30 |
| 21.0$``$22.0 | 0.78 | 3.04 | 134 |
| 22.0$``$23.0 | 0.81 | 3.16 | 328 |
| 23.0$``$24.0 | 0.74 | 2.89 | 165 |
| 24.0$``$25.0 | 0.60 | 2.34 | 4 |
| 25.0$``$26.0 | 0.43 | 1.68 | 0 |
| 26.0$``$27.0 | 0.28 | 1.09 | 0 |
Since the bright stars (\[F160W\] $`<`$ 20.0) in the NIC2 image have colors consistent with Galactic dwarfs, it seems likely that Cohen model under-predicts the contamination since it predicts only 1.27 stars in the range \[F160W\]=17.0–20.0. The Cohen model has not been validated at these faint magnitudes and therefore we feel justified in applying an arbitrarily scaling to the Cohen predictions. A maximum-likelihood analysis shows that scaling up the predicted counts ($`N_p`$) by a factor of 3.9 produces the best fit to the observed counts ($`N_{obs}`$) with a 95% confidence interval of $`2.18.3`$. With this scale factor we predict the contamination in fainter magnitude bins as given in Table 4. Clearly, this implies that Galactic contamination is insignificant for all bins fainter than \[F160W\]=20.0.
### 6.4 The Intermediate-Age Population
The presence of an intermediate-age population of $`5`$ Gyr in the halo of NGC 5128 was proposed by SMW96 based on their detection of $`200`$ stars brighter than the TRGB. We look for the presence of intermediate-age stars as postulated by SMW96 by comparing the IR data with the 2 Gyr old isochrone from BC00 with metallicity \[Fe/H\]=$`1.7`$, $`0.7`$, $`0.4`$, $`+0.0`$, and $`+0.4`$, which include the carbon AGB stars. Figure The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope shows that the presence of an intermediate-age population would be revealed by asymptotic giant branch (AGB) stars $`1`$ mag above the TRGB. This effect is also demonstrated by moving the LMC and SMC AGB stars with ages between $`13`$ Gyr and the younger stars with ages between $`40120`$ Myr to the distance of NGC 5128 (Frogel et al., 1990). The AGB is composed of M-type (oxygen-rich) stars and, for young- and intermediate-age populations, of C-type (carbon-rich) stars which occupy the bright end of the AGB.
As Figure The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope shows, the LMC and SMC AGB stars belonging to an intermediate-age population occupy a part of the CMD that is also occupied by the brightest IR stars in the halo of NGC 5128. We find that $`10`$% of the stars resolved in our NICMOS images and appearing in our CMD are brighter than the TRGB of a 12 Gyr old population, assuming a metallicity of \[Fe/H\]=$`+0.0`$ (this brightest isochrone TRGB terminates at \[F160W\]=21.2 or $`M_{H_{BB}}=6.6`$; BC00). The IR data suggest the presence of an intermediate-age population, also seen in the WFPC2 observations of SMW96. A preliminary match between the IR and the SMW96 optical data shows that our intermediate-age population consists of the same stars that make up the intermediate-age population of SMW96 (Marleau et al., 2000). For their WFPC2 field, located 18$``$ .32 south of the galaxy center, HHP99 claim that the halo is composed almost entirely of old stars (at most 300 stars, representing $`3`$% of their sample, belong to an intermediate-age population). As the fraction of intermediate-age stars detected in the WFPC2 SMW96 field, in agreement with our NICMOS field, is larger by at least a factor of $`3`$ at half the radial distance from the galaxy center, it is suggestive of a radial gradient in the stellar population in the halo (the younger population being more centrally concentrated). With the limited area of our NICMOS field and hence the small number statistics, it is not possible to accurately estimate the metallicity mean and spread associated with the intermediate-age population only. Assuming an old halo population, the double metallicity peak in the globular cluster systems of NGC 5128 (Harris et al., 1992) and the halo stars in the HPP99 field at \[Fe/H\]$`1.1`$ (barely noticeable for the halo stars) and $`0.3`$ is well within our estimated range in metallicities.
If one excludes the brightest star at \[F160W\]=17.15 which is most certainly due to Galactic contamination, the colors and magnitudes of the other four stars with \[F160W\] $`<`$ 20.0 are close to being consistent with a very young ($`40120`$ Myr) population (see Figure The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope). But since their colors are closer to the late type K$``$M dwarfs which dominate Galactic contamination, we conclude that there is no strong evidence for a very young population of stars in our halo field.
## 7 Conclusion
We have presented the first IR CMD for the halo of a giant elliptical galaxy. Assuming a distance to NGC 5128 of 3.5 Mpc, we have detected a discontinuity in the luminosity function at \[F160W\]$``$20.0 and have measured IR magnitudes and colors for stars in the halo of NGC 5128 to \[F160W\]=23.8 (50% completeness limit). We are confident that we are not confused by crowding to \[F160W\]$`23.5`$ based on careful analysis of artificial-stars tests. The weighted average of the mean color of our giant branch above our 50% completeness limit is \[F110W\]$``$\[F160W\]=1.22$`\pm 0.08`$ ($`(JH)_{CIT}=0.78`$) with a dispersion of 0.19 mag. From our artificial-star experiments we have determined that there is a real spread in color in our CMD. By comparing our data with star cluster giant branches and theoretical isochrones, we were able to constrain the metallicity spread associated with this real color spread. Assuming an old population, we find that, in the halo field of NGC 5128 we surveyed, stars have metallicities ranging from roughly 1% of solar at the blue end of the color spread to roughly solar at the red end, with a mean of \[Fe/H\]=$`0.76`$ and a dispersion of 0.44 dex.
We assert that the five brightest stars above the SMW96 determination of the TRGB are most probably due to Galactic contamination. We found that the majority of stars above the TRGB of an old population belong to an intermediate-age population ($`2`$ Gyr). The presence of an intermediate-age population in the halo of NGC 5128 is consistent with the findings of SMW96. We conclude from our analysis that the IR data are consistent with the halo of NGC 5128 being composed of at least two age populations, a population with ages $`2`$ Gyr and an old population. Assuming an old population, we find that the stars have a wide range of metallicities. Future work will combine the WFPC2 CHIP-3 observations of SMW96 and our NICMOS data to examine the multi-color (F555W, F814W, F110W, and F160W) properties of stars in the halo of NGC 5128 and try to reconstruct the galaxy’s formation history.
F.R.M. would like to thank Nial Tanvir for providing some of the software for the completeness tests and Rachel Johnson for her instructions on aperture photometry corrections. Many thanks to Jay Anderson for the use of his code with which he initially derived for us the WFPC2(SMW96) and NICMOS star matching. We are grateful to Joan Najita and Patricia Royle for private communications while dealing with cosmic rays persistence effects and data reduction related issues. Finally, we would like to mention that Martin Cohen kindly provided us with the star count calculations from his model. This research was supported by the HST NASA grant STScI GO-07852.02-96A. F.R.M. acknowledges an IoA observational astronomy rolling grant from PPARC, ref. no. PPA/G/O/1997/00793.
## Appendix A Magnitude System Transformations
The filters used for the NICMOS observations differ substantially from ground-based IR systems. The central wavelength of F110W is close to that of $`J`$-band (1.3 $`\mu `$m), but this filter is almost twice as wide as $`J`$ and transmits from $`0.81.4\mu `$m. The F160W is analogous to ground-based $`H`$, but extends an additional 0.1 $`\mu `$m blueward. To facilitate comparison of our NICMOS data with previous ground-based IR observations, and with theoretical isochrones, we have derived color transformations between F110W and F160W and their closest ground-based counterparts, $`J`$ and $`H`$.
We need to compare our data with stellar cluster photometry (all of which is on the CIT system) and with the isochrones of Bertelli et al. (1994) which are presented on the homogenized system proposed by Bessell and Brett (1988). They present an empirical transformation from CIT to their system equal to:
$$(JH)_{BB}=0.002+1.098(JH)_{CIT}.$$
Therefore, we concentrate on calculating the transformation between the BB and NICMOS system. This transformation must be deduced theoretically because only a small number of calibration stars (five) were observed both by NICMOS (see NICMOS calibration web page at STScI) and from the ground (Persson, Frogel and Aaronson, 1979). Eventually, when more NICMOS data become public it will be possible to replace this theoretical transformation with an empirical one.
We have calculated synthetic magnitudes for stars in the Pickles (1998) spectral library in the NICMOS and BB systems. Since we have adopted a Vega-based system all colors are relative to that of an A0V star, i.e., if $`\eta _\lambda `$ is the system efficiency, taken from Bessell and Brett (1988) and from the HST Data Handbook version 3.1 March 1998 respectively, and $`f_\lambda `$ is the stellar flux, the corresponding synthetic magnitude $`m_\lambda `$ is:
$$m_\lambda =2.5log\left(\frac{_0^{\mathrm{}}f_\lambda \eta _\lambda 𝑑\lambda }{_0^{\mathrm{}}f(A0V)_\lambda \eta _\lambda 𝑑\lambda }\right).$$
Figure The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope shows an example of synthetic Vega-based colors for the Pickles library stars on the BB and NICMOS systems. The spectral types of the giants are O8III to M10III and O5V to M6V for dwarfs; both normal, and a few metal-poor and metal-rich stars have been included. This figure shows that there is a significant color difference between \[F110W\]$``$\[F160W\] and $`(JH)_{BB}`$. A simple linear fit,
$$[F110W][F160W]=0.058+1.484(JH)_{BB},$$
and
$$[F160W]=0.019+0.095(JH)_{BB}+H_{BB},$$
describes the transformation accurately. There is some scatter about the straight line for for the latest M giants. Given the mean colors of the stars in NGC 5128 we expect any systematic error to be less than 0.1 mag in locating the transformed isochrones. An error of 0.1 mag in \[F110W\]$``$\[F160W\] corresponds to a 0.2 dex error in metallicity.
The colors for five stars observed as part of the NICMOS photometric calibration campaign with magnitudes measured in both the NICMOS and Las Campanas Observatory (LCO) (Persson, Frogel and Aaronson, 1979) systems are also plotted in Figure The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope. The LCO magnitudes of these stars were first transformed to the BB system. Comparison of the data with the synthetic photometry in Figure The Nature of the Halo Population of NGC 5128 Resolved with NICMOS on the Hubble Space Telescope confirms that the theoretical colors are satisfactory and that the transformation can be made reliably for normal stars. Some of Persson’s stars are highly reddened and fall off the linear trend. However, nearly all of our stars have \[F110W\]$``$\[F160W\] $`<`$ 1.5, and so the quadratic term that emerges for very red stars is not relevant.
|
warning/0005/gr-qc0005130.html
|
ar5iv
|
text
|
# Diffusion of the electromagnetic energy due to the backscattering off Schwarzschild geometry.
## I Introduction.
Backscattering is a phenomenon that prevents waves from being transmitted exclusively along null cones. That aspect of waves propagation has been investigated for a long time for various wave equations (see, for instance, ). It has been established that solutions of the Klein-Gordon with nonuniform coefficients generically do exhibit backscattering . This topic has been investigated in general relativity since early 1960’s (, , , ); a comprehensive bibliography can be found in . The propagation of electromagnetic waves and of the resulting tails has been studied in the early seventies ( and ) and recently by Ching et al. in the context of the Schwarzschild spacetime and by Hod in the context of the Kerr spacetime. The backscattering effect can be understood as the result of waves propagation in a nonuniform medium with a varying refraction index .
In has been assessed a classical aspect of the phenomenon that was not previously studied - the energy diffusion through null cones - in the example of a spherically symmetric massless scalar field propagating in the Schwarzschild geometry. The novel aspect of that work was a compact estimate on the magnitude of the backscattered energy in terms of the energy of initial data.
This paper is dedicated to the investigation of propagation of electromagnetic fields in a background Schwarzschild spacetime. Similar to the main attention is focused on obtaining bounds on the backscattered fraction of the radiation energy, in terms of initial data. ¿From the notional point of view the present paper parallels , with three notable exceptions. First, the crucial technical points of the former work could have been applied only to spherically symmetric fields. In order to overcome this difficulty, the electromagnetic fields have to be split, with the extraction of a known part which defines initial data. Then the standard expansion in terms of vector spherical harmonics leads a problem that can be tackled with methods applied earlier in . Second, an energy inequality is proven. Third, this paper shows that the energy diffusion depends on the frequency of the radiation. An example of a dipole radiation allows one to characterize this quantitatively. The magnitude of the backscattering can be characterized as the ratio of the backscattered energy versus the initial energy of outgoing waves. This is vanishingly small in the short wave regime but it can be quite significant in the long part of the radiation spectrum. This kind of dependence on the frequency can be expected to hold also for higher multipoles. The scale is essentially set by the gravitational radius of the gravity source. All results of this paper hold true for any material sources of the Schwarzschild geometry - including stars, white dwarves, neutron stars and black holes, although the effects can really matter only in the two latter classes of objects.
The order of the remaining parts of this paper is following. The next section defines notation, basic equations and a decomposition of the electromagnetic potential. The subsequent sections of this work deal only with dipole radiation. In section III is derived an energy estimate. Section IV is dedicated to the derivation of a bound, depending on the initial energy, of the backscattered part of the potential. Section V is devoted to the derivation of useful estimates of a pair of null-line integrals. In Section VI the equations are formulated in the language of characteristics. Previously found restrictions on the backscattered part of the potential allow one to estimate radiation intensities. Section VII brings an improved estimate of the backscattered potential, again basing on the method of characteristics. The next Section proves the main results - a bound on that fraction of the energy that can diffuse due to the backscatter off the Schwarzschild geometry curvature. Section IX shows that in the case of short-wave radiation the dipole radiation backscatter is negligible. In contrast, in the long-wave regime the effect can be significant. Section X discusses how the effect depends on a distance and evaluates the exactness of the obtained criteria. The last Section presents a short summary and conclusions.
## II Formalism
Spherically symmetric geometry outside matter is given by a Schwarzschildean geometry line element,
$$ds^2=(1\frac{2m}{R})dt^2+\frac{1}{1\frac{2m}{R}}dR^2+R^2d\mathrm{\Omega }^2,$$
(1)
where $`t`$ is a time coordinate, $`R`$ is a radial coordinate that coincides with the areal radius and $`d\mathrm{\Omega }^2=d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2`$ is the line element on the unit sphere, $`0\varphi <2\pi `$ and $`0\theta \pi `$.
As it concerns the electromagnetic fields, it is convenient to assume that the scalar component of the electromagnetic potential vanishes while the vector potential satisfies the Coulomb gauge condition. Using a multipole expansion of the electromagnetic vector potential in terms of vector spherical harmonics one obtains
$$(_0^2+_r^{}^2)\mathrm{\Psi }_l=(1\frac{2m}{R})\frac{l(l+1)}{R^2}\mathrm{\Psi }_l.$$
(2)
$`\mathrm{\Psi }`$’s should be essentially two-index functions, $`\mathrm{\Psi }_{lM}`$, (where $`M`$ is the projection of the angular momentum), but since the evolution equation is $`\varphi `$ independent, the index $`M`$ is suppressed. The variable $`r^{}R+2m\mathrm{ln}(\frac{R}{2m}1)`$ is the Regge-Wheeler tortoise coordinate. The backreaction exerted by the electromagnetic field onto the metric has been neglected in the present analysis. That is readily justified for any gravitational sources other than black holes. In the case of a black hole this approximation holds true some distance away from its horizon .
Consider a set of functions of the form
$$\stackrel{~}{\mathrm{\Psi }}_l(t,r^{})=\underset{s=0}{\overset{l}{}}\frac{\mathrm{\Psi }_{ls}(r^{}t)}{R^s}$$
(3)
where the functions $`\mathrm{\Psi }_{ls}`$ are given by the recurrence relations
$`_r^{}\mathrm{\Psi }_{l1}={\displaystyle \frac{l(l+1)}{2}}\mathrm{\Psi }_{l0}`$ (4)
$`_r^{}\mathrm{\Psi }_{l(s+1)}={\displaystyle \frac{1}{2(s+1)}}\left[\left(s(s+1)l(l+1)\right)\mathrm{\Psi }_{ls}2m(s^21)\mathrm{\Psi }_{l(s1)}\right].`$ (5)
In is shown a dipole solution of this type . In the Minkowski space-time (m=0) $`\stackrel{~}{\mathrm{\Psi }}_l`$ solves (2); it represents a purely outgoing electromagnetic radiation.
Let a function $`\stackrel{~}{\mathrm{\Psi }}_l`$ be given by (3) and (5) and assume that (for space-like sections with $`t0`$) its support is compact and located entirely in the vacuum region outside some radius $`a>2m`$, i. e., outside the Schwarzschild radius. Let initial data of a solution $`\mathrm{\Psi }_l`$ of (2) concide with $`\stackrel{~}{\mathrm{\Psi }}_l`$ at $`t=0`$. Thus initially $`\mathrm{\Psi }_l`$ is a purely outgoing partial wave. It should be noted that the assumption that initial data are (initially) purely outgoing is made in this paper only for the sake of clear presentation. The propagation of electromagnetic waves is a linear process as far as the backreaction can be neglected. Therefore the propagation of the initially outgoing radiation (or even of a fraction of the outgoing radiation) is independent of whether or not the ingoing radiation is present.
It will be convenient to decompose the sought solution $`\mathrm{\Psi }_l(r^{},t)`$ into the known part $`\stackrel{~}{\mathrm{\Psi }}_l`$ and an unknown function $`\delta _l`$
$$\mathrm{\Psi }_l=\stackrel{~}{\mathrm{\Psi }}_l+\delta _l.$$
(6)
Initially $`\delta _l=_0\delta _l=0`$. A similar splitting is done in , who then seek a series expansion of $`\delta _l`$. This will be avoided in the paper, in favour of finding a number of estimates of $`\delta _l`$ that would provide the needed information about the backscattered part of the radiation.
In the rest of this paper will be considered only dipole radiation, $`\mathrm{\Psi }_1`$. Consequently, all angular momentum related subscripts will be omitted.
## III An energy estimate
The dipole term constitutes the most important part of the electromagnetic radiation. Assume dipole-type initial data
$$\stackrel{~}{\mathrm{\Psi }}(x(R))=_r^{}f(x(R))+\frac{f(x(R))}{R(r^{})},$$
(7)
with the initial support $`(a,\mathrm{})`$) of a $`C^2`$-differentiable $`f`$ and $`x(R)r^{}(R)r^{}(a)`$. The differentiability of $`f`$ guarantees that the initial energy density is continuous and vanishes on the boundary $`a`$.
Lemma 1. Define $`I_{a,ϵ}(R)`$
$`I_{a,ϵ}(R){\displaystyle _a^R}𝑑r{\displaystyle \frac{f^2(x(r))}{r^{4+2ϵ}}},`$ (8)
and
$`\beta _a(R){\displaystyle _a^R}𝑑r{\displaystyle \frac{\stackrel{~}{\mathrm{\Psi }}^2(x(r))}{r^2}},`$ (9)
where $`0<2ϵ<1`$. Then for $`a>2m(1+1/\sqrt{1+2ϵ})`$ the following inequality holds
$`I_{a,ϵ}(R){\displaystyle \frac{\beta _a(R)}{ϵa^{2ϵ}}}{\displaystyle \frac{1\left(\frac{a}{R}\right)^{2ϵ}}{(1+2ϵ)(1\frac{2m}{a})^2\frac{4m^2}{a^2}}}.`$ (10)
Remark. The integral $`\beta _a(R)`$ is bounded above by the electromagnetic energy $`E_R(t)/(4\pi )`$ defined later. Therefore
$$I_{a,ϵ}(R)\frac{E_a(R,t)}{4\pi ϵa^{2ϵ}}\frac{1\left(\frac{a}{R}\right)^{2ϵ}}{(1+2ϵ)(1\frac{2m}{a})^2\frac{4m^2}{a^2}}.$$
Proof.
Notice that
$`{\displaystyle \frac{f^2\left(x(r)\right)}{r^2}}=\left[{\displaystyle _a^r}𝑑s_s{\displaystyle \frac{f(s)}{s}}\right]^2=`$ (11)
$`\left[{\displaystyle _a^r}𝑑s{\displaystyle \frac{1}{1\frac{2m}{s}}}\left({\displaystyle \frac{\stackrel{~}{\mathrm{\Psi }}}{s}}+{\displaystyle \frac{2mf}{s^3}}\right)\right]^2`$ (12)
$`2\left[{\displaystyle _a^r}ds{\displaystyle \frac{\stackrel{~}{\mathrm{\Psi }}}{s(1\frac{2m}{s})}}\right]^2+2\left[{\displaystyle _a^r}ds{\displaystyle \frac{2mf}{s^3(1\frac{2m}{s})}}\right)]^2;`$ (13)
the second inequality follows from $`(A+B)^22A^2+2B^2`$. The factor $`1/(12m/s)`$ that appears in the integrands can be bounded above by $`1/\eta _a`$, where
$$\eta _a1\frac{2m}{a}.$$
(14)
Subsequently, the use of the Schwarz inequality and simple integrations yield
$`{\displaystyle \frac{f^2\left(x(r)\right)}{r^2}}`$ (15)
$`{\displaystyle \frac{2\beta _a(r)(ra)}{\eta _a^2}}+{\displaystyle \frac{8m^2I_{a,ϵ}(r)}{\eta _a^2(12ϵ)}}\left({\displaystyle \frac{1}{a^{12ϵ}}}{\displaystyle \frac{1}{r^{12ϵ}}}\right).`$ (16)
The insertion of (16) into the integral of (8) gives
$`I_{a,ϵ}(R)`$ (17)
$`{\displaystyle _a^R}𝑑r{\displaystyle \frac{1}{r^{2+2ϵ}}}\left[{\displaystyle \frac{2\beta _a(r)(ra)}{\eta _a^2}}+{\displaystyle \frac{8m^2I_{a,ϵ}(r)}{\eta _a^2(12ϵ)}}\left({\displaystyle \frac{1}{a^{12ϵ}}}{\displaystyle \frac{1}{r^{12ϵ}}}\right)\right].`$ (18)
$`\beta _a(r)`$ and $`I_{a,ϵ}(r)`$ are nondecreasing functions, therefore taking them in front of the appropriate integrals would not make the corresponding terms smaller. Straightforward integration of the obtained expressions yields
$`I_{a,ϵ}(R){\displaystyle \frac{\beta _a(R)}{\eta _a^2a^{2ϵ}}}\left({\displaystyle \frac{1}{2ϵ}}\left(1\left({\displaystyle \frac{a}{R}}\right)^{2ϵ}\right)+{\displaystyle \frac{1}{1+2ϵ}}\left(1+\left({\displaystyle \frac{a}{R}}\right)^{1+2ϵ}\right)\right)+`$ (19)
$`{\displaystyle \frac{2I_{a,ϵ}(R)}{12ϵ}}\left({\displaystyle \frac{2m}{a2m}}\right)^2\left({\displaystyle \frac{1}{1+2ϵ}}\left(1\left({\displaystyle \frac{a}{R}}\right)^{1+2ϵ}\right){\displaystyle \frac{1}{2}}\left(1\left({\displaystyle \frac{a}{R}}\right)^2\right)\right)`$ (20)
One should note that the expression inside the bracket in the first line can be estimated as follows
$$\frac{1}{2ϵ}\left(1\left(\frac{a}{R}\right)^{2ϵ}\right)+\frac{1}{1+2ϵ}\left(1+\left(\frac{a}{R}\right)^{1+2ϵ}\right)\frac{1}{ϵ(1+2ϵ)}\left(1\left(\frac{a}{R}\right)^{2ϵ}\right),$$
while the expression inside the bracket in the second line is bounded above by
$$\frac{12ϵ}{2(1+2ϵ)}\left(1\left(\frac{a}{R}\right)^{1+2ϵ}\right).$$
Eq. (20) can be now written as
$`I_{a,ϵ}(R){\displaystyle \frac{\beta _a(R)\left(1\left(\frac{a}{R}\right)^{2ϵ}\right)}{\eta _a^2a^{2ϵ}ϵ(1+2ϵ)}}+\left({\displaystyle \frac{2m}{a2m}}\right)^2{\displaystyle \frac{I_{a,ϵ}(r)}{1+2ϵ}}.`$ (21)
Rearranging Eq. (21) so that the two terms with $`I_{a,ϵ}(R)`$ are on the left hand side, one obtains
$$I_{a,ϵ}(R)\times \left(1\left(\frac{2m}{a2m}\right)^2\frac{1}{1+2ϵ}\right)\frac{\beta _a(R)\left(1\left(\frac{a}{R}\right)^{2ϵ}\right)}{\eta _a^2a^{2ϵ}ϵ(1+2ϵ)};$$
(22)
this gives the postulated bound of $`I_{a,ϵ}(R)`$ if $`a>2m(1+1/\sqrt{1+2ϵ})`$, as assumed above.
The obtained formula is not exact, but with the appropriate choice of $`f`$ and $`ϵ`$ the error is small. Take for instance $`f=C`$ within $`(a+a_1,bb_1)`$, $`a_1,b_1<<a`$, $`b>>a`$ (which obviously means that $`ba>>a_1,b_1`$), and let $`f`$ be smoothly joined to zero outside $`(a,b)`$ by some intermediary functions. Under those conditions, a direct calculation gives
$$\frac{I_{a,ϵ}(b)}{\beta _a(b)}\frac{3}{(3+2ϵ)a^{2ϵ}};$$
(23)
as compared with $`\frac{1}{ϵ(1+ϵ)a^{2ϵ}}`$ that follows from (22). If $`ϵ1/2`$, then the exact result differs by less than 25 percent from the bound of (22). Later on will be used the value $`ϵ=1/8`$ (which appears to be more economical in subsequent calculations), in which case the above estimate deteriorates significantly. The exact value of $`\frac{I_{a,ϵ}(b)}{\beta _a(b)}`$ is then roughly 15 % of that predicted by (22).
## IV Estimating $`\delta `$.
$`\delta `$ is initially zero and its evolution is governed by the following equation
$$(_0^2+_r^{}^2)\delta =(1\frac{2m}{R})\left[\frac{2}{R^2}\delta +\frac{6mf}{R^4}\right].$$
(24)
Define $`\stackrel{~}{\mathrm{\Gamma }}_{(R,t)}`$ \- a null geodesic that originates at $`(R,t)`$ and is directed outward. If a point lies on the initial hypersurface, then $`\stackrel{~}{\mathrm{\Gamma }}_{(R,0)}\stackrel{~}{\mathrm{\Gamma }}_R`$. By $`\stackrel{~}{\mathrm{\Gamma }}_{(R_0,t_0),(R,t)}`$ will be understood a segment of $`\stackrel{~}{\mathrm{\Gamma }}_{(R_0,t_0)}`$ ending at $`(R,t)`$.
Later will be needed a following bound.
Theorem 2. Let the support of initial data be $`(a,b)`$, $`b\mathrm{}`$ and let $`\stackrel{~}{\mathrm{\Gamma }}_{R_0,(R,t)}`$ be the outgoing null geodesic from $`(R_0,t=0)`$ to $`(R,t)`$. Then
$$\frac{|\delta (R)|}{R}mC_1\sqrt{\beta _a(b)}\frac{1}{a^ϵ\sqrt{R\eta __R}}\left(\frac{1}{R_0^{1ϵ}}\frac{1}{R^{1ϵ}}\right),$$
(25)
where
$$C_1\frac{6\sqrt{2}}{\eta _a^{3/2}(1ϵ)}\sqrt{\frac{1}{ϵ\left[(1+2ϵ)\eta _a^2\frac{4m^2}{a^2}\right]}}$$
(26)
and $`\eta __R=12m/R`$.
Proof. Define an energy $`H(R,t)`$ of the field $`\delta `$ contained in the exterior of a sphere of a radius $`R`$,
$$H(R,t)=_R^{\mathrm{}}𝑑r\left(\frac{(_0\delta )^2}{1\frac{2m}{r}}+(1\frac{2m}{r})(_r\delta )^2+(\delta )^2\frac{2}{r^2}\right).$$
(27)
One can easily show that
$`(_t+_r^{})H(R,t)=`$ (28)
$`(1{\displaystyle \frac{2m}{R}})\left[(1{\displaystyle \frac{2m}{R}})\left({\displaystyle \frac{_0\delta }{(1\frac{2m}{R})}}+_R\delta \right)^2+{\displaystyle \frac{2}{R^2}}\delta ^2\right]12m{\displaystyle _R^{\mathrm{}}}𝑑r_0\delta {\displaystyle \frac{f}{r^4}}`$ (29)
$`12m{\displaystyle _R^{\mathrm{}}}𝑑r_0\delta {\displaystyle \frac{f}{r^4}}`$ (30)
the inequality follows from omission of the nonpositive boundary term. The right hand side can be bounded further by
$$12m\left[_R^{\mathrm{}}𝑑r(_0\delta )^2\right]^{1/2}\left[_R^{\mathrm{}}𝑑r\frac{f^2}{r^8}\right]^{1/2},$$
(31)
due to the Schwarz inequality. That in turn can be bounded by
$$\frac{12m}{R^{2ϵ}\sqrt{\eta _R}}H^{1/2}\left[_R^{\mathrm{}}𝑑r\frac{f^2}{r^{4+2ϵ}}\right]^{1/2}.$$
(32)
The second integral in (32) can not increase along outgoing null directions, therefore is bounded by initial values, $`_R^{\mathrm{}}𝑑r\frac{f^2}{r^{4+2ϵ}}_{R_0}^{\mathrm{}}𝑑r\frac{f^2}{r^{4+2ϵ}}I_{R_0,ϵ}(R)`$. Since $`I_{R_0,ϵ}(\mathrm{})I_{a,ϵ}(\mathrm{})`$, one arrives at
$$(_t+_r^{})H(R,t)^{1/2}6\sqrt{I_{a,ϵ}(\mathrm{})}\frac{m}{\sqrt{\eta _a}R^{2ϵ}}.$$
(33)
The integration of (33) along $`\stackrel{~}{\mathrm{\Gamma }}_{R_0,(R,t)}`$ yields, replacing $`I_{a,ϵ}(\mathrm{})`$ by its bound expressed in (10),
$`\sqrt{H(R,t)}{\displaystyle \frac{6m}{a^ϵ}}{\displaystyle \frac{\sqrt{\beta _a(\mathrm{})}}{\eta _a^{3/2}(1ϵ)}}\sqrt{{\displaystyle \frac{1}{ϵ\left((1+2ϵ)\eta _a^2\frac{4m^2}{a^2}\right)}}}\left({\displaystyle \frac{1}{R_0^{1ϵ}}}{\displaystyle \frac{1}{R^{1ϵ}}}\right)=`$ (34)
$`={\displaystyle \frac{mC_1}{\sqrt{2}a^ϵ}}\left({\displaystyle \frac{1}{R_0^{1ϵ}}}{\displaystyle \frac{1}{R^{1ϵ}}}\right).`$ (35)
Notice that initially $`\delta `$ vanishes and that its propagation is ruled by a hyperbolic equation. Thence at any finite time $`t`$ the support of $`\delta `$ is bounded. Therefore
$`{\displaystyle \frac{|\delta (R)|}{R}}=|{\displaystyle _{\mathrm{}}^R}_r{\displaystyle \frac{\delta (r)}{r}}|`$ (36)
$`\left({\displaystyle _{\mathrm{}}^R}{\displaystyle \frac{1}{r^2}}\right)^{1/2}\left({\displaystyle _{\mathrm{}}^R}(_r\delta {\displaystyle \frac{\delta }{r}})^2\right)^{1/2}`$ (37)
$`{\displaystyle \frac{1}{\sqrt{R2m}}}(2H)^{1/2}(R)`$ (38)
Inequalities (35) and (38) yield the bound of Theorem 2 in the case when $`b=\mathrm{}`$.
Let the initial data be of finite support $`(a,b)`$. Define a region $`\mathrm{\Omega }_b`$ consisting of points $`(Rb,t)`$ acausal to $`(b,t=0)`$. The energy $`H(b(t),t)`$ obviously vanishes for any point $`(b(t),t)`$ located inside $`\mathrm{\Omega }_b`$. In this case the inequality of Theorem 2 can be stated as follows:
$$\frac{|\delta (R)|}{R}mC_1\sqrt{\beta _a(b)}\frac{\sqrt{1(\frac{a}{b})^{2ϵ}}}{a^ϵ\sqrt{R\eta _R}}\left(\frac{1}{R_0^{1ϵ}}\frac{1}{R^{1ϵ}}\right),$$
(39)
In what follows it will be assumed that the initial data have compact support located in an annular region $`(a,b)`$.
## V Estimates of two (null) line integrals
In analogy with $`\stackrel{~}{\mathrm{\Gamma }}_{(R,t)}`$ defined earlier, let $`\mathrm{\Gamma }_{(R,t)}`$ be a null ingoing geodesic that originates at $`(R,t)`$. $`\mathrm{\Gamma }_{(R,t=0)}`$ will be shortened to $`\mathrm{\Gamma }_R`$. A segment of $`\mathrm{\Gamma }_{(R_1,t_1)}`$ connecting $`(R_1,t_1)`$ with $`(R,t)`$ ($`t_1<t,R_1>R`$) will be denoted as $`\mathrm{\Gamma }_{(R_1,t_1),(R,t)}`$.
Let a point $`(R,t)`$ be an intersection of an ingoing null geodesic $`\mathrm{\Gamma }_{R_1}`$ with an outgoing null geodesic $`\stackrel{~}{\mathrm{\Gamma }}_a`$. Let $`(r,\tau )`$, $`rR`$ be a point of $`\mathrm{\Gamma }_{R_1,(R,t)}`$ and define $`(R_0(r),t=0)`$ as a point of the initial hypersurface such that $`\stackrel{~}{\mathrm{\Gamma }}_{R_0}\mathrm{\Gamma }_{R_1}=(r,\tau )`$. Fixing $`a`$ and $`R_1`$, one can view $`R_0`$ as a function of $`r`$; obviously $`R_0(R)=a`$ while $`R_0(R_1)=R_1`$. On the other hand, fixing only $`a`$ and viewing $`R_1`$ as a function of $`R`$ one has $`R_1(a)=a`$; this will be used in the forthcoming proof.
One can prove
Lemma 3. Under above conditions and if $`R_1b`$, the line integral along a null segment geodesic $`\mathrm{\Gamma }_{R_1,(R,t)}`$ is bounded above,
$$_R^{R_1}𝑑r\frac{rR_0}{R_0r\sqrt{1\frac{2m}{R}}}\frac{1}{2}\left(\mathrm{ln}\frac{R}{b}+\mathrm{ln}(\frac{b2m}{a2m})\right).$$
(40)
Lemma 4. Under the above condition but with the initial point $`(R_1,s)`$ ($`s>0`$) of the null geodesic segment $`\mathrm{\Gamma }_{(R_1,s),(R,t)}`$ lying on $`\stackrel{~}{\mathrm{\Gamma }}_b`$ (Fig. 2), one can prove
$$_R^{R_1}dr\frac{rR_0}{R_0r\sqrt{1\frac{2m}{R}}}\frac{1}{2}\mathrm{ln}(\frac{b2m}{a2m})).$$
(41)
Proof of Lemma 3.
Let $`r`$ be a radial coordinate of a point lying on the intersection of $`\stackrel{~}{\mathrm{\Gamma }}_{R_0}`$ and $`\mathrm{\Gamma }_{R_1}`$ (Fig. 1). One finds that the areal distances of three points $`(R_0(r),0)`$, $`(r,\tau )`$ and $`(R_1,0)`$ satisfy following
$$R_0(r)=2rR_1+2m\mathrm{ln}\left(\frac{(r2m)^2}{(R_0(r)2m)(R_12m)}\right).$$
(42)
That implies
$$dR_0=2\frac{1\frac{2m}{R_0}}{1\frac{2m}{r}}dr.$$
(43)
Replacing $`r`$ by $`R_0`$ in the integral of (40) one obtains
$`{\displaystyle _R^{R_1}}𝑑r{\displaystyle \frac{rR_0}{R_0r\sqrt{1\frac{2m}{r}}}}={\displaystyle _R^{R_1}}{\displaystyle \frac{dr}{\sqrt{1\frac{2m}{r}}}}\left({\displaystyle \frac{1}{R_0}}{\displaystyle \frac{1}{r}}\right)=`$ (44)
$`{\displaystyle \frac{1}{2}}{\displaystyle _a^{R_1}}𝑑R_0{\displaystyle \frac{\sqrt{1\frac{2m}{r}}}{R_02m}}{\displaystyle _R^{R_1}}{\displaystyle \frac{dr}{r\sqrt{1\frac{2m}{r}}}}`$ (45)
$`{\displaystyle \frac{1}{2}}\mathrm{ln}({\displaystyle \frac{R_12m}{a2m}})\mathrm{ln}{\displaystyle \frac{R_1}{R}}=`$ (46)
$`\mathrm{ln}{\displaystyle \frac{R}{\sqrt{aR_1}}}+{\displaystyle \frac{1}{2}}\mathrm{ln}{\displaystyle \frac{1\frac{2m}{R_1}}{1\frac{2m}{a}}}.`$ (47)
Next, one can show that $`R_12Ra`$. Indeed, assuming that $`a`$ is fixed, one has from (43) $`\frac{dR_1}{dR}2`$; since the initial condition is $`R_1(a)=a`$, the conclusion follows.
Taking into account $`R_12Ra`$ one gets
$$\mathrm{ln}\frac{R}{\sqrt{aR_1}}\mathrm{ln}\frac{R}{\sqrt{a(2Ra)}}\mathrm{ln}\sqrt{\frac{R}{a}}.$$
(48)
Replacing $`R_1`$ by $`b`$ in the last term of (47) and inserting (48), one arrives at the first of conjectured inequalities, (40).
In order to prove Lemma 4 one should start from relation between areal distances of four points $`(r,\tau )`$, $`(R_0(r),0)`$, $`(R,t)`$ and $`(a,0)`$ (see Fig. 1):
$$2\left(rR+2m\mathrm{ln}\frac{r2m}{R2m}\right)=R_0a+2m\mathrm{ln}\frac{R_02m}{a2m};$$
(49)
The variable $`r`$ ranges from $`R_1>b`$ to $`R`$. Fixing $`a`$ and $`R`$, one again obtains
$$dR_0=2\frac{1\frac{2m}{R_0}}{1\frac{2m}{r}}dr.$$
(50)
A straightforward calulation, in which $`dr`$ is replaced by $`dR_0`$, shows that
$$_R^{R_1}dr\frac{rR_0}{R_0r\sqrt{1\frac{2m}{R}}}\frac{1}{2}\mathrm{ln}(\frac{b2m}{a2m})\mathrm{ln}\frac{R_1}{R}).$$
(51)
Since $`R_1R`$, one immediately obtains (41).
## VI An estimate of the amplitudes backscattered inward
Define the intensity of the backscattered radiation that is directed inward
$$h_{}(R,t)=\frac{1}{1\frac{2m}{R}}(_0+_r^{})\delta .$$
(52)
Eq. (24) reads now
$$(_0+_r^{})\left((1\frac{2m}{R})h_{}\right)=(1\frac{2m}{R})\left[\frac{2}{R^2}\delta +\frac{6mf}{R^4}\right].$$
(53)
The integral form of (53) reads,
$$(1\frac{2m}{R})h_{}(R,t)=_{R_1}^R𝑑r\left[\frac{2}{r^2}\delta +\frac{6mf}{r^4}\right]+h_{}(R_1,s);$$
(54)
here the integration contour coincides with a null ingoing geodesics $`\mathrm{\Gamma }_{(R_1,s),(R,t)}`$. $`(R_1,s)`$ lies on the initial hypersurface ($`s=0`$) if $`R_1b`$; thus $`h_{}(R_1,s=0)=0`$, since the initial data are entirely outgoing. If $`R_1>b`$ then $`(R_1,s)\stackrel{~}{\mathrm{\Gamma }}_{b,(R_1,s)}`$; also in this case $`h_{}(R_1,s)=0`$, because $`\stackrel{~}{\mathrm{\Gamma }}_{b,R_1}`$ constitutes the outer boundary of the outgoing impulse. In either case the radiation amplitude satisfies the integral equation
$$(1\frac{2m}{R})h_{}(R,t)=_{R_1}^R𝑑r\left[\frac{2}{r^2}\delta +\frac{6mf}{r^4}\right].$$
(55)
The second term is bounded above by
$`{\displaystyle _R^{R_1}}𝑑r{\displaystyle \frac{6m|f|}{r^4}}`$ (56)
$`6m\left({\displaystyle _R^{R_1}}𝑑r{\displaystyle \frac{f^2}{r^{4+2ϵ}}}\right)^{1/2}\left({\displaystyle _R^{R_1}}𝑑r{\displaystyle \frac{1}{r^{42ϵ}}}\right)^{1/2}=`$ (57)
$`{\displaystyle \frac{6m}{\sqrt{32ϵ}}}\left({\displaystyle _R^{R_1}}𝑑r{\displaystyle \frac{f^2}{r^{4+2ϵ}}}\right)^{1/2}{\displaystyle \frac{1}{R^{(3/2)ϵ}}}\sqrt{1{\displaystyle \frac{R^{32ϵ}}{R_1^{32ϵ}}}}`$ (58)
$`{\displaystyle \frac{6m}{\sqrt{32ϵ}}}\left({\displaystyle _R^{R_1}}𝑑r{\displaystyle \frac{f^2}{r^{4+2ϵ}}}\right)^{1/2}{\displaystyle \frac{1}{R^{(3/2)ϵ}}}\sqrt{1{\displaystyle \frac{a^{32ϵ}}{b^{32ϵ}}}}`$ (59)
where the second line follows from the Schwartz inequality and the last inequality is due to the fact that $`R/R_1a/b`$ (Appendix A).
In order to find the integral from the last line of of (59) it is useful to project it onto the initial data surface, along outgoing null geodesics $`\stackrel{~}{\mathrm{\Gamma }}_{R_0,(r,\tau )}`$. Notice that $`dR_0=2\frac{1\frac{2m}{R_0(r)}}{(1\frac{2m}{r})}dr`$ \- see (43). The $`f^2/r^{4+2ϵ}`$ term cannot decrease during this projection. One arrives at
$$_R^{R_1}𝑑r\frac{f^2(r)}{r^{4+2ϵ}}_a^{R_1}𝑑R_0\frac{1\frac{2m}{r}}{2(1\frac{2m}{R_0(r)})}\frac{f^2(R_0)}{R_0^{4+2ϵ}}\frac{I_{a,ϵ}(R_1)}{2\eta _a}.$$
(60)
Inserting the energy estimate of Lemma 1 into (60) one gets finally
$$_R^{R_1}𝑑r\frac{6m|f|}{r^4}\sqrt{\beta _a(b)}\frac{mC_2}{a^ϵR^{3/2ϵ}}\sqrt{1(\frac{a}{b})^{2ϵ}}.$$
(61)
Here the constant $`C_2`$ is given by
$$C_2=\frac{3\sqrt{2\left(1\frac{a^{32ϵ}}{b^{32ϵ}}\right)}}{\eta _a^{3/2}\sqrt{ϵ[(1+2ϵ)\eta _a^2\frac{4m^2}{a^2}](32ϵ)}}.$$
(62)
The $`\delta `$-related term of (55) is bounded, due to (25), by
$$\frac{2mC_1}{a^ϵ}\sqrt{\beta _a(b)}\sqrt{1(\frac{a}{b})^{2ϵ}}_R^{R_1}𝑑r\frac{1}{\eta _Rr^{3/2}}\left(\frac{1}{R_0^{1ϵ}}\frac{1}{r^{1ϵ}}\right)$$
(63)
Here $`rR_0`$ and $`r,R_0\stackrel{~}{\mathrm{\Gamma }}_{R_0,(r,\tau )}`$. Thus $`1/(r^ϵR_0^{1ϵ})1/R_0`$. Therefore the expression of (63) is bounded above by
$$\frac{2mC_1}{a^ϵ}\sqrt{\beta _a(b)}\frac{\sqrt{1(\frac{a}{b})^{2ϵ}}}{R^{3/2ϵ}}_R^{R_1}𝑑r\frac{1}{\eta _R}\left(\frac{1}{R_0}\frac{1}{r}\right).$$
(64)
Results of Lemma 3 and 4 lead now to a pair of estimates. If $`R_1b`$ then
$$2_R^{R_1}𝑑r\frac{|\delta |}{r^2}mC_1\sqrt{\beta _a(b)}\frac{\sqrt{1(\frac{a}{b})^{2ϵ}}}{a^ϵR^{3/2ϵ}}\left(\mathrm{ln}\frac{R}{a}+\mathrm{ln}\frac{\eta _b}{\eta _a}\right)$$
(65)
and if $`R_1>b`$ (in which case $`(R,t)\mathrm{\Gamma }_{(R_1,s)}`$) then
$$2_R^{R_1}𝑑r\frac{|\delta |}{r^2}mC_1\sqrt{\beta _a(b)}\frac{\sqrt{1(\frac{a}{b})^{2ϵ}}}{a^ϵR^{3/2ϵ}}\left(\mathrm{ln}\frac{b}{a}+\mathrm{ln}\frac{\eta _b}{\eta _a}\right).$$
(66)
In summary, the radiation amplitude is bounded above by
$`(1{\displaystyle \frac{2m}{R}})|h_{}(R,t)|`$ (67)
$`{\displaystyle \frac{C_3}{a^ϵR^{3/2ϵ}}}\left[C_4+C_1\mathrm{ln}{\displaystyle \frac{b\mathrm{\Theta }(R_1(R)b)+R\mathrm{\Theta }(R_1(R)+b)}{a}}\right],`$ (68)
where $`\mathrm{\Theta }(bR_1)=0`$ if $`bR_1<0`$ and $`\mathrm{\Theta }(bR_1)=1`$ if $`bR_10`$. The constants $`C_3`$ and $`C_4`$ are defined by
$`C_3m\sqrt{\beta _a(b)}\sqrt{1({\displaystyle \frac{a}{b}})^{2ϵ}},`$ (69)
$`C_4C_2+C_1\mathrm{ln}{\displaystyle \frac{\eta _b}{\eta _a}}.`$ (70)
## VII Refining the bound on $`\delta `$
Equation (52) can be written in the integral form
$$\delta (R,t)=_{R_0}^{(R,t)}𝑑rh_{}+\delta (R_0),$$
(71)
where the integration contour coincides with $`\stackrel{~}{\mathrm{\Gamma }}_{(R_0),(R,t)}`$ and $`R_0`$ is a point of the initial Cauchy slice defined earlier. Since initially $`\delta `$ vanishes, one has $`\delta (R,t)=_{R_0}^R𝑑rh_{}`$. It becomes clear in Sec. VIII that one needs to bound $`\delta (R,t)`$ only along $`\stackrel{~}{\mathrm{\Gamma }}_{a,\mathrm{}}`$; in what follows is always meant this situation. Define $`R(b)(R:R_1(R)=b)`$ (see Fig. 2). Inserting the bound of (68) (but notice that (68) bounds $`\eta _r|h(r)|`$, not $`|h(r)|`$ itself), one obtains
$`|\delta (R)|`$ (72)
$`{\displaystyle \frac{C_3}{\eta _aa^ϵ}}[C_4{\displaystyle _a^R}{\displaystyle \frac{dr}{r^{3/2ϵ}}}+C_1\mathrm{\Theta }(R(b)R){\displaystyle _a^R}dr{\displaystyle \frac{\mathrm{ln}\frac{r}{a}}{r^{3/2ϵ}}}+`$ (73)
$`C_1\mathrm{\Theta }(RR(b))({\displaystyle _a^{R(b)}}dr{\displaystyle \frac{\mathrm{ln}\frac{r}{a}}{r^{3/2ϵ}}}+\mathrm{ln}{\displaystyle \frac{b}{a}}{\displaystyle _{R(b)}^R}{\displaystyle \frac{dr}{r^{3/2ϵ}}})].`$ (74)
The integrand of the second integral is nonnegative, therefore extending the integration up to $`R(b)`$ can give only a bigger quantity. Thus one gets, after elementary integration,
$`|\delta (R)|`$ (75)
$`{\displaystyle \frac{2C_3}{\eta _a(12ϵ)a^{1/2}}}[C_4(1\left({\displaystyle \frac{a}{R}}\right)^{1/2ϵ})++{\displaystyle \frac{2C_1}{12ϵ}}(1\left({\displaystyle \frac{a}{R(b)}}\right)^{1/2ϵ})+`$ (76)
$`C_1\left({\displaystyle \frac{a}{R(b)}}\right)^{1/2ϵ}(\mathrm{ln}{\displaystyle \frac{R(b)}{a}}+\mathrm{\Theta }(RR(b))\mathrm{ln}{\displaystyle \frac{b}{a}}(1\left({\displaystyle \frac{R(b)}{R}}\right)^{1/2ϵ}))]`$ (77)
Dropping out the negative term $`\mathrm{\Theta }(RR(b))\mathrm{ln}\frac{b}{a}\left(\frac{R(b)}{R}\right)^{1/2ϵ}`$ and taking into account that $`\mathrm{ln}\frac{R(b)}{a}+\mathrm{\Theta }(RR(b))\mathrm{ln}\frac{b}{a}\mathrm{ln}\frac{b}{R(b)}`$, one arrives at
$`|\delta (R)|`$ (78)
$`{\displaystyle \frac{2C_3}{\eta _a(12ϵ)a^{1/2}}}[C_4(1\left({\displaystyle \frac{a}{R}}\right)^{1/2ϵ})+`$ (79)
$`{\displaystyle \frac{2C_1}{12ϵ}}(1\left({\displaystyle \frac{a}{R(b)}}\right)^{1/2ϵ})+C_1\left({\displaystyle \frac{a}{R(b)}}\right)^{1/2ϵ}\mathrm{ln}{\displaystyle \frac{b}{R(b)}}]`$ (80)
Define
$$\kappa (ba)/a.$$
(81)
One can show (see Appendix B) that
$$\frac{a+b}{2}m\kappa R(b)\frac{a+b}{2};$$
(82)
the equality is achieved in the Minkowski space-time (m=0). Since $`b=a+a\kappa `$, one has $`R(b)a+\eta _a\kappa /2`$ or, defining
$$\alpha \frac{\eta _a}{2},$$
(83)
$`R(b)a+\alpha \kappa `$. The insertion of the above into (80) yields
$`|\delta (R)|{\displaystyle \frac{C_3}{\eta _a(12ϵ)a^{1/2}}}\left[C_4\left(1\left({\displaystyle \frac{a}{R}}\right)^{1/2ϵ}\right)+C_5\right],`$ (84)
where
$$C_5C_1\frac{\mathrm{ln}\frac{1+\kappa }{1+\alpha \kappa }}{\left(1+\alpha \kappa \right)^{1/2ϵ}}+\frac{2C_1}{12ϵ}\left(1\frac{1}{(1+\kappa /2)^{1/2ϵ}}\right).$$
(85)
This estimate gives a better control over the asymptotic behaviour of $`\delta `$ than the former one, (38), by a factor $`1/\sqrt{R}`$. In particular, now $`\delta ^2/R^2`$ is known to be integrable. This integrability will be exploited in the next section.
## VIII Bounding the radiation energy loss
The energy $`E_R(t)`$ of the electromagnetic field $`\mathrm{\Psi }`$ contained in the exterior of a sphere of a radius $`R`$ reads
$$E_R(t)=2\pi _R^{\mathrm{}}𝑑r\left(\frac{(_0\mathrm{\Psi })^2}{1\frac{2m}{r}}+(1\frac{2m}{r})(_r\mathrm{\Psi })^2+\frac{2(\mathrm{\Psi })^2}{r^2}\right).$$
(86)
Let the initial data be as specified hitherto, $`\mathrm{\Psi }(t=0)=\stackrel{~}{\mathrm{\Psi }}`$ and $`_0\mathrm{\Psi }(t=0)=_0\stackrel{~}{\mathrm{\Psi }}`$ for some $`\stackrel{~}{\mathrm{\Psi }}`$; thus they vanish outside an annular region $`(a,b)`$. Define $`E_a^b=E_a(0)`$ as the energy of the initial pulse. If initial configuration is purely outgoing, then $`_t\mathrm{\Psi }=_r^{}\stackrel{~}{\mathrm{\Psi }}+f_r^{}(1/R)`$.
Let an outgoing null cone $`C_a`$ originate from $`(a,0)`$. In the Minkowski spacetime the outgoing radiation contained outside $`C_a`$ does not leak inward and its energy remains constant. In a curved spacetime, however, some energy will be lost from the main stream due to the diffusion of the radiation $`h_{}`$ through $`C_a`$. Most of the backscattered radiation will fall onto the center of the gravitational attraction. The forthcoming theorem gives a bound on the amount of diffused energy.
Theorem 3. Under the above assumptions, the fraction of the diffused energy $`\delta E_a/E_a^b`$ satisfies the inequality
$`{\displaystyle \frac{\delta E_a}{E_a^b}}\left({\displaystyle \frac{2m}{a}}\right)^2{\displaystyle \frac{1\frac{1}{(1+\kappa )^{2ϵ}}}{\eta _a^2}}\times `$ (87)
$`[{\displaystyle \frac{C_4^2}{(1ϵ)}}({\displaystyle \frac{\eta _a}{16}}+{\displaystyle \frac{2ϵ^2}{(12ϵ)^2(32ϵ)}})+{\displaystyle \frac{C_5^2}{(12ϵ)^2}}+C_6+{\displaystyle \frac{2C_4C_5}{(12ϵ)(32ϵ)}}],`$ (88)
where $`C_1`$ \- $`C_5`$ have been defined earlier and
$`C_6={\displaystyle \frac{\eta _aC_1^2}{16(1ϵ)}}\left({\displaystyle \frac{\mathrm{ln}^2\frac{1+\kappa }{1+\alpha \kappa }}{(1+\alpha \kappa )^{22ϵ}}}+{\displaystyle \frac{1\frac{1}{(1+\kappa /2)^{22ϵ}}}{(1ϵ)^2}}\right)+`$ (89)
$`{\displaystyle \frac{C_1C_4\eta _a}{16(1ϵ)}}\left({\displaystyle \frac{(2+\frac{C_1\mathrm{ln}(1+\kappa /2)}{C_4(1ϵ)})\mathrm{ln}\frac{1+\kappa }{1+\alpha \kappa }}{(1+\alpha \kappa )^{22ϵ}}}+{\displaystyle \frac{1\frac{1}{(1+\kappa /2)^{22ϵ}}}{(1ϵ)}}\right).`$ (90)
Proof. The rate of the energy change along $`C_a`$ is given by
$`(_0+_r^{})E_a=`$ (91)
$`2\pi (1{\displaystyle \frac{2m}{R}})[(1{\displaystyle \frac{2m}{R}})({\displaystyle \frac{_0\mathrm{\Psi }}{1\frac{2m}{R}}}+_R\mathrm{\Psi })^2+{\displaystyle \frac{2}{R^2}}\mathrm{\Psi }^2]=`$ (92)
$`2\pi (1{\displaystyle \frac{2m}{R}})[(1{\displaystyle \frac{2m}{R}})(h_{}{\displaystyle \frac{f}{R^2}})^2+{\displaystyle \frac{2}{R^2}}(\stackrel{~}{\mathrm{\Psi }}+\delta )^2].`$ (93)
The functions $`f`$ and $`\stackrel{~}{\mathrm{\Psi }}`$ are assumed to vanish on the null cone $`C_a`$. Therefore $`\mathrm{\Psi }=\delta `$, $`_R\mathrm{\Psi }=_R\delta `$ and $`_t\mathrm{\Psi }=_t\delta `$ on $`C_a`$. In such a case the rate of the energy change reads
$`(_0+_r^{})E_a=`$ (94)
$`2\pi (1{\displaystyle \frac{2m}{R}})[((1{\displaystyle \frac{2m}{R}})h_{}^2+{\displaystyle \frac{2\delta ^2}{R^2}}].`$ (95)
The energy loss is equal to a line integral along $`\stackrel{~}{\mathrm{\Gamma }}_a`$,
$`\delta E_aE_aE_{\mathrm{}}=`$ (96)
$`2\pi {\displaystyle _a^{\mathrm{}}}dr[(1{\displaystyle \frac{2m}{r}})h_{}^2+{\displaystyle \frac{2\delta ^2}{r^2}}].`$ (97)
The derivation of (88) requires the use of estimates (68) and (84). The calculation of the $`\delta `$\- related part of the right hand side of (97) is straightforward and it yields
$`4\pi {\displaystyle _a^{\mathrm{}}}𝑑r{\displaystyle \frac{\delta ^2}{r^2}}`$ (98)
$`4\pi \beta _a(b)\left({\displaystyle \frac{2m}{a}}\right)^2{\displaystyle \frac{1\frac{1}{(1+\kappa )^{2ϵ}}}{\eta _a^2(12ϵ)^2}}\times \left[C_4^2{\displaystyle \frac{2ϵ^2}{(32ϵ)(1ϵ)}}+C_5^2+2C_4C_5{\displaystyle \frac{12ϵ}{32ϵ}}\right].`$ (99)
In order to bound the contribution coming from the backscattered radiation amplitude $`h_{}`$ one needs the estimate (68). A straightforward calculation shows that
$`2\pi {\displaystyle _a^{\mathrm{}}}𝑑r(1{\displaystyle \frac{2m}{r}})h_{}^2`$ (100)
$`\pi \beta _a(b)\left({\displaystyle \frac{2m}{a}}\right)^2{\displaystyle \frac{1\frac{1}{(1+\kappa )^{2ϵ}}}{2\eta _a^2(1ϵ)}}\times `$ (101)
$`[{\displaystyle \frac{C_4^2}{2}}+C_1C_4(y\mathrm{ln}{\displaystyle \frac{b}{R(b)}}+{\displaystyle \frac{1}{2(1ϵ)}}(1y))+`$ (102)
$`C_1^2({\displaystyle \frac{y}{2}}\mathrm{ln}^2{\displaystyle \frac{b}{R(b)}}{\displaystyle \frac{y}{2(1ϵ)}}\mathrm{ln}{\displaystyle \frac{R(b)}{a}}(12\mathrm{ln}{\displaystyle \frac{b}{R(b)}})+{\displaystyle \frac{1}{4(1ϵ)^2}}(1y)],`$ (103)
where $`y(a/R(b))^{22ϵ}`$. Neglecting the negative term with $`\mathrm{ln}\mathrm{}`$ and using the bounds of Appendix B on $`b/R(b)`$, one arrives at a bound that in conjuction with (99) proves Theorem 3.
Remark. The above estimate depends on the parameter $`ϵ`$, that should be chosen in such a way as to optimize the bound. The exact value of the optimal $`ϵ`$ depends on $`a`$ and $`\kappa `$, but the value $`ϵ=1/8`$ is proven to yield satisfactory estimates.
## IX Dependence of backscatter on the frequency of waves
The coefficients $`C_4`$ \- $`C_6`$ appearing in Theorem 3 change with $`\kappa `$ but remain finite in the whole $`(0,\mathrm{})`$ range of possible values of $`\kappa =(ba)/a`$. In the case when the support of the initial radiation is very narrow, i. e., $`\kappa <<1`$, then the coefficient $`\frac{1\frac{1}{(1+\kappa )^{2ϵ}}}{\eta _a^2}\kappa `$. I such a case one obtains that
$$\frac{\delta E_a}{E_a^b}C\left(\frac{2m}{a}\right)^2\kappa ,$$
(104)
where $`C`$ is a constant. In the limit $`\kappa 0`$ the ratio $`\frac{\delta E_a}{E_a^b}`$ becomes 0; the backscattering is negligible in the case of initial pulses of electromagnetic energy that are very narrow. And conversely, the bound becomes bigger with increase of the width of the radiation pulse. The physical meaning of that can be deduced with the help of the Fourier transform theory. The similarity theorem () states that compression of the support of a function corresponds to expansion of the frequency scale. If a support of initial data is made narrow, then the wavelengths scale of the pulse extends in the direction of short lengths. Therefore the message behind the obtained results must be that high frequency radiation is essentially unhindered by the effect of backscattering and that long waves can be backscattered.
This dependence of the backscattering on the wave length has been in fact observed in the numerical investigation of the propagation of pulses of scalar massless fields . In this case halving of the length has led to a similar decrease of the fraction of the diffused energy.
In the case of a black hole or a neutron star the scale is set essentially by the Schwarzschild radius $`R_S=2m`$; waves with lengths much shorter than $`R_S`$ are not backscattered, while waves of lengths $`R_S`$ can reveal quite a strong effect. Moreover, one can show that the $`(2m/R)^2`$ dependence of the effect implies that most of the energy diffusion occurs in regions that are not very far (as compared to the Schwarzschild radius) from the center.
In order to exemplify the above remarks, consider the diffusion effect in following two cases. Assume the same location $`a=4R_S`$, of both radiative dipoles and
i) $`\kappa =1/8`$ (i. e., the fundamental wavelength $`R_S`$) for a pulse I;
ii) $`\kappa =1/128`$ (i. e., the fundamental wavelength $`R_S/8`$) for the pulse II.
In the calculation that is reported below, $`ϵ`$ is chosen to be 1/8, in accordance with the remark ending the preceding section. Then in the case I one obtains $`\delta E_a/E_a^b<0.37`$, while in the case II (of shorter waves) one gets $`\delta E_a/E_a^b<0.004`$.
The evolution equation (53) can be written in another form as
$$(_0+_r^{})\left((1\frac{2m}{R})h_+\right)=(1\frac{2m}{R})[\frac{2}{R^2}\delta +\frac{6mf}{R^4})].$$
(105)
where
$$h_+(R,t)=\frac{1}{1\frac{2m}{R}}(_0+_r^{})\left(\delta +f\right).$$
(106)
is the intensity of the outgoing part of the radiation. The inequality (16) can be written as follows, applying Lemma 1 and the remark following it,
$$\frac{|f(R)|}{R^{3/2}}C\sqrt{E_a^b}\left(1\left(\frac{a}{b}\right)^{2ϵ}\right);$$
(107)
here $`C`$ is some constant. The integration of (105) along a null geodesic $`\stackrel{~}{\mathrm{\Gamma }}_a`$ yields now
$$(1\frac{2m}{R})h_+(R(t),t)(1\frac{2m}{a})h_+(R(0),t=0)\frac{1}{a^{3/2}}\left(1\left(\frac{a}{b}\right)^{2ϵ}\right),$$
(108)
where the proportionality constant depends only on $`ϵ`$, $`2m/a`$ and the initial energy $`E_a^b`$. Fixing the energy $`E_a^b`$, one notices that in the regime $`(ba)/a<<1`$ the right hand side of (108) is essentially zero. Thus the product $`(1\frac{2m}{R})h_+`$ is constant. In this case one clearly sees the manifestation of the redshift - the rescaling of the amplitude $`h_+`$
$`h_+(\mathrm{})=\eta _ah_+(a,t=0).`$ (109)
See also a discussion of that fact in a massless scalar field theory.
## X Distance dependence of energy diffusion and sharpness of the estimates.
The bound of Theorem 3 depends on the source location - it contains, among other factors, a square of the factor $`2m/a`$. Thus the bounds in question decrease with the increase of $`a`$. The dependence on the distance can actually be much stronger. In order to see this, consider the dipole radiation II of ii), described in the preceding section, but located at $`a=4m`$ (instead of $`a=8m`$, as assumed formerly). One obtains that now $`\delta E_a/E_a^b0.77`$, instead of the former bound $`0.001`$. Numerical results concerning the propagation of massless scalar fields also show that the backscattered energy decreases rapidly with the increasing distance .
It is of interest to establish how accurate the final estimate is. Most of the inequalities derived in this paper are sharp, in the sense, that one can find examples that saturate them. Thus, for instance, the two null-line integrals of Section V are estimated sharply (the inequalities saturate in Minkowski space-time). Similarly results of Appendices A and B are also exact; again, the inequalities become equalities in Minkowski space-time. The energy estimate of Section III is not sharp; but the ”loss of sharpness”, to say, can be less than 25 % (see the final remark in Sec. III). The main source of unsharpness is the omission of negative terms in a bound on $`\delta `$ (Sec. VII) and in bounds of diffused energy in Sec. VIII; but that becomes insignificant with the decrease of $`\kappa `$, i. e., when the width of the pulse becomes small in comparison to the Schwarzschild radius. On the other hand, the combination of two exact steps can be associated with some loss in the accuracy.
Taking this into account, it is quite likely, that in the case of sources characterized by $`\kappa <1`$ the bound in question gives an order of the diffused energy. On the other hand, if $`\kappa >>1`$, then the bound of Theorem 3 becomes very inaccurate, with
$`{\displaystyle \frac{\delta E_a}{E_a^b}}C\left({\displaystyle \frac{2m}{a}}\right)^2,`$ (111)
where $`C10^4`$. As mentioned before, the main source of unsharpness is the omission of negative terms in a bound on $`\delta `$ (Sec. VII) and in bounds of Sec. VIII. A more accurate treatment would significantly improve the estimates in the long wave regime.
## XI Discussion
The main result of this paper, Theorem 3, states that the dipole energy diffusion due to the backscattering depends on the square of $`2m/a`$, where $`m`$ is the mass of the gravitational source and $`a`$ is a location of the pulse of radiation. Sections IX and X show that the high-frequency radiation essentially is not backscattered, but that the low-frequency radiation can manifest a significant diffusion effect. The last statement is best described in terms of the dimensionless parameter $`\stackrel{~}{\kappa }R_S/\lambda `$, where $`\lambda `$ is the fundamental radiation length. If $`\stackrel{~}{\kappa }>>1`$ then the backscattering is negligible, but if $`\stackrel{~}{\kappa }1`$ then it can be significant. The above results demonstrate that the effect becomes negligible at distances much bigger than the Schwarzschild radius of a central mass. That rules out most stars as objects that can induce observable backscattering effects. For a star of a solar type and $`\lambda R_S`$, for instance, the ratio $`\frac{\delta E_a}{E_a^b}`$ can be at most $`10^{20}`$. In the case of white dwarves and $`\lambda R_S`$ the bound (111) gives $`\frac{\delta E_a}{E_a^b}<10^8`$. For long-wave radiation the bounds are bigger - the effect even looks as marginally relevant, for white dwarves, when $`\frac{\delta E_a}{E_a^b}10^2`$. However a sharper estimate would lower that by several orders.
On the other hand two astrophysical compact objects, neutron stars and (most likely) black holes, are not excluded as objects of interest.
The backscattering would damp the total luminosity produced in accretion disks that exist in vicinities of compact objects, but since the most efficient regions of the disks are located at a distance of (at least) several Schwarzschild radii, the effect would be probably weak. More relevant can be ”echoes” - aftermaths of violent flashy eruptions, produced by a part of radiation reflected from a close vicinity of a horizon of a black hole. Numerical calculations done in the massless scalar fields propagation suggest that the amplitude of the reflected radiation can constitute up to 20 % of the incident one, assuming that the length of the wave is comparable to the Schwarzschild radius of a black hole.
The results of this section can be in principle generalized into the case of higher order multipoles. The key point would consist in showing analogues of the energy estimates of Sec. III that would bound the higher multipole moments. That should lead to a variant of Theorem 3 valid under reservations similar to those expressed earlier.
An analysis similar to that of the present paper can be repeated also in the case of a weak gravitational radiation produced in disks rotating around Schwarzschildean black holes. The conclusions concerning the fraction of the diffused energy can be similar.
Acknowledgements. This work has been suported in part by the KBN grant 2 PO3B 010 16. The author is grateful to Niall O’ Murchadha for reading of the manuscript, many discussions and valuable comments. Thanks are also due to Peter Aichelburg for a useful remark and to Krzysztof Roszkowski for his help in editing of this paper.
## XII Appendix A
Lemma A. Let $`(R,t),R_1\mathrm{\Gamma }_{R_1,(R,t)}`$, $`R_1R`$ and $`R>4m`$. Then
$$\frac{R}{R_1}\frac{a}{b}.$$
(112)
Proof.
There are two separate cases that need to be considered.
i) If $`R_1`$ lies on the initial hypersurface then $`R_1b`$ and $`Ra`$ and the inequality follows immediately.
ii) If $`(R_1,s)\stackrel{~}{\mathrm{\Gamma }}_{b,(R_1,s)}`$. In this case one has
$$2\left(R_1R+2m\mathrm{ln}\frac{R_12m}{R2m}\right)=ba+2m\mathrm{ln}\frac{b2m}{a2m}$$
(113)
Define $`XR_1R`$. One finds from (113) that
$$\frac{d}{dR}\mathrm{ln}\frac{X}{R}=\frac{\frac{2m}{R}}{1\frac{2m}{R}}\frac{1}{R+X}\frac{1}{R}0$$
(114)
provided that $`R4m`$. Thus $`X/R`$ is a non-increasing function which means that $`R/R_1`$ is a non-decreasing function and $`R/R_1R(b)/b`$. Since $`R(b)a`$, one arrives at the postulated inequality.
## XIII Appendix B
Lemma B. Define $`\kappa (ba)/a0`$. Define $`(R(b),t)`$ as the the intersection point of $`\mathrm{\Gamma }_b`$ and $`\stackrel{~}{\mathrm{\Gamma }}_a`$. Then
$$\frac{a+b}{2}m\kappa R(b)\frac{a+b}{2}.$$
(115)
Proof. The relation (42) (see the main text), with $`R_1=b`$, $`r=R(b)`$ and $`R_0=a`$, can be written as ,
$$a=2R(b)b+2m\mathrm{ln}\left(\frac{(R(b)2m)^2}{(a2m)(b2m)}\right).$$
(116)
We will treat (116) as a relation between $`b`$ and $`R(b)`$, with fixed $`a`$. Obviously $`R(b)=b=a`$ when $`b=a`$. One easily finds that
$$_bR(b)=\frac{1}{2}\frac{\eta _{_{R(b)}}}{\eta _b}.$$
(117)
Notice that $`R(b)b`$. Thus $`_bR(b)1/2`$. On the other hand $`R(b)a`$. Thus $`_bR(b)(1/2)\eta _{_{R(b)}}(1/2)\eta _a`$. The use of those two bounds on $`_bR(b)`$ and the initial condition $`R(a)=a`$ immediately imply the Lemma.
|
warning/0005/astro-ph0005075.html
|
ar5iv
|
text
|
# Cosmic ray acceleration at relativistic and ultrarelativistic shock waves
## Chapter 1 Introduction
Ever since the discovery of radiation which comes from cosmos by Hess in 1912 and christened by Millikan in 1925 as ‘cosmic rays’, physicists and astronomers have speculated upon their origin. Fermi (1949) made the first serious attempt at explaining the power law nature of the cosmic ray spectrum. He noted that a particle could increase its energy at collisions against magnetic field irregularities. In his model cosmic rays interact with galactic molecular clouds that move randomly. Particles increase their energy in head-on collisions which are more frequent than overtaking collisions when they loose energy. The process is known as the second-order Fermi acceleration because the mean particle momentum gain $`\mathrm{\Delta }p/p`$ in one interaction is proportional to $`(V/v)^2`$, $`V`$ is the root-mean-square velocity of a cloud and $`v`$ is the particle velocity, considered below to be comparable to the speed of light – c. Presently the second-order Fermi acceleration is considered in plasma where the magnetic field fluctuations play a role of the Fermi ‘clouds’.
Nonrelativistic shocks. The concept that shock waves accelerate particles in a mechanism similar to the one described by Fermi (1949) appeared in four seminal papers: Krymski (1977), Axford et al. (1977), Bell (1978a,b) and Blandford & Ostriker (1978). The idea was foreshadowed by Hoyle (1960) who postulated that shocks could efficiently accelerate particles but without specifying a mechanism. Parker (1958) and Hudson (1965, 1967) attempted to obtain such mechanism based on pairs of converging shocks and, most notably, Schatzman (1963) constructed a theory based on perpendicular shocks. Contrary to the original mechanism in the convergent shock flow pattern particles interact with the flowing plasma only like in head-on collisions. Mean momentum gain in such interaction is proportional to $`U_1/c`$ ($`U_1`$ is the shock velocity) and hence the process is known as the first-order Fermi acceleration. Efficiency of the first-order relative to the second-order Fermi acceleration equals roughly $`U_1/V_A`$, where $`V_A`$ is the Alfvén speed in plasma (cf. Ostrowski & Schlickeiser 1996).
A shock wave, or briefly a shock, can be described as a sharp transition layer which propagates through plasma with a velocity exceeding the speed of sound and changes its state through the compression. The thickness of the layer is determined by the physical process responsible for thermodynamic parameters transfer from incoming plasma, upstream of the shock, to flowing away plasma, downstream of the shock. In tenuous plasma the transfer proceeds through collective electromagnetic effects and the shock width is of the order of the gyroradius of a thermal ion. In the acceleration process we will consider relativistic particles which move with speeds close to the speed of light and have a gyroradii much larger than thermal ions and consequently they see the shock as a discontinuity.
The acceleration processes in nonrelativistic – $`U_1c`$ – shocks yield power-law particles momentum spectra, $`f(p)p^\alpha `$, with a very simple formula for the spectral index of accelerated particles
$$\alpha =\frac{3r}{r1},$$
$`(1.1)`$
where
$$r=\frac{\gamma _a+1}{\gamma _a1+2M^2}$$
$`(1.2)`$
is a shock compression ratio, $`M`$ is the shock Mach number and $`\gamma _a`$ is the plasma adiabatic index. For a strong shock, $`M\mathrm{}`$, propagating in a nonrelativistic plasma with $`\gamma _a=5/3`$ we have $`r4`$ and $`\alpha 4+`$. This is encouragingly close to the index of 4.3 inferred for the source of the galactic cosmic rays. Similarly, the acceleration time expressed by a simple diffusive formula is discussed in Section 3.
Relativistic shocks. A consistent method to tackle the problem of first-order Fermi acceleration in relativistic shock waves was conceived by Kirk & Schneider (1987a; see also Kirk 1988). They assumed a parallel shock geometry and that particles are subject to pitch-angle scattering on each side of the shock. By extending the diffusion approximation to higher order terms in the anisotropy of the particle distribution, they obtained solutions to a kinetic equation of the Fokker–Planck type with the isotropic form of pitch angle diffusion coefficient. Since pitch-angle scattering conserves the particle momentum in the fluid frame, the energy spectrum is obtained by matching the solutions at the shock. Their $`Q_L`$ method yielded a particle energy spectral index for strong nonrelativistic shocks as $`\sigma 2.0`$ – where $`\sigma \alpha 2`$ – in agreement with previous results. For relativistic shocks with realistic compression of Heavens & Drury (1988), the method produced particle spectra with $`\sigma `$ slightly smaller than 2 provided the Lorentz factor of the shock $`\gamma 5`$, and slightly larger at higher $`\gamma `$. The authors derived also an angular distribution function at the shock as measured in the upstream and the downstream fluid frame. In the upstream fluid frame the distribution is strongly peaked even for a mildly relativistic case of $`U_1=0.3c`$.
Next, Kirk & Schneider (1988) extended the analysis by involving both diffusion and large-angle scattering in particle pitch angle. They discovered that – in relativistic shock waves – the presence of large angle scattering can substantially modify the spectrum of accelerated particles. An extension of the Kirk & Schneider’s (1987a) approach to more general conditions in the shock was given by Heavens & Drury (1988) who took into consideration the fluid dynamics of relativistic shock waves. They also noted that the resulting particle spectral indices depend on the perturbations spectrum near the shock in contrast to the nonrelativistic case.
Kirk & Heavens (1989) considered the acceleration process in shocks with magnetic fields oblique to the shock normal (see also Ballard & Heavens 1991) by extending the method of Kirk & Schneider (1987a). Oblique shock fronts may be conveniently classified into two categories: subluminal and superluminal. In the former ones it is possible to find a Lorentz transformation to a frame of reference in which the electric field is zero in both the upstream and the downstream regions, and the shock front is stationary. In this frame, called the de Hoffman-Teller frame (de Hoffman & Teller 1950), the energy of a particle remains constant provided it does not suffer scattering. Superluminal shocks, however, do not admit a transformation to such a frame of reference. They correspond to shock fronts in which the point of intersection of the front with a magnetic field line moves at a speed greater than c. Kirk & Heavens used the de Hoffman-Teller frame to consider the subluminal shocks. They showed, contrary to nonrelativistic results again, that such shocks led to flatter spectra than parallel ones approaching the value $`\sigma 1.0`$ when the shock velocity along the magnetic field $`U_Bc`$. Their work relied on the assumption of adiabatic invariant $`p_{}^2/B`$ conservation for particles interacting with the shock, which restricted considerations to the case of a nearly uniform magnetic field upstream and downstream of the shock.
A different approach to particle acceleration was presented by Begelman & Kirk (1990) who noted that in relativistic shocks most field configurations lead to super-luminal conditions for the acceleration process. In such conditions, particles are accelerated in a single shock transmission by drifting parallel to the electric field present in the shock. Begelman & Kirk showed that there is more efficient acceleration in relativistic conditions than that predicted by a simple adiabatic theory.
The acceleration process in the presence of finite amplitude magnetic field perturbations was considered by Ostrowski (1991; 1993) and Ballard & Heavens (1992). Ostrowski considered a particle acceleration process in the relativistic shocks with oblique magnetic fields in the presence of field perturbations, where the assumption $`p_{}^2/B=\mathrm{const}`$ was no longer valid. To derive particle spectral indices he used a method of particle Monte Carlo simulations and noted that the spectral index was not a monotonic function of the perturbation amplitude enabling the occurrence of steeper spectra than those for the limits of small and large perturbations. It was also revealed that conditions leading to very flat spectra involve an energetic particle density jump at the shock. The acceleration process in the case of a perpendicular shock shows a transition between the compressive acceleration described by Begelman & Kirk (1990) and, for larger perturbations, the regime allowing for formation of the wide range power-law spectrum. The Ostrowski (1991) method was based on the ‘mean field + perturbation’ decomposition of magnetic field, i.e. a particle is considered to propagate in the mean field along its undisturbed ‘adiabatic’ trajectory, while the magnetic field inhomogeneities are allowed for by perturbing the trajectory parameters in finite time-steps. The simulations of Ostrowski (1993) were based on the numerical integration of the particle equations of motion in a perturbed magnetic field. Finite-amplitude field perturbations were described with analytic formulae as a superposition of static sinusoidal waves.
The analogous simulations by Ballard & Heavens (1992) for highly disordered background magnetic fields show systematically steeper spectra in comparison to the above results, as discussed by Ostrowski (1993). In terms of the Lorentz factor of the shock Ballard & Heavens found a rough relation $`\alpha (3\gamma +1)/8`$ that is valid up to $`\gamma 5`$. They checked their results considering different power-law fluctuations spectra for the magnetic field and stated that differences between the resulting particle spectra were quite small.
The particle spectrum formation in the presence of non-linear coupling of accelerated particles to the plasma flow has been commented by Ostrowski (1994).
The shock waves propagating with relativistic velocities rise also interesting questions concerning the cosmic ray acceleration time scale, $`T_{acc}`$. Until our results published in 1996 (Bednarz & Ostrowski 1996 - see chapter 3) there was only somewhat superficial information available about that problem. A simple comparison to the nonrelativistic formula based on numerical simulations shows that $`T_{acc}`$ relatively decreases with increasing shock velocity for parallel (Quenby & Lieu 1989; Ellison et al. 1990) and oblique (Takahara & Terasawa 1990; Newman et al. 1992; Lieu et al. 1994; Quenby & Drolias 1995; Naito & Takahara 1995) shocks. However, the numerical approaches used there, based on assuming the particle isotropization at each scattering, neglect or underestimate a significant factor controlling the acceleration process – the particle anisotropy. Ellison et al. (1990) and Naito & Takahara (1995) included also derivations applying the pitch angle diffusion approach. The calculations of Ellison et al. for parallel shocks show similar results to the ones they obtained with large amplitude scattering. In their computations for the shock with velocity $`0.98c`$ the acceleration time scale is reduced on a factor $`3`$ with respect to the nonrelativistic formula. Naito & Takahara considered shocks with oblique magnetic fields. They confirmed the reduction of the acceleration time scale with increasing inclination of the magnetic field derived earlier for nonrelativistic shocks (Ostrowski 1988). However, their approach neglected the effects of particle cross field diffusion and assumed the adiabatic invariant conservation at particle interactions with the shock. These two simplifications limit their results to the cases with small amplitude turbulence near the shock<sup>1</sup><sup>1</sup>1One should note that the spatial distributions near the shock derived by these authors (their figures 1 and 2) do not show a particle density jump proved to exist in oblique relativistic shocks by Ostrowski (1991). It is also implicitly present in analytic derivations of Kirk & Heavens (1989).. One should also note that comparing some of the mentioned papers the derived time scales to the nonrelativistic expression does not have any clear physical meaning when dealing with relativistic shocks.
In the present paper we use pitch angle diffusion approximation for particle transport in the acceleration process. Let us note that some earlier derivations of the acceleration time scale were based on the numerical simulations involving particle scattering at point like scattering centers isotropizing the particle momentum at each scattering, the so called large angle scattering model. This approach does not provide a proper description for the acceleration processes in shock waves moving with velocities comparable to the particle velocity because it removes particle anisotropy and changes the factors related to it. Moreover, against arguments presented in some papers such scattering patterns can not be realized in turbulent magnetic fields near relativistic shocks, where most particles active in the acceleration process are able to diffuse only a short distance below a few particle gyroradii off the shock<sup>2</sup><sup>2</sup>2 However, for the nonrelativistic shock velocity and particles much above the injection energy such approximations can be safely used (cf. Jones & Ellison 1991).. Such distances are most often insufficient to allow for big particle pitch-angle changes. In shocks with oblique magnetic fields such large angle scattering patterns can substantially change the shape of the accelerated particle spectrum with respect to the pitch angle diffusion model. Additionally, as an individual particle interaction with the shock can involve a few revolutions along the magnetic field, the usually assumed adiabatic invariant conservation, $`p_{}^2/B=\mathrm{const}`$, cannot be valid for short inter-scattering intervals.
Ultrarelativistic shocks. The acceleration mechanism described in section 4.1 is quite different from that in the nonrelativistic and mildly relativistic regime so that we distinguish a class of ultrarelativistic shocks if their Lorentz factors $`\gamma 1`$. The condition $`\gamma 1`$ implies also some simplifications that allow to consider ultrarelativistic shocks as a separate class. First, the magnetic field inclination downstream of the shock is, in practice, always perpendicular to the shock normal as one can derive from Eq. 2.14. Similarly, we can approximate in Eq. 2.13 the ratio of the value the magnetic field downstream of the shock to upstream as $`B_2/B_1\sqrt{8}\gamma \mathrm{sin}\psi _1`$. A turbulence downstream of the shock could amplify this value and for example assuming equipartition with the thermal pressure downstream, one obtains $`B_2/B_1(c/V_A)\gamma `$. Moreover, independently of the plasma composition (proton-electron or electron-positron) the shock velocity relative to the downstream medium is $`U_2=c/3`$ in the limit of large $`\gamma `$.
The ultrarelativistic shocks are characterized by large anisotropy of particle momentum distribution near the shock that was presented in Bednarz & Ostrowski (1998, see Figs. 4.4 - 4.7 below). The values of two main parameters describing the acceleration process, namely the energy spectral index and the acceleration time, are independent of shock conditions if fluctuations upstream of the shock ensure the acceleration process to be effective. They tend to 2.2 (spectral index, Bednarz & Ostrowski 1998; also Bednarz & Ostrowski 1997a,b) and 1.0 $`r_g/c`$ (acceleration time, Bednarz 1998, 1999). The rough analytical calculations of Gallant & Achterberg (1999) are consistent with the Bednarz & Ostrowski (1998) paper and Gallant et al. (1998) confirm the value of spectral index for the specific condition of the extremely disordered magnetic field downstream of the shock.
Ultrarelativistic shocks are considered as sources of cosmic rays with energies exceeding $`10^{20}`$ eV and several papers suggested that gamma ray bursts (GRBs) could be sources of these particles (cf. Waxman 1995, Vietri 1995). Vietri (1995) argued that in the Fermi-type acceleration at an ultrarelativistic shock, a particle could have an relative energy gain $`\gamma ^2`$ per shock crossing cycle. Gallant & Achterberg (1999) showed that particles with initial momenta isotropically distributed upstream of the shock gain $`\gamma ^2`$ energy, but only at the first interaction of the shock. They also showed that for parameters typical of the millisecond pulsars in the neutron star binaries observed in our Galaxy, the gamma ray burst blast wave would decelerate within the pulsar wind bubble, yielding an energy spectrum $`\sigma 2`$ for the boosted particles. Moreover, this spectrum would typically extend over the energy region $`10^{18.5}10^{20}`$ eV, i.e. precisely where the ultra-high-energy cosmic rays (UHECR) component is observed. Bednarz (1999) suggested that such extremely energetic particles could be produced by reflections of the shock directly in GRBs.
Relativistic shocks in astrophysical objects. Results presented further in the theses could be applied in models of some galactic and extragalactic objects. One of those are active galactic nuclei where a central black hole ejects plasma in the form of relativistic jets. A few tens of blazars has been detected in GeV $`\gamma `$-rays by the EGRET detector (von Montigny et. al. 1995). It is widely believed that the $`\gamma `$-ray production in blazars is strictly related to the existence of relativistic jets because many of them show superluminal motions (Vermeulen & Cohen 1994). Jiang et al. (1998) applied the Königl inhomogeneous jet model (Blandford & Königl 1979; Königl 1981) to a sample of quasars and BL Lacs objects and found the Lorentz factors of jets to be a significant part the ultrarelativistic ones. In unified schemes for active galactic nuclei the Fanaroff-Riley type II (FR II) radio sources are formed by AGNs, similarly to blazars, but jets are ejected at higher angles to the line of sight. Evidence that they are relativistic even on tens or hundreds kiloparsec scales suggest that the hotspots in these sources are the downstream regions of relativistic shocks.
The recent finding of microquasars in our Galaxy, a class of objects that mimic – on scales million of times smaller – the properties of quasars opened new possibilities to study physical processes in accretion disks of black holes. The observations of Mirabell & Rodriguez (1994), Tingay et al. (1995), and Hjellming & Rupen (1995) confirm the existence of relativistic flows related to these objects, and it is expected that they form relativistic shocks in the interstellar medium.
A relativistic wind of magnetized electron-positron plasma blowing from a pulsar with the flow Lorentz factor of $`10^6`$ is expected to form a termination shock (e.g. Kennel & Coroniti 1984; Gallant & Arons 1994 and Chiueh et al. 1998). Non-thermal radiation apparently seen in the class of such objects – plerions – suggests the existence of acceleration processes inside the nebula. The Crab Nebula as the young and energetic source is the best plerion to study it. Recent optical observations of Crab Nebula using the Hubble Space Telescope and also the X-ray observations of ROSAT (cf. Hester et al. 1995) show a fascinating structure of jets, a torus of X-ray emission and complexes of sharp wisps. $`\gamma `$ ray observations of the Crab Nebula exhibit the existence of extremely energetic electrons near the pulsar (cf. de Jager et al. 1996). The electron energy is a few magnitudes larger than that in the blowing wind so an acceleration mechanism has to take place near the pulsar. Gallant & Arons (1994) proposed a mechanism where electrons gain their energy from electromagnetic waves generated by gyrating ions. The mechanism tries to explain wisps at the distance of 10” from the pulsar but a knot found at 0.7” (cf. Hester et al. 1995) is not explained in the model. We expect that acceleration mechanism presented by Bednarz & Ostrowski (1998) and Bednarz (1999) is able to account for the generation of such energetic electrons at if the ultrarelativistic shock formed near the pulsar.
Observations carried out by the Burst and Transient Source Experiment show that GRBs originate from cosmological sources (Meegan et al. 1992 and Dermer 1992). Identification of the host galaxy for the GRB 971214 (Kulkarni et al. 1998) and several other bursts causes there is little doubt now that some, and most likely all GRBs are cosmological. These phenomena are surely related to ultrarelativistic shocks with $`\gamma >10^2`$ (cf. Woods & Loeb 1995). The power-law form of the spectrum often observed at high photon energies suggests the existence of nonthermal populations of energetic particles. Bednarz & Ostrowski (1998, see chapter 4 below) showed that such shocks are able to accelerate charged particles and values of their energy spectral indices converge to $`\sigma =2.2`$ when $`\gamma \mathrm{}`$ and/or the magnetic turbulence amplitude grows.
Below, we will present our results on relativistic shock acceleration published in a series of papers Bednarz & Ostrowski (1996, 1998, 1999) and Bednarz (1999). In the next chapter we discuss our numerical simulations and problems with their application to relativistic shock conditions. Then, in chapter 3, the acceleration time scales in mildly relativistic shocks are derived for a number of magnetic field configurations. Chapter 4 is devoted to ultrarelativistic shocks. We show convergence of the particle energy spectral index to the asymptotic value $`\sigma _{\mathrm{}}2.2`$ for $`\gamma \mathrm{}`$. We also discuss particle reflections from large $`\gamma `$ shocks providing a limit for models involving GRBs as sources of UHECR. The acceleration time scale is also derived. In the last chapter 5 a short summary is presented.
## Chapter 2 Numerical simulations
In order to consider the role of particle anisotropic distributions and different configurations of the magnetic field in shocks the present work is based on the small angle particle momentum scattering approach described by Ostrowski (1991). It enables us to model effects of cross-field diffusion, important in shocks with oblique magnetic fields. Let us note (cf. Ostrowski 1993) that this code allows for a reasonable description of particle transport in the presence of large amplitude magnetic field perturbations also.
Some earlier derivations of the acceleration time scale were based on the numerical simulations involving particle scattering at point like scattering centers isotropizing the particle momentum at each scattering. This approach does not provide a proper description for the acceleration processes in shock waves moving with velocities comparable to the particle velocity because it removes particle anisotropy and changes factors related to it. Moreover, against arguments presented in some papers, such scattering pattern can not be realized in turbulent magnetic fields near relativistic shocks, where most particles active in the acceleration process are able to diffuse only a short distance below a few particle gyroradii off the shock<sup>1</sup><sup>1</sup>1 However, for the nonrelativistic shock velocity and particles much above the injection energy such approximation can be safely used (cf. Jones & Ellison 1991).. Such distances are often insufficient to allow for big particle pitch-angle changes occurring with the point-like scattering centers which isotropize particle momentum at each scattering. In shocks with oblique magnetic fields such scattering pattern can substantially change the shape of the accelerated particle spectrum with respect to the pitch angle diffusion model. Additionally, as an individual particle interaction with the shock can involve a few revolutions along the magnetic field, the usually assumed adiabatic invariant conservation, $`p_{}^2/B=\mathrm{const}`$, cannot be valid for short inter-scattering intervals.
Below, the light velocity is used as the velocity unit, $`c=1`$. As the considered particles are ultrarelativistic ones, $`p=E`$, we often put the particle momentum for its energy. In the shock we label all upstream (downstream) quantities with the subscript ‘1’ (‘2’). The quantities are given in their respective plasma rest frames but subscripts U or D mean that a parameter is measured in upstream plasma rest frame or downstream plasma rest frame, respectively.
The shock normal rest frame is the one with the plasma velocity normal to the shock, both upstream and downstream the shock (cf. Begelman & Kirk 1990). The acceleration time scales in relativistic shocks (chapter 3), $`T_{acc}`$, are always given in this particular frame in units of the upstream gyroradius divided by c but the downstream plasma rest frame quantities are used (chapter 4) for the case of ultrarelativistic ones, $`t_{acc}`$.
Here we affix a gyroradius with the index ‘$`g`$’ when it is a value given for the local uniform (tantamount to mean or homogeneous) magnetic field component. Index ‘$`e`$’ means the effective field including the field perturbations (see Eq. 2.15).
Let us denote parallel diffusion coefficient as $`\kappa _{}`$ and perpendicular diffusion coefficient as $`\kappa _{}`$. Moreover, we will sometimes use shortcuts $`\tau \kappa _{}/\kappa _{}`$ and $`\lambda \mathrm{log}_{10}(\kappa _{}/\kappa _{})`$.
If it will not cause ambiguity we will use symbol $`\psi `$ to designate the magnetic field inclination to the shock normal upstream of the shock, instead of $`\psi _1`$, and the Lorentz factor of the shock as seen upstream of the shock as $`\gamma `$, instead of $`\gamma _1`$. For the same magnetic field fluctuation patterns upstream and downstream of the shock we will use symbols without indices for these patterns.
### 2.1 Acceleration time scale in nonrelativistic versus relativistic shock waves
In the case of a nonrelativistic shock wave, with velocity $`U_11`$, the acceleration time scale can be defined as
$$T_{acc}\frac{E}{\frac{\overline{\mathrm{\Delta }E}}{\mathrm{\Delta }t}},$$
$`(2.1)`$
where $`\overline{\mathrm{\Delta }E}`$ is the mean energy gain at particle interaction with the shock and $`\mathrm{\Delta }t`$ is the mean time between successive interactions. One can use mean values here because any substantial increase of particle momentum requires a large number of shock-particle interactions and the successive interactions are only very weakly correlated with each other. The respective expression for $`T_{acc}`$ in parallel shocks,
$$T_{acc}^0=\frac{3}{U_1U_2}\left\{\frac{\kappa _1}{U_1}+\frac{\kappa _2}{U_2}\right\},$$
$`(2.2)`$
where $`\kappa _i`$ is the respective particle spatial diffusion coefficient, has been discussed by Lagage & Cesarsky (1983). Ostrowski (1988) provided the analogous scale for shocks with oblique magnetic fields and small amplitude magnetic field perturbations. It can be written in the form
$$T_{acc}^\psi =\frac{3}{U_1U_2}\left\{\frac{\kappa _{n,1}}{U_1\sqrt{\frac{\kappa _{n,1}}{\kappa _{,1}\mathrm{cos}^2\psi _1}}}+\frac{\kappa _{n,2}}{U_2\sqrt{\frac{\kappa _{n,2}}{\kappa _{,2}\mathrm{cos}^2\psi _2}}}\right\},$$
$`(2.3)`$
where the index $`n`$ denotes quantities normal to the shock, the index $``$ those parallel to the magnetic field, $`\psi `$ is an angle between the magnetic field and the shock normal and $`U_1/\mathrm{cos}\psi _1c`$ is assumed. The terms $`\sqrt{\kappa _n/(\kappa _{}\mathrm{cos}^2\psi )}`$ represent a ratio of the mean normal velocity of a particle to such velocity in the absence of cross-field diffusion. One may note that for negligible cross-field diffusion the expression (2.3) coincides with (2.2) if we put $`\kappa _{n,i}`$ for $`\kappa _i`$ ($`i`$ = $`1`$, $`2`$). The case of oblique shock with finite amplitude field perturbations has not been adequately discussed yet, but we expect the respective acceleration scale to be between the values given by the above formulae for $`T_{acc}^0`$ and $`T_{acc}^\psi `$. The influence of the particle escape boundary on the acceleration time scale and the particle spectrum is discussed by Ostrowski & Schlickeiser (1996).
If the shock velocity becomes relativistic, the particle energy change at a single interaction with the shock can be comparable, or even larger than the original energy. Moreover, after interaction with the shock, the upstream particles with small initial angles between its momenta and the mean magnetic field have a larger chance to travel far away from the shock. On average, such particles spend longer times and are able to change its pitch angles substantially until the next hits at the shock. Then, larger pitch angles allow for particle reflections with large energy gains or for transmissions downstream (cf. Ostrowski 1991, Lucek & Bell 1994). Therefore, correlations of the times between successive interactions, $`\mathrm{\Delta }t_{diff}`$, the energy gains at these interactions, $`\mathrm{\Delta }E`$, and possibly the probability of particle escape occur. As an example, in Fig. 2.1 we map the number of particle interactions with the shock in coordinates ($`\mathrm{\Delta }t_{diff}`$, $`\mathrm{\Delta }E`$). A cut of the presented surface at any particular value of $`\mathrm{\Delta }t_{diff}`$ gives the distribution of energy gains for particles who have spent this time since the last interaction with the shock. A general trend seen on the map for increasing $`\mathrm{\Delta }t_{diff}`$ is the growing value of $`\mathrm{\Delta }E/E`$ for the distribution maximum. Because of these correlations are accompanied with the large energy gains $`\mathrm{\Delta }EE`$, we propose a different approach to the derivation of the acceleration time scale with respect to the one used for nonrelativistic shocks. Usually the acceleration time scale is applied for the derivation of the highest energies occurring in the particle spectrum, characterized by its cut-off energy, $`E_c`$. Thus we use this energy scale to define the acceleration time scale as
$$T_{acc}^{(c)}\frac{E_c}{\dot{E}_c},$$
$`(2.4)`$
where $`\dot{E}_cdE_c/dt`$. The rate of the cut-off energy increase is a well-defined quantity and the time scale (2.4) has a clear physical interpretation. The above definition does not require any limit for the energy gains of individual particles and all possible correlations are automatically included here. From the meaning of the definition (2.4) it follows that $`T_{acc}^{(c)}`$ is somewhat shorter than the respective scale at the same energy for later times required for the respective part of the spectrum to become a pure power-law (cf. Ostrowski & Schlickeiser 1996). One should also note that in relativistic shocks the time scale depends on the reference frame we use for its measurement. In the present paper the acceleration time scales are given in the respective normal shock rest frame. However, the applied time units $`r_{e,1}/c`$ are defined with the use of the upstream gyration time.
### 2.2 Transport of particles
To derive particle trajectories in a disturbed magnetic field one should, in general, integrate full equation of motion along these trajectories (see summary in Decker 1988 and Ostrowski 1988). However, for slightly inhomogeneous fields it was proposed a ‘quasi-linear’ approximation for analytical calculations (e.g. Jokipii 1971) consisting of distinguishing between two factors determining a particle’s trajectory: the ‘adiabatic’ undisturbed motion in the mean field $`\stackrel{}{B_0}`$, and perturbations to this trajectory derived by averaging the effect of magnetic field perturbations $`\delta \stackrel{}{B}=\stackrel{}{B}\stackrel{}{B_0}`$ along the trajectory. As a result the description of particle transport in terms of the Fokker-Planck equation includes the diffusive term in the pitch angle $`\vartheta `$, where $`\vartheta \mathrm{}(\stackrel{}{p},\stackrel{}{B_0})`$, which describes trajectory perturbations and all quantities are averaged over the phase angle along the trajectory $`\phi `$. In the case of efficient particle scattering maintaining the particle distribution function $`f(\stackrel{}{r},p,\vartheta ,t)`$ is very nearly isotropic, the equation can be reduced to the spatial diffusion equation. Concerning the pitch angle diffusion, all information on the particle scattering process is contained in the pitch angle diffusion coefficient $`D_\vartheta =\mathrm{\Delta }\vartheta ^2/(2\mathrm{\Delta }t)`$ where $`\mathrm{\Delta }\vartheta `$ ($`1`$) is the change of $`\vartheta `$ during an individual ‘scattering act’, and $`\mathrm{}`$ denotes taking the average (see Chandrasekhar 1943). Perturbing force acts at the particle trajectory in a continuous way, and the notion of the scattering act may be introduced by summing up all changes to the orbit over some time $`\mathrm{\Delta }t`$, long enough for the corresponding pitch angle changes to be uncorrelated in the successive scattering acts. In applications, usually a process of diffusion in parameter $`\mu \mathrm{cos}\vartheta `$ with the corresponding diffusion coefficient, $`D_\mu =D_\vartheta (1\mu ^2)`$ is considered. The relation of the above Fokker-Planck approach to the general situation also involving large-angle scattering was discussed by Kirk & Schneider (1988). Our numerical approach resembles the one applied by Kirk & Schneider (1987b).
We restrict our consideration to the test-particle approximation, in which it is assumed that particles are scattered by scattering centers in the fluid but have no effect either on the fluid velocity or on the density of scattering centers. Between two successive scatterings, the particle is assumed to proceed along the undisturbed path in the mean field. We will furthermore assume that the scattering centers are frozen into the fluid. The assumption implicates that $`|\stackrel{}{p}|=\mathrm{const}`$. We introduce discrete uncorrelated perturbations of the particle’s direction \[i.e. perturbations in $`\mathrm{\Delta }\vartheta `$ (or $`\mathrm{\Delta }\mu `$) and $`\mathrm{\Delta }\phi `$\], in finite time steps, $`\mathrm{\Delta }t`$. Thus, all particle momentum vectors can be represented as points on the sphere of constant $`|\stackrel{}{p}|`$ parameterized with two angles $`\vartheta `$ (or $`\mu `$) and $`\phi `$ (Fig. 2.2). All our calculations were performed in the respective local plasma rest frame. In any such frame the electric field vanishes and particle energies are conserved.
If the distribution of the particle orientation $`\mathrm{\Omega }(\vartheta ,\phi )`$ at the sphere (Fig. 2.2) maintain the same form at any point on that sphere independently of the local coordinate lines then the scattering process is not affected by the orientation. In the case one gets ‘isotropic’ diffusion coefficient $`D_\vartheta `$ that is independent of $`\vartheta `$. Let us denote the scattering amplitude (angle between the original orientation $`\mathrm{\Omega }`$ and the one after scattering $`\stackrel{~}{\mathrm{\Omega }}`$) as $`\mathrm{\Delta }\mathrm{\Omega }`$, and the angle between the meridian, $`\phi =\mathrm{const}`$, and the great circle connecting $`\mathrm{\Omega }`$ with $`\stackrel{~}{\mathrm{\Omega }}`$ as $`\beta `$. For consistency we demand that, in the limit of small scattering amplitudes, the considered scattering model should lead to an isotropic diffusion coefficient. The considered scattering probability distribution, $`F=F(\mathrm{\Delta }\mathrm{\Omega },\beta )`$ must satisfy the ‘elliptical’ symmetry: $`F(\mathrm{\Delta }\mathrm{\Omega },\beta )`$=$`F(\mathrm{\Delta }\mathrm{\Omega },\pi \beta )`$=$`F(\mathrm{\Delta }\mathrm{\Omega },\beta )`$. Kirk & Schneider (1987b) used the distribution $`F(\mathrm{\Delta }\mathrm{\Omega },\beta )`$ derived from the heat conduction equation. It was equivalent to assuming that the diffusive character of particle trajectories is also preserved at the limit $`\mathrm{\Delta }t0`$. However, in a general case, one should assume a form determined by considered form of the magnetic field perturbations. The simple choice is to take $`F=F(\mathrm{\Delta }\mathrm{\Omega })`$ which does not distinguish any direction in space. Unless one has any particular pattern of field perturbations it is the most natural choice and we will restrict ourselves to such distributions below. In particular, we take it in a normalized form
$$F(\mathrm{\Delta }\mathrm{\Omega })=\{\begin{array}{cc}(1\mathrm{cos}\mathrm{\Delta }\mathrm{\Omega }_{max})^1\mathrm{sin}\mathrm{\Delta }\mathrm{\Omega }\hfill & (\mathrm{\Delta }\mathrm{\Omega }\mathrm{\Delta }\mathrm{\Omega }_{max})\hfill \\ 0\hfill & (\mathrm{\Delta }\mathrm{\Omega }>\mathrm{\Delta }\mathrm{\Omega }_{max}\text{)}\hfill \end{array}$$
$`(2.5)`$
which ensures an equal probability of reaching any unit surface element of the sphere within the range $`\mathrm{\Delta }\mathrm{\Omega }_{max}`$ from the original position. One should note that this model scattering is no longer a symmetric one in $`\mu `$. The fact is visible after averaging the spherical triangle relation (Eq. 2.11) over $`\beta `$, for a given $`\mathrm{\Delta }\mathrm{\Omega }`$, the mean change of $`\mu `$ is $`\mathrm{\Delta }\mu =\mu (\mathrm{cos}\mathrm{\Delta }\mathrm{\Omega }1)`$. The anisotropy results from the projection of the circle $`\mathrm{\Delta }\mathrm{\Omega }=\mathrm{const}`$ in the spherical coordinates ($`\mu ,\phi `$) and, for constant $`\mathrm{\Delta }\mathrm{\Omega }_{max}`$, does not lead to any actual particle anisotropy. In the Fokker-Planck equation
$$\frac{f}{t}+v\mu \frac{f}{z}=\frac{}{\mu }\left(\frac{\mathrm{\Delta }\mu }{\mathrm{\Delta }t}f\right)+\frac{1}{2}\frac{^2}{\mu ^2}\left(\frac{\mathrm{\Delta }\mu ^2}{\mathrm{\Delta }t}f\right),$$
$`(2.6)`$
where $`ff(z,\mu ,t)`$ is the particle distribution function presented in simplified form with a spatial coordinate $`z`$ along the magnetic field, any homogeneous stationary solution must be the isotropic one. Thus the consistency condition for the Fokker-Planck coefficients is
$$\mathrm{\Delta }\mu +\frac{1}{2}\frac{}{\mu }\mathrm{\Delta }\mu ^2=0,$$
$`(2.7)`$
and the above mentioned anisotropy is compensated for by a gradient of the diffusion coefficient $`D_\mu `$.
From the distribution (2.5), for $`\mathrm{\Delta }\mathrm{\Omega }_{max}1`$, one obtains the relation between mean values $`2\mathrm{\Delta }\mathrm{\Omega }^24\mathrm{\Delta }\vartheta ^2\mathrm{\Delta }\mathrm{\Omega }_{max}^2`$. Using the definition $`D_\vartheta \mathrm{\Delta }\vartheta ^2/(2\mathrm{\Delta }t)`$, we obtain
$$\mathrm{\Delta }\mathrm{\Omega }_{max}^2=8D_\vartheta \mathrm{\Delta }t.$$
$`(2.8)`$
A great number of reasonable distributions of $`\mathrm{\Delta }t`$ could be proposed for any value of $`\mathrm{\Delta }t`$, which may be interpreted as representative for different perturbations spectra. In the limit of infinitesimal scattering amplitude the physical picture is not sensitive to the particular choice of this distribution. However, for higher amplitudes and anisotropic particle distributions this selection may qualitatively affect the simulation results and should be done with great care (cf. Kirk & Schneider 1987b, 1988). For the simulation of large-amplitude scattering one can use equation (2.8) only in a formal manner. Now, the factor $`8D_\vartheta `$ still provides the relation between $`\mathrm{\Delta }\mathrm{\Omega }_{max}`$ and $`\mathrm{\Delta }t`$, but for scatterings of small amplitude it has the additional property of being 8 times the pitch angle diffusion coefficient. Let us also note that in our method we make use of the concept of a mean field and assume that the scatterings are not correlated. For highly perturbed magnetic fields both assumptions may be of limited validity.
Based on the above model one is able to construct an algorithm for the derivation of the scattering momentum orientation ($`\mu _{new},\phi _{new}`$) from the original one ($`\mu ,\phi `$), after an individual scattering act. Let us denote two independent random values from the range ($`0,1`$) as $`R_1`$ and $`R_2`$. Using equation (2.5) we can generate the value for $`\mathrm{\Delta }\mathrm{\Omega }`$ as
$$\mathrm{cos}(\mathrm{\Delta }\mathrm{\Omega })=1(1\mathrm{cos}\mathrm{\Delta }\mathrm{\Omega }_{max})R_1$$
$`(2.9)`$
and the orientation angle $`\beta `$
$$\beta =2\pi R_2.$$
$`(2.10)`$
The new value for pitch angle cosine is derived from the spherical triangle $`\mathrm{\Omega }\mathrm{\Pi }\stackrel{~}{\mathrm{\Omega }}`$ of Fig 2.2 as
$$\mu _{new}=\mu \mathrm{cos}\mathrm{\Delta }\mathrm{\Omega }+\sqrt{1\mu ^2}\mathrm{sin}\mathrm{\Delta }\mathrm{\Omega }\mathrm{cos}\beta .$$
$`(2.11)`$
In equation (2.11) an exact value for $`\mu _{new}`$ is obtained, and one can consider high-amplitude scattering ($`\mathrm{\Delta }\mathrm{\Omega }1`$) as well. However, as was mentioned previously, one should consider carefully the meaning of the diffusion coefficient in this case. For instance, in simulating a diffusion perpendicular to the field one should also account for the possibility of phase perturbation along the trajectory. In our approach the considered spherical triangle yields
$$\phi _{new}=\phi +\mathrm{arctan}\left(\frac{\mathrm{sin}\mathrm{\Delta }\mathrm{\Omega }\mathrm{sin}\beta }{\mathrm{cos}\mathrm{\Delta }\mathrm{\Omega }\sqrt{1\mu ^2}\mathrm{sin}\mathrm{\Delta }\mathrm{\Omega }\mu \mathrm{cos}\beta }\right)+\pi H(\mu \mu _{new}\mathrm{cos}\mathrm{\Delta }\mathrm{\Omega }),$$
$`(2.12)`$
where $`H(x)`$ is the Heaviside step function.
### 2.3 Magnetic field
In the present discussion we consider the role of the mean magnetic field configuration and the amount of particle scattering. In order to avoid effects of varying shock compression due to the presence of different magnetic field configurations we take the field as a trace one without any dynamical effects on the plasma flow. The shock compression, as seen in the shock normal rest frame, $`r=U_1/U_2`$, is derived from the approximate formulae presented by Heavens & Drury (1988). For illustration of the results, in the present theses we consider the shock waves propagating in the cold electron-proton plasma. For the mean magnetic field $`B_1`$ taken in the upstream plasma rest frame and inclined at the angle $`\psi _1`$ with respect to the shock normal we derive its downstream value and inclination, $`B_2`$ and $`\psi _2`$, with the use of jump conditions presented for relativistic shocks by e.g. Appl & Camenzind (1988)
$$B_2=B_1\sqrt{\mathrm{cos}^2\psi _1+R^2\mathrm{sin}^2\psi _1},$$
$`(2.13)`$
$$\mathrm{tan}\psi _2=R\mathrm{tan}\psi _1,$$
$`(2.14)`$
where $`R=r\gamma _1/\gamma _2`$ and the Lorentz factors $`\gamma _i1/\sqrt{1U_i^2}`$ ($`i`$ = $`1`$, $`2`$). These formulae are valid for both sub- and super-luminal magnetic field configurations.
We model particle trajectory perturbations by introducing small-angle random momentum scattering along the mean-field trajectory (cf. Ostrowski 1991). The particle momentum scattering distribution is uniform within a cone wide at $`\mathrm{\Delta }\mathrm{\Omega }`$, along the original momentum direction. The presented simulations for mildly relativistic shocks use a constant value of $`\mathrm{\Delta }\mathrm{\Omega }=0.173`$ ($`=10^{}`$). Scattering events are at discrete instants, equally spaced in time as measured in the units of the respective $`r_{g,i}/c`$ ($`i`$ = $`1`$, $`2`$). The increasing perturbation amplitude is reproduced in simulations by decreasing the time period $`\mathrm{\Delta }t`$ between the successive scatterings.
In ultrarelativistic shock waves efficient particle scattering with a very small $`\mathrm{\Delta }\mathrm{\Omega }`$ requires derivation of a large number of scattering acts and the respective numerical code becomes extremely time-consuming. In order to overcome this difficulty we propose a hybrid approach involving ‘very small’ $`\mathrm{\Delta }\mathrm{\Omega }_C`$ ($`0.5\gamma ^1`$) close to the shock, where the scattering details play a role, and much larger scattering amplitude $`\mathrm{\Delta }\mathrm{\Omega }_F=9^{}`$ to describe particle diffusion further away from the shock. The respective scaling of the scattering time $`\mathrm{\Delta }t`$ is performed in both cases ($`\mathrm{\Delta }\mathrm{\Omega }_C^2/\mathrm{\Delta }t_C=\mathrm{\Delta }\mathrm{\Omega }_F^2/\mathrm{\Delta }t_F`$) to yield the same turbulence amplitudes measured by the values of the cross-field diffusion coefficient, $`\kappa _{}`$ and the parallel diffusion coefficient $`\kappa _{}`$. For a few instances we checked the validity of this approach by reproducing the results for the small $`\mathrm{\Delta }\mathrm{\Omega }_C`$ everywhere.
For simplicity, except sections 4.1 and 4.3, we use the same scattering pattern ($`\mathrm{\Delta }\mathrm{\Omega }`$ and $`\mathrm{\Delta }t`$ in units of $`r_g/c`$) upstream and downstream the shock, leading to the same values of $`\kappa _{}/\kappa _{}`$ in these regions (see, however, Ostrowski 1993). One should note that the particle momentum scattering due to the presence of the turbulent magnetic field is equivalent to the effective magnetic field larger than the respective uniform mean component, $`B_1`$ or $`B_2`$. In our model, the effective field can be estimated as
$$B_{e,i}=B_i\sqrt{1+\left(0.67\frac{\mathrm{\Delta }\mathrm{\Omega }}{\mathrm{\Delta }t}\right)^2}(i=1,2).$$
$`(2.15)`$
It is the lower limit for the actual field since the amount of power in perturbations with wave-lengths smaller than $`c\mathrm{\Delta }t`$ cannot be considered within such a simple model. The amount of energy in magnetic turbulence with the waves shorter than $`c\mathrm{\Delta }t`$ is required to be small because the presented estimate assumes the particle momentum perturbation in $`\mathrm{\Delta }t`$ occurs on the uniform effective perturbing field. To compare the scattering processes with different $`\mathrm{\Delta }t`$ one has to neglect the unknown factor of the ratio of the averaged actual magnetic field to the estimated value (like the one in Eq. 2.15). Let us note that this factor, as well as the notion of the effective field were not considered earlier.
For relativistic shocks the derived acceleration time scales are presented in units of the formal diffusive scale $`T_04(\kappa _{n,1}/U_1+\kappa _{n,2}/U_2)/c`$ or in units of $`r_{e,1}/c`$ , in the shock normal rest frame but for ultrarelativistic ones in units of $`r_{g,2}/c`$ in the downstream plasma rest frame.
### 2.4 Fitting the spectrum and the acceleration time
Our numerical calculations involve particles with momenta systematically increasing over several orders of magnitude. In order to avoid any energy dependent systematic effect we consider the situation with all spatial and time scales – defined by the diffusion coefficient, the mean time between scatterings and the shock velocity – to be proportional to the particle gyroradius, $`r_g=p/(eB)`$, i.e. to its momentum.
For a chosen shock velocity and the magnetic field configuration we inject particles in the shock at some initial momentum $`p_0`$ and follow their phase-space trajectories. We assume the constant particle injection to continue in time after the initial time $`t_0=0`$. A particle is excluded from simulations if it escapes through the free-escape boundary placed far downstream of the shock or reaches the energy larger than the assumed upper limit. These particles are replaced with the ones arising from splitting the remaining high-weight particles, preserving their physical parameters (cf. Kirk & Schneider 1987b; Ostrowski 1991). Particles that exist longer than the time upper limit for simulations are excluded from simulations without replacing. Here we put the boundary at the distance $`6\kappa _{2,n}/U_2+4r_{g,2}`$ for relativistic and $`4r_{g,2}`$ for ultrarelativistic shocks. We checked by simulations that any further increase of this distance does not influence the results in any noticeable way. For every shock crossing, the particle weight factor multiplied by the inverse of the particle velocity normal to the shock ($``$ particle density) is added to the respective time and momentum bin of the spectrum as measured in the shock normal rest frame. As one considers a continuous injection in all instants after $`t_0`$, in order to obtain the particle spectrum at some time $`t_j>t_0`$ one has to add to particle density in a bin $`p_i`$ at $`t_j`$ the densities in this momentum bin for all the earlier times. The resulting particle spectra are represented as power-law functions with the squared exponential cut-off in momentum
$$f(p,t)=Ap^\alpha e^{\left(\frac{p}{p_c}\right)^2}.$$
$`(2.16)`$
In this formula three parameters are to be fitted: the normalization constant $`A`$, the spectral index for the stationary solution $`\alpha `$, and the momentum cut-off $`p_c`$ (Fig. 2.3). Any simulated spectrum evolved in time by increasing the width of its power-law section and thus the best fit of this power-law was possible with the use of the final spectrum at maximum time. Therefore, in the simulations we used the last spectrum to fit parameters $`A`$ and $`\alpha `$. Next, for any earlier spectrum, these parameters were assumed to be constant and we were fitting only the cut-off momentum, $`p_c`$. For each fit we used 20 last points of the spectrum preceding the point where particle density fell below $`0.16`$ of $`Ap^\alpha `$. The number of $`0.16`$ was chosen experimentally in order to obtain the best fits to the cut-off region of the spectrum. As the distribution (2.16) represents only an approximation to the actual particle distribution, there was no reason to use points corresponding to lower densities of lesser statistical significance.
In the simulations, due to our proportional momentum scaling of the respective quantities, the derived acceleration time scale (2.4) must be also proportional to $`p`$, and thus to $`r_g(p)/c`$. Therefore, this time scale measured in units of $`r_g(p_c)/c`$ (or $`r_e(p_c)/c`$) is momentum independent and can be easily scaled to any momentum. The parameter $`T_r`$ gives the value of the acceleration time scale in units of $`r_{e,1}/c`$, i.e. $`T_{acc}^{(c)}=T_rr_{e,1}/c`$. The value of $`T_{acc,i}^{(c)}`$ at a particular time $`t_i`$ is derived from the respective values of $`p_{c,i}`$:
$$T_{acc,i}^{(c)}=\frac{p_{c,i}}{\frac{p_{c,i}p_{c,i1}}{t_it_{i1}}},$$
$`(2.17)`$
where we consider the advanced phase of acceleration ($`p_{c,i}p_0`$). As in our simulations $`p_ct`$ the condition $`(p_{c,i}p_{c,i1})/p_{c,i}1`$ is not required to hold in equation (2.17). Therefore, with all scales proportional to the particle momentum, the formula (2.17) reduces to $`T_{acc,i}^{(c)}=t_i`$ and the parameter $`T_r`$ tends to a constant (Fig. 2.4). The extension of the simulated spectra over several decades in particle energy allows to avoid problems with the initial conditions and decrease the relative error of the derived time scale by averaging over a larger number of instantaneous $`T_{acc,i}`$.
## Chapter 3 The acceleration time scale in relativistic shock waves
For a given relativistic shock velocity particle anisotropy in the shock depends on the mean magnetic field inclination to the shock normal and the form of turbulent field. Below, we describe the results of simulations performed in order to understand the time dependence of the acceleration process in various conditions. In order to do that we consider shock waves propagating with velocities $`U_1`$ = $`0.3`$, $`0.5`$, $`0.7`$ and $`0.9`$ of the velocity of light and the magnetic field inclination: $`\psi _1`$ = $`1^{}`$, $`25.8^{}`$, $`45.6^{}`$, $`60^{}`$, $`72.5^{}`$, $`84.3^{}`$ and $`89^{}`$. The first one is for a parallel shock, the last two ones are for perpendicular super-luminal shock with all velocities $`U_1`$. The intermediate values define luminal shocks ($`U_1/\mathrm{cos}\psi _1=1.0`$) at the successive velocities considered, respectively $`U_1`$ = $`0.9`$, $`0.7`$, $`0.5`$ and $`0.3`$.
In all these cases we investigate the role of varying magnitude of turbulence characterized here by the value of $`\mathrm{\Delta }t`$ or by the ratio of the diffusion coefficient across the mean field and that along the field, $`\kappa _{}/\kappa _{}`$ . The relation between these parameters for $`\mathrm{\Delta }\mathrm{\Omega }=10^{}`$ is presented in Fig. 3.1 where at $`\mathrm{\Delta }t>0.01`$ the presented relation has the power-law form $`\kappa _{}/\kappa _{}=6.310^5(\mathrm{\Delta }t)^2`$.
### 3.1 Parallel shocks
The most simple case for discussion of the first-order Fermi acceleration is a shock wave with parallel configuration of the mean magnetic field. As an example we consider the shock with negligible field inclination $`\psi _1=1^{}`$. For such a shock, the present simulations confirm the expected relation of decreasing the acceleration time scale with increasing the shock velocity and the amplitude of trajectory perturbations (Fig. 3.2). One should note at the upper panel of the figure that for short $`\mathrm{\Delta }t`$ the presented time scales decrease more and more slowly with decreasing $`\mathrm{\Delta }t`$. It is due to the fact that starting from some value of $`\mathrm{\Delta }t`$ we reach conditions of nearly isotropic diffusion, $`\kappa _{}\kappa _{}`$ and further decreasing of the time delay between scatterings decreases the acceleration time in much the same proportion as the time unit $`r_{e,1}/c`$ used to measure it (cf. Eq. 2.15, Fig. 3.1). In the lower panel of Fig. 3.2 the diffusive time scale $`T_0`$ ($`4(\kappa _{n,1}/U_1+\kappa _{n,2}/U_2)/c`$) is used as the time unit. The minute differences between the successive curves reflect the statistical fluctuations arising during simulations. Without such fluctuations all curves should coincide. The one sigma fit errors of $`T_{acc}^{(c)}`$ are indicated near the respective points. One should note that for increasing the shock velocity the acceleration time scale decreases with respect to the diffusive time scale.
### 3.2 Variation of $`T_{acc}^{(c)}`$ with magnetic field inclination
In order to compare the acceleration time scales for different magnetic field inclinations $`\psi _1`$ we performed simulations assuming a constant scattering parameters upstream and downstream yielding the same ratio of $`\kappa _{}/\kappa _{}`$ in these regions. However, due to shock compression the particle gyration period is shorter downstream than upstream in proportion to the mean magnetic field compression (Eq. 2.13). In Fig. 3.3 we present the values of the acceleration time scale derived in such conditions at different $`\psi _1`$. For super-luminal shocks the results are presented for the cases allowing for particle power-law energy spectra, i.e. when the cross-field diffusion is sufficiently effective. Actually, the spectra with inclinations $`\alpha <10.0`$ are only included.
In general, the acceleration time scale decreases with increasing field inclination, reaching in some cases the values comparable, or even smaller than the particle upstream gyroperiod (6.28 in our units of $`r_{e,1}/c`$). The trend can be reversed for intermediate wave amplitudes when the magnetic field configuration changes into the luminal and super-luminal one. Such changes are accompanied with the steepening of the spectrum (see below). The acceleration rate at different scattering amplitudes changes with $`\psi _1`$ in a way that at different inclinations the minimum acceleration times occur at different perturbation amplitudes (different $`\mathrm{\Delta }t`$).
An important feature of the acceleration process in relativistic shocks should be mentioned at this point. The variations of $`T_{acc}^{(c)}`$ in oblique shocks are accompanied by changes of the particle spectrum inclination (cf. Kirk & Heavens 1989; Ostrowski 1991). In Fig. 3.4, the curves at ($`T_{acc}^{(c)}`$, $`\alpha `$) plane represent the results for decreasing the scattering amplitude expressed with parameter $`\mathrm{\Delta }t`$, and joined with lines for the same magnetic field inclination $`\psi _1`$ . For parallel shocks the changes in $`T_{acc}^{(c)}`$ do not lead to any variation of the spectral index. However, for oblique sub-luminal ($`\psi _1`$ = $`25.8^{}`$, $`45.6^{}`$) and luminal ($`\psi _1`$ = $`60^{}`$) shocks a non-monotonic behavior is seen. The trend in changing $`T_{acc}^{(c)}`$ and $`\alpha `$ observed at smaller perturbation amplitudes (larger $`\mathrm{\Delta }t`$) is reversed at larger amplitudes when the substantial cross-field diffusion is possible.
For oblique shocks (cf. Fig. 3.5) we observe an analogous reduction of the acceleration time scale as that reported by Naito & Takahara (1995) with the pitch angle diffusion model allowing for a more rapid acceleration than the large angle scattering model. Of course this agreement is broken for short $`\mathrm{\Delta }t`$, where the cross field diffusion can not be neglected and the particle magnetic momentum is not conserved at interactions with the shock.
### 3.3 Variation of $`T_{acc}^{(c)}`$ with varying turbulence levels
In a parallel shock the acceleration time scale reduces with the increased turbulence level in it’s neighborhood. This phenomenon, well known for nonrelativistic shocks (cf. Lagage & Cesarsky 1983), is confirmed here for relativistic shock velocities (Fig. 3.2). In general, there are two main reasons for this change. The first one is a simple reduction of the diffusion time of particles outside the shock due to shorter intervals between scatterings analogous to the decrease observed in nonrelativistic shocks. However, the increased amount of scattering influences also the acceleration process due to changing (decreasing) the particle anisotropy at the shock and thus,
modifying the mean energy gain of particles interacting with the shock discontinuity. Additionally, in oblique shocks the upstream-downstream transmission probability may increase. One should note that the present approach is not able to describe fully the effect of decreasing anisotropy with the small amplitude random scattering model applied. It is due to the fact that correlations between the successive modifications of a trajectory (a sequence of small angle scattering acts in this paper) in a single MHD wave cannot be accurately modeled within the simplified approach used. A more exact approach requires integration along the particle trajectories in realistic configurations of the magnetic field. However, the comparison of the present simplified method to the one involving such an integration shows a reasonably good agreement (Ostrowski 1993) suggesting that averaging over realistic trajectories is equivalent in some way to such averaging within our random scattering approach.
In shocks with oblique magnetic fields a non-monotonic change of the acceleration time scale with the amount of scattering along the particle trajectory is observed (Fig. 3.6, see also Fig. 3.3; cf. Ostrowski 1991 for the spectral index). Increasing the amount of turbulence up to some critical amplitude decreases the diffusion time along the magnetic field and thus $`T_{acc}^{(c)}`$. However, as the mean diffusion time outside the shock is related to the normal diffusion coefficient<sup>1</sup><sup>1</sup>1One should note that for the relativistic shocks, due to particle anisotropy, the respective relation may be not so simple as that given in equation (2.2) for nonrelativistic shocks. $`\kappa _n`$ ($`\kappa _{n,i}=\kappa _{,i}\mathrm{cos}^2\psi _i+\kappa _{,i}\mathrm{sin}^2\psi _i`$, $`i`$ $`=`$ $`1`$, $`2`$), the increasing $`\kappa _{}`$ will lead, for large scattering amplitudes to
longer $`T_{acc}`$ in units of $`r_e/c`$. In the units of $`T_0`$ the acceleration time depends only weakly on the turbulence level and shows a small maximum for the minimum at the presented figure. For super-luminal shocks one can note the absence of data points corresponding to low turbulence levels, where the power-law spectrum cannot be formed or it is extremely steep. In these excluded cases, the upstream population of energetic particles is only compressed at the shock with the characteristic upstream time of $`r_{e,1}/U_1`$ (cf. Begelman & Kirk 1990; Ostrowski 1993).
## Chapter 4 Ultrarelativistic shock waves
In the present chapter we discussed several aspects of the first order acceleration process active at ultrarelativistic ($`\gamma 1`$) shock waves. These results are partly published in Bednarz & Ostrowski (1998, 1999) and in Bednarz (1999).
Below the downstream magnetic field is derived for the relativistic shock with the compression $`R`$ obtained with the formulas of Heavens & Drury (1988) for a cold ($`e`$, $`p`$) plasma – $`R3.6`$ for our smallest value of $`\gamma =3`$ and tends to $`R=3`$ for $`\gamma >>1`$, as measured in the shock rest frame.
### 4.1 Acceleration mechanism
A particle crossing the shock to upstream medium has a momentum vector nearly parallel to the shock normal. Then the particle momentum changes its inclination in two ways by: 1) scattering in an inhomogeneous magnetic field and 2) smooth variation in a homogeneous field component. Hereafter, the mean deflection angle in these two cases will be denoted by $`\mathrm{\Delta }\mathrm{\Omega }_S`$ and $`\mathrm{\Delta }\mathrm{\Omega }_H`$, respectively. The first process is a diffusive one and the second depends on time linearly. That means that with increasing shock velocity, keeping other parameters constant, $`\mathrm{\Delta }\mathrm{\Omega }_S`$ decreases slower as a square root of time in comparison with $`\mathrm{\Delta }\mathrm{\Omega }_H`$. The Lorentz transformation shows that with $`\gamma 1`$ even a tiny angular deviation in the upstream plasma rest frame can lead to a large angular deviation in the downstream plasma rest frame. Let us denote a particle phase by $`\phi `$ and the angle between momentum and a magnetic field vector by $`\vartheta `$ both measured in the downstream plasma rest frame. Values of these parameters at the moment when a particle crosses the shock downstream determine if it is able to reach the shock again in the case of neglected magnetic field fluctuations downstream of the shock. In fact a motion in the homogeneous magnetic field carries a particle in such a way that in most cases it cannot reach the shock again. The magnetic field fluctuations perturbing the momentum direction lead to broadening the ($`\phi ,\vartheta `$) range that allows particles to reach the shock again. Thus, as we show below for efficient scattering, when $`\mathrm{\Delta }\mathrm{\Omega }_H`$ becomes unimportant in comparison to $`\mathrm{\Delta }\mathrm{\Omega }_S`$, the spectral index and the acceleration time reach their asymptotic values. The discussed relation between $`\mathrm{\Delta }\mathrm{\Omega }_H`$ and $`\mathrm{\Delta }\mathrm{\Omega }_S`$ is reproduced in our simulations and presented in Fig. 4.1. There are shown 11 points from $`\gamma `$ = 100 to 320 and three other for $`\gamma `$ = 640, 1280, 2560. The expected linear dependence of these quantities can be noticed.
### 4.2 Energy spectra
Particle spectral indices were derived for different mean magnetic field configurations, measured by the magnetic field inclination $`\psi `$ with respect to the shock normal in the upstream plasma rest frame and for different amounts of turbulence measured by $`\kappa _{}/\kappa _{}`$. In the simulations we considered a few configurations of the upstream magnetic field with inclinations with respect to the shock normal being $`\psi `$ = $`0^{}`$, $`10^{}`$, $`20^{}`$, $`30^{}`$, $`60^{}`$ and $`90^{}`$. The first case represents the parallel shock, the second is for the oblique shock - subluminal at $`\gamma =3`$ and a superluminal one at larger $`\gamma `$, and the larger $`\psi `$ are for superluminal perpendicular shocks for all considered velocities. We applied the same patterns upstream and downstream of the shock for the fluctuation levels $`\lambda =5.34,4.39,3.44,2.49,1.56,0.67,0.16,0.00`$ ($`\lambda \mathrm{log}_{10}\kappa _{}/\kappa _{}`$).
In successive panels in Fig. 4.2 the energy spectral indices, $`\sigma `$, for varying $`\psi `$ and $`\kappa _{}/\kappa _{}`$ are presented. For a parallel shock ($`\psi =0^{}`$) the amount of scattering does not influence the spectral index and for the growing $`\gamma `$ it approaches $`\sigma _{\mathrm{}}2.2`$. One may note that essentially the same limiting value was anticipated for the large-$`\gamma `$ parallel shocks by Heavens & Drury (1988). Let us remember that the results for $`\psi =10^{}`$ are for superluminal shocks if $`\gamma >5.75`$. In this case, when we go from the ‘slow’ $`\gamma =3`$ shocks to higher $`\gamma `$ ones, at first the spectrum inclination increases ($`\sigma `$ grows) but at large $`\gamma `$ the spectrum flattens to approach the asymptotic value close to 2.2. The spectrum steepening phase is more pronounced for small amplitude perturbations (small $`\kappa _{}/\kappa _{}`$), but even at very low turbulence levels the final range of the spectrum flattening is observed. For larger $`\psi `$ the situation does not change considerably, but the phase of spectrum steepening is wider involving larger values of $`\sigma `$ and starting at smaller velocities below the lower limit of our considerations (there may be no such range involving the steepening phase if the required velocity is below the sound velocity).
The spectral indices for different magnetic field inclinations, but for the same value of $`\mathrm{log}_{10}(\kappa _{}/\kappa _{})=3.44`$ are presented at Fig. 4.3. The large spectral indices occurring in the steepening phase are usually interpreted as a spectrum cutoff. In this case the main factor increasing the particle energy density is a nonadiabatic compression in the shock (Begelman & Kirk 1990).
The particle angular distributions $`F(\mu )`$ in the $`\gamma 1`$ shocks can be extremely anisotropic when considered in the upstream plasma rest frame. However, when presented in the shock rest frame the distribution is always ‘mildly’ anisotropic. This feature is illustrated in Fig. 4.4 when $`\gamma `$ equals 3 or 27 (note that in Figs. 4.4 - 4.7 the area below each curve is normalized to 100). In the simulations we observed an interesting phenomenon accompanying previously discussed spectrum convergence to the limiting inclination: spectra close to the limit exhibit similar angular distributions at the shock as measured in the shock rest frame (Fig. 4.5).
Again, this feature is independent of the background conditions, and the difference between the actual angular distribution and the limiting one reflects the difference between the spectral index $`\sigma `$ and $`\sigma _{\mathrm{}}`$ (cf. Fig. 4.6). For parallel shocks with $`\gamma 9`$ where the spectral index is essentially constant $`\sigma =\sigma _{\mathrm{}}`$ this distribution is independent of the value of $`\gamma `$ and the perturbation amplitude $`\kappa _{}/\kappa _{}`$ (Fig. 4.7).
For large $`\gamma `$ shocks we observe the convergence of the derived energy spectral indices to the value $`\sigma _{\mathrm{}}2.2`$ independent of the background conditions. This unexpected result provides a strong constraint for the acceleration process in large $`\gamma `$ shocks and it can be quantitatively explained with arguments presented at the beginning of this chapter. Our interesting finding do not fully explained with such simple arguments is of the belief that the resulting spectral index is the same for oblique and parallel shocks. Our derivations are limited to the test particle approach. However, as the obtained spectra are characterized with $`\sigma >2.0`$ any nonlinear back reaction effects are not expected to affect the acceleration process within the spectrum high energy tail with $`\sigma \sigma _{\mathrm{}}`$.
### 4.3 The acceleration time scale
In the following simulations we consider shocks with $`\gamma =`$ 20, 40, 80, 160, 320, magnetic field inclinations $`\psi =`$ $`15^{}`$, $`30^{}`$, $`45^{}`$, $`60^{}`$, $`75^{}`$, $`90^{}`$ and downstream values of magnetic field fluctuations $`\tau _2=0`$, $`1.010^3`$, $`1.110^2`$, $`0.11`$, $`0.69`$.
Simulations prove that fluctuations upstream of the shock (measured by $`\tau _1`$, $`\tau \kappa _{}/\kappa _{}`$) and downstream of the shock (measured by $`\tau _2`$) influence the acceleration process independently. The minimum fluctuations upstream of the shock needed to run the acceleration process efficiently tend to zero when $`\gamma \mathrm{}`$. We checked by simulations with different $`\tau _2`$ that its value does not influence the spectral index considerably for any given $`\tau _1`$ with exception of only the injection phase of the upstream isotropic distribution.
Thus, as a first case we consider downstream conditions without magnetic field fluctuations. By simple data inspection (cf. Fig. 4.8) we look for minimum $`\tau _1`$ where the spectral index reaches its limit of 2.2 and we apply this value in further simulations. The relation between $`\tau _1`$, $`\gamma `$ and $`\psi `$ can be roughly fitted with the equation $`\tau _1=0.25\gamma ^{1.2}\psi `$ in the considered range of shock parameters. We repeated simulations for a number of cases with different $`\gamma `$ and $`\psi `$ and $`\tau _20`$. The obtained results are in good agreement with the ones derived from the above equation up to $`\tau _2=0.11`$.
Values of the acceleration time $`t_{acc}`$ for three amplitudes of magnetic field fluctuations downstream of the shock are presented in Fig. 4.9. In the figure one can see the lack of change of $`t_{acc}`$ with $`\psi `$ but it slowly decreases to the asymptotic value with $`\gamma `$. In the simulations we have observed the tendency of $`t_{acc}`$ to grow when $`\sigma `$ increases up to 2.3-2.4 and no further change if magnetic field fluctuations upstream of the shock grow. For $`\tau _20.11`$ the asymptotic value of the acceleration time is close to $`r_g/c`$. It occurs that $`r_g/c`$ is a good unit provided that the homogeneous magnetic field dominates the randomly component. Unfortunately, when this condition fails the meaning of $`t_{acc}`$ becomes unclear in the simulations then. For this reason we will not discuss further the case of $`\tau _2`$=0.69.
Approximate calculations of Gallant & Achterberg (1999) showed that $`t_U^U/t_U^D1`$, where $`t_U^U`$ is the particle mean residence time upstream of the shock (upper index) as measured in the upstream plasma rest frame (lower index), and D in $`t_U^D`$ stands for the downstream residence time. However, they were not able to consider the anisotropic particle momentum distribution and our results in Fig. 4.10 transformed to the upstream plasma rest frame with $`t_U^U/t_U^D`$ within the range $`0.010.1`$ are more adequate for real situations. Additionally, the above authors applied an extremely irregular magnetic field upstream of the shock represented by randomly oriented magnetic cells with field amplitude $`B`$ and they measured time in the upstream unit of $`r_g(B)/c`$. As a result they obtained that $`t_U^U/t_U^D`$ could be much larger than 1 in the case.
Just before the spectral index reaches its minimal value (cf. Fig. 4.8) $`\mathrm{\Delta }\mathrm{\Omega }_S`$ stabilizes near the limit which value does not further depend on the magnetic field inclination as is seen in Fig. 4.11. Momentum vectors of particles crossing downstream of the shock have similar distributions as measured in the downstream plasma rest frame if $`\mathrm{\Delta }\mathrm{\Omega }_S`$ approaches the maximum value. Then, it follows that parameters we consider below depend only on $`\tau _2`$.
For growing $`\tau _2`$ ($`\tau _2=0,1.010^3,1.110^2,0.11`$) <sup>1</sup><sup>1</sup>1Below, we provide the respective series of a given simulated parameter for this sequence of $`\tau _2`$. the acceleration time is constant and accompanied by a slow increase of the mean energy gain in one cycle downstream-upstream-downstream $`\mathrm{\Delta }E/E_D=0.89,\mathrm{\hspace{0.17em}0.94},\mathrm{\hspace{0.17em}1.0},\mathrm{\hspace{0.17em}1.1}`$, and a slight decrease of the fraction of particles that reach the shock again after crossing it downstream, $`\mathrm{\Delta }n/n=0.51,\mathrm{\hspace{0.17em}0.50},\mathrm{\hspace{0.17em}0.48},\mathrm{\hspace{0.17em}0.44}`$. Simultaneously the mean time a particle spends downstream of the shock grows as $`t_D^D=0.96,\mathrm{\hspace{0.17em}1.0},\mathrm{\hspace{0.17em}1.2},\mathrm{\hspace{0.17em}1.35}`$. Time that a particle spends upstream of the shock can be neglected in this rest frame as is visible in Fig. 4.10. It implies, approximately,
$$t_{acc}=t_D^D/\mathrm{\Delta }E/E_D$$
$`(4.1)`$
if one neglects correlations between these quantities (cf. Fig. 4.12). Similarly we can roughly estimate the value of the energy spectral index of accelerated particles as
$$\sigma 1\mathrm{ln}(\mathrm{\Delta }n/n)/\mathrm{ln}(\mathrm{\Delta }E/E_D+1).$$
$`(4.2)`$
### 4.4 The acceleration through particle reflection
Particles with an initial momentum $`p_0`$ taken as the momentum unit, $`p_0=1`$, were injected at the distance of $`2r_g`$ ($`r_g`$ \- particle gyroradius) upstream of the shock front. For all particles we derived their trajectories until crossing the shock downstream, and then upstream, or were advected with the downstream plasma, to reach a distance of $`4r_g`$ downstream of the shock. For each single particle interaction with the shock the particle momentum vector was recorded so we were able to consider angular and energy distributions of such particles. We considered shocks with Lorentz factors $`\gamma =10`$, $`160`$ and $`320`$. For each shock we discussed the acceleration processes in conditions with the magnetic field inclinations $`\psi =0^{}`$, $`10^{}`$, $`70^{}`$ and with 16 values for the turbulence amplitude measured by the ratio $`\tau `$ of the cross-field diffusion coefficient $`\kappa _{}`$ to the parallel diffusion coefficient $`\kappa _{}`$. The applied values of $`\tau `$ were taken from the range of ($`3.210^6`$, $`0.95`$), approximately uniformly distributed in $`\mathrm{log}\tau `$ . In each simulation run we derived trajectories of $`510^4`$ particles with the initial momenta isotropically distributed in the upstream rest frame.
In the downstream plasma rest frame the shock moves with velocity $`c/3`$. This velocity is comparable to the particle velocity $`c`$. Therefore, from all particles crossing the shock downstream only the ones with particular momentum orientations will interact with the shock again; the remaining particles will be caught in the downstream plasma flow and advected far from the shock front. In the simulations we considered this process quantitatively. However, let us first present a simple illustration.
Large compression ratios occurring in ultrarelativistic shocks, as measured between the upstream and downstream plasma rest frames, lead for nearly all oblique upstream magnetic field configurations to the quasi-perpendicular configurations downstream of the shock. Thus, let us consider for this illustrative example a shock with a non-perturbed perpendicular downstream magnetic field distribution. Particle crossing the shock downstream with inclination to the magnetic field $`\vartheta `$ and the phase $`\phi `$ – both measured in the downstream plasma rest frame, $`\phi =\pi /2`$ for a particle velocity normal to the shock and directed downstream – will be able to cross the shock upstream only if the equation
$$\frac{c}{3}t=r_\mathrm{g}\left[\mathrm{cos}(\phi +\omega _\mathrm{g}t)\mathrm{cos}\phi \right]$$
$`(4.3)`$
has a solution at positive time $`t`$. Here $`r_\mathrm{g}=\frac{pc}{eB}\mathrm{sin}\vartheta `$ is the particle gyroradius, $`\omega _\mathrm{g}=\frac{eB}{p}`$ is the gyration frequency, and other symbols have the usual meaning.
An angular range in the space ($`\vartheta `$, $`\phi `$) enabling particles crossing the shock downstream to reach the shock again can be characterized for illustration by three values of $`\vartheta `$. Particles with $`\mathrm{sin}\vartheta =1`$ are able to reach the shock again if $`\phi `$(1.96, 3.48), with $`\mathrm{sin}\vartheta =0.5`$ if $`\phi `$(2.96, 3.87) and with $`\mathrm{sin}\vartheta =1/3`$ only for $`\phi =4.71`$. That means that all particles with $`\phi `$ smaller than 1.96 (Fig. 4.13) are not able to reach the shock again if fluctuations of the magnetic field downstream of the shock are not present.
For perturbed magnetic fields some downstream trajectories starting in the ($`\vartheta `$, $`\phi `$) plane outside the reflection range can be scattered toward the shock to cross it upstream. We prove it by simulations presented in Fig. 4.14. One may observe that increasing the perturbation amplitude leads to an increased number of reflected particles reaching $`13`$% in the limit of $`\tau =1`$. For large magnetic field fluctuations the mean relative energy gains of reflected particles are close to $`1.2\gamma ^2`$ for the shock Lorentz factors considered. One may note a variation of the energy gain with growing $`\tau `$. The points resulting from simulations for the smallest values of $`\tau `$ were not included in Fig. 4.15 because of the small number of reflected particles (cf. Fig. 4.14).
## Chapter 5 Summary
We performed Monte Carlo simulations for shock waves with parallel and oblique (either, sub-luminal and super-luminal) magnetic field configurations with different amounts of scattering along particle trajectories. Field perturbations with amplitudes ranging from very small ones up to $`\delta BB`$ are considered.
In chapter 2 (cf. Bednarz & Ostrowski 1996) we demonstrate the existence of correlation between particle energy gains and its diffusive times. The analogous correlation is expected for the probability of particle escape downstream the shock. Therefore, for defining the acceleration time scale we use the rate of change of the spectrum cut-off momentum which accommodate all such correlations.
Acceleration times scales in relativistic shocks are discussed in chapter 3 (cf. Bednarz & Ostrowski 1996). In parallel shocks $`T_{acc}^{(c)}`$ diminishes with the growing perturbation amplitude and the shock velocity. However, it is approximately constant for the increasing turbulence level if we use the respective diffusive time scale as the time unit. Another feature discovered in oblique shocks is that due to the cross-field diffusion $`T_{acc}^{(c)}`$ can change with $`\delta B`$ in a non-monotonic way. The acceleration process leading to the power-law particle spectrum in a super-luminal shock is possible only in the presence of large amplitude turbulence. Then, the shorter acceleration times occur when the perturbations’ amplitudes are smaller and the respective spectra steeper. We discussed the coupling between the acceleration time scale and the particle spectral index in oblique shock waves with various field inclinations and revealed a possibility for non-monotonic relations of these quantities. The shortest acceleration time scales seen in the simulations are below the particle gyroperiod upstream of the shock. These times do not require the ultrarelativistic shock velocities, but may occur in mildly relativistic ones with the quasi-perpendicular magnetic field configuration. One should note that due to the larger magnetic field downstream of the shock in this short time the particle trajectory can follow a few revolutions near the shock with only a short section of each one penetrating the upstream region.
The presented estimates of the acceleration time scale provide an interesting possibility for modeling shock waves in the conditions where the electron spectrum cut-off energy is determined by the balance of gains and losses. If one is able to derive the respective acceleration rate from the knowledge of the energy loss process and the particle spectral index is also known then both these values provide constraints for the acceleration process which could be further used to reduce the parameter space available for the considered shock wave (cf. Fig. 3.4).
In chapter 4 (cf. Bednarz & Ostrowski 1998, 1999; Bednarz 1999) we discussed the acceleration mechanism that holds in ultrarelativistic shocks. We discovered convergence of spectral indices to the universal asymptotic value $`\sigma _{\mathrm{}}2.2`$ and we considered high particle anisotropy that accompanies particle acceleration. The simulations yielded that acceleration time derived from formula applied in ultrarelativistic shocks approximately equals the time derived from the formula neglecting correlations (4.1) and the constant value of 1.0 can be used for $`\tau 0.11`$.
The presented results are to be applied in models of GRB sources involving ultrarelativistic shock waves (cf. Bednarz 1999). One should note that the mean downstream plasma proton energies can reach there several tens of GeV (cf. Paczyński & Xu 1994) and the lower limit of the considered cosmic ray energies has to be larger than this scale. For shocks propagating in (e<sup>-</sup>, e<sup>+</sup>) plasma the involved thermal energies are lower, $`\gamma `$ MeV. These estimates provide the respective lower limits for the accelerated cosmic ray particles. For the physical conditions considered in GRB sources the acceleration process can provide particles with much larger energies limited only by the condition that the energy loss processes (radiative, or due to escape) are ineffective in the downstream gyroperiod time scale. We note a striking coincidence of our limiting spectral index with the value derived for energetic electrons from gamma-burst afterglow observations. Waxman (1997) used a fireball model of GRBs and showed from the functional dependence of the flux on time and frequency that $`\sigma =2.3\pm 0.1`$ in the afterglow of GRB 970228. Galama et al. (1998) made two independent measurements of the electron spectrum index in the afterglow of GRB 970508 which was very close to $`2.2`$.
In the end we have shown that efficiency of ‘$`\gamma ^2`$’ reflections in ultrarelativistic shock waves strongly depends on fluctuations of magnetic field downstream of the shock. In the most favorable conditions with high amplitude turbulence downstream the shock the reflection efficiency is a factor of 10 or more smaller than the values assumed by other authors. Moreover, due to the magnetic field compression at the shock we do not expect the required large values of $`\kappa _{}/\kappa _{}`$ to occur behind the shock (cf. a different approach of Medvedev & Loeb 1999). Therefore, with the actual efficiency of 1 - 10 % there is an additional difficulty for models postulating UHE particle acceleration at GRB shocks (cf. Gallant & Achterberg 1999). Let us note, however, that the mean downstream trajectory of the reflected particle involves only a fraction of its gyroperiod. Thus the presence of compressive long waves in this region leading to non-random trajectory perturbations could modify our estimates.
## Chapter 6 Bibliography
Appl, S., Camenzind, M., 1988, Astron. Astrophys. 206, 258
Axford W. I., Leer E., Skadron G., 1977, in Proc. 15th Int. Cosmic Ray Conf. 11, p. 132, Plovdiv
Ballard K. R., Heavens A. F., 1991, Mon. Not. R. Astron. Soc. 251, 438
Ballard K. R., Heavens A. F., 1992, Mon. Not. R. Astron. Soc. 259, 89
Bednarz J., 1998, in Proc. 16th European Cosmic Ray Symposium, p. 287, Alcalá de Henares
Bednarz J., 1999, submitted
Bednarz J., Ostrowski M., 1996, Mon. Not. R. Astron. Soc. 283, 447
Bednarz J., Ostrowski M., 1997a, in Proc. Conf. Relativistic Jets in AGN’s, p. 158, Cracow
Bednarz J., Ostrowski M., 1997b, in Proc. 25th Int. Cosmic Ray Conf. 4, p. 385, Durban
Bednarz J., Ostrowski M., 1998, Phys. Rev. Lett. 80, 3911
Bednarz J., Ostrowski M., 1999, Mon. Not. R. Astron. Soc. 310, L11
Begelman M. C., Kirk J. G., 1990, Astrophys. J. 353, 66
Bell A. R., 1978a, Mon. Not. R. Astron. Soc. 182, 147
Bell A. R., 1978b, Mon. Not. R. Astron. Soc. 182, 443
Blandford R. D., Königl A., 1979, Astrophys. J. 232, 34
Blandford R. D., Ostriker J. P., 1978, Astrophys. J. 221, L29
Decker R. B., 1988, Space Sci. Rev. 48, 195
de Jager O. C., Harding A. K., Michelson P. F., Nel H. J., Nolan P. L., Sreekumar P., Thompson D. J., 1996, Astrophys. J. 457, 253
de Hoffman F., Teller E., 1950, Phys. Rev. 80, 692
Dermer C. D., 1992, Phys. Rev. Lett. 68, 1799
Chandrasekhar S., 1943, Rev. Mod. Phys. 15, 1
Chiueh T., Li Z., Begelman M. C., 1998, Astrophys. J. 505, 835
Ellison D. C., Jones F. C., Reynolds S. P., 1990, Astrophys. J. 360, 702
Fermi E., 1949, Phys. Rev. 75, 1169
Galama T. J., Wijers R. A. M. J., Bremer M., Groot P. J., Strom R. G., de Bruyn A. G., Kouveliotou C., Robinson C. R., van Paradijs J., 1998, Astrophys. J. 500, L101
Gallant Y. A., Achterberg A., 1999, Mon. Not. R. Astron. Soc. 305, L6
Gallant Y. A., Achterberg A., Kirk J. G., 1998, in Proc. 16th European Cosmic Ray Symposium, p. 371, Alcalá de Henares
Gallant Y. A., Arons J., 1994, Astrophys. J. 435, 230
Heavens A., Drury L’O. C., 1988, Mon. Not. R. Astron. Soc. 235, 997
Hester J. J. et al., 1995, Astrophys. J. 448, 240
Hjellming R. M., Rupen M. P., 1995, Nature 375, 464
Hoyle F., 1960, Mon. Not. R. Astron. Soc. 120, 338
Hudson P. D., 1965, Mon. Not. R. Astron. Soc. 131, 23
Hudson P. D., 1967, Mon. Not. R. Astron. Soc. 137, 205
Jiang D. R., Cao X., Hong X., 1998, Astrophys. J. 494, 139
Jokipii J. R., 1971, Rev. Geophys. Space Phys. 9, 27
Jones F. C., Ellison D. C., 1991, Space Sci. Rev. 58, 259
Kennel C. F., Coroniti F. V., 1984, Astrophys. J. 283, 694
Kirk J. G., 1988, Habilitation Theses, preprint No. 345, Max-Planck-Institut für Astrophysik, Garching
Kirk J. G., Heavens A., 1989, Mon. Not. R. Astron. Soc. 239, 995
Kirk J. G., Schneider P., 1988, Astron. Astrophys. 201, 177
Kirk J. G., Schneider P., 1987a, Astrophys. J. 315, 425
Kirk J. G., Schneider P., 1987b, Astrophys. J. 322, 256
Königl A., 1981, Astrophys. J. 243, 700
Krymsky G. F., 1977, Dokl. Akad. Nauk SSSR 234, 1306, (Engl. transl. Sov. Phys.-Dokl. 23, 327)
Kulkarni S. R. et al., 1998, Nature 393, 35
Lagage P. O., Cesarsky C., 1983, Astron. Astrophys. 125, 249
Lieu R., Quenby J. J., Drolias B., Naidu K., 1994, Astrophys. J. 421, 211
Lucek S. G., Bell A. R., 1994, Mon. Not. R. Astron. Soc. 268, 581
Medvedev M. V., Loeb A., 1999, Astrophys. J. submitted (astro-ph/9904363)
Meegan C. A., Fishman G. J., Wilson R. B., Paciesas W. S., Pendleton G. N., Ho- rack J. M., Brock M. N., Kouveliotou C., 1992, Nature 355, 143
Mirabel I. F., Rodriguez, L. F., 1994, Nature, 371, 46
Naito T., Takahara F., 1995 Mon. Not. R. Astron. Soc. 275, 1077
Newman P. L., Moussas X., Quenby J. J., Valdes-Galicia J. F., Theodossiou-Ekateri- nidi Z., 1992, Astron. Astrophys. 255, 443
Paczyński B., Xu G., 1994, Astrophys. J. 427, 708
Parker E. N., 1958, Phys. Rev. 109, 1328
Quenby J. J., Drolias B., 1995, in Proc. 24th Int. Cosmic Ray Conf. 3, p. 261, Rome
Quenby J. J., Lieu R., 1989, Nature 342, 654
Ostrowski M., 1988, Mon. Not. R. Astron. Soc. 233, 257
Ostrowski M., 1991, Mon. Not. R. Astron. Soc. 249, 551
Ostrowski M., 1993, Mon. Not. R. Astron. Soc. 264, 248
Ostrowski M., 1994, Comments Astrophys. 17, 207
Ostrowski M., Schlickeiser R., 1996, Sol. Phys. 167, 381
Schatzman E., 1963, Ann. Astrophys. 26, 234
Takahara F., Terasawa T., 1990, in Proc. ICRR Int. Symp. on ”Astrophysical Aspects of the Most Energetic Cosmic Rays”, Kofu
Tingay S. J. et al., 1995, Nature 374, 141
Vermeulen R. C., Cohen M. H., 1994, Astrophys. J. 430, 467
Vietri M., 1995, Astrophys. J. 453, 883
von Montigny C. et al., 1995, Astrophys. J. 440, 525
Waxman E., 1995, Phys. Rev. Lett. 75, 386
Waxman E., 1997, Astrophys. J., 485, L5
Woods E., Loeb A., 1995, Astrophys. J. 453, 583
|
warning/0005/gr-qc0005068.html
|
ar5iv
|
text
|
# SPACETIMES ADMITTING A 3-PARAMETER SIMILARITY GROUP
## 1 Introduction
Spacetimes admitting (local) groups of (local) isometries have been widely studied, especially in those cases where the dimension of the Lie algebra of Killing vectors fields (KV) is high or where there is a non-trivial isotropy subgroup. On the other hand, Hall and Steele have investigated those spacetimes admitting an r-dimensional Lie algebra of homothetic vector fields (HVF) (which gives rise to an r-parameter group of similarities); solving the problem completely in those cases where $`r6`$ and also giving some general results when $`r5`$. The purpose of this paper is, up to a certain extent, to complement the study carried out in the above reference, especially in the case $`r=3`$, which could be of interest in Cosmology (see, for instance and references cited therein).
Throughout this paper $`(M,g)`$ will denote a spacetime: $`M`$ then being a (smooth) Hausdorff, simply connected 4-dimensional manifold, and $`g`$ a (smooth) Lorentz metric with signature (-+++). A semicolon will denote a covariant derivative with respect to the metric connection associated with $`g`$, and a comma will denote a partial derivative as usual. A global vector field $`X`$ on $`M`$ is called homothetic if either of the two equivalent conditions
$`_Xg_{ab}X_{a;b}+X_{b;a}=2\lambda g_{ab}`$
$`X_{a;b}=\lambda g_{ab}+F_{ab}(F_{ab}=F_{ba})`$ (1)
holds on a local chart, where $`\lambda `$ is a constant on $`M`$, $`F`$ is the homothetic bivector, and $``$ denotes the Lie derivative operator. If $`\lambda 0`$, $`X`$ is called proper homothetic and it can always be scaled so as to have $`\lambda =1`$, if $`\lambda =0`$ then $`X`$ is a KV on $`M`$. For a geometrical interpretation of (1) we refer the reader to .
A necessary condition that $`X`$ be homothetic is
$$X_{}^{a}{}_{;bc}{}^{}=R_{}^{a}{}_{bcd}{}^{}X^d$$
(2)
where $`R_{}^{a}{}_{bcd}{}^{}`$ are the components of the Riemann tensor in a coordinate chart; thus, an HVF is a particular case of affine collineation and therefore it will satisfy
$$_XR_{}^{a}{}_{bcd}{}^{}=_XR_{ab}=_XC_{}^{a}{}_{bcd}{}^{}=0$$
(3)
where $`R_{ab}`$ ($`R_{}^{c}{}_{acb}{}^{}`$) and $`C_{}^{a}{}_{bcd}{}^{}`$ stand, respectively, for the components of the Ricci and the Conformal Weyl tensor. Also, from the Einstein’s Field equations (EFE), it follows
$$_XT_{ab}=0$$
(4)
where $`T_{ab}`$ is the energy momentum tensor representing the material content of the spacetime ($`M,g`$). It can easily be shown that whenever a proper HVF exists in a Lie algebra of HVF’s $`_r`$, this necessarily contains an $`(r1)`$-dimensional Lie subalgebra of KV $`𝒢_{r1}`$; therefore one can always choose a basis for $`_r`$ in such a way that it contains at most one proper HVF, the $`r1`$ remaining ones thus being KV’s. If these vector fields in the basis of $`_r`$ are all complete vector fields, then $`_r`$ gives rise in a well known way to a Lie group of homotheties; otherwise, it gives rise to a local group of local homothetic transformations of $`M`$ and, although the usual concepts of isotropy and orbits still hold, a little more care is required. We shall not go into the details here, but refer the reader to for further information on this particular issue. It is also immediate to see from (1) that the Lie bracket of a proper HVF and a KV is a KV. Further information on the isotropy structure as well as on the fixed point structure of $`H_r`$ can be found in (see also ).
In this paper we shall be concerned with spacetimes admitting a 3-parameter group of homotheties acting on non-null orbits, providing a classification of all possible Lie algebra structures (in terms of the Bianchi type of $`_3`$), and giving in each case the form of the metric as well as that of the proper HVF and the two KV’s, in terms of local coordinates. We shall also present a few selected examples of spacetimes, satisfying the above properties, which can represent perfect fluid cosmological models, as well as vacuum solutions.
The paper is organized as follows: Section 2 contains a brief summary of results on groups of homotheties and spacetimes admitting them. In section 3 we present the Bianchi classification of Lie algebras $`_3`$ along with some remarks on the topology of the orbits of the corresponding Lie subalgebras $`𝒢_2`$. The general form of the metric is provided in each case. Finally, section 4, contains a few selected examples of perfect fluid spacetimes admitting a maximal $`H_3`$; some of them we believe are new and, whenever this is possible, we relate our results to those already existing in the literature.
## 2 Preliminary Results
In this section we provide (without proof) some general results regarding spacetimes admitting Lie groups of HVF. In what is to follow we shall be assuming that a proper HVF exists; $`r`$ will then denote the dimension of the Lie algebra of HVF $`_r`$ and $`𝒢_{r1}`$ will be its associated Killing subalgebra. Most of the following results, together with their proofs, can be found in ref. .
1. The orbits of $`_r`$ and $`𝒢_{r1}`$ can only coincide if they are 4-dimensional or 3-dimensional and null. Thus, the case in which we will be interested mainly ($`r=3`$, non-null orbits) corresponds to a transitive action of the homothety group (and therefore, the dimensions of the orbits of $`_3`$ and $`𝒢_2`$ will be 3 and 2 respectively).
2. If $`r=11`$ then $`M`$ is flat. The cases $`r=10`$, $`9`$ are impossible, as it follows from consideration of the dimension of $`𝒢_{r1}`$. $`r=8`$ corresponds to $`M`$ being a conformally flat, homogeneous generalized plane wave . The case $`r=7`$ implies that $`M`$ is a type $`N`$, homogeneous or a conformally flat non-homogeneous generalized plane wave or one of the special Robertson-Walker spacetimes (or their equivalent, with Segre type $`\{(1,11)1\}`$ for the Ricci tensor). $`r=6`$ implies that $`M`$ is a type $`N`$, non-homogeneous, generalized plane wave. In the $`r=5`$ case, the associated $`𝒢_4`$ subalgebra has necessarily 3-dimensional non-null orbits, the Petrov type being $`D`$, $`N`$ or $`O`$ for timelike Killing orbits, and $`D`$ or $`O`$ for spacelike ones.
3. If $`r=4`$ and a multiply transitive action is assumed, then $`_4`$ and $`𝒢_3`$ have respectively 3-dimensional and 2-dimensional orbits.
4. Spacetimes admitting a 4-parameter group of homotheties acting transitively on $`M`$ were studied by Rosquist and Jantzen (see also where a thorough study of vacuum Bianchi $`I`$ solutions admitting a proper HVF is carried out).
5. The case $`r=3`$ has an associated Killing subalgebra $`𝒢_2`$ and the respective dimensions of their orbits are 3 and 2 (see remark (1) above). In this case one can classify the Lie algebras $`_3`$ according to their Bianchi type (see for example ); the only possible types being those corresponding to soluble groups ($`I`$ to $`VII`$ in the previous reference), as it follows from the fact that $`_3`$ must contain a 2-dimensional subalgebra $`𝒢_2`$, which in all cases but one, turns out to be abelian. In this case (abelian $`𝒢_2`$), there are only two different topologies possible for the (non-null) orbits $`V_2`$; namely: $`V_2`$ diffeomorphic to $`^2`$, and $`V_2`$ diffeomorphic to $`𝒮^1\times `$; and it follows that in the latter case the only Bianchi type possible for $`_3`$ is $`I`$; as for the case $`V_2^2`$, all seven types can, in principle, occur.
## 3 Bianchi types of $`_3`$
The purpose of this section is to analyze the possible Bianchi types of $`_3`$, giving in each case the coordinate forms of the proper HVF and the metric tensor. We shall restrict ourselves to the case of non-null orbits, and furthermore we shall assume (as is customary) that the Killing orbits $`V_2`$ admit orthogonal 2-surfaces. We shall denote the KV’s spanning $`𝒢_2`$ by $`\xi `$ and $`\eta `$, and the proper HVF in the basis of $`_3`$ as $`X`$; also we shall treat separately the case where $`𝒢_2`$ is abelian from that where $`𝒢_2`$ is non-abelian.
### 3.1 Case $`𝒢_2`$ abelian
Under the above assumptions, the possible Bianchi types of $`_3`$ containing an abelian $`𝒢_2`$ are :
$`(I)`$ $`[\xi ,\eta ]=[\xi ,X]=[\eta ,X]=0`$
$`(II)`$ $`[\xi ,\eta ]=[\xi ,X]=0[\eta ,X]=\xi `$
$`(III)`$ $`[\xi ,\eta ]=0[\xi ,X]=\xi [\eta ,X]=0`$
$`(IV)`$ $`[\xi ,\eta ]=0[\xi ,X]=\xi [\eta ,X]=\xi +\eta `$
$`(V)`$ $`[\xi ,\eta ]=0[\xi ,X]=\xi [\eta ,X]=\eta `$
$`(VI)`$ $`[\xi ,\eta ]=0[\xi ,X]=\xi [\eta ,X]=q\eta `$
$`(VII)`$ $`[\xi ,\eta ]=0[\xi ,X]=\eta [\eta ,X]=\xi +q\eta (q^2<4)`$
Assume now that the Killing orbits $`V_2`$ are spacelike and diffeomorphic to $`^2`$; since $`\xi `$ and $`\eta `$ commute, we locally have
$$\xi =\frac{}{x}\eta =\frac{}{y}$$
(5)
taking now two more coordinates, $`t`$ and $`z`$, it follows from the assumption that the Killing orbits $`V_2`$ admit orthogonal surfaces, that the line element associated to the metric $`g`$ can be written as
$$ds^2=\mathrm{\Psi }^2\{dt^2+dz^2+s^2dy^2+b^2(Pdy+dx)^2\}$$
(6)
where $`\mathrm{\Psi }`$, $`s`$, $`b`$ and $`P`$ are all functions of $`t`$ and $`z`$ alone, their functional dependence on these coordinates to be determined (to some extent) by the HVF $`X`$ in each case ($`IVII`$).
For timelike Killing orbits, also diffeomorphic to $`^2`$, one (locally) has:
$$\xi =\frac{}{t}\eta =\frac{}{z}$$
(7)
and the metric would then read:
$$ds^2=\mathrm{\Psi }^2\{dx^2+dy^2+s^2dz^2b^2(dt+Pdz)^2\}$$
(8)
where $`\mathrm{\Psi }`$, $`s`$, $`b`$ and $`P`$ are now functions of $`x`$ and $`y`$, coordinates on the surfaces orthogonal to the Killing orbits.
The case $`V_2𝒮^1\times `$, either spacelike (cylindrical symmetry) or timelike (stationary axially symmetric metric) is easily obtained from (5) and (6) (respectively (7) and (8)) by simply changing $`y`$ (respectively $`z`$) to $`\phi `$, angular coordinate (and then one has to impose the regularity condition on the axis, ref. p.192; in order to ensure that $`\phi `$ has the standard periodicity $`2\pi `$). It is precisely the fact that the axis of rotation is an (invariant) submanifold of the spacetime manifold, what implies that -if no other CKV exists on $`M`$\- then $`_3`$ must be abelian (Bianchi type $`I`$) .
In what is to follow, and for the sake of simplicity, we shall assume that the Killing orbits $`V_2`$ are spacelike and diffeomorphic to $`^2`$, so that the forms (5) and (6) will hold for the KV’s and the line element respectively. The case of timelike Killing orbits can be formally obtained by changing $`(x,y)`$ into $`(t,z)`$ in the expression (5) for the KV’s $`\xi `$ and $`\eta `$; and
$$tixyzPiP$$
(9)
in (6), to get the line element in this case.
Now, assume $`X`$ is a proper-HVF satisfying (1) (with $`\lambda =1`$); in a coordinate chart it will have an expression of the form
$$X=X^a(x^b)_a$$
(10)
Specializing now the equation (1) to the HVF (10) and the metric (6) one sees that, assuming non-null homothetic orbits $`V_3`$, it is always possible to perform a coordinate change in the 2-spaces orthogonal to the Killing orbits; such that preserves the form of the metric and brings the HVF (10) to one of the two following forms:
$$X=_t+X^x(x,y)_x+X^y(x,y)_y$$
(11)
or
$$X=_z+X^x(x,y)_x+X^y(x,y)_y$$
(12)
Assuming the form (11) for the proper HVF $`X`$ (i.e.: 3-dimensional timelike homothetic orbits), equation (1), along with each particular Lie algebra structure, yields for every different Bianchi type ($`I`$) to ($`VII`$) the following forms for $`X`$ and the functions $`\mathrm{\Psi }`$, $`s`$, $`b`$ and $`P`$ appearing in (6),
$`(I)`$ $`X=_t\mathrm{\Psi }=e^tf(z)s=\widehat{s}(z)b=\widehat{b}(z)P=\widehat{p}(z)`$ (13)
$`(II)`$ $`X=_t+y_x\mathrm{\Psi }=e^tf(z)s=\widehat{s}(z)b=\widehat{b}(z)P=\widehat{p}(z)t`$ (14)
$`(III)`$ $`X=_t+x_x\mathrm{\Psi }=e^tf(z)s=\widehat{s}(z)b=e^t\widehat{b}(z)P=e^t\widehat{p}(z)`$ (15)
$`(IV)`$ $`X=_t+(x+y)_x+y_y\mathrm{\Psi }=e^tf(z)`$
$`s=e^t\widehat{s}(z)b=e^t\widehat{b}(z)P=\widehat{p}(z)t`$
$`(V)`$ $`X=_t+x_x+y_y\mathrm{\Psi }=e^tf(z)`$
$`s=e^t\widehat{s}(z)b=e^t\widehat{b}(z)P=\widehat{p}(z)`$
$`(VI)`$ $`X=_t+x_x+qy_y(q0,1)\mathrm{\Psi }=e^tf(z)`$
$`s=e^{qt}\widehat{s}(z)b=e^t\widehat{b}(z)P=e^{(1q)t}\widehat{p}(z)`$
$`(VII)`$ $`X=_ty_x+(x+qy)_y(q^2<4)\mathrm{\Psi }=e^tf(z)`$ (19)
$$s=\frac{e^{\frac{q}{2}t}\frac{\sqrt{4q^2}}{2}a(z)}{(\sqrt{a(z)^2+c(z)^2+g(z)^2}+c(z)\mathrm{cos}(\sqrt{4q^2}t)+g(z)\mathrm{sin}(\sqrt{4q^2}t))^{\frac{1}{2}}}$$
$$b=e^{\frac{q}{2}t}(\sqrt{a(z)^2+c(z)^2+g(z)^2}+c(z)\mathrm{cos}(\sqrt{4q^2}t)+g(z)\mathrm{sin}(\sqrt{4q^2}t))^{\frac{1}{2}}$$
$$P=\frac{q}{2}+\frac{\sqrt{4q^2}}{2}\frac{g(z)\mathrm{cos}(\sqrt{4q^2}t)+c(z)\mathrm{sin}(\sqrt{4q^2}t)}{\sqrt{a(z)^2+c(z)^2+g(z)^2}+c(z)\mathrm{cos}(\sqrt{4q^2}t)+g(z)\mathrm{sin}(\sqrt{4q^2}t)}$$
If $`P=`$constant two mutually orthogonal KV’s exist in $`𝒢_2`$ and one can always, by means of a linear change of coordinates in the Killing orbits $`V_2`$, bring $`P`$ to zero. It is worth noticing that this is not possible for families ($`II`$) and ($`IV`$). Note that in case $`VII`$, $`P=`$constant implies the existence of a third KV tangent to the Killing orbits $`V_2`$, therefore these orbits are of constant curvature and the Lie algebra of HVF is 4-dimensional.
The other ansatz for the proper HVF $`X`$, (12) corresponds to the homothetic orbits being spacelike, and would yield similar results to those above but with the role of the coordinates $`t`$ and $`z`$ reversed.
As for the case $`V_2𝒮^1\times `$, and from the previous remarks on this issue, it follows that the metric would be that given in (6) (or (8)) exchanging $`y`$ (or $`z`$) for $`\phi `$, and with $`\mathrm{\Psi }`$, $`s`$, $`b`$ and $`P`$ being those given in (13) (or their equivalent under the substitution (9) in the case of timelike Killing orbits).
### 3.2 Case $`𝒢_2`$ non-abelian
The Lie algebra structure in this case is
$$[\xi ,\eta ]=\xi [\xi ,X]=[\eta ,X]=0$$
(20)
Assuming that the theorem of Bilyalov and Defrise-Carter holds (for a precise statement of the conditions under which this happens, see ), there exists a (smooth) function $`\sigma =\sigma (x^c)`$ such that $`\xi `$, $`\eta `$ and $`X`$ span a 3-dimensional Lie algebra of KV’s in a spacetime ($`M,\widehat{g}`$) where $`\widehat{g}=e^{2\sigma }g`$; i.e.: ($`M,\widehat{g}`$) is a Bianchi type $`III`$ spacetime. One can now adapt coordinates to $`\xi `$, $`\eta `$ and $`X`$ in ($`M,\widehat{g}`$) as follows:
$$\xi =\frac{}{x^1}\eta =A\frac{}{x^3}+B\frac{}{x^1}+C\frac{}{x^2}X=\frac{}{x^3}$$
(21)
where $`A`$, $`B`$ and $`C`$ are functions of $`x^\alpha `$, $`\alpha =1,2,3`$. The commutation relations (20) imply: $`A=A(x^2)`$, $`B=x^1+B_0(x^2)`$ and $`C=C(x^2)`$; and one can then always carry out a change of coordinates $`(x^\alpha )(x^\alpha ^{})`$ so as to write $`\xi `$, $`\eta `$ and $`X`$ as:
$$\xi =\frac{}{x^1^{}}\eta =x^1^{}\frac{}{x^1^{}}+\frac{}{x^2^{}}X=\frac{}{x^3^{}}$$
(22)
dropping now the primes and choosing a new coordinate $`x^4`$; it follows that the metric $`\widehat{g}`$ can be written as
$$\widehat{g}_{ab}=\left(\begin{array}{cccc}e^{2x^2}a_{11}& e^{x^2}a_{12}& e^{x^2}a_{13}& 0\\ e^{x^2}a_{12}& a_{22}& a_{23}& 0\\ e^{x^2}a_{13}& a_{23}& a_{33}& 0\\ 0& 0& 0& ϵ\end{array}\right)$$
(23)
where $`ϵ=\pm 1`$ and $`a_{\alpha \beta }=a_{\alpha \beta }(x^4)`$. It is now immediate to find out the general form for the metric $`g`$; for from $`g=e^{2\sigma }\widehat{g}`$ along with $`_Xg=2g`$ it readily follows that $`\sigma =x^3+\sigma _0(x^4)`$ (since $`_\xi g=_\eta g=0`$); redefining now the coordinate $`x^4`$ one has
$$g_{ab}=e^{2x^3}\left(\begin{array}{cccc}e^{2x^2}A_{11}& e^{x^2}A_{12}& e^{x^2}A_{13}& 0\\ e^{x^2}A_{12}& A_{22}& A_{23}& 0\\ e^{x^2}A_{13}& A_{23}& A_{33}& 0\\ 0& 0& 0& ϵ\end{array}\right)$$
(24)
where again $`A_{\alpha \beta }=A_{\alpha \beta }(x^4)`$ and $`ϵ=\pm 1`$. The case $`ϵ=+1`$ corresponds to the 3-dimensional homothetic orbits being timelike, whereas $`ϵ=1`$ corresponds to spacelike homothetic orbits.
## 4 Examples
The aim of this section is to provide a sample of physically significant spacetimes admitting an $`H_3`$ on non-null orbits as maximal group of similarity.
### 4.1 Perfect Fluid Spacetimes
A perfect fluid spacetime ($`M,g`$) satisfies the EFE’s for an energy-momentum tensor of the form:
$$T_{ab}=(\mu +p)u_au_b+pg_{ab}$$
(25)
where $`u^a`$ is the velocity flow of the fluid ($`u^au_a=1`$), $`\mu `$ is a positive function representing the energy density as measured by an observer comoving with the fluid, and $`p`$ represents the pressure, usually satisfying an equation of state of the form $`p=p(\mu )`$ (barotropic equation of state; the fluid is then said to be isentropic; i.e.: zero density of entropy production). If a KV $`𝒳`$ exists in the spacetime, one has :
$$_𝒳u_a=_𝒳\mu =_𝒳p=0$$
(26)
and the existence of a proper HVF $`X`$ implies in turn :
$$_Xu_a=u_a_X\mu =2\mu p=(\gamma 1)\mu $$
(27)
where $`\gamma [1,2]`$ in order to comply with the energy conditions . Henceforth, all of the examples we shall present will correspond to the case of spacelike Killing orbits diffeomorphic to $`^2`$ and the choice (11) for the proper HVF. Thus, in the adapted coordinate system set up in section 3.1 (see for example (6)) we shall have:
$$\mu =e^{2t}\widehat{\mu }(z)u_t=e^tf^1\mathrm{cosh}\alpha (z)u_z=e^tf^1\mathrm{sinh}\alpha (z)$$
(28)
where $`\alpha (z)`$ is a function to be determined via the field equations. The components of the 4-velocity on the Killing orbits are zero ($`u_x=u_y=0`$) also as a consequence of the field equations. As for the remaining cases (timelike Killing orbits, Killing orbits diffeomorphic to a cylinder and/or the ansatz (12) for the proper HVF $`X`$), see remarks in the previous section concerning this. Nevertheless, no correspondence will exist -in general- between solutions obtained in all those various cases and those we will present here (the only similarity being the general form of the metrics under the changes of coordinates suggested in section 3, as it was already pointed out there).
Going back to the current case, the generic form of the spacetime metric will be that given in (6) with the functions $`\mathrm{\Psi }`$, $`s`$, $`b`$ and $`P`$ that appear in (13)-(19), depending on the particular Bianchi type we are interested in. Furthermore, the generic forms of the fluid 4-velocity and of the energy density will be those given by (28). From the expression of the velocity of the fluid, it is easy to see that, in general, it is non-geodesic, expanding and shearing; its vorticity being zero since $`u^a`$ is orthogonal to the Killing orbits.
At this point it is convenient to split up our study into two cases:
#### 4.1.1 The fluid flow velocity is tangent to the homothetic orbits.
In this case, and since we chose the coordinates $`t`$, $`x`$ and $`y`$ adapted to the homothetic orbits; it follows that $`u_z=0`$, and therefore the fluid is comoving (for this particular choice of coordinates). One then has
$$u_t=e^tf^1(z)\dot{u}_z=\frac{f^{}}{f}$$
(29)
$$\theta =e^tf(k+3)$$
(30)
where a dash indicates differentiation with respect to $`z`$, $`\theta `$ stands for the expansion of the fluid ($`\theta u_{}^{a}{}_{;a}{}^{}`$); $`k`$ is defined as $`sb=e^{kt}\widehat{s}(z)\widehat{b}(z)`$, and the remaining components of the 4-velocity $`u_a`$ and the acceleration $`\dot{u}_a`$ are zero.
From the contracted Bianchi identities it follows
$$\gamma =\frac{2}{k+3}$$
(31)
Now, from the classification of $`_3`$ into Bianchi types, we see that for families $`I`$ and $`II`$ one has $`\gamma =\frac{2}{3}`$. Such a value for $`\gamma `$ lies out of the interval permitted ($`\gamma [1,2]`$); nevertheless, it is physically significant since matter becomes attractive for $`\gamma >\frac{2}{3}`$; therefore the value $`\frac{2}{3}`$ may be of interest in inflationary models . For family $`III`$ one has $`k=1`$ and hence $`\gamma =1`$; i.e.: $`p=0`$ (dust). Families $`IV`$ and $`V`$ correspond to $`k=2`$, i.e.: $`\gamma =2`$; that is: $`p=\mu `$ (stiff matter). For family $`VI`$, $`k=(q+1)`$ (with $`q0,1`$) and thus $`\gamma =2/(2q)`$; which taking into account the permitted values of $`\gamma `$, implies that $`q(0,1)`$. For family $`VII`$, $`k=q`$ (with $`q^2<4`$) and therefore $`\gamma =2/(3q)`$.
Furthermore, it is possible to see from the field equations, that the families $`(III)`$ and $`(V)`$ admit no solutions of this type with $`\mu 0`$.
The case when $`P=0`$ (i.e.: $`𝒢_2`$ contains two mutually orthogonal KV’s) and $`u^a`$ is tangent to the homothetic orbits, has been thoroughly studied by Wainwright and collaborators in a series of papers dedicated to investigate the role of self-similarity in Cosmology. They interpret these self-similar models as asymptotic states (at late times) of more general inhomogeneous cosmological models; since they are precisely those corresponding to the equilibrium points of the EFE’s, written as an autonomous system, for orthogonally transitive $`G_2`$ Cosmologies.
From our remarks above, it follows that their solutions must be of the type $`VI`$ (type $`I`$ is ruled out since $`\gamma [1,2]`$, types $`III`$ and $`V`$ cannot admit solutions with $`\mu 0`$; and type $`VII`$ together with $`P=0`$ implies the existence of a further KV tangent to the Killing orbits $`V_2`$ and the metric would then admit a non-transitive group $`H_4`$ of homotheties).
Eardley studied the case of 3-dimensional spacelike homothetic orbits.
#### 4.1.2 The fluid flow is not tangent to the homothetic orbits (tilted case).
In this case $`u_z0`$ for our particular choice of coordinates, and consequently the expressions for the acceleration $`\dot{u}_a`$ and the expansion $`\theta `$ of the fluid become more complicated, as well as the field equations. In orthogonal transitive abelian $`G_2`$ models, it is possible, though, to perform a change of coordinates in the $`t`$, $`z`$ plane so as to bring the 4-velocity of the fluid to a comoving form, preserving the diagonal form of the induced metric there , as a consequence, the field equations can be written in a much simpler form. Since most of the solutions of these characteristics appearing in the literature are given in those coordinates, we found convenient to translate our results (6) and (13-19) to them. Obviously, the form of the proper HVF $`X`$ will change and the coordinates will no longer be adapted to the homothetic orbits; thus, we next give the equivalents of (6) and (13)-(19) in the new coordinates. Following ; the metric can now be written as:
$$ds^2=A^2dt^2+B^2dz^2+r\{f(dx+wdy)^2+f^1dy^2\}$$
(32)
where $`A`$, $`B`$, $`r`$, $`f`$ and $`w`$ are functions of $`t`$ and $`z`$ and the two (commuting) KV’s are $`\xi =\frac{}{x}`$ and $`\eta =\frac{}{y}`$ (same as before). The 4-velocity of the fluid is now:
$$u=A^1\frac{}{t},\mathrm{or}\mathrm{equivalently}u_a=(A,0,0,0)$$
(33)
we can now use the remaining coordinate freedom in the $`t`$, $`z`$ plane ($`tm(t)`$, $`zn(z)`$) to bring the (non-null) proper HVF $`X`$ satisfying ($`IVII`$) to either of the following three forms:
$`(i)`$ $`X=_t+X^x(x,y)_x+X^y(x,y)_y`$ (34)
$`(ii)`$ $`X=_z+X^x(x,y)_x+X^y(x,y)_y`$ (35)
$`(iii)`$ $`X=_t+_z+X^x(x,y)_x+X^y(x,y)_y`$ (36)
$`X^x(x,y)`$ and $`X^y(x,y)`$ being linear functions of the coordinates $`x`$ and $`y`$, to be determined for each particular algebraic Bianchi type $`I`$ to $`VII`$. Notice that ($`i`$) corresponds to $`u`$ being tangent to the homothetic orbits, and therefore it has been dealt with above. ($`ii`$) corresponds to the orbits of the homothety group being spacelike (and also $`u^aX_a=0`$); and this is the case studied by Eardley (in this case, and since $`X`$ and $`u`$ are mutually orthogonal it follows $`\gamma =2`$; i.e.: $`p=\mu `$ stiff matter). Finally, ($`iii`$) is precisely the case we are currently interested in; namely $`u`$ not tangent to the homothetic orbits.
Specializing now the equation (1) to the metric (32) and the HVF $`X`$ given by (36), we get for the metric functions
$$A^2=e^{t+z}\widehat{A}^2(tz)B^2=e^{t+z}\widehat{B}^2(tz)$$
(37)
in all seven types; and:
$`(I)`$ $`r=e^{t+z}\widehat{r}(tz)f=\widehat{f}(tz)w=\widehat{w}(tz)`$ (38)
$`(II)`$ $`r=e^{t+z}\widehat{r}(tz)f=\widehat{f}(tz)w=\widehat{w}(tz){\displaystyle \frac{t+z}{2}}`$ (39)
$`(III)`$ $`r=e^{\frac{t+z}{2}}\widehat{r}(tz)f=e^{\frac{t+z}{2}}\widehat{f}(tz)w=e^{\frac{t+z}{2}}\widehat{w}(tz)`$ (40)
$`(IV)`$ $`r=\widehat{r}(tz)f=\widehat{f}(tz)w=\widehat{w}(tz){\displaystyle \frac{t+z}{2}}`$ (41)
$`(V)`$ $`r=\widehat{r}(tz)f=\widehat{f}(tz)w=\widehat{w}(tz)`$ (42)
$`(VI)`$ $`r=e^{\frac{1q}{2}(t+z)}\widehat{r}(tz)f=e^{\frac{1q}{2}(t+z)}\widehat{f}(tz)w=e^{\frac{1q}{2}(t+z)}\widehat{w}(tz)`$ (43)
$`(VII)`$ $`r=e^{\frac{2q}{2}(t+z)}\sqrt{{\displaystyle \frac{4q^2}{4}}}\widehat{r}(tz)`$ (44)
$$f=\frac{\sqrt{\widehat{r}^2+\widehat{b}^2+\widehat{c}^2}+\widehat{b}\mathrm{cos}(\sqrt{\frac{4q^2}{4}}(t+z))+\widehat{c}\mathrm{sin}(\sqrt{\frac{4q^2}{4}}(t+z))}{\sqrt{\frac{4q^2}{4}}\widehat{r}(tz)}$$
$$w=\frac{q}{2}+\sqrt{\frac{4q^2}{4}}\frac{\widehat{b}\mathrm{sin}(\sqrt{\frac{4q^2}{4}}(t+z))\widehat{c}\mathrm{cos}(\sqrt{\frac{4q^2}{4}}(t+z))}{\sqrt{\widehat{r}^2+\widehat{b}^2+\widehat{c}^2}+\widehat{b}\mathrm{cos}(\sqrt{\frac{4q^2}{4}}(t+z))+\widehat{c}\mathrm{sin}(\sqrt{\frac{4q^2}{4}}(t+z))}$$
where $`\widehat{b}`$ and $`\widehat{c}`$ are both functions of ($`tz`$).
It is interesting to notice that all diagonal ($`w=0`$), perfect fluid solutions of the form (32) (i.e.: admitting an orthogonally transitive abelian $`G_2`$ with flat spacelike orbits) and such that the functions $`A`$, $`B`$, $`r`$ and $`f`$ are separable in the variables $`t`$ and $`z`$ are already known . Note that the only Bianchi types which can contain diagonal metrics such that, for them, the maximal isometry group is the abelian $`G_2`$ generated by $`\xi `$ and $`\eta `$, are the types $`I`$, $`III`$, $`V`$ and $`VI`$ (families $`II`$ and $`IV`$ do not contain diagonal metrics, and the diagonal, type $`VII`$ case admits a further KV). For them, the metric functions are all of the form
$$F=e^{a(t+z)}\varphi (tz),a=constant$$
(45)
and it is immediate to prove that $`F`$ is separable in $`t`$ and $`z`$ if and only if $`\varphi `$ is of the form:
$$\varphi =Ce^{k(tz)},C,k=constants$$
(46)
We next present a few solutions which have been obtained for the Bianchi types $`I`$, $`III`$ and $`V`$ assuming $`w=0`$ (diagonal), but which are not separable in the above sense.
Type $`I`$
$$ds^2=\frac{e^{t+z}}{f_{o}^{}{}_{}{}^{2}|1e^{2(tz)}|^{c^2}}\{e^{tz}dt^2+e^{(tz)}dz^2\}+e^{2t}dx^2+e^{2z}dy^2$$
(47)
$$\mu =\frac{c^2f_{o}^{}{}_{}{}^{2}}{e^{2t}}|1e^{2(tz)}|^{c^2},p=\mu $$
Type $`III`$
$$ds^2=e^{t+z}k^2|1e^{(tz)}|^\beta \{e^{tz}dt^2+dz^2\}+e^{t+z}dx^2+dy^2$$
(48)
$$\mu =\frac{1\beta }{4k^2e^{2t}|1e^{(tz)}|^\beta },p=\mu $$
For $`\beta =0`$ in the above solution, the fluid has geodesic flow ($`\dot{u}_a=0`$) and the metric admits a further spacelike KV which is not tangent to the Killing orbits $`V_2`$; the solution being therefore a special type of spatially homogeneous Bianchi cosmological model. For $`\beta 0`$ the fluid is non-geodesic and the metric admits no further KV’s.
Type $`V`$
$$ds^2=\frac{e^{t+z}}{f_{o}^{}{}_{}{}^{2}}\left\{dt^2\frac{c^2\phi ^2}{1c^2\phi ^2}+\frac{dz^2}{1c^2\phi ^2}\right\}+\phi ^2dx^2+dy^2$$
(49)
$$\mu =\frac{1c^2\phi ^2}{2c\phi ^2}\frac{f_{o}^{}{}_{}{}^{2}}{e^{t+z}},p=\mu $$
where $`\phi `$ is a function of $`tz`$ given implicitly by:
$$c(tz)=\mathrm{ln}\phi \frac{c^2}{2}\phi ^2$$
(50)
Notice that all these solutions have a stiff matter equation of state ($`p=\mu `$) and therefore can be derived from vacuum solutions (also admitting an abelian $`G_2`$) using a method proposed by Wainwright et al.
### 4.2 Vacuum Solutions
From our previous developments (see (13) and (3.1)) it is possible to find all vacuum solutions corresponding to types $`I`$ and $`V`$ in our classification. Solving the vacuum field equations for them we get, respectively:
Type $`I`$
$`ds^2`$ $`=`$ $`e^{2t}\{dt^2+dz^2+{\displaystyle \frac{(e^{2z}(\alpha ^2+c^2)e^{2z})^2}{\alpha ^2e^{2z}+(e^zce^z)^2}}dy^2+`$
$`+`$ $`(\alpha ^2e^{2z}+(e^zce^z)^2)({\displaystyle \frac{2\alpha }{\alpha ^2e^{2z}+(e^zce^z)^2}}dy+dx)^2\}`$
where $`\alpha `$, $`c`$ and $`\beta `$ are constants.
Type $`V`$
$$ds^2=e^{2t}e^{\alpha z}\{dt^2+dz^2\}+dy^2+dx^2$$
(52)
where $`\alpha `$ is a constant. It is immediate to see that this is (locally) Minkowski spacetime and it can be brought to the standard form by means of the following coordinate change:
$$\widehat{t}=\frac{1}{2}(\frac{e^{(1+\alpha )(t+z)}}{1+\alpha }+\frac{e^{(1\alpha )(tz)}}{1\alpha }),\widehat{z}=\frac{1}{2}(\frac{e^{(1+\alpha )(t+z)}}{1+\alpha }\frac{e^{(1\alpha )(tz)}}{1\alpha })$$
(53)
Acknowledgments
The authors would like to thank Drs J.M.M. Senovilla (Universitat de Barcelona) and G.S. Hall (University of Aberdeen) for many helpful discussions. Financial support from DGICYT Research project PB 91-0335 is also acknowledged.
|
warning/0005/astro-ph0005587.html
|
ar5iv
|
text
|
# The Sun’s acoustic asphericity and magnetic fields in the solar convection zone
## 1 Introduction
Rotational splittings of solar oscillation frequencies have been successfully utilized to infer the rotation rate in the solar interior. To first order, rotation affects only the splitting coefficients which represent odd terms in the azimuthal order $`m`$ of the global resonant modes. The even terms in these splitting coefficients, which reflect the Sun’s effective acoustic asphericity, can arise from second order effects contributed both by the rotation and magnetic field as also from latitudinal temperature variations. Since the rotation rate can be inferred using the odd splitting coefficients, the inferred profile can be used to estimate the second order effects. These can then be subtracted from the observed even coefficients to estimate the magnetic field strength (Gough & Thompson 1990) or other latitudinal variations in sound propagation speed. The distortion introduced by rotation can be compared with the measured oblateness at the solar surface.
The even coefficients of splittings are fairly small, and no definitive results have so far been obtained regarding the magnetic field strength in the solar interior. With the good quality data now becoming available from GONG (Global Oscillation Network Group) and MDI (Michelson Doppler Imager) projects, it is desirable to investigate the possibility of inferring the strength of magnetic field in solar interior. There should also be some shift in the mean frequency for each $`n,\mathrm{}`$ multiplet due to second order effects from rotation and magnetic field, which can also be estimated. It is difficult to measure this frequency shift from observed data as it is hard to separate it from the effects of other uncertainties in the spherical structure of the Sun. Nevertheless, these frequency shifts can affect the helioseismic inferences and it would be interesting to estimate their effect.
## 2 The technique
The frequencies of solar oscillations can be expressed in terms of the splitting coefficients:
$$\nu _{n,\mathrm{},m}=\nu _{n,\mathrm{}}+\underset{j=1}{\overset{J_{\mathrm{max}}}{}}a_j^{n,\mathrm{}}𝒫_j^{\mathrm{}}(m),(J_{\mathrm{max}}2\mathrm{})$$
(1)
where $`𝒫_j^{\mathrm{}}(m)`$ are orthogonal polynomials of degree $`j`$ in $`m`$ (Ritzwoller & Lavely 1991; Schou, Christensen-Dalsgaard & Thompson 1994). The odd coefficients $`a_1,a_3,a_5,\mathrm{}`$ can be used to infer the rotation rate in the solar interior, while the even coefficients arise basically from second order effects due to rotation and magnetic field. Since forces due to rotation or magnetic field in the solar interior are smaller by about 5 orders of magnitude as compared to gravitational forces, it is possible to apply a perturbative treatment to calculate their contribution to frequency splittings. In this approach, we estimate the effects of rotation and magnetic field on the frequencies but without explicitly constructing a model of a rotating, magnetic star.
We adopt the formulation due to Gough & Thompson (1990), with the difference that we include perturbation in the gravitational potential and also assume differential rotation in the interior, though the symmetry axis of magnetic field is taken to coincide with rotation axis.
In an inertial frame the oscillation equations can be formally written as
$$\xi \xi \xi +\rho \omega ^2\xi \xi \xi =\omega \xi \xi \xi +𝒩\xi \xi \xi +\xi \xi \xi ,$$
(2)
where
$`\xi \xi \xi `$ $`=`$ $`(\rho c_s^2\xi \xi \xi +\xi \xi \xi p)(\xi \xi \xi +\xi \xi \xi \mathrm{ln}\rho )p`$ (3)
$`\rho G\left({\displaystyle \frac{(\rho \xi \xi \xi )}{|𝐫𝐫^{}|}d^3𝐫^{}}\right),`$
$`\xi \xi \xi `$ $`=`$ $`2i\rho 𝐯\xi \xi \xi ,`$ (4)
$`𝒩\xi \xi \xi `$ $`=`$ $`\rho \xi \xi \xi (𝐯𝐯)+\rho (𝐯)^2\xi \xi \xi ,`$ (5)
$``$ $`=`$ $`{\displaystyle \frac{1}{4\pi }}({\displaystyle \frac{(\rho \xi \xi \xi )}{\rho }}(\times 𝐁)\times 𝐁+(\times 𝐁_1)\times 𝐁`$ (6)
$`+(\times 𝐁)\times 𝐁_1).`$
Here $`𝐁_1=\times (\xi \xi \xi \times 𝐁)`$ is the linearized Eulerian perturbation to magnetic field, $`𝐁`$, $`\xi `$$`\xi `$$`\xi `$ is the displacement eigenfunction, $`𝐯=\mathrm{\Omega }\mathrm{\Omega }\mathrm{\Omega }\times 𝐫`$ is the velocity due to rotation, and $`p,\rho ,c_s`$ are respectively, pressure, density and sound speed in the equilibrium state.
In the presence of rotation and magnetic field the equilibrium state will naturally undergo a distortion that needs to be included in the calculations. To account for this deformation we consider a transformation to map each point $`𝐫`$ in the distorted star to a point $`𝐱`$ in the spherical volume occupied by the undistorted star by a transformation
$$x=(1+h_\mathrm{\Omega }(𝐫)+h_B(𝐫))r,$$
(7)
where the functions $`h_\mathrm{\Omega }(𝐫)`$ and $`h_B(𝐫)`$ which depend on the rotation and magnetic field respectively, are to be determined by solving the equations for equilibrium in a distorted star (Gough & Thompson 1990). This will give us the perturbation to a nonrotating spherically symmetric solar model and the extent of distortion at the surface may be compared with observed values. Here, $`x`$ is chosen so that $`x=R`$ can be regarded as the distorted solar surface, where $`R`$ is the radial distance of the outermost layer included in the solar model. Similarly, various equilibrium quantities are also expressed in the form
$$\rho (r)=\rho _0(x)+\rho _\mathrm{\Omega }(𝐱)+\rho _B(𝐱).$$
(8)
In all these expansions higher order terms have been neglected.
We consider the terms on the right hand side of Eq. 2, as perturbations to basic equations for linear adiabatic oscillations for non-magnetic and non-rotating star. Rotation introduces a first order perturbation through $``$ which gives the odd splitting coefficients, while magnetic field can only give rise to even terms in $`m`$ and contributes to the even splitting coefficients. The distortion from a spherically symmetric equilibrium state also introduces even order terms. The relative magnitude of contributions from rotation and magnetic field will, of course, depend on the rotation rate and magnetic field strength. For the solar case we know that odd splitting coefficients arising from the first order effect of rotation are much larger than the even coefficients and we therefore expect the magnetic field to make a comparatively smaller contribution. We must therefore include the effect of rotation to second order, while magnetic field and distortion effects need be retained only to first non-vanishing terms. The first order perturbation arising in frequencies on account of rotation also introduces a perturbation to eigenfunctions which will give a second order contribution. We can formally express the frequency and eigenfunction as
$$\omega =\omega _0+\omega _1+\omega _2,\xi \xi \xi =\xi \xi \xi _0+\xi \xi \xi _1.$$
(9)
Retaining terms to second order, we get
$`_0(\xi \xi \xi _0+\xi \xi \xi _1)+_\mathrm{\Omega }\xi \xi \xi _0+_B\xi \xi \xi _0+\rho _0(\omega _0^2+2\omega _0\omega _1)(\xi \xi \xi _0+\xi \xi \xi _1)`$
$`+\rho _0(\omega _1^2+2\omega _0\omega _2)\xi \xi \xi _0+\rho _\mathrm{\Omega }\omega _0^2\xi \xi \xi _0+\rho _B\omega _0^2\xi \xi \xi _0`$
$`=\omega _0(\xi \xi \xi _0+\xi \xi \xi _1)+\omega _1\xi \xi \xi _0+𝒩\xi \xi \xi _0+\xi \xi \xi _0.`$ (10)
Here, $`_\mathrm{\Omega }`$ and $`_B`$ are the perturbations to $``$ arising from distortion of equilibrium state due to rotation and magnetic field respectively. Taking the scalar product with $`\xi \xi \xi _0^{}`$ and integrating over the entire volume, we recover
$`2\omega _0\rho _0\xi \xi \xi _0^{}\xi \xi \xi _0\omega _2=\xi \xi \xi _0^{}(𝒩_\mathrm{\Omega }\rho _\mathrm{\Omega }\omega _0^2)\xi \xi \xi _0`$
$`+\xi \xi \xi _0^{}(_B\rho _B\omega _0^2)\xi \xi \xi _0\omega _1^2\rho _0\xi \xi \xi _0^{}\xi \xi \xi _0`$
$`2\omega _0\omega _1\rho _0\xi \xi \xi _0^{}\xi \xi \xi _1+\omega _1\xi \xi \xi _0^{}\xi \xi \xi _0+\omega _0\xi \xi \xi _0^{}\xi \xi \xi _1,`$ (11)
where the angular brackets denote
$$f(x,\theta ,\varphi )=_{x<R}f(x,\theta ,\varphi )x^2\mathrm{sin}\theta dxd\theta d\varphi $$
(12)
The first order correction to frequency is given by
$$\omega _1=\frac{\xi \xi \xi _0^{}\xi \xi \xi _0}{2\rho _0\xi \xi \xi _0^{}\xi \xi \xi _0},$$
(13)
while perturbation to the eigenfunction may be calculated using
$$\xi \xi \xi _1+\rho _0\omega _0^2\xi \xi \xi _1=2\rho _0\omega _0\omega _1\xi \xi \xi _0+\omega _0\xi \xi \xi _0.$$
(14)
The observed odd splitting coefficients can be used to infer the rotation rate inside the Sun (Thompson et al. 1996; Schou et al. 1998). We approximate this rotation rate using the first three terms in the expansion of the angular velocity,
$$\mathrm{\Omega }(r,\theta )=\mathrm{\Omega }_0(r)+\mathrm{\Omega }_2(r)\mathrm{cos}^2\theta +\mathrm{\Omega }_4(r)\mathrm{cos}^4\theta ,$$
(15)
where $`\theta `$ is the colatitude. This rotation rate is then used to compute the second order rotational contribution to frequency splitting, which may be subtracted from the observed splittings to obtain the residual which may be due to magnetic field, any other velocity field or asphericity in solar structure.
In the present analysis we use only the toroidal magnetic field, taken to be of the form,
$$𝐁=[0,0,a(r)\frac{dP_k}{d\theta }\genfrac{}{}{0pt}{}{(\mathrm{cos}\theta )}{}],$$
(16)
with the axis of symmetry coinciding with the rotation axis. Here $`P_k(x)`$ is the Legendre polynomial of degree $`k`$. The Lorentz force due to a field of this form can be written as
$$𝐅=\rho (r)\underset{\lambda =0}{\overset{k}{}}[f_{r\lambda }(r)P_{2\lambda }(\mathrm{cos}\theta ),f_{\theta \lambda }(r)\frac{dP_{2\lambda }}{d\theta },0].$$
(17)
Each of this term can be treated separately and the results can be combined to yield the net effect.
We calculate the second order frequency shift due to rotation and magnetic field for each value of $`m`$ and then use Eq. 1 to obtain the corresponding splitting coefficients. These can then be compared with observed coefficients from GONG (Hill et al. 1996) or MDI (Rhodes et al. 1997) data. To evaluate the angular integrals we use the following recursion relations
$`\mathrm{cos}\theta Y_{\mathrm{}}^m`$ $`=`$ $`C_{\mathrm{}}^mY_{\mathrm{}+1}^m+C_\mathrm{}1^mY_\mathrm{}1^m,`$ (18)
$`\mathrm{sin}\theta {\displaystyle \frac{Y_{\mathrm{}}^m}{\theta }}`$ $`=`$ $`\mathrm{}C_{\mathrm{}}^mY_{\mathrm{}+1}^m(\mathrm{}+1)C_\mathrm{}1^mY_\mathrm{}1^m,`$ (19)
where
$$C_{\mathrm{}}^m=\sqrt{\frac{(\mathrm{}+1+m)(\mathrm{}+1m)}{(2\mathrm{}+1)(2\mathrm{}+3)}}.$$
(20)
Since we have used only the first two terms in the expansion of rotation rate as a function of latitude, we restrict to calculation of the splitting coefficients $`a_2`$ and $`a_4`$ in this work.
## 3 Results
We use the rotation rate inferred from the GONG data for the months 4–14 (Antia, Basu & Chitre 1998) to estimate the second order frequency shift and the corresponding splitting coefficients $`a_2`$ and $`a_4`$, as outlined in the previous section. We incorporate all the second-order contributions arising from rotation, including those from the distortion of equilibrium state and the perturbation to the eigenfunctions. Although there may be some variation in rotation rate with time, the estimated variation is very small and its effect on the inferred splitting coefficient would be much smaller than the errors in observed values.
### 3.1 Shift in the mean frequency
In principle, the shift in the mean frequency arising from second order effects of rotation can be calculated with the help of the prescription outlined in the previous section, by taking the spherically symmetric component of the perturbing force ($`\lambda =0`$ term in Eq. 17). However, this will also change the mass, radius and luminosity of the solar model. The change may be smaller than the errors in observed radius or luminosity, but it may tend to give a different estimate for modified frequency compared to what will be obtained if the observed constraints on mass, radius and luminosity were to be exactly applied. Hence, for obtaining a consistent estimate of the effect of distortion, we construct a spherically symmetric solar model with correct mass, radius and luminosity by modifying the effective acceleration due to gravity, $`g`$ to account for the spherically symmetric component of forces due to rotation. The difference in frequency of this model in relation to a standard, non-rotating model would give the frequency shift due to distortion. All the other second-order rotational terms are added to this shift, to obtain the total shift in frequency due to rotation which is displayed in Fig. 1. This figure includes all modes with $`0.5<\nu <4.5`$ mHz and $`\mathrm{}250`$. The corrections to mean frequencies due to general relativistic effects as discussed towards the end of this subsection, are also shown in the figure.
This relative frequency shift, which is less than $`10^5`$, is nonetheless comparable to the estimated errors in the observed frequencies and the correction should, in principle, be applied while doing inversions (e.g. Gough et al. 1996) for the Sun’s spherical structure. In order to estimate the error introduced by neglect of this effect, we can carry out an inversion for sound speed and density in the solar interior using this frequency shift due to rotation as the frequency difference and the results are shown in Fig. 2. The inversions are performed using a regularized least squares inversion technique (Antia 1996). The resulting $`\delta c_s^2/c_s^2`$ and $`\delta \rho /\rho `$ are almost an order of magnitude less than the estimated errors in inversions.
As an aside, we note that the internal rotation rate from Antia et al. (1998) adopted in our study was obtained assuming a spherically symmetric background state for the Sun, as is usual for inversions for the solar rotation. We realise that both the mean frequencies of solar oscillations and the rotational splittings will be modified by departures in the equilibrium solar model from spherical symmetry, as discussed in this paper. In order to estimate the resulting shift in rotational splittings we would need to calculate the third order terms in perturbation expansion of Gough & Thompson (1990). We have not included these terms in our analysis, but we expect that their contribution would have the same relative magnitude of $`10^5`$ as that found for the shift in mean frequencies. This is clearly, much smaller than the estimated errors in splitting coefficients in current helioseismic data sets. Therefore we do not expect the rotational splittings and hence the inverted rotation rate to be significantly affected by this higher order effect.
It may be noted that mean frequencies of f-modes get diminished by up to 15 nHz on account of the effect of rotation. Since rotation effectively reduces the acceleration due to gravity $`g`$, this leads to a decrease in the frequencies of f-modes. The relative change in f-mode frequencies is shown in Fig. 3. If this effect is taken into account the estimated solar radius using f-mode frequencies (Schou et al. 1997, Antia 1998) would effectively be decreased by about 4 km. This is again much less than the systematic errors in estimated radius, though the decrease is larger than the statistical errors (Tripathy & Antia 1999).
It is interesting to note that apart from second order effects of rotation, there would also be corrections to the frequencies arising from general relativity. The relativistic effect can be measured by $`Gm(r)/(rc^2)`$, where $`G`$ is the gravitational constant, $`m(r)`$ is the mass contained within spherical shell of radius $`r`$, and $`c`$ the speed of light. Fig. 4 shows this ratio in a solar model and it can be seen that it is comparable to the ratio of centrifugal to gravitational forces. It is possible to calculate a solar model using Oppenheimer-Volkoff equation of relativistic stellar structure instead of the standard equation of hydrostatic equilibrium:
$$\frac{dp}{dr}=\frac{G(\rho +p/c^2)(m+4\pi r^2p/c^2)}{r^2(12Gm/rc^2)}.$$
(21)
It is clear that general relativistic effect would be of opposite sign to that due to rotation, as rotation effectively reduces the acceleration due to gravity, $`g`$, while the relativistic correction tends to increase it. Thus there is a partial cancelation between the two effects. It is possible to calculate the change in solar models due to the relativistic effect, although a detailed calculation of frequencies using relativistic stellar oscillations equations would require considerable effort and is beyond the scope of the present work. To a first approximation we may calculate the effect by using the normal equations of stellar oscillations with gravity modified according to Eq. 21. Such a calculation shows that the effect of relativity more or less cancels the frequency shift due to rotation for low degree modes. The frequency shift due to general relativity are also shown in Fig. 1. If this frequency shift is added to the contribution arising from rotation then the effect on helioseismic inversion is significantly reduced in the solar core as can be seen from Fig. 2 (compare the thick and thin lines).
### 3.2 Oblateness due to rotation
During the course of computing the splitting coefficients, it is necessary to calculate the deformation induced by rotation as outlined by Gough and Thompson (1990). This deformation may be compared with the observed oblateness at the solar surface. The surface amplitudes of the $`P_2(\mathrm{cos}\theta )`$ and $`P_4(\mathrm{cos}\theta )`$ components of deformation are found to be $`5.84\times 10^6`$ and $`6.2\times 10^7`$ respectively, which are consistent with the estimates obtained by Armstrong & Kuhn (1999). These can be compared with measured values of $`(5.44\pm 0.46)\times 10^6`$ and $`(1.48\pm 0.58)\times 10^6`$ respectively, from MDI measurement during 1997 (Kuhn et al. 1998). Kuhn et al. (1998) find a large temporal variation in the $`P_4`$ component, but it is not clear if the variation is statistically significant. It can be seen that the measured values of solar oblateness are reasonably close to those expected from rotational distortion. There may be some residual arising from other effects, like magnetic field or other asphericities. The contribution from magnetic field is indeed expected to vary with solar cycle and may account for the variation in $`P_4`$ component, if the variation is in fact real.
It is also possible to estimate the global parameters for the Sun, like angular momentum, rotational kinetic energy and gravitational quadrupole and hexadecapole moments due to rotational distortion (Pijpers 1998) and the results are summarized below:
$`\text{Moment of Inertia, }I=7.11\times 10^{53}\mathrm{gm}\mathrm{cm}^2,`$ (22)
$`\text{Angular Momentum, }H=1.91\times 10^{48}\mathrm{gm}\mathrm{cm}^2\mathrm{s}^1,`$ (23)
$`\text{Kinetic Energy, }T=2.57\times 10^{42}\mathrm{gm}\mathrm{cm}^2\mathrm{s}^2,`$ (24)
$`\text{Quadrupole Moment, }J_2=2.18\times 10^7,`$ (25)
$`\text{Hexadecapole Moment, }J_4=4.64\times 10^9,`$ (26)
which are consistent with estimates of Pijpers (1998), who obtained his estimates by working in terms of kernels for the various quantities. The value of $`J_2`$ will yield a precession of the perihelion of planet Mercury by about 0.03 arcsec/century, which is small enough to maintain consistency of the general theory of relativity.
### 3.3 Second order splitting due to rotation
The contribution to splitting coefficients $`a_2`$ and $`a_4`$ due to rotation is shown in Fig. 5. This contribution needs to be subtracted from the observed splitting coefficients for obtaining the residual contribution which may arise from effects due to magnetic field, other velocity fields or asphericity in solar structure. Since the errors in individual splitting coefficients are too large to give significant differences, we average over 30 neighbouring modes in $`w=\nu /(\mathrm{}+1/2)`$ and the corresponding results are shown in Fig. 6. There is reasonable agreement between the GONG data for months 4–14 (23 August 1995 to 21 September 1996) and MDI data for the first 360 days of its operation (1 May 1996 to 25 April 1997). It is well known that the even splitting coefficients vary with solar activity cycle (Libbrecht & Woodard 1990; Dziembowski et al. 1998; Howe, Komm & Hill 1999) and there may not be agreement between observations taken at different epochs. But in the present case there is considerable overlap in period and the observations are near the minimum phase of solar activity, when these coefficients are not expected to vary significantly.
The difference between the observed splitting coefficients and the estimated contribution from rotation is significant for modes with turning points in the convection zone. For modes penetrating more deeply, the errors are larger and the difference is probably not significant.
### 3.4 Splitting due to magnetic field near the base of the convection zone
There have been some suggestions that a significant toroidal magnetic field may be concentrated in a layer around the base of the convection zone (Dziembowski & Goode 1992). We therefore first investigate splittings that are expected from such a field by assuming the magnetic field to be given by Eq. 16 with
$$a(r)=\{\begin{array}{cc}\sqrt{8\pi p_0\beta _0}(1(\frac{rr_0}{d})^2)\hfill & \text{if }|rr_0|d\hfill \\ 0\hfill & \text{otherwise}\hfill \end{array}$$
(27)
where $`p_0`$ is the gas pressure, $`\beta _0`$ is a constant giving the ratio of magnetic to gas pressure, $`r_0`$ and $`d`$ are constants defining the mean position and thickness of layer where the field is concentrated. Fig. 7 shows the splitting coefficients resulting from a toroidal magnetic field of this form concentrated at the base of the convection zone ($`r_0=0.713R_{}`$ and $`d=0.02R_{}`$). The splitting shown in this and subsequent figures includes both the direct and distortion contributions as defined by Gough and Thompson (1990).
The coefficients $`a_2`$ and $`a_4`$ from a toroidal magnetic field concentrated near the base of the convection zone have a characteristic signature for modes with turning point near the base of the convection zone; it should be possible to detect such a signal in the observed splittings if a strong enough magnetic field is indeed present in these layers. The computed splittings, particularly for the deeply penetrating modes in Fig. 7 show a great spread, which is characteristic of the splittings arising from a thin magnetic layer. We return below to the use that can potentially be made of this signature. In the present study, however, we choose to average over neighbouring modes, as discussed in Section 3.3, which suppresses this spread. Our rationale is that the errors in the real data are too large for the spread to be visibly distinguished from noise in the measured splittings at present. Thus we take averages over neighbouring modes and compare the residual after removing the contribution due to rotation with the expected splitting from the magnetic field and the results are shown in Fig. 8. Note that even after averaging a clear signature of the magnetic field is seen in the splitting coefficients. Since we are comparing the average over the same set of modes for the observed splittings and computed splittings for magnetic field, we should be able to get some estimate of magnetic field if a strong enough field does indeed exist. From Fig. 8 it can be seen that there is no clear signature of any feature near the base of the convection zone in the observed splittings, and hence we can only set an upper limit on the magnetic field in this layer. This will of course, depend on the thickness of the magnetic layer. Since there is no clear signature of any signal near the base of the convection zone, for quantitative purpose we take the difference between the lowest and highest point in the range $`0.6<r_t/R_{}<0.8`$ in observed splitting coefficients. For $`a_2`$ this difference is 8.7 nHz for MDI and 7.0 nHz for GONG data, while computed splittings with $`\beta _0=10^4`$ show a difference of 12.6 nHz for a half-thickness of $`0.02R_{}`$. Thus, we can put an upper limit of $`0.7\times 10^4`$ on $`\beta _0`$ which corresponds to a magnetic field strength of 300 kG for a layer of half-thickness $`0.02R_{}`$ near the base of the convection zone. Similar analysis for splitting coefficient $`a_4`$ yields a slightly larger upper limit of 400 kG. These limiting values are close to what was obtained by Basu (1997) using a similar technique and is also consistent with the value independently inferred by D’Silva & Choudhuri (1993). Note, this limit roughly increases as $`1/\sqrt{d}`$, and clearly, if the thickness of this region is smaller, the upper limit would be larger. It should be noted that the tachocline, where the rotation rate undergoes a transition from differential rotation in the convection zone to a solid-body like rotation in the radiative interior may have a thickness as small as $`0.01R_{}`$ (Basu 1997; Antia, Basu & Chitre 1998). With this thickness the upper limit on magnetic field would naturally be increased.
There is the possibility of distinguishing seismologically between magnetic layers of different thicknesses by using modes that penetrate well beneath the magnetic layer. A thin layer will induce a signature in the $`a_2`$ and $`a_4`$ coefficients which is periodic in mode frequency (Gough & Thompson 1988; Vorontsov 1988; Thompson 1988), in much the same way that the rather sharp transition near the base of the convective envelope produces a periodic signature in the mean frequencies (e.g., Gough 1990; Basu, Antia & Narasimha 1994; Monteiro, Christensen-Dalsgaard & Thompson 1994). Indeed it is this signature which is largely responsible for the vertical spread of points for modes with turning points at radii $`r0.6R`$ in Fig. 7. Basu (1997) attempted to use this oscillatory signal to obtain an upper limit on magnetic field near the base of convection zone (see also Gough & Thompson 1988). Fig. 9 shows $`\mathrm{}a_2`$ for modes with $`\mathrm{}10`$ for magnetic field concentrated near the base of the convection zone, with two different values of $`d`$. It is clear that the amplitude of oscillatory signal varies significantly with $`d`$. However, the observed splitting coefficients for low values of $`\mathrm{}`$ have large errors and it is difficult to extract the small oscillatory signal from these.
### 3.5 Field in the upper convection zone
Having considered a magnetic field at the base of the convection zone, where theory suggests a field might be stored, we consider where else the data might indicate the presence of magnetic field. There is no signature for the presence of significant magnetic field in the radiative interior, since the averaged residual splitting after correcting for rotation seem to be consistent with zero. However, within the convection zone there is some significant residual splitting, which could be due to the effect of a magnetic field. An inspection of these residuals indicates the existence of a peak around $`r=0.96R_{}`$, and indeed, if it is due solely to magnetic field, the field may be distributed around this depth ($`28000`$ km). It may be noted that this is approximately the depth to which shear layer seen in rotation profile extends (Antia, Basu & Chitre 1998; Schou et al. 1998).
We now attempt to estimate splittings due to the field concentrated in this region. Fig. 10 shows the splittings due to a few magnetic field configurations which are concentrated in the upper part of the convection zone. A comparison of these with the observed splittings indicates that there may be an azimuthal magnetic field with $`\beta <10^4`$ (i.e., $`B20000`$ G), with peak around $`r=0.96R_{}`$.
The possible existence of a magnetized layer with field of order 20 kG located around $`r=0.96R_{}`$ is, indeed, a significant inference drawn from our analysis. The physical interpretation for the origin of such a moderately strong magnetic field at this depth below the Sun’s surface is naturally a challenging task for theories of solar dynamo to accommodate. It may be useful to recall here that the numerical simulations of the Sun’s outer convection zone (Nordlund 1999) indicate a major presence of downward moving plumes. It is conceivable that these downdrafts could gather the turbulent magnetic field in the sub-surface layers and carry them to depths in the convective envelope until some sort of equipartition is reached. Interestingly, the density, $`\rho `$ at a depth of 25–30 Mm is upwards of $`4\times 10^3`$ gm cm<sup>-3</sup>, while the downward velocity for the plumes is of order 500 m s<sup>-1</sup>. The dynamical pressure of the plumes, $`\rho v^210^7`$ dyne cm<sup>-2</sup>, then becomes comparable with the magnetic pressure, $`B^2/8\pi `$, corresponding to a field strength of 20–30 kG. It is, therefore, tempting to envisage the formation of such a magnetized layer by the pounding of the downdrafts which tend to concentrate the field at depths where the equipartition of the kind outlined above is approached.
In this study we have assumed a smooth toroidal magnetic field, but in practice we do not expect such a field inside the convection zone. Turbulence may be expected to randomize the magnetic field and such a field may not be expected to produce any significant distortion in the equilibrium state. The direct effect of magnetic field will still be felt though the contribution would be different. Thus our results may be treated as indicating the order of magnitude of field that may be expected if the observed splitting coefficients are indeed due to the magnetic field. If the field is concentrated in flux tubes which occupy only a small fraction of the volume then the required magnetic field could be correspondingly larger. If we assume that the flux tubes occupy a fraction $`f`$ of the total volume, the magnetic field strength should increase by $`1/\sqrt{f}`$. If we consider only direct contribution to the splittings then it turns out that $`a_2`$ is always negative for all toroidal field configurations that we tried and hence such a contribution is not likely to explain the observed splittings. But a different magnetic field configuration, e.g., poloidal field might produce $`a_2`$ with the required sign using only direct contribution. The order-of-magnitude splitting caused by a magnetic field is $`\mathrm{}a_2/\nu \beta _0v_A^2/c_s^2`$, where $`v_A`$ is the Alfvén speed. We therefore regard it as unlikely that a different magnetic field configuration would produce a markedly different answer for the field strength required to account for the observed signal in $`a_2`$ and $`a_4`$.
A nonmagnetic latitudinally-dependent perturbation to the wave propagation speed might be responsible for the signal we detect (cf., Gough & Zweibel 1995). Once again we may expect a perturbation of order $`10^4`$ located in the region around $`r=0.96R_{}`$ to yield the observed splittings. Gough et al. (1996) inferred a perturbation of that magnitude, of unspecified origin, from earlier GONG data. A temperature variation of order 10K, suitably confined, might conceivably produce a similar signature. In fact, Kuhn (e.g., Kuhn 1996) has argued that the thermal shadow of belts of magnetic flux near the bottom of the convection zone can have a significant effect on the even a-coefficients. But Kuhn’s models show the largest temperature perturbations occurring in the very superficial superadiabatic layer, at a depth of a small fraction of one per cent of the solar radius. Such a perturbation alone would be consistent with the f-modes having a small residual splitting, but would not explain the apparent overturning of the p-mode splittings at $`r_t0.96R`$. A magnetic field at some depth below the surface may explain both aspects. We certainly do not rule out the possibility that some nonmagnetic asphericity, which we have not considered in detail in this study, may account for some of the observed splittings.
## 4 Conclusions
Second order correction to mean frequencies due to rotation is comparable to the error estimates in the observed frequencies. The error in helioseismic inversion introduced by the frequency shift due to rotation is $`10^4`$, which is much smaller than the estimated errors in inversions. Further, a part of this frequency shift is expected to be nullified by the general relativistic effects. The shift in f-mode frequencies due to rotation can reduce the estimated solar radius by 4 km. The distortion due to rotation can yield surface oblateness of $`5.8\times 10^6`$ and $`6.2\times 10^7`$ in the $`P_2(\mathrm{cos}\theta )`$ and $`P_4(\mathrm{cos}\theta )`$ components, respectively. This is in reasonable agreement with observed oblateness at the solar surface (Kuhn et al. 1998) and it appears that most of the observed distortion is accounted by the seismically inferred rotation rate in solar interior. The quadrupole moment $`J_2=2.18\times 10^7`$ resulting from rotational distortion is small enough to maintain consistency of the general theory of relativity.
After subtracting the estimated contribution from rotation to the splitting coefficients $`a_2`$ and $`a_4`$ from the observed splittings, there is a small residual which is statistically significant in the convection zone. This could arise from a magnetic field. From the magnitude of residual in observed splittings we can tentatively conclude that magnetic field with $`\beta 10^4`$ may be present in the upper part of the convection zone. This corresponds to an azimuthal magnetic field of $`20`$ kG around $`r=0.96R_{}`$. However, we cannot rule out the possibility that this signal in splitting coefficients may arise from some aspherical perturbation to the temperature field. This would be practically indistinguishable from the effect of a magnetic field using just the mode frequencies (Gough & Zweibel 1995); but complementary analyses such as time-distance helioseismology might be able to distinguish them, since the local direct effect of a magnetic field on the waves is anisotropic, whereas that of a temperature perturbation is not.
A toroidal magnetic field that is concentrated near the base of the convection zone gives a characteristic pattern in the splittings for modes with lower turning point in that region. Since no such signal is seen in observed frequencies, we can put an upper limit of about 300 kG on the strength of the magnetic field in this region.
###### Acknowledgements.
We thank J.-P. Zahn for useful comments, and J. Schou for providing MDI splittings. This work utilizes data obtained by the Global Oscillation Network Group (GONG) project, managed by the National Solar Observatory, a Division of the National Optical Astronomy Observatories, which is operated by AURA, Inc. under a cooperative agreement with the National Science Foundation. The data were acquired by instruments operated by the Big Bear Solar Observatory, High Altitude Observatory, Learmonth Solar Observatory, Udaipur Solar Observatory, Instituto de Astrofísico de Canarias, and Cerro Tololo Interamerican Observatory. This work also utilizes data from the Solar Oscillations Investigation / Michelson Doppler Imager (SOI/MDI) on the Solar and Heliospheric Observatory (SOHO). SOHO is a project of international cooperation between ESA and NASA. MJT acknowledges the support of the UK Particle Physics and Astronomy Research Council.
|
warning/0005/hep-th0005121.html
|
ar5iv
|
text
|
# UPR-887T, OUTP-99-03P Five–Brane BPS States in Heterotic M–Theory
## 1 Introduction:
In a series of papers , it was shown that when Hořava–Witten theory is compactified on Calabi–Yau threefolds, a “brane Universe” theory of particle physics, called Heterotic M–Theory, naturally emerges. This Universe consists of a five–dimensional bulk space which is bounded by two four–dimensional BPS three–branes, one at one end of the fifth dimension and one at the other. One of these branes contains observed matter, such as quarks and leptons, as well as their gauge interactions, either as a grand unified theory or as the standard model. This is called the “observable sector”. The other brane can be chosen to contain pure gauge fields and their associated gauginos. Gaugino condensates on this brane may provide a source of supersymmetry breaking in the theory. This brane is called the “hidden sector”.
As discussed in , the $`E_8`$ gauge group of Horǎva–Witten theory can be spontaneously broken to a GUT group or the standard model gauge group on the observable brane by the appearance of a semi–stable holomorphic vector bundle with structure group $`G`$ on the associated Calabi–Yau threefold. This bundle acts as a non–zero vacuum expectation value, breaking the $`E_8`$ gauge group to the commutant $`H`$ of $`G`$. Commutant $`H`$ acts as the gauge group on the observable brane. The physical requirements that $`H`$ be a realistic GUT or standard model gauge group, that the theory be anomaly free and that there be three families of quarks and leptons, puts strong constraints on the allowed vector bundle. In general, there will be a solution if and only if one further requires the appearance of five–branes, wrapped on holomorphic curves in the associated Calabi–Yau threefold, located in the bulk space.
As shown in , the exact structure of the holomorphic curve on which these bulk five–branes are wrapped can be computed from the anomaly cancellation condition. In this paper, we will take the Calabi–Yau threefold, $`X`$, to be an elliptic fibration over a base surface, $`B`$. In this case, it was shown in that the holomorphic curve on which the bulk five–branes wrap generically has both a base and a fiber component. However, as discussed in , there are always regions of the associated moduli space where such a curve decomposes into independent curves, at least one of which is a pure fiber. We denote this pure fiber by $`𝒞_2`$. When wrapped on this component, the manifold of a five–brane is of the form $`𝒞_2\times M_4`$. As is well known, as long as the elliptic fiber is smooth, the $`M_4`$ worldvolume spectrum consists of an $`N=4`$ Abelian vector supermultiplet, which is broken to an $`N=1`$ Abelian vector supermultiplet plus uncharged chiral supermultiplets by higher dimensional operators in the effective theory.
However, not all fibers of an elliptic fibration are smooth. There are a number of ways, classified by Kodaira , in which the associated torii can degenerate. The locus of all points in the base $`B`$ over which the fiber is degenerate forms a divisor, called the discriminant curve. The structure of this discriminant curve depends on the precise theory one is considering and, in general, can be quite intricate. This curve generically has smooth sections, cusps, and both normal and more complicated intersections. The Kodaira type of degeneration of the fiber can change substantially from place to place over the discriminant curve. If we choose some point on the discriminant locus, and wrap the five–brane over the associated degenerate elliptic fiber, then new physics emerges. If the point chosen is in a smooth part of the curve, then the worldvolume theory on $`M_4`$ has $`N=2`$ supersymmetry at low energy. If one chooses the point at a singular locus of the discriminant curve, then the supersymmetry at low energies may be further reduced to $`N=1`$. In either case, new, massless states emerge, in addition to the “standard” worldvolume multiplets mentioned above. These new states arise from the fluctuations of M membranes stretching between the vanishing cycles of the torus. As the cycles go to zero, massless states emerge.
The purpose of this paper is threefold. First, we present, for specificity, a quasi–realistic, three–family $`SU(5)`$ grand unified theory within the context of Heterotic M-Theory. This theory corresponds to an explicit semi–stable holomorphic vector bundle with structure group $`G=SU(5)`$ over an elliptically fibered Calabi–Yau threefold with base $`B=\widehat{𝔽}_3`$. We discuss the bulk five–branes and compute the class of the holomorphic curve over which they are wrapped. The moduli space of this class is then presented and it is shown that a region of this space corresponds to a single five–brane wrapped on a pure elliptic fiber $`𝒞_2`$. This work is presented in Section $`2`$. In Section $`3`$, we give the general theory for explicitly computing the discriminant curves of elliptically fibered Calabi–Yau threefolds. We apply these methods to determine the discriminant curves associated with the specific Calabi–Yau threefold with base $`B=\widehat{𝔽}_3`$ presented in Section $`2`$. We show that there are several possible curves, each with smooth sections, cusps, and normal and tangential self–intersections. We compute the Kodaira type fiber degeneracy over all points of the discriminant curve. In the following section, Section $`4`$, we explicitly compute the $`M_4`$ worldvolume BPS states that arise when a five–brane is wrapped on the degenerate fibers over the smooth parts of the discriminant curves. This reduces the problem from one of two complex moduli to one modulus, and allows the application of standard Kodaira theory. Using the theory of string junctions developed in , we present the spectrum of BPS states, and the associated $`N=2`$ hyper- and vector multiplets, for fibers of each Kodaira type occuring in the explicit theory of Section $`2`$. The computation of light states over cusps and points of self–intersection, being inherently more intricate, will be presented in future publications . Finally, in the Appendix, we outline those aspects of string junction theory necessary for the calculations in this paper.
One can only speculate at this point about the possible physical role of the BPS multiplets that arise in this manner. They will appear in the brane world scenario as “exotic” charged matter living on the worldvolumes of hidden sector three–branes. These new multiplets could be involved in new mechanisms of supersymmetry breakdown , might be relevant to cosmology , such as the dark matter in the Universe and so on. We will return to these issues elsewhere.
## 2 A Three Generation GUT Theory:
In this section, we will construct a quasi–realistic particle physics theory with three generations of quarks and leptons and a grand unified gauge group SU(5). This is carried out within the framework of Heterotic M–Theory compactified on an elliptically fibered Calabi–Yau threefold X. The requirement that X be a Calabi–Yau manifold means that $`c_1(TX)=0`$, which, in turn, implies that the base B of the elliptic fibration is restricted to be a del Pezzo surface $`d_i`$, a Hirzebruch surface $`𝔽_r`$, an Enriques surface $`𝐄`$ or certain “blown–up” $`𝔽_r`$ surfaces. For specificity, in this paper we will consider Calabi–Yau spaces with the base restricted to be of the latter type, that is, a blown–up Hirzebruch surface. As will become clear in the next section, we make this choice because elliptic fibrations over such a base have non–trivial discriminant curves involving not only different Kodaira type fibers, but also both normal and tangential crossing points. Again, for concreteness we will construct a model over a specific blown–up Hirzebruch surface, although our results apply generally.
Consider the Hirzebruch surface $`𝔽_3`$. This is a ruled surface which is a natural fibration of $`𝔽_3^1`$. We denote the fiber class of this fibration by $``$. This class has vanishing self–intersection. In addition, there is a unique curve of self–intersection $`3`$ which we denote by $`𝒮`$. These two classes have a single point of intersection since $`𝒮=1`$. We now modify this surface by blowing up a point on the curve $``$ which is not the point of intersection. The blow–up at that point introduces a new exceptional class, which we denote as $`𝒢`$. This is the surface that we will use as the base of our elliptically fibered Calabi–Yau threefold and we denote it by
$$B=\widehat{𝔽}_3$$
(2.1)
It is convenient to introduce $``$, where $`+𝒢=`$. The three classes $`𝒮`$, $``$ and $`𝒢`$ are each effective classes and together they form a basis of $`H_2(\widehat{𝔽}_3,)`$. Furthermore, they generate the Mori cone of effective classes. The intersection numbers of these three classes are given by
$$𝒮𝒮=3,=1,𝒢𝒢=1$$
(2.2)
and
$$𝒮=1,𝒮𝒢=0,𝒢=1$$
(2.3)
The first and second Chern classes of $`\widehat{𝔽}_3`$ can be written as
$$c_1(\widehat{𝔽}_3)=2𝒮+5+4𝒢,c_2(\widehat{𝔽}_3)=5$$
(2.4)
respectively. Having specified the Chern classes of the base, we note that the Chern classes of the tangent bundle, $`TX`$ of $`X`$ can now be computed. Since $`X`$ is a Calabi–Yau threefold, $`c_1(TX)=0`$. However, $`c_2(TX)`$ and $`c_3(TX)`$ are found to be non-vanishing functions of $`c_1(\widehat{𝔽}_3)`$ and $`c_2(\widehat{𝔽}_3)`$. The exact expression for $`c_2(TX)`$ is given in .
We now want to specify a stable, holomorphic $`SU(n)`$ vector bundle over the elliptically fibered Calabi–Yau threefold X with $`B=\widehat{𝔽}_3`$. For specificity, we will demand that the grand unification group be given by
$$H=SU(5)$$
(2.5)
which then requires that we choose the structure group of the vector bundle to be
$$G=SU(5)$$
(2.6)
Hence, $`n=5`$. Having chosen $`n`$, the class of the spectral cover $`𝒞`$ of the vector bundle is given by specifying a curve $`\eta `$ in the base $`B=\widehat{𝔽}_3`$. This curve must be effective to ensure that the spectral cover is effective, as it must be. In addition, $`\eta `$ must be an irreducible curve so that the associated vector bundle will be stable. Since $`𝒮`$, $``$ and $`𝒢`$ are a basis of $`H_2(\widehat{𝔽}_3,)`$, we can always write
$$\eta =a𝒮+b+c𝒢$$
(2.7)
where $`a`$, $`b`$ and $`c`$ are integers. Recalling that the classes $`𝒮`$, $``$ and $`𝒢`$ generate the Mori cone, it follows that the condition for $`\eta `$ to be an effective class is simply that
$$a0,b0,c0$$
(2.8)
The conditions for the irreducibility of $`\eta `$ are a little more intricate to derive. Here we simply state the result, which is that either
$$a0,b=c=0b0,a=c=0c0,a=b=0$$
(2.9)
or
$$ba,bc,abc$$
(2.10)
Having chosen the spectral cover $`𝒞`$ by giving $`n`$ and $`\eta `$, it is now necessary to specify the spectral line bundle $`𝒩`$ over $`𝒞`$. Generically, the allowed line bundles are indexed by a rational number $`\lambda `$. For $`n`$ odd, which is the case in our theory, this parameter must satisfy $`\lambda +\frac{1}{2}`$. For specificity, we will choose
$$\lambda =\frac{1}{2}$$
(2.11)
Having specified $`𝒞`$ and $`𝒩`$ subject to the above conditions, a stable, holomorphic $`SU(5)`$ vector bundle $`V`$ can be constructed using the Fourier–Mukai transformation
$$(𝒞,𝒩)V$$
(2.12)
We need not discuss $`V`$ here, other than to say that its properties can be calculated from the above data. In particular, in addition to its vanishing first Chern class, $`c_1(V)=0`$, the second and third Chern classes, $`c_2(V)`$ and $`c_3(V)`$ respectively, can be computed and are found to be functions of $`n`$, $`\eta `$ and $`\lambda `$. The exact expressions are given in .
As discussed in , the physical requirement that the theory be anomaly free leads to the topological expression
$$W=c_2(TX)c_2(V)$$
(2.13)
$`W`$ is a class that is interpreted as being a holomorphic curve in the Calabi–Yau threefold $`X`$ on which five–branes, located in the five–dimensional bulk space, are wrapped. Using the exact expressions for $`c_2(TX)`$ and $`c_2(V)`$, we can compute this five–brane class explicitly. We find that
$$W=W_B\sigma +a_fF$$
(2.14)
where $`\sigma `$ is the class of the zero section of the elliptic fibration $`X`$ and $`F`$ is the generic class of its fiber. Furthermore,
$$W_B=12c_1(B)\eta $$
(2.15)
and
$$a_f=c_2(B)+(11+\frac{n^3n}{24})c_1(B)^2\frac{3n}{2\lambda }(\lambda ^2\frac{1}{4})$$
(2.16)
For the above specific theory, both $`W_B`$ and $`a_f`$ can be computed and are given by
$$W_B=(24a)𝒮+(60b)+(48c)𝒢$$
(2.17)
and
$$a_f=117$$
(2.18)
respectively, where we have used equations (2.4), (2.7) and (2.11), as well as $`n=5`$, $`\lambda =\frac{1}{2}`$ and equations (2.2), (2.3). As discussed in , the class $`W`$ is further constrained by the requirement that it be an effective class in $`H_2(X,)`$. It was shown that this will be the case if and only if $`W_B`$ is an effective class in $`H_2(\widehat{𝔽}_3,)`$ and $`a_f`$ is a non–negative integer. It follows from (2.18) that the last condition is satisfied. We see from equation (2.17) that $`W_B`$ and, therefore, $`W`$ will be an effective curve if and only if
$$24a,60b,48c$$
(2.19)
In addition to the anomaly cancellation condition (2.13), there is another property required of any realistic theory of particle physics, that is, that the number of quark and lepton generations be 3. As discussed in , the number of quark and lepton generations is given by $`N_{gen}=\frac{c_3(V)}{2}`$. Using the expression for $`c_3(V)`$, it was shown that the three family condition imposes the further constraint that
$$\lambda (W_B^2(24n)W_Bc_1(B)+12(12n)c_1(B)^2)=3$$
(2.20)
For the above specific theory, using equations (2.4) and (2.17), as well as $`n=5`$ and the intersections in (2.2), (2.3), this condition becomes
$$3a^2+2abb^2+2bcc^2+5(abc)=6$$
(2.21)
Therefore, to get a realistic $`SU(5)`$ grand unified theory with three families of quarks and leptons for elliptically fibered Calabi–Yau threefolds with base $`B=\widehat{𝔽}_3`$, we must find a curve $`\eta `$ of form (2.7) that satisfies the conditions (2.8),(2.9) or (2.10), (2.19) and (2.21) simultaneously.
The solution of these constraints is not entirely trivial. For the purposes of this paper, we will only give the simplest solution. We find that it is impossible to find any solutions for curves $`\eta `$ satisfying the conditions in equation (2.9). We must, therefore, impose equation (2.10). As an ansatz, we try for solutions with $`a0`$ and $`b=c0`$. Under these conditions, we find a solution with
$$a=6,b=c=42$$
(2.22)
corresponding to the curve
$$\eta =6𝒮+42+42𝒢$$
(2.23)
and the five–brane class $`W=W_B\sigma +117F`$, where
$$W_B=18𝒮+18+6𝒢$$
(2.24)
The bulk space five–branes of this specific theory, described by the class $`W=W_B\sigma +117F`$ with $`W_B`$ given in (2.24), are the objects of interest in this paper.
To analyze the physical structure of these five–branes, it is useful to first construct their moduli space. The moduli spaces of M-theory five–branes wrapped on holomorphic curves in elliptically fibered Calabi–Yau threefolds were constructed in . Here, we will simply apply the results of to the specific theory discussed above. We find that the moduli space of these five–branes is given by
$$(W_B\sigma +117F)=\underset{n=0}{\overset{117}{}}_0(W_B\sigma +nF)\times ((117n)F)$$
(2.25)
where $`W_B`$ is given in (2.24) and
$$((117n)F)=(\widehat{𝔽}_3\times _a^1)^{117n}/_{117n}$$
(2.26)
$`_a^1`$ is the moduli space of the translation/axion multiplet associated with the fifth direction of the bulk space. The components $`_0(W_B\sigma +nF)`$ of the moduli space can also be explicitly computed. This construction was presented in , where specific examples were given. It was shown that these components, which can be rather complicated, depend sensitively on the choice of the base surface of the elliptic fibration and on the explicit form of the curve $`W_B`$ . However, in this paper, it is not necessary to know the explicit form of these components of moduli space, and we will not discuss them further.
We are particularly interested in the component of moduli space where all the fibers are associated with the base curve $`W_B`$ except for one, that is, when
$$n=116$$
(2.27)
The relevant component of moduli space is then
$$_0(W_B\sigma +(116)F)\times (F)$$
(2.28)
where
$$(F)=\widehat{𝔽}_3\times _a^1$$
(2.29)
Physically, this region of moduli space describes a situation where all the five–branes are wrapped on the complicated curve $`W_B\sigma +116F`$ except for one, which is wrapped on a vertical curve described by the pure fiber class $`F`$. The moduli of this single five–brane, specified by the space (2.29), are independent of the all other moduli. Therefore, this single five–brane can be discussed entirely by itself, without reference to the rest of the five–brane class. In the remainder of this paper, we will consider only this single five–brane wrapped on a curve described by the fiber class $`F`$. If, furthermore, we hold this five–brane fixed at a point $`y_0`$ in the fifth direction of the bulk space, the moduli space is reduced to
$$(F)=\widehat{𝔽}_3$$
(2.30)
The physical properties of this single five–brane are then completely determined by the behavior of the elliptic fiber $`F`$ on which the five–brane is wrapped as one moves around the base space $`B=\widehat{𝔽}_3`$.
The fiber over a generic point in $`B=\widehat{𝔽}_3`$ is a smooth torus. Therefore, the worldvolume fields of a five–brane wrapped on a generic fiber consist of the standard ones associated with the self–dual anti–symmetric tensor multiplet. However, as is well known, the torus fibers can become singular at specific points in the base surface. The locus of points where the fibers degenerate is a divisor of the base called the discriminant locus. That is, the discriminant locus is a smooth curve in $`B=\widehat{𝔽}_3`$. Let us choose a point on the discriminant curve and wrap a five–brane around the degenerate torus fibered over that point. Then, as is well known, in addition to the usual fields associated with the self–dual anti–symmetric tensor multiplet, there are also massless BPS multiplets that appear on the five–brane worldvolume. The properties of these new states are directly related to the “singularity structure” of the elliptic fiber at that point, which has been classified by Kodaira and will be discussed below. The ensemble of these new states form a conformal field theory whose construction will be one of the goals of this paper.
Furthermore, the structure of the torus degeneration, that is, the Kodaira type of the elliptic fiber, can change as one considers different points on the discriminant curve. Therefore, the conformal field theory that appears when a five–brane is wrapped on a degenerate torus over one point of the discriminant curve, can be very different from that which occurs when it is wrapped on a degenerate fiber over a different point on the curve. In this paper, we will present the general theory for determining the exotic multiplets and the associated conformal field theories on the five–brane worldvolume. It is clear that the starting point of our analysis must be the construction of the allowed discriminant curves in the base $`B=\widehat{𝔽}_3`$, and the computation of the exact Kodaira structure of the fiber degeneration. This will be carried out in the next section.
## 3 Constructing The Discriminant Curves:
A simple representation of an elliptic curve is given in the projective space $`^2`$ by the Weierstrass equation
$$zy^2=4x^3g_2xz^2g_3z^3$$
(3.1)
where $`(x,y,z)`$ are the homogeneous coordinates of $`^2`$ and $`g_2`$, $`g_3`$ are constants. The origin of the elliptic curve is located at $`(x,y,z)=(0,1,0)`$. The torus described by (3.1) can become degenerate if one of its cycles shrinks to zero. Such singular behavior is characterized by the vanishing of the discriminant
$$\mathrm{\Delta }=g_2^327g_3^2$$
(3.2)
Equation (3.1) can also represent an elliptically fibered threefold, $`W`$, if the coefficients $`g_2`$ and $`g_3`$ in the Weierstrass equation are functions over a base surface, $`B`$. Clearly $`W`$ constructed in this way has a zero section $`\sigma `$ and we denote by $``$ the co–normal bundle over $`\sigma `$. The correct way to express this fibration globally is to replace the projective plane $`^2`$ by a fourfold $`^2`$-bundle $`PB`$ where $`P=(𝒪_B^2^3)`$. The notation $`(M)`$ stands for the projectivization of a vector bundle $`M`$. There is a hyperplane line bundle $`𝒪_P(1)`$ on $`P`$ which corresponds to the divisor $`(^2^3)P`$ and the coordinates $`x,y,z`$ are sections of $`𝒪_P(1)^2,𝒪_P(1)^3`$ and $`𝒪_P(1)`$ respectively. Equation (3.1) can now be interpreted as the affine form of a global equation on $`P`$ involving the sections $`x,y,z`$, as long as we require $`g_2`$ and $`g_3`$ to be sections of appropriate line bundles over the base $`B`$. It follows from (3.1) that
$$g_2^4,g_3^6$$
(3.3)
The symbol “$``$” simply means “section of”. The zero locus of equation (3.1) defines an elliptically fibered threefold hypersurface of $`P`$, which is called a Weierstrass fibration over the base $`B`$ and which we will denote by $`W`$. In addition, note that the discriminant defined in equation (3.2) is a section of the line bundle
$$\mathrm{\Delta }^{12}$$
(3.4)
over the base $`B`$. The zero locus of the discriminant section specifies a divisor of $`B`$, the discriminant curve. An elliptic fiber in $`W`$ is degenerate if and only if it lies over a point in the discriminant curve.
Let us now demand that $`W`$ be a Calabi–Yau threefold. It follows that we must require $`c_1(TW)=0`$. It can be shown that this, in turn, implies
$$=K_B^1$$
(3.5)
where $`K_B`$ is the canonical bundle on the base, $`B`$. Condition (3.5) is rather strong and, as stated earlier, restricts the allowed base spaces of an elliptically fibered Calabi–Yau threefold to be rational (del Pezzo, Hirzebruch, as well as certain blow–ups of Hirzebruch surfaces) and Enriques surfaces (see for example ). Henceforth, we will only discuss Weierstrass fibrations that are, in addition, Calabi–Yau threefolds. Using the fact that $`K_B^1=𝒪_B(c_1(B))`$, it follows from $`g_2`$ and $`g_3`$ are sections of the line bundles
$$g_2𝒪_B(4c_1(B)),g_3𝒪_B(6c_1(B))$$
(3.6)
This places constraints on the sections $`g_2`$ and $`g_3`$ that we will return to below. In addition, note from (3.4) and (3.5) that the discriminant is a section of the line bundle
$$\mathrm{\Delta }𝒪_B(12c_1(B))$$
(3.7)
This condition implies that the discriminant curve in $`B`$ must lie in the class $`12c_1(B)`$, a strong restriction on allowed discriminant curves.
Generally, Weierstrass fibrations can be singular. These singularities, however, can be removed by “blowing–up” the singular points, producing a smooth elliptically fibered threefold that we will denote by $`X`$. Despite the fact that $`c_1(TW)=0`$, the first Chern class of the blown–up fibration $`X`$ need not vanish and, hence, $`X`$ need not be a Calabi–Yau threefold. It is possible to enforce the condition $`c_1(TX)=0`$, but only at the cost of putting additional constraints on the sections $`g_2`$ and $`g_3`$. In this paper, we will always demand that $`X`$ is a Calabi–Yau threefold and, hence, that $`g_2`$ and $`g_3`$ are suitably constrained. In this paper, for simplicity, we will restrict the discussion to blow–ups such that each fiber of the induced fibration $`XB`$ will still be a complex curve. This then places further constraints on the allowed sections $`g_2`$ and $`g_3`$. Finally, and conversely, it can be shown that any smooth elliptically fibered Calabi–Yau threefold $`X`$ with a zero section can be obtained as a resolution of the singularities of a Weierstrass fibration.
To make these statements concrete, we now explicitly compute the discriminant curve for an elliptically fibered Calabi–Yau threefold with base $`B=𝔽_3`$.
### Example 1: $`B=𝔽_3`$
Consider a smooth elliptically fibered Calabi–Yau threefold $`X`$ with zero section $`\sigma `$ and base $`B=𝔽_3`$. As discussed previously, there are two effective classes $``$ and $`𝒮`$ that together form a basis of $`H_2(𝔽_3,)`$ with intersection numbers
$$=0,𝒮𝒮=3,𝒮=1$$
(3.8)
The first and second Chern classes of $`B=𝔽_3`$ are given by
$$c_1(𝔽_3)=2𝒮+5,c_2(𝔽_3)=4$$
(3.9)
$`X`$ is the resolution of a Weierstrass fibration over the base $`B=𝔽_3`$ with sections $`g_2`$, $`g_3`$ satisfying (3.6) and discriminant $`\mathrm{\Delta }`$ defined by (3.2) and satisfying (3.7). Let us first consider the consequences of (3.6). To do this, we must discuss several important assertions. The first of these, Claim I, is the following.
* The zero locus of any section $`g_2𝒪_B(4c_1(𝔽_3))`$ and $`g_3𝒪_B(6c_1(𝔽_3))`$ must necessarily vanish along $`𝒮`$ with order at least 2.
The proof of this assertion is outlined as follows. First consider $`g_2𝒪_B(4c_1(𝔽_3))`$. Note that one can always write
$$𝒪_B(4c_1(𝔽_3))=𝒪_B(2𝒮)𝒪_B(4c_1(𝔽_3)2𝒮)$$
(3.10)
To prove Claim I, we need to show that any section $`g_2𝒪_B(4c_1(𝔽_3))`$ can be written in the form
$$g_2=s_{2𝒮}G_2$$
(3.11)
where
$$s_{2𝒮}𝒪_B(2𝒮),G_2𝒪_B(4c_1(𝔽_3)2𝒮)$$
(3.12)
Denote the vector space of all sections of $`𝒪_B(4c_1(𝔽_3))`$ and $`𝒪_B(4c_1(𝔽_3)2𝒮)`$ by $`H^0(𝔽_3,4c_1(𝔽_3))`$ and $`H^0(𝔽_3,4c_1(𝔽_3)2𝒮)`$ respectively. The decomposition (3.11) will be satisfied if the dimensions of these two vector spaces are identical, which we now proceed to show. To do this, consider the exact sequence
$$0H^0(𝔽_3,4c_1(𝔽_3)𝒮)H^0(𝔽_3,4c_1(𝔽_3))H^0(𝒮,4c_1(𝔽_3)|_𝒮)$$
(3.13)
Note, however, that
$$H^0(𝒮,(4c_1(𝔽_3))|_𝒮)=H^0(^1,4c_1(𝔽_3)𝒮)=H^0(^1,4)=0$$
(3.14)
where we have used (3.8) and (3.9). Hence, dim$`H^0(𝔽_3,4c_1(𝔽_3)𝒮)=`$ dim$`H^0(𝔽_3,4c_1(𝔽_3))`$. Similarly, consider the exact sequence
$$0H^0(𝔽_3,4c_1(𝔽_3)2𝒮)H^0(𝔽_3,4c_1(𝔽_3)𝒮)H^0(𝒮,(4c_1(𝔽_3)𝒮)|_𝒮)$$
(3.15)
and the relation
$$H^0(𝒮,(4c_1(𝔽_3)𝒮)|_𝒮)=H^0(^1,(4c_1(𝔽_3)𝒮)𝒮)=H^0(^1,1)=0$$
(3.16)
Therefore, dim$`H^0(𝔽_3,4c_1(𝔽_3)2𝒮)=`$ dim$`H^0(𝔽_3,4c_1(𝔽_3)𝒮)`$. Combining with the above relation implies
$$dimH^0(𝔽_3,4c_1(𝔽_3)2𝒮)=dimH^0(𝔽_3,4c_1(𝔽_3))$$
(3.17)
This establishes the claim for the section $`g_2𝒪_B(4c_1(𝔽_3))`$. The proof is similar for the section $`g_3𝒪_B(6c_1(𝔽_3))`$ and establishes that
$$g_3=s_{2𝒮}G_3$$
(3.18)
where
$$s_{2𝒮}𝒪_B(2𝒮),G_3𝒪_B(6c_1(𝔽_3)2𝒮)$$
(3.19)
Simply put, Claim I says that if in local coordinates u,v on $`B=𝔽_3`$ the curve $`𝒮`$ is given by $`u=0`$, then, near $`𝒮`$, $`g_2`$ and $`g_3`$ must be of the form
$$g_2=u^2G_2,g_3=u^2G_3$$
(3.20)
However, it does not specify the form of $`G_2`$ and $`G_3`$, which may or may not vanish on $`𝒮`$. The properties of $`G_2`$ and $`G_3`$ are further refined in a second claim, Claim II, which we state without proof.
* The divisors $`4c_1(𝔽_3)2𝒮`$ and $`6c_1(𝔽_3)2𝒮`$ are “very ample”.
For a divisor $`D`$ in a space $`B`$ to be very ample means that a) for any point $`pB`$ there must be a curve in the class of $`D`$ that passes through $`p`$ and b) for any two disjoint points $`p,qB`$ there must be a curve in the class of $`D`$ that passes through $`p`$ but not $`q`$, and a curve that passes through $`q`$ but not $`p`$. Claim II tells us that there must exist sections $`g_2`$ and $`g_3`$ that vanish at exactly order 2 along $`𝒮`$, that is, for which $`G_2(0,v)`$ and $`G_3(0,v)`$ are non-vanishing.
Using these results, we now turn to the discriminant curve $`\mathrm{\Delta }`$ defined in (3.2) and the consequences of (3.7). The third assertion, Claim $`III`$, is the following.
* The zero locus of the discriminant $`\mathrm{\Delta }=g_2^327g_3^2𝒪_B(12c_1(𝔽_3))`$ must necessarily vanish along $`𝒮`$ with order exactly 4. Furthermore, the zero locus of $`\mathrm{\Delta }`$ also vanishes along the curve described by $`\mathrm{\Sigma }=20𝒮+60`$. The curves $`𝒮`$ and $`\mathrm{\Sigma }`$ do not intersect, that is, $`𝒮\mathrm{\Sigma }=0`$.
The proof of Claim III is as follows. First note from the definition of the discriminant section $`\mathrm{\Delta }=g_2^327g_3^2`$ and (3.11), (3.18) that
$$\mathrm{\Delta }=s_{2𝒮}^2\mathrm{\Delta }_{12c_1(B)4𝒮}$$
(3.21)
where
$$\mathrm{\Delta }_{12c_1(B)4𝒮}=s_{2𝒮}G_2^327G_3^2$$
(3.22)
¿From Claim I we know that the zero locus of $`\mathrm{\Delta }`$ vanishes along $`𝒮`$ with order at least 4. A in Claim I, we can show that the zero locus of $`\mathrm{\Delta }`$ vanishes along $`𝒮`$ with order exactly 4. Simply put, this implies, in terms of the local coordinates u,v on $`B=𝔽_3`$, that near $`𝒮`$ the discriminant $`\mathrm{\Delta }`$ must be of the form
$$\mathrm{\Delta }=u^4(27G_3(0,v)^3)$$
(3.23)
where the factor $`27G_3(0,v)^3`$ is non-vanishing. Now, note from (3.9) that
$$12c_1(𝔽_3)4𝒮=20𝒮+60$$
(3.24)
We conclude that $`\mathrm{\Delta }`$ must also vanish along a curve of the form $`\mathrm{\Sigma }=20𝒮+60`$. Finally, using equation (3.8) it is easy to see that $`𝒮\mathrm{\Sigma }=0`$. This establishes Claim III. Finally, we prove a fourth assertion, Claim IV.
* $`𝒮`$ is a smooth rational curve in $`B=𝔽_3`$. A torus over any point on the curve $`𝒮`$ degenerates as Kodaira type $`IV`$. The curve $`\mathrm{\Sigma }`$ in $`B=𝔽_3`$ is smooth except at 200 cusps. The fibers degenerate as Kodaira type $`I_1`$ over the smooth parts of the curve and as Kodaira type $`II`$ over each of the cusps.
To prove Claim IV, first consider the curve $`𝒮`$. It follows from (3.20), Claim II and (3.23) that, in local coordinates near $`𝒮`$, sections $`g_2`$, $`g_3`$ and $`\mathrm{\Delta }`$ have the form
$$g_2=u^2G_2,g_3=u^2G_3,\mathrm{\Delta }=u^4(27G_3(0,v)^3)$$
(3.25)
Using Table 1, we find that any torus over a point on the curve $`𝒮`$ degenerates as a Kodaira fiber of Type $`IV`$. Now consider the curve specified by $`\mathrm{\Sigma }`$. There are two ways in which the section $`\mathrm{\Delta }_{12c_1(B)4𝒮}=s_{2𝒮}G_2^327G_3^2`$ can vanish on $`\mathrm{\Sigma }`$. First, $`G_2`$ and $`G_3`$ can vanish simultaneously. We see from (3.12) and (3.19) that this will occur precisely at the points of intersection of the divisors $`4c_1(𝔽_3)2𝒮`$ and $`6c_1(𝔽_3)2𝒮`$. Since these divisors, by Claim II, are very ample, then, $`G_2`$ and $`G_3`$ can be chosen so that, near each such intersection point, they are local coordinates of the base. Denoting $`G_2=s`$ and $`G_3=t`$, then the discriminant curve is described by the equation
$$f(s,t)s^327t^2=0$$
(3.26)
where the function $`f(s,t)`$, which represents $`u^2`$ evaluated in $`s,t`$ coordinates, is non-vanishing in the neighborhood of the curve. One can check that this curve has a cusp at the point $`(s,t)=(0,0)`$. That is, each intersection point corresponds to a cusp of the curve $`\mathrm{\Sigma }`$. It follows that in the neighborhood of a cusp, the Weierstrass equation (3.1) for the elliptic fiber can be written as
$$y^2=4x^3f(s,t)sxf(s,t)t$$
(3.27)
where we have used the affine coordinate $`z=1`$. This defines a smooth threefold. Exactly at the cusp, this equation becomes
$$y^2=4x^3$$
(3.28)
which is well known to describe a fiber of Kodaira type $`II`$. Clearly, the total number of such cusps, $`N_{cusps}`$, of the curve $`\mathrm{\Sigma }`$ is given by
$$N_{cusps}=(4c_1(𝔽_3)2𝒮)(6c_1(𝔽_3)2𝒮)=200$$
(3.29)
where we have used expressions (3.8) and (3.9). Away from these 200 cusps, the curve $`\mathrm{\Sigma }`$ is smooth. Along the smooth part of the curve, we can introduce new local coordinates $`S=f(s,t)s^327t^2`$ and an independent variable $`T`$ such that $`\mathrm{\Sigma }`$ is defined by $`S=0`$. We see then that the sections $`g_2`$, $`g_3`$ and $`\mathrm{\Delta }`$, in a neighborhood of the smooth part of the curve $`\mathrm{\Sigma }`$, have the form
$$g_2=S^0𝒢_2,g_3=S^0𝒢_3,\mathrm{\Delta }=S^1f^2$$
(3.30)
where $`𝒢_2,𝒢_2`$ and $`f`$ are functions of $`S`$ and $`T`$ which do not vanish on $`\mathrm{\Sigma }`$. It follows from Table 1 that the torus over any point on the smooth part of the curve $`\mathrm{\Sigma }`$ degenerates as Kodaira type $`I_1`$. This completes the proof of Claim IV.
| Kodaira type | A-D-E | monodromy | N,L,K |
| --- | --- | --- | --- |
| $`I_n`$ | $`A_{n1}`$ | $`\left(\begin{array}{cc}1& n\\ 0& 1\end{array}\right)`$ | $`N=n,L=0,K=0`$ |
| $`II`$ | | $`\left(\begin{array}{cc}1& 1\\ 1& 0\end{array}\right)`$ | $`N=2,L>0,K=1`$ |
| $`III`$ | $`A_1`$ | $`\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$ | $`N=3,L=1,K>1`$ |
| $`IV`$ | $`A_2`$ | $`\left(\begin{array}{cc}0& 1\\ 1& 1\end{array}\right)`$ | $`N=4,L>1,K=2`$ |
| $`I_0^{}`$ | $`D_4`$ | $`\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)`$ | $`N=6,L>1,K>2`$ |
| $`I_n^{}`$ | $`D_{n+4}`$ | $`\left(\begin{array}{cc}1& n\\ 0& 1\end{array}\right)`$ | $`N=6+n,L=2,K=3`$ |
| $`IV^{}`$ | $`E_6`$ | $`\left(\begin{array}{cc}1& 1\\ 1& 0\end{array}\right)`$ | $`N=8,L>2,K=4`$ |
| $`III^{}`$ | $`E_7`$ | $`\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$ | $`N=9,L=3,K>4`$ |
| $`II^{}`$ | $`E_8`$ | $`\left(\begin{array}{cc}0& 1\\ 1& 1\end{array}\right)`$ | $`N=10,L>3,K=4`$ |
Table 1: The integers $`N`$, $`L`$ and $`K`$ characterize the behavior of $`\mathrm{\Delta }`$, $`g_2`$ and $`g_3`$ near the discriminant locus $`u=0`$; $`\mathrm{\Delta }=u^Na,g_2=u^Lb`$ , and $`g_3=u^Kc`$ .
To conclude, the discriminant curve for an elliptically fibered Calabi–Yau threefold with base $`B=𝔽_3`$ has the following structure. First,
$$\mathrm{\Delta }=𝒮\mathrm{\Sigma }$$
(3.31)
where
$$\mathrm{\Sigma }=20𝒮+60$$
(3.32)
and
$$𝒮\mathrm{\Sigma }=0$$
(3.33)
The component curve $`𝒮`$ is smooth, whereas the other component curve $`\mathrm{\Sigma }`$ is smooth except at 200 points, where it degenerates into cusps. We first list the Kodaira type for the fibers over the smooth parts of these curves.
* $`𝒮`$– Kodaira type $`IV`$
* $`\mathrm{\Sigma }`$– Kodaira type $`I_1`$
The Kodaira type over the cusp points of $`\mathrm{\Sigma }`$ are
* $`\mathrm{\Sigma }`$– cusps– Kodaira type $`II`$
This is shown pictorially in Figure $`1`$.
We now apply this technology to the case of interest in this paper, an elliptically fibered Calabi–Yau threefold with the base $`B=\widehat{𝔽}_3`$ introduced above. Since the computations are similar to the ones just described, we will only present the results. (These examples were first constructed torically in ).
### Example 2: $`B=\widehat{𝔽}_3`$
Consider a smooth elliptically fibered Calabi–Yau threefold $`X`$ with zero section $`\sigma `$ and base $`B=\widehat{𝔽}_3`$. As discussed previously, there are three effective classes $`𝒮`$, $``$ and $`𝒢`$ that together form a basis of $`H_2(\widehat{𝔽}_3,)`$ with intersection numbers
$$𝒮𝒮=3,=1,𝒢𝒢=1$$
(3.34)
and
$$𝒮=1,𝒮𝒢=0,𝒢=1$$
(3.35)
The first and second Chern classes of $`B=\widehat{𝔽}_3`$ are given by
$$c_1(\widehat{𝔽}_3)=2𝒮+5+4𝒢,c_2(\widehat{𝔽}_3)=5$$
(3.36)
It is helpful in the following to define the class
$$\sigma =𝒮+3+2𝒢$$
(3.37)
$`X`$ is the resolution of a Weierstrass fibration over the base $`B=\widehat{𝔽}_3`$ with sections $`g_2`$, $`g_3`$ satisfying (3.6) and discriminant $`\mathrm{\Delta }`$ defined by (3.2) and satisfying (3.7). Using this data, one can, as outlined above, compute the discriminant curves. We will simply state the results.
### Case 1:
The discriminant curve is composed of three components
$$\mathrm{\Delta }=𝒮\sigma \mathrm{\Sigma }$$
(3.38)
where
$$\mathrm{\Sigma }=16𝒮+54+44𝒢$$
(3.39)
It follows from (3.34) and (3.35) that
$$𝒮\sigma =0,\sigma \mathrm{\Sigma }=44,𝒮\mathrm{\Sigma }=6$$
(3.40)
The intersection number $`\sigma \mathrm{\Sigma }=44`$ corresponds to 36 intersection points of $`\sigma `$ with $`\mathrm{\Sigma }`$, 28 of which are simple normal crossing, whereas 8 have tangential intersections of multiplicity 2. The 6 intersection points of $`𝒮`$ with $`\mathrm{\Sigma }`$ are all simple normal crossings. The number of cusps of the component curve $`\mathrm{\Sigma }`$ (see ) is
$$N_{cusps}=150$$
(3.41)
Away from the intersection points and the cusps, the discriminant curve is smooth. The Kodaira types for the smooth parts of the component curves $`\sigma `$ and $`\mathrm{\Sigma }`$ are given by
* $`\sigma `$– Kodaira type $`I_2`$
* $`\mathrm{\Sigma }`$– Kodaira type $`I_1`$
* $`𝒮`$– Kodaira type $`I_0^{}`$
The Kodaira type of the $`\sigma \mathrm{\Sigma }`$ and $`𝒮\mathrm{\Sigma }`$ intersection points are
* $`\sigma \mathrm{\Sigma }`$– 28 simple normal crossings– Kodaira type $`I_3`$
* $`\sigma \mathrm{\Sigma }`$– 8 tangential crossings– Kodaira type $`III`$
* $`𝒮\mathrm{\Sigma }`$– 6 simple normal crossings– Not Kodaira
Finally, over each of the cusp points of curve $`\mathrm{\Sigma }`$ we find
* $`\mathrm{\Sigma }`$– cusps– Kodaira type $`II`$
This is shown pictorially in Figure 2.
We now present three more discriminant curves. In all of these, the discriminant is composed of the three component curves
$$\mathrm{\Delta }=𝒮\sigma \mathrm{\Sigma }$$
(3.42)
where
$$𝒮\sigma =0,\sigma \mathrm{\Sigma }=44$$
(3.43)
The intersection number $`\sigma \mathrm{\Sigma }=44`$ corresponds to 36 intersection points of $`\sigma `$ with $`\mathrm{\Sigma }`$, 28 of which correspond to simple normal crossing whereas 8 have tangential intersections of multiplicity 2. The component curve $`\mathrm{\Sigma }`$ always has a finite number of cusps. Away from the intersection points and the cusps, the discriminant curve is smooth. The Kodaira type for the smooth parts of the component curves $`\sigma `$ and $`\mathrm{\Sigma }`$ are in all cases given by
* $`\sigma `$– Kodaira type $`I_2`$
* $`\mathrm{\Sigma }`$– Kodaira type $`I_1`$
Similarly, in all cases the Kodaira type of the $`\sigma \mathrm{\Sigma }`$ intersection points are
* $`\sigma \mathrm{\Sigma }`$– 28 simple normal crossings– Kodaira type $`I_3`$
* $`\sigma \mathrm{\Sigma }`$– 8 tangential crossings– Kodaira type $`III`$
The final shared characteristic is that over the cusp points of the curve $`\mathrm{\Sigma }`$ we always have
* $`\mathrm{\Sigma }`$– cusps– Kodaira type $`II`$
### Case 2:
The component curve $`\mathrm{\Sigma }`$ is given by
$$\mathrm{\Sigma }=14𝒮+54+44𝒢$$
(3.44)
with intersection number
$$𝒮\mathrm{\Sigma }=12$$
(3.45)
This corresponds to 6 intersection points of $`𝒮`$ with $`\mathrm{\Sigma }`$, each a tangential crossing with multiplicity 2. In addition, the number of cusps of the curve $`\mathrm{\Sigma }`$ is
$$N_{cusps}=142$$
(3.46)
Over the smooth parts of component curve $`𝒮`$ we find
* $`𝒮`$– Kodaira type $`IV^{}`$
The Kodaira type over the remaining intersection points are
* $`𝒮\mathrm{\Sigma }`$– 6 tangential crossing– Not Kodaira
### Case 3:
The component curve $`\mathrm{\Sigma }`$ is given by
$$\mathrm{\Sigma }=18𝒮+54+44𝒢$$
(3.47)
with intersection number given by
$$𝒮\mathrm{\Sigma }=0$$
(3.48)
In addition, the number of cusps of the curve $`\mathrm{\Sigma }`$ is
$$N_{cusps}=152$$
(3.49)
The Kodaira type for the fibers over the smooth part of the component curve $`𝒮`$ is
* $`𝒮`$– Kodaira type $`IV`$
### Case 4:
The component curve $`\mathrm{\Sigma }`$ is given by
$$\mathrm{\Sigma }=13𝒮+54+44𝒢$$
(3.50)
The remaining intersection number is now given by
$$𝒮\mathrm{\Sigma }=15$$
(3.51)
The intersection number $`𝒮\mathrm{\Sigma }=15`$ corresponds to 5 intersection points of $`𝒮`$ with $`\mathrm{\Sigma }`$, all of which are tangential crossings of multiplicity 3. In addition, the number of cusps of the curve $`\mathrm{\Sigma }`$ is
$$N_{cusps}=137$$
(3.52)
The Kodaira type for the fibers over the smooth part of $`𝒮`$ is
* $`𝒮`$– Kodaira type $`III^{}`$
The Kodaira type over the remaining intersection points are
* $`𝒮\mathrm{\Sigma }`$– 5 tangential crossings– Kodaira type $`III^{}`$
To conclude, we have presented a general formalism for constructing the discriminant curves of elliptically fibered Calabi–Yau threefolds. This formalism was applied to a specific Calabi–Yau threefold with a blown–up Hirzebruch base $`B=\widehat{𝔽}_3`$. Each of the discriminant curves was shown to be composed of smooth sections, as well as cusps, simple normal crossing points and tangential crossings. In this paper, we will be concerned with five–branes wrapped on elliptic fibers near the smooth parts of the discriminant curves only. We leave the discussion of five–branes near the cusps, simple normal crossings and tangential crossing points to another publication.
## 4 Monodromy, Massless States and Brane Junctions:
In this section, we will compute the massless spectrum on the worldvolume of five–branes wrapped on elliptic fibers near the smooth parts of discriminant curves. This is accomplished by demonstrating the equivalence, in this context, of degenerating M membranes to string junctions, and using the junction lattice techniques developed in . For concreteness, we present our results within the context of the three–family, $`SU(5)`$ model compactified on a Calabi–Yau threefold with blown–up Hirzebruch base $`B=\widehat{𝔽}_3`$ presented above. This formalism, however, is applicable to any theory compactified on elliptically fibered Calabi–Yau threefolds.
We begin by briefly reviewing the worldvolume theory of an M–theory five–brane wrapped on a smooth elliptic fiber $`𝒞_2`$ located far from the discriminant locus. First, consider a five–brane in flat space–time. The massless worldvolume degrees of freedom then form a six–dimensional $`(2,0)`$ supersymmetric tensor multiplet. The bosonic content of this multiplet consists of five scalar fields, $`X^i`$, $`i=1,..,5`$ describing translations in the transverse directions together with a two–form $``$ whose field strength $`H=d`$ is self–dual, $`H=H`$. When the five–brane is wrapped on a smooth elliptic fiber of a Calabi–Yau threefold, four of the scalars become moduli for the location of the fiber, that is, they are coordinates of the base surface $`B`$. We denote these four fields by $`u_1,u_2,v_1`$, and $`v_2`$. The fifth scalar, which we denote by $`y`$, continues to parameterize translations in the remaining transverse direction. The behavior of the self–dual field strength is a little more complicated. The cohomology classes of a smooth elliptic fiber are those of a real two–torus, that is, $`H^0`$, $`H^1`$ and $`H^2`$ with dimensions one, two and one respectively. We denote a basis of $`H^0`$ by $`1`$, of $`H^1`$ by the two harmonic one–forms $`\lambda _c`$, $`c=1,2`$ and of $`H^2`$ by the harmonic volume form $`\mathrm{\Omega }`$. If we decompose the five–brane worldvolume as $`𝒞_2\times M_4`$, then the field strength $`H`$ can be expanded as
$$H=da\mathrm{\Omega }+F^c\lambda _c+h$$
(4.1)
where the four–dimensional fields are a scalar “axion” $`a`$, two $`U(1)`$ vector fields $`F^c=dA^c`$ and a three–form field strength $`h=db`$. However, not all of these fields are independent because of the self–duality condition $`H=H`$. Applying this condition leads to the constraints
$$F^2=F^1,h=da$$
(4.2)
We conclude that the two–form $``$ decomposes into an axionic scalar field $`a`$ and a single $`U(1)`$ gauge connection $`A_\mu `$. Putting everything together, we find that for a five–brane compactified on a smooth elliptic fiber $`𝒞_2`$, the bosonic worldvolume fields on $`M_4`$ are $`A_\mu ,a,y,u_1,u_2,v_1,v_2`$. These fields, along with their fermionic superpartners, form an $`N=4`$ vector supermultiplet. The low energy worldvolume field theory on $`M_4`$ exhibits $`N=4`$ supersymmetry, but this is broken to $`N=1`$ supersymmetry by higher dimension operators. The $`N=4`$ vector supermultiplet decomposes under $`N=1`$ supersymmetry into an Abelian Yang–Mills multiplet with bosonic field $`A_\mu `$, and three chiral multiplets with bosonic fields $`Y=y+ia`$, $`u=u_1+iu_2`$ and $`v=v_1+iv_2`$ respectively. Note that none of these chiral supermultiplets is charged. Before continuing, we record the fact that the first homology group of a smooth elliptic fiber, $`H_1(𝒞_2,)`$, has dimension two and a basis of one–cycles $`\omega _c`$, $`c=1,2`$. It follows that, in this basis, any cycle in $`H_1(𝒞_2,)`$ is specified by $`(p,q)`$, where $`p,q`$.
We see from the examples given in the previous section that elliptic fibers above different points on the smooth parts of discriminant curves can have very different degeneration characteristics. The simplest, which, for example, occurs at the smooth parts of the $`\mathrm{\Sigma }`$ component curve, is associated with Kodaira type $`I_1`$. We will discuss this case first.
### Kodaira Type $`I_1`$:
In the neighborhood of a smooth point on the discriminant curve, one can always define two special complex coordinates, a coordinate $`u`$ transverse to the curve and a coordinate $`v`$ along the curve. Note that previously these coordinates were designated by various symbols, such as $`S`$ and $`T`$, but, henceforth, we will call the complex coordinates near any component of the discriminant curve $`u`$ and $`v`$. Now pick an arbitrary point on the smooth part of the discriminant curve with a fiber of Kodaira type $`I_1`$, and choose the origin of the $`u,v`$ coordinates to be at that point. Recall from (3.30) that, in the neighborhood of this point, the associated sections have the form
$$g_2=u^0𝒢_2,g_3=u^0𝒢_3,\mathrm{\Delta }=u^1f^2$$
(4.3)
where $`𝒢_2`$, $`𝒢_3`$ and $`f`$, to leading order in $`u`$, are non–zero functions of $`v`$ only. We see from (3.2) that the first two functions must satisfy
$$𝒢_2^327𝒢_3^2=0$$
(4.4)
It follows from (3.1) and (4.3) that, at the origin, the Weierstrass form of the elliptic fiber becomes
$$y^2=4x^3𝒢_2x𝒢_3$$
(4.5)
where we have used affine coordinates with $`z=1`$. This is the standard Weierstrass representation of an elliptic fiber of Kodaira type $`I_1`$.
In the neighborhood of a point on the smooth part of the discriminant curve, the symmetry between the complex moduli $`u`$ and $`v`$, that exists far from the discriminant locus, is broken. Hence, near such a point, the $`N=4`$ supersymmetry of the low energy $`M_4`$ worldvolume theory of the wrapped five–brane is reduced to $`N=2`$ supersymmetry. The $`N=4`$ vector supermultiplet discussed above then decomposes into two $`N=2`$ supermultiplets, an Abelian Yang-Mills multiplet with bosonic field content $`A_\mu `$, $`u`$ and a hypermultiplet with the bosonic fields $`Y`$, $`v`$. As is the case far from the discriminant curve, the supersymmetry is broken to $`N=1`$ supersymmetry by higher dimension operators.
To discuss the monodromy, we must restrict to the fibration over a curve that intersects the discriminant locus at the origin. The “generic” curves are transverse to the discriminant locus, that is, their intersection multiplicity with the discriminant locus is $`1`$. Surfaces over curves with multiplicity greater than $`1`$ are singular at the discriminant point, while those over generic curves are smooth. Therefore, we will choose the intersection surface to be the fibration over a generic curve. Note that the path $`𝒫`$ defined by $`v=0`$, which has coordinate $`u`$, is tranverse to the discriminant curve and, hence, is generic. For specificity, we will take the intersecting surface to be the elliptic fibration over $`𝒫`$, which we denote by $`𝒯`$. Restricted to this surface, the discriminant curve appears as a point in the one–dimensional complex base $`𝒫`$. Furthermore, the forms of $`g_2`$, $`g_3`$, and $`\mathrm{\Delta }`$, as well as the Weierstrass form of the elliptic fiber, restricted to this twofold remain those given in (4.3) and (4.5) with $`𝒢_2`$, $`𝒢_3`$, and $`f`$ evaluated at $`v=0`$. It follows from Table $`1`$ that, within $`𝒯`$, the degeneration of the fiber near the discriminant point remains Kodaira type $`I_1`$. Since the problem has now been reduced to an elliptic twofold over a one–dimensional base, we can apply standard Kodaira theory to analyze the monodromy.
Before proceeding, it is important to discuss the structure of the elliptic fibration $`𝒯`$ over the generic base curve $`𝒫`$. $`𝒯`$ will be a $`K3`$ twofold if and only if $`𝒫𝒫=0`$, which guarantees $`c_1(𝒯)=0`$ and $`𝒫`$ is rational, which implies that $`H_1(𝒯,)=0`$. For the specific threefold base $`B=\widehat{𝔽}_3`$ and the smooth parts of the discriminant curves discussed in this paper, both of these properties are satisfied. However, they need not be true for general threefold bases and discriminant curves, and have to be checked on a case by case basis. In general, $`𝒯`$ need not be a $`K3`$ surface.
The $`SL(2,)`$ monodromy transformation for a type $`I_1`$ fiber can be found from Table $`1`$, and is given by
$$A=\left(\begin{array}{cc}1& 1\\ 0& 1\end{array}\right)$$
(4.6)
This transformation acts on the elements of $`H_1(𝒞_2,)`$ and has, up to multiplication by non–zero integers, a single eigenvector
$$\stackrel{}{𝒱}=(1,0)$$
(4.7)
with eigenvalue $`+1`$. The meaning of this result is the following. Consider an elliptic fiber $`𝒞_2`$ over a point, $`z`$, near, but not at, the discriminant locus. Then the one–cycle, $`\stackrel{}{𝒱}=(1,0)`$, is the unique cycle (up to multiplication) in $`H_1(𝒞_2,)`$ that contracts to zero as the base point of the fiber is moved to the discriminant. That is, $`\stackrel{}{𝒱}=(1,0)`$ is the unique “vanishing cycle” (see for example ).
The physical implications of this arise as follows. Fix a curve in the base from point $`z`$ to the discriminant locus. Now consider a membrane in $`𝒯`$ bounded by $`\stackrel{}{𝒱}=(1,0)`$ in the elliptic fiber over $`z`$, whose intersection with the fiber over each point on this curve is the vanishing cycle. Note that this membrane “ends” on the degenerate cycle of the $`I_1`$ fiber over the discriminant locus. A $`𝒯`$–membrane of this type is shown in Figure $`3`$. One can show that membranes constructed in this manner have representatives which are holomorphic two-cycles in $`𝒯`$, albeit with respect to a different complex structure than the usual one compatible with the elliptic fibration. The class of such membranes lies in the relative homology $`H_2(𝒯/F,)`$ and will be denoted by $`\stackrel{}{v}`$. If $`𝒯=K3`$, then the self–intersection number of any class $`\stackrel{}{v}`$ containing a holomorphic two–cycle is given by
$$\stackrel{}{v}\stackrel{}{v}=2g2+b$$
(4.8)
where $`g`$ is the genus and $`b`$ is the number of boundaries. In our case, $`g=0,b=1`$ and, hence
$$\stackrel{}{v}\stackrel{}{v}=1$$
(4.9)
When $`𝒯`$ is not a $`K3`$ surface, the situation requires further analysis. However, we are able to show that for a general surface $`𝒯`$ self–intersection (4.9) continues to hold. This is done as follows. Consider a second membrane $`\stackrel{}{v}^{}`$ homologous to $`\stackrel{}{v}`$. These membranes intersect only at the collapsed cycle in the Kodaira fiber over the discriminant point and at the $`(1,0)`$ cycle in the fiber $`𝒞_2`$ over $`z`$. Using the fact that the self–intersection of the $`(1,0)`$ cycle is zero, it follows that the intersection number of $`\stackrel{}{v}^{}`$ with $`\stackrel{}{v}`$ and, hence, the self–intersection of $`\stackrel{}{v}`$ is $`\pm 1`$. Since this result is topological, we see from (4.9) that for any surface $`𝒯`$, not just $`K3`$, $`\stackrel{}{v}\stackrel{}{v}=1`$.
As long as the five–brane is wrapped on an elliptic fiber reasonably far from the discriminant, the fluctuations of this membrane are massive and can be ignored in the low energy effective theory. However, as the elliptic fiber approaches the discriminant point, these fluctuations become less and less heavy, finally becoming exactly massless when the five–brane is wrapped on the degenerate $`I_1`$ Kodaira fiber. Therefore, at the discriminant locus we expect light BPS states to enter the $`M_4`$ low energy theory of the wrapped five–brane.
To compute these states, first note that any cycle of the form $`(p,0)`$, where $`p`$, is an eigenvector of the monodromy $`A`$, not just $`(1,0)`$. The class of the associated membrane is denoted by
$$\stackrel{}{J}=n\stackrel{}{v}$$
(4.10)
where $`n`$. The vanishing cycle associated with this class is given by
$$(p,q)=n(1,0)$$
(4.11)
and, hence, $`p=n`$ and $`q=0`$. It follows from (4.9) that the self–intersection number is
$$\stackrel{}{J}\stackrel{}{J}=n^2$$
(4.12)
At this point, it is very useful to note from Figure $`3`$ that the projection of such a membrane into the base produces a string indexed by $`\stackrel{}{v}`$, integer $`n`$ and $`(p,q)`$. These are the fundamental objects that compose string junction lattices , which we review briefly in the Appendix. Furthermore, the self–intersection number of class $`\stackrel{}{v}`$ given in (4.9) coincides with the norm of the string $`\stackrel{}{v}`$ defined in . Therefore, computations from the points of view of the $`𝒯`$–membrane and the string junction lattice coincide. We will, therefore, use either method interchangeably. Note that, in general, membranes with boundary on the $`(p,q)`$ cycle within a wrapped five–brane give rise to states with electric and magnetic charges $`p`$ and $`q`$. As discussed in the Appendix, the condition for the associated state to be BPS saturated is that
$$\stackrel{}{J}\stackrel{}{J}2+gcd(p,q)$$
(4.13)
where “gcd” stands for the greatest common positive divisor of $`p`$ and $`q`$. In the case under consideration, since $`q=0`$, this condition becomes
$$\stackrel{}{J}\stackrel{}{J}|n|2$$
(4.14)
Comparing this with (4.12), it is clear that the state will be BPS saturated if and only if $`n=\pm 1`$. This CPT conjugate pair of allowed BPS states, $`([1],[1])`$, combine to form a single, stable $`N=2`$ hypermultiplet $`\mathrm{\Phi }_{(1,0)}`$, of electric charge $`Q_e=1`$ and vanishing magnetic charge, in the $`M_4`$ worldvolume theory of the wrapped five–brane. The non–BPS states are either unstable or massive.
We conclude that, near a point on a smooth part of the discriminant curve with elliptic fiber of Kodaira type $`I_1`$, the $`M_4`$ worldvolume theory of a wrapped five–brane has $`N=2`$ supersymmetry at low energy. In addition to the “standard” Abelian Yang–Mills supermultiplet with gauge connection $`A_\mu `$ and an uncharged hypermultiplet, the $`I_1`$ degeneracy of the elliptic fiber produces a light BPS hypermultiplet
* $`\mathrm{\Phi }_{(1,0)}`$ $`Q_e=1`$, $`Q_m=0`$
This electric charge couples to the gauge connection $`A_\mu `$.
The next simplest possibility, which, for example, occurs at the smooth parts of the $`\sigma `$ component curve, is associated with Kodaira type $`I_2`$. We will discuss this case in the next subsection.
### Kodaira Type $`I_2`$:
Pick any point on the smooth part of the discriminant curve with a fiber of Kodaira type $`I_2`$, and choose the origin of the $`u,v`$ coordinates to be at that point. In the neighborhood of this point, the associated sections have the form
$$g_2=u^0𝒢_2,g_3=u^0𝒢_3,\mathrm{\Delta }=u^2f^2$$
(4.15)
where $`𝒢_2`$, $`𝒢_3`$, and $`f`$, to leading order in $`u`$, are non–zero functions of $`v`$ only (see (3.30)). We see from (3.2) that the first two functions must satisfy (4.4). It follows from (3.1) and (4.15) that, at the origin, the Weierstrass form for the elliptic fiber is again given by (4.5). This is the standard Weierstrass representation of an elliptic fiber of Kodaira type $`I_2`$. As discussed above, near such a point, the $`N=4`$ supersymmetry of the low energy $`M_4`$ worldvolume theory of the wrapped five–brane is reduced to $`N=2`$ supersymmetry. This is broken to $`N=1`$ supersymmetry by higher dimension operators.
As above, we can restrict the discussion to the elliptic fibration $`𝒯`$ over the path $`𝒫`$ with coordinate $`u`$ that is transverse to the discriminant curve . The discriminant curve now appears as a point in the one–dimensional complex base $`𝒫`$. Furthermore, the forms of $`g_2`$, $`g_3`$ and $`\mathrm{\Delta }`$, as well as the Weierstrass form of the elliptic fiber, restricted to this twofold remains those given in (4.15) and (4.5) with $`𝒢_2`$, $`𝒢_3`$ and $`f`$ evaluated at $`v=0`$. Then it follows from Table $`1`$ that, within $`𝒯`$, the degeneracy of the fiber near the discriminant point remains Kodaira type $`I_2`$.
The $`SL(2,)`$ monodromy transformation for a type $`I_2`$ fiber can be found from Table $`1`$, and is given by
$$_{I_2}=\left(\begin{array}{cc}1& 2\\ 0& 1\end{array}\right).$$
(4.16)
This transformation acts on the elements of $`H_1(𝒞_2,)`$ and has, up to multiplication by a non–zero integer, the single eigenvector (4.7) with eigenvalue $`+1`$. At this point, however, the situation diverges substantially from that of the Kodaira type $`I_1`$ case above. This is signaled by the appearance of the 2 in the monodromy matrix, which implies that it can be decomposed as
$$_{I_2}=\left(\begin{array}{cc}1& 1\\ 0& 1\end{array}\right)\left(\begin{array}{cc}1& 1\\ 0& 1\end{array}\right)=AA$$
(4.17)
It follows that there are now two copies of the eigenvector $`\stackrel{}{𝒱}=(1,0)`$ that are relevant to the problem. One way to explore this issue is to locally deform the Weierstrass equation and, hence. the discriminant in such a way that its locus, the single point at the origin with Kodaira type $`I_2`$, is split into two nearby discriminant loci, which we label $`1`$ and $`2`$, each with Kodaira type $`I_1`$. Then, the monodromy at each of these points is simply given by (4.6), with a single eigenvector (4.7). A deformation of the Weierstrass equation that will accomplish this is given by
$$y^2=4x^3𝒢_2x𝒢_3+u(u+ϵ)b$$
(4.18)
Here, $`b`$ is the coefficient of the quadratic term in the $`u`$ expansion of $`g_3`$ and $`ϵ`$ is a complex deformation parameter. Note, in passing, that under such a deformation a $`K3`$ twofold remains $`K3`$. The meaning of this result is the following. Consider an elliptic fiber $`𝒞_2`$ over a point, $`z`$, near, but not at, the two discriminant loci. Then the one–cycle $`\stackrel{}{𝒱}=(1,0)`$ is the unique cycle in $`H_1(𝒞_2,)`$ that contracts to zero as the fiber is moved to either of the two discriminant points. That is, $`\stackrel{}{𝒱}`$ is the unique vanishing cycle associated with both discriminant loci.
The physical implications arise from considering a $`𝒯`$–membrane whose boundary in the elliptic fiber over $`z`$ is $`\stackrel{}{𝒱}=(1,0)`$ and which “ends” on the $`I_1`$ fiber over either point $`1`$ or $`2`$ in the discriminant locus. Membranes of this type have a very specific structure. Recall that in the $`I_1`$ case above there was only one type of membrane, which ended on the unique discriminant point. In contrast, in the $`I_2`$ case there are two types of membranes, one type ending on the $`I_1`$ fiber over discriminant point $`1`$ and the other type ending on the $`I_1`$ fiber over discriminant point $`2`$. We will denote the classes associated with discriminant points $`1`$ and $`2`$ by $`\stackrel{}{v_1}`$ and $`\stackrel{}{v_2}`$ respectively. Independently of whether or not $`𝒯`$ is a $`K3`$ surface, we can show that the intersections of these membrane classes are given by
$$\stackrel{}{v_1}\stackrel{}{v_1}=\stackrel{}{v_2}\stackrel{}{v_2}=1,\stackrel{}{v_1}\stackrel{}{v_2}=0$$
(4.19)
The proof for the self–intersections was presented in the section on $`I_1`$ and does not change here. Using the fact that $`\stackrel{}{v}_1`$ and $`\stackrel{}{v}_2`$ intersect only on the $`(1,0)`$ cycle in the fiber $`𝒞_2`$ over $`z`$, and that the self–intersection of this cycle vanishes, it follows that $`\stackrel{}{v}_1\stackrel{}{v}_2=0`$. A generic membrane of this type, projected into the base, is the string junction shown in Figure $`4`$. We note that the membrane intersection numbers in (4.19) are identical to the string junction lattice intersection form presented in (6.3) and (6.4) in the Appendix. Therefore, once again, we can use the membrane and string junction lattice formalism interchangeably. As the elliptic fiber approaches either of the discriminant points, the membrane fluctuations become less and less heavy, becoming exactly massless when the five–brane is wrapped on the degenerate $`I_1`$ Kodaira fiber over point $`1`$ or $`2`$. Therefore, at either of the two discriminant loci, we expect light BPS states to enter the $`M_4`$ low energy theory of the wrapped five–brane.
To compute these states, we first note that any cycle of the form $`(p,0)`$, where $`p`$, is an eigenvector of the monodromy $`_{I_2}`$, not just $`(1,0)`$. The class of the associated membrane is denoted by
$$\stackrel{}{J}=n_1\stackrel{}{v_1}+n_2\stackrel{}{v_2}$$
(4.20)
where $`n_1,n_2`$. The boundary cycle associated with this class is given by
$$(p,q)=n_1(1,0)+n_2(1,0)$$
(4.21)
and, hence, $`p=n_1+n_2`$ and $`q=0`$. It follows from (4.19) that the self–intersection number is
$$\stackrel{}{J}\stackrel{}{J}=n_1^2n_2^2$$
(4.22)
However, recall that (4.13) is the condition for the associated state to be BPS saturated. In the case under consideration, since $`q=0`$, this condition becomes
$$\stackrel{}{J}\stackrel{}{J}|n_1+n_2|2$$
(4.23)
Comparing this with (4.22), it is clear that a state will be BPS saturated if and only if the pair $`n_1,n_2`$ takes the values $`\pm 1,0`$ or $`0,\pm 1`$. The two BPS states described by $`([1,0],[1,0])`$, combine to form a single, stable $`N=2`$ hypermultiplet, $`\mathrm{\Phi }_{(1,0)}^1`$, of electric charge $`Q_e=1`$ and vanishing magnetic charge, in the $`M_4`$ worldvolume theory of the wrapped five–brane. Similarly, the states associated with $`([0,1],[0,1])`$ form a second, stable $`N=2`$ hypermultiplet, $`\mathrm{\Phi }_{(1,0)}^2`$, also with $`Q_e=1`$ and $`Q_m=0`$. The non-BPS states are either unstable or massive.
Further important information can be extracted by rewriting the membrane class $`\stackrel{}{J}`$ in (4.20) as follows. Define classes
$$\stackrel{}{w}_p=\frac{1}{2}(v_1+v_2),\stackrel{}{w}=\frac{1}{2}(v_1v_2)$$
(4.24)
Then, in terms of these classes, $`\stackrel{}{J}`$ can be written as
$$\stackrel{}{J}=p\stackrel{}{w}_p+a\stackrel{}{w}$$
(4.25)
where $`p=n_1+n_2`$ and $`a=n_1n_2`$. It is useful to proceed one step further and write this equation as
$$\stackrel{}{J}=p\stackrel{}{w}_p+aC^1\stackrel{}{\alpha }$$
(4.26)
where $`C^1=\frac{1}{2}`$ and $`\stackrel{}{\alpha }=2\stackrel{}{w}`$. Note that
$$\stackrel{}{\alpha }=v_1v_2$$
(4.27)
which is associated with the cycle $`(p,q)=(0,0)`$ and, hence, corresponds to no boundary cycle at all. Such classes are described by curves from discriminant point $`1`$ to discriminant point $`2`$, as shown in Figure $`5`$. Furthermore, it follows from (4.19) that
$$\stackrel{}{\alpha }\stackrel{}{\alpha }=2$$
(4.28)
We see, by comparing equations (4.26) and (4.27) with (6.6) and (6.7) in the Appendix, that class $`\stackrel{}{\alpha }`$ corresponds to the simple root of the Lie algebra of $`SU(2)`$.
Note that, written as a pair $`[n_1,n_2]`$, class $`\stackrel{}{\alpha }=[1,1]`$. It follows that the BPS states that make up the hypermultiplet $`\mathrm{\Phi }_{(1,0)}^1`$, specified by $`[1,0]`$ and $`[1,0]`$, are related to the BPS states, $`[0,1]`$ and $`[0,1]`$, that make up the hypermultiplet $`\mathrm{\Phi }_{(1,0)}^2`$ through addition and subtraction of the root vector $`\stackrel{}{\alpha }`$. That is
$$[1,0]\stackrel{}{\alpha }=[0,1],[1,0]+\stackrel{}{\alpha }=[0,1]$$
(4.29)
Further addition and subtraction of the root $`\stackrel{}{\alpha }`$ leads to unstable or massive non–BPS states. It follows that hypermultiplets $`\mathrm{\Phi }_{(1,0)}^1`$ and $`\mathrm{\Phi }_{(1,0)}^2`$ form a doublet representation of $`SU(2)`$. The appearance of an $`SU(2)`$ global group can be read off directly from the A-D-E column of Table $`1`$. For Kodaira type $`I_2`$, the A-D-E classification is $`A_1`$, which corresponds to the group $`SU(2)`$.
We conclude that, near a point on a smooth part of the discriminant curve with elliptic fiber of Kodaira type $`I_2`$, the $`M_4`$ worldvolume theory of a wrapped five–brane has $`N=2`$ supersymmetry at low energy. In addition to the “standard” Abelian Yang–Mills supermultiplet with gauge connection $`A_\mu `$ and an uncharged hypermultiplet, the $`I_2`$ degeneracy of the elliptic fiber produces an $`SU(2)`$ doublet $`\mathrm{𝟐}`$ of light BPS hypermultiplets
* $`\mathrm{\Phi }_{(1,0)}^i`$ $`Q_e=1,Q_m=0`$
where $`i=1,2`$. The electric charge couples to the gauge connection $`A_\mu `$.
We now turn our attention to more a more complicated situation, specifically that of a Kodaira type $`IV`$ fiber occuring at the smooth part of the discriminant curve. Such fibers are found, for example, over the smooth part of the component curve $`𝒮`$ in Case $`3`$ above. The analysis of this example contains most of the elements necessary to compute the light BPS states for a fiber of any Kodaira type over a smooth part of the discriminant curve.
### Kodaira Type IV:
Pick any point on the smooth part of the discriminant curve with a fiber of Kodaira type $`IV`$, and choose the origin of the $`u,v`$ coordinates to be at that point. In the neighborhood of this point, the associated sections have the form
$$g_2=u^L𝒢_2,g_3=u^2𝒢_3,\mathrm{\Delta }=u^4f^2$$
(4.30)
where $`L>1`$ and $`𝒢_2`$, $`𝒢_3`$ and $`f`$, to leading order in $`u`$, are non–vanishing functions of $`v`$ only. Further analysis of the explicit curve $`𝒮`$ in Case $`3`$ above fixes the value of coefficient $`L`$ to be $`L=2`$. We see from (3.2) that, near the discriminant locus,
$$\mathrm{\Delta }=g_2^327g_3^2=u^4(27𝒢_3^2)$$
(4.31)
which vanishes at the locus without any further restrictions on $`𝒢_2`$ and $`𝒢_3`$. It follows from (3.1) and (4.30) that, at the origin, the Weierstrass form of the elliptic fiber becomes
$$y^2=4x^3u^2𝒢_2xu^2𝒢_3$$
(4.32)
where we have used affine coordinates with $`z=1`$. This is the Weierstrass representation of an elliptic fiber of Kodaira type $`IV`$. Near such a point, the $`N=4`$ supersymmetry of the low energy $`M_4`$ worldvolume theory of the wrapped five–brane is reduced to $`N=2`$ supersymmetry. This is broken to $`N=1`$ supersymmetry by higher dimensional operators.
As above, we can restrict the discussion to the elliptic fibration over the path $`𝒫`$ with coordinate $`u`$ that is transverse to the discriminant curve . The discriminant curve now appears as a point in the one–dimensional complex base $`𝒫`$. The forms of $`g_2`$, $`g_3`$ and $`\mathrm{\Delta }`$, as well as the Weierstrass form of the elliptic fiber, restricted to this twofold remain those given in (4.30) and (4.32) with $`𝒢_2`$, $`𝒢_3`$ and $`f`$ evaluated at $`v=0`$. It follows from Table $`1`$ that, within $`𝒯`$, the degeneracy of the fiber near the discriminant point remains that of Kodaira type $`IV`$.
The $`SL(2,)`$ monodromy transformation for a type $`IV`$ fiber can be found from Table $`1`$. Here, for technical reasons, it is convenient to use a monodromy matrix conjugate to the one listed in Table $`1`$ given by
$$_{IV}=\left(\begin{array}{cc}1& 1\\ 1& 0\end{array}\right)$$
(4.33)
This transformation acts on the elements of $`H_1(𝒞_2,)`$. However, unlike the previous cases of Kodaira type $`I_1`$ and $`I_2`$, $`_{IV}`$ has no real eigenvector and, therefore, there is no obvious associated vanishing cycle. An indication as how to proceed is given by the fact that $`_{IV}`$ can be decomposed as
$$_{IV}=\left(\begin{array}{cc}1& 1\\ 0& 1\end{array}\right)\left(\begin{array}{cc}1& 1\\ 0& 1\end{array}\right)\left(\begin{array}{cc}1& 1\\ 0& 1\end{array}\right)\left(\begin{array}{cc}2& 1\\ 1& 0\end{array}\right)=AAAB$$
(4.34)
Following the technique employed in the $`I_2`$ case, we proceed by deforming the discriminant curve in such a way that its locus, the single point at the origin with Kodaira type $`IV`$, is split into four nearby discriminant points, which we label $`1`$,$`2`$,$`3`$ and $`4`$. This corresponds to a relevant deformation of the Weierstrass representation given by
$$y^2=4x^3u^2𝒢_2xu^2𝒢_3+ϵ$$
(4.35)
where $`ϵ`$ is a constant deformation parameter. The first three points each have a Kodaira fiber of type $`I_1`$ with the monodromy given in (4.6) and, up to multiplication by a non–zero integer, the single eigenvector (4.7). The fourth point, however, has monodromy
$$B=\left(\begin{array}{cc}2& 1\\ 1& 0\end{array}\right)$$
(4.36)
with, up to multiplication by a non–zero integer, the eigenvector
$$\stackrel{}{\stackrel{~}{𝒱}}=(1,1)$$
(4.37)
At first glance, it is not obvious what Kodaira type one is finding in this last case. However, it is easy to show that $`A`$ is, in fact, conjugate to $`B`$ and, hence, the fiber at point $`4`$ is also of Kodaira type $`I_1`$. The meaning of these results is the following. Consider an elliptic fiber $`𝒞_2`$ over a point, $`z`$, near, but not at, the four discriminant loci. Then the one–cycle $`\stackrel{}{𝒱}=(1,0)`$ is the unique cycle in $`H_1(𝒞_2,)`$ that contracts to zero as the fiber is moved to each of the discriminant points $`1`$, $`2`$, and $`3`$ without encircling $`4`$. That is, $`\stackrel{}{𝒱}`$ is the unique vanishing cycle associated with the first three discriminant loci. Similarly, the one–cycle $`\stackrel{}{\stackrel{~}{𝒱}}=(1,1)`$ is the unique cycle in $`H_1(𝒞_2,)`$ that contracts to zero as the fiber is moved to the discriminant point $`4`$ without encircling $`1,2`$ or $`3`$. That is, $`\stackrel{}{\stackrel{~}{𝒱}}`$ is the unique vanishing cycle associated with the fourth discriminant locus.
The physical implications of this arise as follows. $`𝒯`$– membranes that are bounded by the vanishing cycle $`\stackrel{}{𝒱}=(1,0)`$ in the elliptic fiber over $`z`$ are of a very specific structure. There are three membranes of this type, each “ending” on the $`I_1`$ fiber over discriminant points $`1`$,$`2`$ or $`3`$ respectively. Similarly, membranes that are bounded by the vanishing cycle $`\stackrel{}{\stackrel{~}{𝒱}}=(1,1)`$ have a very specific structure. These membranes “end” on the $`I_1`$ fiber over the discriminant point $`4`$. We will denote the $`𝒯`$–membrane classes associated with discriminant points $`1`$, $`2`$,$`3`$ and $`4`$ by $`\stackrel{}{v_1}`$,$`\stackrel{}{v_2}`$, $`\stackrel{}{v_3}`$ and $`\stackrel{}{v_4}`$ respectively. A generic membrane class of this type, projected into the base, is the string junction shown in Figure $`6`$. Note from the Appendix that the intersections of these membrane classes are given by
$$\stackrel{}{v_1}\stackrel{}{v_1}=\stackrel{}{v_2}\stackrel{}{v_2}=\stackrel{}{v_3}\stackrel{}{v_3}=\stackrel{}{v_4}\stackrel{}{v_4}=1$$
(4.38)
$$\stackrel{}{v_1}\stackrel{}{v_2}=\stackrel{}{v_1}\stackrel{}{v_3}=\stackrel{}{v_2}\stackrel{}{v_3}=0$$
(4.39)
$$\stackrel{}{v_1}\stackrel{}{v_4}=\stackrel{}{v_2}\stackrel{}{v_4}=\stackrel{}{v_3}\stackrel{}{v_4}=\frac{1}{2}$$
(4.40)
As the elliptic fiber approaches any one of the discriminant points, the membrane fluctuations become less and less heavy, becoming exactly massless when the five–brane is wrapped on the degenerate $`I_1`$ Kodaira fiber over point $`1`$,$`2`$,$`3`$ or $`4`$. Therefore, at any of the four discriminant loci, we expect light BPS states to enter the $`M_4`$ low energy theory of the wrapped five–brane.
To compute these states, we first note, by analogy with the above discussion, that the classes of the associated membranes are given by
$$\stackrel{}{J}=\mathrm{\Sigma }_{i=1}^4n_i\stackrel{}{v_i}$$
(4.41)
where $`n_i`$, $`i=1,..,4`$ are arbitrary integers. Using (4.7) and (4.37), it follows that the associated boundary cycles in $`H_1(𝒞_2,)`$ are given by
$$(p,q)=\mathrm{\Sigma }_{i=1}^3n_i(1,0)+n_4(1,1)$$
(4.42)
and, hence, $`p=\mathrm{\Sigma }_{i=1}^4n_i`$ and $`q=n_4`$. From (4.39), we find that the self–intersection number of class $`\stackrel{}{J}`$ is
$$\stackrel{}{J}\stackrel{}{J}=\mathrm{\Sigma }_{i=1}^4n_i^2(n_1+n_2+n_3)n_4$$
(4.43)
However, recall that (4.13) is the condition for the associated state to be BPS saturated. Comparing this to (4.43), it was shown in that a state will be BPS saturated if and only if $`n_1,n_2,n_3,n_4`$ take the values
| $`(n_1,n_2,n_3,n_4)`$ | $`(p,q)`$ |
| --- | --- |
| $`(1,0,0,0),(0,1,0,0),(0,0,1,0)`$ | $`(1,0)`$ |
| $`(1,0,0,0),(0,1,0,0),(0,0,1,0)`$ | $`(1,0)`$ |
| $`(1,0,0,1),(0,1,0,1),(0,0,1,1)`$ | $`(0,1)`$ |
| $`(1,0,0,1),(0,1,0,1),(0,0,1,1)`$ | $`(0,1)`$ |
| $`(0,1,1,1),(1,0,1,1),(1,1,0,1)`$ | $`(1,1)`$ |
| $`(0,1,1,1),(1,0,1,1),(1,1,0,1)`$ | $`(1,1)`$ |
| $`(1,1,1,2)`$ | $`(1,2)`$ |
| $`(1,1,1,2)`$ | $`(1,2)`$ |
| $`(0,0,0,1)`$ | $`(1,1)`$ |
| $`(0,0,0,1)`$ | $`(1,1)`$ |
| $`(1,1,1,1)`$ | $`(2,1)`$ |
| $`(1,1,1,1)`$ | $`(2,1)`$ |
Table 2: The BPS states associated with a fiber of Kodaira type $`IV`$.
The six BPS states in the first row combine to form a triplet of $`N=2`$ hypermultiplets, each of electric charge $`Q_e=1`$ and vanishing magnetic charge, in the $`M_4`$ worldvolume theory of the wrapped five–brane. We denote these hypermultiplets by $`\mathrm{\Phi }_{(1,0)}^i`$, where $`i=1,2,3`$. Similarly, the second and third rows each correspond to a triplet of $`N=2`$ hypermultiplets, which we denote by $`\mathrm{\Phi }_{(0,1)}^i`$ and $`\mathrm{\Phi }_{(1,1)}^i`$ with $`i=1,2,3`$ respectively. As indicated by the notation, the multiplet $`\mathrm{\Phi }_{(0,1)}^i`$ has vanishing electric charge and magnetic charge $`Q_m=1`$, whereas $`\mathrm{\Phi }_{(1,1)}^i`$ carries both electric charge $`Q_e=1`$ and magnetic charge $`Q_m=1`$. The two BPS states in the fourth row combine to form a single, $`N=2`$ hypermultiplet, $`\mathrm{\Psi }_{(1,2)}`$ with $`Q_e=1`$ and $`Q_m=2`$. Finally, the last two rows each correspond to a single $`N=2`$ hypermultiplet, denoted by $`\mathrm{\Psi }_{(1,1)}`$ and $`\mathrm{\Psi }_{(2,1)}`$, with $`Q_e=1`$, $`Q_m=1`$ and $`Q_e=2`$, $`Q_m=1`$ respectively. The non–BPS states are either unstable or massive. When the deformation that split the $`IV`$ fiber is undone, all these BPS multiplets become simultaneously massless as the fivebrane is moved towards the $`IV`$ fiber. These states have mutually non-local $`(p,q)`$ charges, meaning that they can not be made simultaneously purely electric by an $`SL(2,)`$ transformation. This leads to a very exotic low energy theory without a local Lagrangian description. Such theories were obtained originally, in a different context, in . When the fivebrane wraps the type $`IV`$ Kodaira fiber, the low energy theory flows to an interacting fixed point in the infrared.
Further important information can be extracted by rewriting the membrane class $`\stackrel{}{J}`$ as follows. Define classes
$$\stackrel{}{w}_p=\frac{1}{3}(v_1+v_2+v_3),\stackrel{}{w}_q=\frac{1}{3}(v_1+v_2+v_3)v_4$$
(4.44)
and
$$\stackrel{}{w}^1=\frac{1}{3}(2v_1v_2v_3),\stackrel{}{w}^2=\frac{1}{3}(v_1+v_22v_3)$$
(4.45)
Then, in terms of these classes, $`\stackrel{}{J}`$ can be written as
$$\stackrel{}{J}=p\stackrel{}{w}_p+q\stackrel{}{w}_q+a_i\stackrel{}{w}^i$$
(4.46)
where $`p`$ and $`q`$ are given in (4.42) and
$$a_1=n_1n_2,a_2=n_2n_3$$
(4.47)
It is useful to proceed one step further and write equation (4.46) as
$$\stackrel{}{J}=p\stackrel{}{w}_p+q\stackrel{}{w}_q+a_iC^{1ij}\stackrel{}{\alpha }_i$$
(4.48)
where $`\stackrel{}{\alpha }_1=2w^1w^2`$, $`\stackrel{}{\alpha }_2=w^1+2w^2`$ and $`C_{ij}=\stackrel{}{\alpha }_i\stackrel{}{\alpha }_j`$. Note that
$$\stackrel{}{\alpha }_1=v_1v_2,\stackrel{}{\alpha }_2=v_2v_3$$
(4.49)
which are both associated with the cycle $`(p,q)=(0,0)`$ and, hence, each corresponds to no boundary cycle at all. Therefore, $`\stackrel{}{\alpha }_1`$ is described by curves from discriminant point $`1`$ to discriminant point $`2`$ and $`\stackrel{}{\alpha }_2`$ is described by curves from discriminant point $`2`$ to discriminant point $`3`$, as shown in Figure $`6`$. Furthermore, it follows from (4.39) that
$$C_{ij}=\left(\begin{array}{cc}2& 1\\ 1& 2\end{array}\right)$$
(4.50)
which is the Cartan matrix of the Lie algebra of $`SU(3)`$. Hence, the classes $`\stackrel{}{\alpha }_1`$ and $`\stackrel{}{\alpha }_2`$ correspond to the simple roots of the $`SU(3)`$. Note that, written as a four–tuplet $`[n_1,n_2,n_3,n_4]`$, class $`\stackrel{}{\alpha }_1=[1,1,0,0]`$ and class $`\stackrel{}{\alpha }_2=[0,1,1,0]`$.
Consider, for example, the six BPS states in the first row of Table $`2`$, and note that
$$[1,0,0,0]\stackrel{}{\alpha }_1=[0,1,0,0],[0,1,0,0]\stackrel{}{\alpha }_2=[0,0,1,0]$$
(4.51)
$$[1,0,0,0]+\stackrel{}{\alpha }_1=[0,1,0,0],[0,1,0,0]+\stackrel{}{\alpha }_2=[0,0,1,0]$$
(4.52)
Further addition or subtraction of the roots leads to non–BPS states that are either unstable or massive. Hence, the pairs of BPS states $`([1,0,0,0],[1,0,0,0])`$, $`([0,1,0,0],[0,1,0,0])`$ and $`([0,0,1,0],[0,0,1,0])`$ combine to form the hypermultiplets $`\mathrm{\Phi }_{(1,0)}^1,\mathrm{\Phi }_{(1,0)}^2`$ and $`\mathrm{\Phi }_{(1,0)}^3`$ respectively. Furthermore, these hypermultiplets form a triplet $`\mathrm{𝟑}`$ representation of $`SU(3)`$. The same is true for the hypermultiplets $`\mathrm{\Phi }_{(0,1)}^i`$ and $`\mathrm{\Phi }_{(1,1)}^i`$ for $`i=1,2,3`$, each of which transforms as an $`SU(3)`$ $`\mathrm{𝟑}`$ representation. Now consider the two BPS states in the fourth row of Table $`2`$. Addition or subtraction of any root leads immediately to non–BPS states which are either unstable of massive. Hence, the pair of BPS states $`([1,1,1,2],[1,1,1,2])`$ combine to form a hypermultiplet $`\mathrm{\Psi }_{(1,2)}`$, which is a singlet under $`SU(3)`$. The same is true for $`\mathrm{\Psi }_{(1,1)}`$ and $`\mathrm{\Psi }_{(2,1)}`$ which are both $`SU(3)`$ singlets. The appearance of the global group $`SU(3)`$ can be read off directly from the A-D-E column of Table $`1`$. For Kodaira type $`IV`$, the A-D-E classification is $`A_2`$, which corresponds to the group $`SU(3)`$.
We conclude that, near a point on a smooth part of the discriminant curve with elliptic fiber of Kodaira type $`IV`$, the $`M_4`$ worldvolume theory of a wrapped five–brane has $`N=2`$ supersymmetry at low energy. In addition to the “standard” Abelian Yang–Mills supermultiplet with gauge connection $`A_\mu `$ and an uncharged hypermultiplet, the $`IV`$ degeneracy of the elliptic fiber produces light BPS hypermultiplets. These fall into $`SU(3)`$ $`\mathrm{𝟑}`$ representations
* $`\mathrm{\Phi }_{(1,0)}^i`$ $`Q_e=1,Q_m=0`$
* $`\mathrm{\Phi }_{(0,1)}^i`$ $`Q_e=0,Q_m=1`$
* $`\mathrm{\Phi }_{(1,1)}^i`$ $`Q_e=1,Q_m=1`$
with $`i=1,2,3`$ and $`SU(3)`$ singlets
* $`\mathrm{\Psi }_{(1,2)}`$ $`Q_e=1,Q_m=0`$
* $`\mathrm{\Psi }_{(1,1)}`$ $`Q_e=1,Q_m=1`$
* $`\mathrm{\Psi }_{(2,1)}`$ $`Q_e=2,Q_m=1`$
The electric charge couples to gauge connection $`A_\mu `$ whereas the magnetic charge couples to $`\stackrel{~}{A_\mu }`$ defined by $`d\stackrel{~}{A_\mu }=F`$.
We end this section by presenting the BPS states associated with each of the remaining Kodaira types, $`I_0^{}`$ ,$`III^{}`$, and $`IV^{}`$, over the smooth parts of the discriminant curve in the $`B=\widehat{𝔽}_3`$ model.
### Kodaira Types $`I_0^{}`$, $`III^{}`$, and $`IV^{}`$:
The A-D-E symmetry algebra associated with a fiber of Kodaira type $`I_0^{}`$ can be read off from Table $`1`$ and is given by $`D_4`$. The associated global symmetry group is $`SO(8)`$. The low energy theory that arises in the neighborhood of an $`I_0^{}`$ fiber is the same as that of an $`N=2`$, $`SU(2)`$ Yang–Mills theory with four quark flavors. The BPS multiplets in the $`I_0^{}`$ case can be easily found by comparing to Yang-Mills theory results , or by using string junctions as in . The results are summarized in the table below.
| $`(p,q)`$ charges | $`SO(8)`$ representation |
| --- | --- |
| $`(2n,2m)`$ | $`\mathrm{𝟏}`$ |
| $`(2n+1,2m)`$ | $`\mathrm{𝟖}_𝐯`$ |
| $`(2n,2m+1)`$ | $`\mathrm{𝟖}_𝐬`$ |
| $`(2n+1,2m+1)`$ | $`\mathrm{𝟖}_𝐜`$ |
Table 3: The BPS multiplets associated with a fiber of Kodaira type $`I_0^{}`$.
Table $`3`$ is to be read as follows. For each $`n`$ and $`m`$, there is an $`N=2`$ multiplet with the $`(p,q)`$ charges listed in the first column, which transforms in the representation of $`SO(8)`$ listed in the last column. In the first row, $`n`$ and $`m`$ are constrained to be relatively prime, whereas in the remaining rows, $`p`$ and $`q`$ must be relatively prime. There are no further constraints. Note that, in analogy to the Kodaira type $`IV`$ case, there are a finite number of representations of $`SO(8)`$ which appear, specifically singlets and octets only. However, unlike the type $`IV`$ case, each of these representations occurs with infinite multiplicity. Of course not all of these multiplets are simultaneously stable, depending on the location of the five–brane in moduli space.
Let us summarize these results in terms of $`N=2`$ supermultiplets. In addition to the “standard” Abelian vector supermultiplet with gauge connection $`A_\mu `$ and an uncharged hypermultiplet, the $`I_0^{}`$ degeneracy of the elliptic fiber produces light BPS hyper- and vector multiplets. The extra vector multiplets are $`SO(8)`$ singlets with $`p=\pm 2`$ and $`q=0`$
* $`V_\pm `$
As the five–brane approaches the discriminant curve, these combine with the usual uncharged Abelian vector multiplet to form massless $`N=2`$ vector multiplets $`V_a`$ which transform as the adjoint representation of an enhanced $`SU(2)`$ gauge group. The remaining states belong to hypermultiplets
* $`\mathrm{\Phi }_{(p,q)}^i`$ $`Q_e=p`$, $`Q_m=q`$
where the charges $`(p,q)`$ and the $`SO(8)`$ representation multiplets $`i`$ are given in Table $`3`$. As the five–brane approaches the discriminant curve, the hypermultiplets transforming as $`\mathrm{𝟖}_𝐯`$ under $`SO(8)`$ with $`q=0`$ combine to form an $`SU(2)`$ doublet $`\mathrm{𝟐}`$, $`Q_A^i`$, where $`i`$ is the index of the $`\mathrm{𝟖}_𝐯`$ representation and $`A=1,2`$. These correspond to hypermultiplets of four $`SU(2)`$ doublet quark flavors.
The A-D-E symmetry algebras associated with fibers of Kodaira type $`III^{}`$ and $`IV^{}`$ can be read off from Table $`1`$. The associated global symmetry groups are the exceptional groups $`E_7`$ and $`E_6`$ respectively. The complete spectrum of BPS multiplets for these two Kodaira types are much harder to determine, for reasons discussed below, and we will present only partial results . First, note from Table $`4`$ in the Appendix that the monodromy associated with a Kodaira fiber of type $`I_0^{}`$ decomposes into type $`I_1`$ monodromies as $`AAAABC`$. Similarly, one sees that Kodaira fibers of type $`III^{}`$ and $`IV^{}`$ decompose as $`AAAAAABCC`$ and $`AAAAABCC`$ respectively. It follows that the $`I_0^{}`$ string junction lattice is a sublattice of the $`III^{}`$ and $`IV^{}`$ junction lattices. This reflects the fact that a fiber of type $`I_0^{}`$ may be obtained from type $`III^{}`$ or type $`IV^{}`$ fibers by deformations of the discriminant curve; one simply moves some of the $`I_1`$ fibers in the decomposition of the $`III^{}`$ or $`IV^{}`$ to infinity, leaving only the $`I_1`$ fibers which make up the $`I_0^{}`$. For instance, the $`IV^{}`$ fiber decomposes into $`I_1`$ fibers of monodromy type AAAAABCC. By moving one $`I_1`$ of type A and another of type $`C`$ to infinity, one is left with $`I_1`$ fibers of monodromy type AAAABC, which is the decomposition of an $`I_0^{}`$ fiber. It follows from this that each of the $`I_0^{}`$ BPS multiplets listed in Table $`3`$ is also a BPS multiplet of both the Kodaira type $`III^{}`$ and type $`IV^{}`$ fibers. However, it turns out that $`I_0^{}`$ multiplets with different $`(p,q)`$ charges transforming in the same representation of $`SO(8)`$, transform as different representations of the exceptional groups of the Kodaira fibers $`III^{}`$ and $`IV^{}`$. For example , the $`\mathrm{𝟖}_𝐯`$ appearing for the $`I_0^{}`$ fiber with charges $`(p,q)=(1,0)`$ is embedded in a $`\mathrm{𝟓𝟔}`$ of $`E_7`$, or a $`\mathrm{𝟐𝟕}`$ of $`E_6`$, while the $`\mathrm{𝟖}_𝐯`$ with charges $`(p,q)=(1,2)`$ is embedded in a $`\mathrm{𝟐𝟕𝟔𝟔𝟒}`$ of $`E_7`$, or a $`\mathrm{𝟑𝟓𝟏}`$ of $`E_6`$. Consequently, although the type $`I_0^{}`$ BPS multiplets with $`(p,q)`$ charges listed in Table $`3`$ are also BPS multiplets of type $`III^{}`$ and type $`IV^{}`$ fibers, they are classified by an infinite number of different $`E_7`$ and $`E_6`$ representations. Although these can be computed on a case by case basis, a simple listing of such multiplets is impossible. In addition, there are BPS multiplets associated with both type $`III^{}`$ and type $`IV^{}`$ fibers that are unrelated to those of the $`I_0^{}`$. These multiplets arise from string junctions involving the type $`A`$ and type $`C`$ $`I_1`$ fibers not contained in the decomposition of $`I_0^{}`$. We will not discuss these states here, referring the interested reader to . Again, one expects BPS states with mutually non-local charges to become simultaneously massless as the five-brane approaches the singular fiber. The low energy theory on the five-brane wrapping the singular fiber flows to an exotic interacting fixed point theory. The fixed point theories with exceptional global symmetries were first discussed, in a different context, in .
## 5 Conclusion:
In this paper, we have presented detailed techniques for computing the discriminant curves of elliptically fibered Calabi–Yau threefolds. These were applied to a specific three–family, $`SU(5)`$ GUT model of particle physics within the context of Heterotic M–Theory. In this theory, and in general, the discriminant curves have an intricate structure, consisting of smooth sections, cusps and tangential and normal self–intersections. The type of degeneration of the elliptic fiber, classified by Kodaira, changes in the different regions of the discriminant curve. In this paper, we discussed how to find the Kodaira type of the fiber singularities and explicitly computed them for the discriminant curves associated with the $`SU(5)`$ GUT model. In Heterotic M–Theory, anomaly cancellation generically requires the existence of five–branes, located in the bulk space, wrapped on a holomorphic curve in the associated Calabi–Yau threefold. We showed that there is always a region of the moduli space of this holomorphic curve that corresponds a single five–brane wrapped on a pure fiber elliptic curve. Since this fiber degenerates as it approaches the discriminant curve, one expects light BPS states to appear in the worldvolume theory of this five–brane. For points on the smooth parts of the discriminant curve, we demonstrated in detail how the M–theory membranes associated with these states are, when projected into the base space, related to string junction lattices. Using string junction techniques, we computed the massless BPS hyper- and vector multiplets for all the Kodaira type degeneracies that occur in our specific $`SU(5)`$ GUT model. The Kodaira theory, as well as the computation of light states, is considerably more intricate at the cusp and self–intersection points of the discriminant curve. These topics will be discussed in a forthcoming publication.
It is important to note that the solutions of the BPS constraint on string junctions only gives a list of the possible BPS states in the spectrum. Determining their stability at different points in moduli space is a dynamical question, and the answer is not completely known for the exceptional Kodaira types. The answer is known for the non–exceptional Kodaira degeneracies, including the types $`I_1`$, $`I_2`$, $`IV`$ and $`I_0^{}`$ discussed in this paper . In the $`I_0^{}`$ case, the states listed in Table $`3`$ all appear in the stable spectrum at different points in moduli space. It follows that the multiplets of the exceptional groups $`E_7`$ and $`E_6`$ into which these states are embedded also exist in the spectrum. However, nothing is known about the stability of the multiplets which decouple upon deforming the exceptional fibers into $`I_0^{}`$.
Acknowledgements
We would like to thank P. Candelas, D. Luest, D. Morrison and B. Zwiebach for useful conversations and discussions. B. Ovrut is supported in part by a Senior Alexander von Humboldt Award, by the DOE under contract No. DE-ACO2-76-ER-03071 and by the University of Pennsylvania Research Foundation Grant. Z. Guralnick is supported in part by the DOE under contract No. DE-ACO2-76-ER-03071. A. Grassi is supported in part by an NSF grant DMS-9706707.
## 6 Appendix A– Review of String Junctions:
This Appendix contains a review of string junctions, and is a brief summary of work found in . There, string junctions were discussed in the context of type IIB string theory on $`^1`$. As described in this paper, membrane junctions in an elliptically fibered surface in M–theory become string junctions when projected into the base.
The equivalence classes of string junctions form a lattice with a quadratic form. To define the junction lattice, one initially splits all Kodaira fibers into a number of type $`I_1`$ fibers by a suitable deformation of the Weierstrass model (a “relevant deformation”). The string junctions lie in the base space. The $`I_1`$ loci are points in the base and the string junctions start from a fixed point $`P`$ away from these loci and may end only at the $`I_1`$ points. Each segment of the junction is an oriented string with charges $`(p,q)`$, which are conserved at branching points. The charges $`(p,q)`$ correspond to a one–cycle in the associated elliptic fiber $`F`$ over the point $`P`$. The segment of a string junction ending at an $`I_1`$ locus has charges which are proportional to the vanishing cycle $`(p_i,q_i)`$ for the $`i`$-th $`I_1`$ fiber. The vanishing cycle is the eigenvector of the monodromy matrix
$$_i=\left(\begin{array}{cc}1p_iq_i& p_i^2\\ q_i^2& 1+p_iq_i\end{array}\right)$$
(6.1)
In a convenient decomposition , any Kodaira fiber may be split into $`I_1`$ fibers on the real axis each with vanishing cycles $`(1,0)`$, $`(1,1)`$, or $`(1,1)`$. The associated monodromy matrices are, using (6.1), given by
$$A=\left(\begin{array}{cc}1& 1\\ 0& 1\end{array}\right),B=\left(\begin{array}{cc}2& 1\\ 1& 0\end{array}\right),C=\left(\begin{array}{cc}0& 1\\ 1& 2\end{array}\right)$$
(6.2)
respectively. The decompositions of the general Kodaira type fibers into $`I_1`$ fibers, from left to right on the real axis, are listed below.
| Kodaira type | decomposition |
| --- | --- |
| $`I_n`$ | $`A^n`$ |
| $`II`$ | $`AB`$ |
| $`III`$ | $`AAB`$ |
| $`IV`$ | $`AAAB`$ |
| $`I_n^{}`$ | $`A^{4+n}BC`$ |
| $`IV^{}`$ | $`AAAAABCC`$ |
| $`III^{}`$ | $`AAAAAABCC`$ |
| $`II^{}`$ | $`AAAAAAABCC`$ |
Table 4: The decomposition of Kodaira fibers into type $`I_1`$ fibers and the associated monodromy.
The basis cycles on the elliptic fiber are not globally defined, and there are branch cuts which may be taken to extend vertically downward from each $`I_1`$ loci. As a segment of the string junction crosses the branch cut of the $`i`$-th $`I_1`$ locus, the $`(p_i,q_i)`$ charges labeling the cycle over this segment are acted on by the monodromy $`_i`$. Such a segment can be pulled across the $`I_1`$ locus so that it no longer crosses the branch cut. However, because of a Hanany-Witten effect, an additional segment appears stretching between the original segment and the $`I_1`$ locus, as illustrated in Figure $`8`$. The charges of the new segment are determined by charge conservation to be $`(p^{},q^{})=_i(p,q)(p,q)`$, which is proportional to the vanishing cycle $`(p_i,q_i)`$. The two junctions related by pulling this segment across the $`I_1`$ locus are equivalent. Thus, the equivalence classes of junctions can be determined by considering only junctions which have no components below the real axis, as in Figure $`8(b)`$.<sup>1</sup><sup>1</sup>1Alternatively: One would like to write the vanishing cycle $`(p,q)`$ for any path $`\mathrm{\Gamma }`$, which starts from $`P`$ and ends at one $`I_1`$ point $`Q`$, as a suitable sum of our basic vanishing cycles. This sum should uniquely identify a class of string junctions in the base. In fact, $`\mathrm{\Gamma }`$ can be decomposed into a product of simple loops around each $`I_1`$ point and a simple path to $`Q`$ (that is a path which is homotopically trivial in the complement of the $`I_1`$ points in the base). The composition of the corresponding monodromy matrices applied to the cycle $`(p,q)`$ is a multiple of the vanishing cycle for the point $`Q`$. Now because $`_i(r,s)(r,s)`$ is always proportional to the vanishing cycle $`(p_i,q_i)`$, for all $`i`$ and $`(r,s)`$, we can show that the vanishing cycle for the path $`\mathrm{\Gamma }`$ has the expected form. The equivalence classes are then labeled by vectors $`\stackrel{}{J}=_in_i\stackrel{}{v}_i`$, where the integers $`n_i`$ indicate that the charges of the segment ending on the $`i`$-th $`I_1`$ locus are $`n_i(p_i,q_i)`$. The dimension of the string junction lattice is equal to the number of $`I_1`$ fibers. The reader is referred to for details on the construction of the quadratic intersection form on the junction lattice. Here, we simply state the results. The intersection form is given by
$$\stackrel{}{v}_i\stackrel{}{v}_i=1$$
(6.3)
$$\stackrel{}{v}_i\stackrel{}{v}_j=\frac{1}{2}(p_iq_jq_ip_j),ij$$
(6.4)
The basis in which an equivalence class $`\stackrel{}{J}`$ is labeled by the integers $`n_i`$ is not the most useful one. There is another basis in which the gauge and global symmetry charges appear explicitly. Viewed as quantum numbers of a BPS state, the electric and magnetic charges of a string junction are $`(p,q)=_in_i(p_i,q_i)`$. The remaining directions in the lattice are related to global symmetry charges. Using the intersection form on the lattice, one can show that the junction lattice contains the weight lattice of a (simply laced) Lie algebra . The simple roots of the Lie algebra correspond to a basis set of string junctions with vanishing charges $`(p,q)=(0,0)`$, and have self–intersection $`2`$. The intersection matrix of the simple root junctions $`\stackrel{}{\alpha }_i`$ is minus the Cartan matrix of the Lie algebra
$$\stackrel{}{\alpha }_i\stackrel{}{\alpha }_j=C_{ij}.$$
(6.5)
When $`I_1`$ fibers coalesce to form another Kodaira fiber, some root junctions have vanishing length. The Lie algebra associated with these simple roots generates a global symmetry of the theory which matches the A-D-E type of the Kodaira fiber. Fundamental weight junctions $`\stackrel{}{w}^i`$ with vanishing $`(p,q)`$ charges are defined by $`\stackrel{}{w}^i\stackrel{}{\alpha }_j=\delta _j^i`$, and are related to the simple root junctions by
$$\stackrel{}{w}^i=C^{1ij}\stackrel{}{\alpha }_j.$$
(6.6)
The fundamental weight junctions are “improper” in the sense that the $`n_i`$ associated with them are not integers. To get a complete basis, one defines another pair of string junctions $`\stackrel{}{w}_p`$ and $`\stackrel{}{w}_q`$. These are orthogonal to the fundamental weight junctions $`\stackrel{}{w}^i`$ and carry total charges $`(p,q)=(1,0)`$ and $`(0,1)`$ respectively. The junctions $`\stackrel{}{w}_p`$ and $`\stackrel{}{w}_q`$ are also improper. Any proper junction $`\stackrel{}{J}`$ can be written as
$$\stackrel{}{J}=p\stackrel{}{w}_p+q\stackrel{}{w}_q+a_i\stackrel{}{w}^i.$$
(6.7)
where the integers $`a_i`$ are the Dynkin labels. The weights in a BPS representation consist of junctions $`\stackrel{}{J}`$ related by the addition of simple roots, and satisfying the BPS condition
$$\stackrel{}{J}\stackrel{}{J}2+gcd(p,q),$$
(6.8)
where $`gcd`$ denotes the greatest common positive divisor of $`p`$ and $`q`$.
|
warning/0005/cond-mat0005120.html
|
ar5iv
|
text
|
# Superfluid Weight vs. Superconducting Temperature based on a U(1) Slave-Boson Approach to the t-J Hamiltonian
## Abstract
Based on an improved U(1) slave-boson approach to the t-J Hamiltonian, we investigate a relationship between the superfluid weight $`n_s/m^{}`$(the superconducting charge carrier density/the effective mass of the charge carrier) and the superconducting temperature $`T_c`$. From the present study we find a linear increase of $`n_s/m^{}`$ with $`T_c`$ with the doping concentration in the underdoped region, a saturation around the optimal doping and a decrease in both $`n_s/m^{}`$ and $`T_c`$ in the overdoped region. Such a trend of the ‘boomerang’ shaped locus in $`n_s/m^{}`$ vs. $`T_c`$ with increasing doping concentration from the underdoped to the heavily overdoped region is predicted to be in complete agreement with muon-spin-relaxation measurements. The boomerang behavior is found to occur in correlation with reduction in the spin singlet pairing(spinon pairing) order in the heavily overdoped region.
The transverse field muon-spin-relaxation($`\mu `$-SR) measurements of the magnetic penetration depth $`\lambda `$ in high $`T_c`$ copper oxide superconductors reveal an universal linear increase of the superfluid weight $`n_s/m^{}`$(the superconducting charge carrier density / the effective mass of the charge carrier) with the superconducting transition temperature $`T_c`$ in the underdoped region, a saturation near the optimal doping and a decrease in both $`n_s/m^{}`$ and $`T_c`$ as the hole doping concentration increases from the underdoped region to the heavily overdoped region. In the underdoped region the observed linear increase in both $`n_s/m^{}`$ and $`T_c`$ scales well with the doped hole concentration $`p`$ while $`n_s/m^{}`$ vs. $`T_c`$ does not scale with $`p`$ near and over the optimal doping. One can conjecture that the observed reduction of $`n_s/m^{}`$ with increasing hole carrier concentration in the overdoped region is attributed to either a decrease in the superconducting charge carrier density $`n_s`$ or to an increase in the effective mass of the charge carrier $`m^{}`$. In the present study we report a study of the observed boomerang behavior in $`n_s/m^{}`$ vs. $`T_c`$ as a function of doped hole concentration by using an improved U(1) slave-boson theory which correctly allows coupling between the spin and charge degrees of freedom. Contrary to our earlier U(1) slave-boson approach, such introduction of both the spin and charge degrees of freedom was found to predict the arch shaped bose condensation temperature as a function of doped hole concentration $`p`$, in agreement with the experimentally observed phase diagram in the plane of $`T_c`$ and $`p`$. For this study we first compute the doping and temperature dependences of the superfluid weight and discuss the cause of reduction in both the superfluid weight and the superconducting temperature in the overdoped region. In addition, a numerical comparison with the SU(2) slave-boson theory is briefly made to discuss a difference.
We write the t-J Hamiltonian,
$`H`$ $`=`$ $`t{\displaystyle \underset{<i,j>}{}}(c_{i\sigma }^{}c_{j\sigma }+c.c.)+J{\displaystyle \underset{<i,j>}{}}(𝐒_i𝐒_j{\displaystyle \frac{1}{4}}n_in_j)\mu _0{\displaystyle \underset{i,\sigma }{}}c_{i\sigma }^{}c_{i\sigma },`$ (1)
where $`𝐒_i𝐒_j\frac{1}{4}n_in_j=\frac{1}{2}(c_i^{}c_j^{}c_i^{}c_j^{})(c_jc_ic_jc_i)`$. Here $`𝐒_i`$ is the electron spin operator at site $`i`$, $`𝐒_i=\frac{1}{2}c_{i\alpha }^{}𝝈_{\alpha \beta }c_{i\beta }`$ with $`𝝈_{\alpha \beta }`$, the Pauli spin matrix element and $`n_i`$, the electron number operator at site $`i`$, $`n_i=c_{i\sigma }^{}c_{i\sigma }`$. Using the single occupancy constraint and thus $`c_{i\sigma }=b_i^{}f_{i\sigma }`$(with $`f_{i\sigma }`$, spinon annihilation operator of electron spin $`\sigma `$ and $`b_i^{}`$, holon creation operator at site $`i`$), the U(1) slave-boson representation of the above t-J Hamiltonian leads to Eq.(4),
$`H`$ $`=`$ $`t{\displaystyle \underset{<i,j>}{}}(f_{i\sigma }^{}f_{j\sigma }b_j^{}b_i+c.c.)`$ (4)
$`{\displaystyle \frac{J}{2}}{\displaystyle \underset{<i,j>}{}}b_ib_jb_j^{}b_i^{}(f_i^{}f_j^{}f_i^{}f_j^{})(f_jf_if_jf_i)`$
$`\mu _0{\displaystyle \underset{i,\sigma }{}}f_{i\sigma }^{}f_{i\sigma }+i{\displaystyle \underset{i}{}}\lambda _i(f_{i\sigma }^{}f_{i\sigma }+b_i^{}b_i1),`$
where $`\lambda _i`$ is the Lagrange multiplier field to enforce the single occupancy constraint at each site. In the SU(2) slave-boson theory, the electron operator is given by $`c_\alpha =\frac{1}{\sqrt{2}}h^{}\psi _\alpha `$ with $`\alpha =1,2`$, where $`\psi _1=\left(\begin{array}{c}f_1\\ f_2^{}\end{array}\right)`$ and $`\psi _2=\left(\begin{array}{c}f_2\\ f_1^{}\end{array}\right)`$ and $`h=\left(\begin{array}{c}b_1\\ b_2\end{array}\right)`$ are respectively the doublets of spinon and holon annihilation operators in the SU(2) theory. The SU(2) slave-boson representation of the above t-J Hamiltonian shows
$`H={\displaystyle \frac{t}{2}}{\displaystyle \underset{<i,j>\sigma }{}}[(f_{\sigma i}^{}f_{\sigma j})(b_{1j}^{}b_{1i}b_{2i}^{}b_{2j})`$ (5)
$`+(f_{\sigma j}^{}f_{\sigma i})(b_{1i}^{}b_{1j}b_{2j}^{}b_{2i})`$ (6)
$`+(f_{2i}f_{1j}f_{1i}f_{2j})(b_{1j}^{}b_{2i}+b_{1i}^{}b_{2j})`$ (7)
$`+(f_{1j}^{}f_{2i}^{}f_{2j}^{}f_{1i}^{})(b_{2i}^{}b_{1j}+b_{2j}^{}b_{1i})]`$ (8)
$`{\displaystyle \frac{J}{2}}{\displaystyle \underset{<i,j>}{}}(1h_i^{}h_i)(1h_j^{}h_j)\times `$ (9)
$`(f_{2i}^{}f_{1j}^{}f_{1i}^{}f_{2j}^{})(f_{1j}f_{2i}f_{2j}f_{1i})\mu _0{\displaystyle \underset{i}{}}h_i^{}h_i`$ (10)
$`{\displaystyle \underset{i}{}}[i\lambda _i^{(1)}(f_{1i}^{}f_{2i}^{}+b_{1i}^{}b_{2i})+i\lambda _i^{(2)}(f_{2i}f_{1i}+b_{2i}^{}b_{1i})`$ (11)
$`+i\lambda _i^{(3)}(f_{1i}^{}f_{1i}f_{2i}f_{2i}^{}+b_{1i}^{}b_{1i}b_{2i}^{}b_{2i})],`$ (12)
where $`\lambda _i^{(1),(2),(3)}`$ are the real Lagrangian multipliers to enforce the local single occupancy constraint in the SU(2) slave-boson representation. It is noted that in both the U(1) and SU(2) approaches above, coupling between spin and charge degrees of freedom is correctly introduced in the Heisenberg term.
After Hubbard Stratonovich transformations for the direct, exchange and pairing channels, we find the total free energy in the functional integral representation, for the U(1) slave-boson theory,
$`F(𝐀)={\displaystyle \frac{1}{\beta }}\mathrm{ln}{\displaystyle D\chi D\mathrm{\Delta }^fD\mathrm{\Delta }^bDbDf}`$ (13)
$`e^{(S^b(𝐀,\chi ,\mathrm{\Delta }^f,\mathrm{\Delta }^b,b)+S^f(\chi ,\mathrm{\Delta }^f,f))},`$ (14)
where $`S^b`$ represents the holon action for the charge degree of freedom,
$`S^b(𝐀,\chi ,\mathrm{\Delta }^f,\mathrm{\Delta }^b,b)={\displaystyle _0^\beta }d\tau [{\displaystyle \underset{i}{}}b^{}(𝐫_i,\tau )(_\tau \mu ^b)b(𝐫_i,\tau )`$ (15)
$`+{\displaystyle \frac{J}{2}}{\displaystyle \underset{<i,j>}{}}|\mathrm{\Delta }_{ij}^f|^2\left[|\mathrm{\Delta }_{ij}^b|^2+p^2\right]`$ (16)
$`t{\displaystyle \underset{<i,j>}{}}e^{iA_{ij}}\chi _{ij}b^{}(𝐫_i,\tau )b(𝐫_j,\tau )+c.c.`$ (17)
$`{\displaystyle \frac{J}{2}}{\displaystyle \underset{<i,j>}{}}|\mathrm{\Delta }_{ij}^f|^2[\mathrm{\Delta }_{ij}^bb^{}(𝐫_i,\tau )b^{}(𝐫_j,\tau )+c.c.]`$ (18)
and $`S^f`$ is the spinon action for the spin degrees of freedom,
$`S^f(\chi ,\mathrm{\Delta }^f,f)={\displaystyle _0^\beta }d\tau [{\displaystyle \underset{i}{}}f_\sigma ^{}(𝐫_i,\tau )(_\tau \mu ^f)f_\sigma (𝐫_i,\tau )`$ (19)
$`+{\displaystyle \frac{J(1p^2)}{2}}{\displaystyle \underset{<i,j>}{}}\left[|\mathrm{\Delta }_{ij}^f|^2+{\displaystyle \frac{1}{2}}|\chi _{ij}|^2\right]`$ (20)
$`{\displaystyle \frac{J}{4}}(1p)^2{\displaystyle \underset{<i,j>}{}}\chi _{ij}f_\sigma ^{}(𝐫_i,\tau )f_\sigma (𝐫_j,\tau )+c.c.`$ (21)
$`{\displaystyle \frac{J}{2}}(1p)^2{\displaystyle \underset{<i,j>}{}}\mathrm{\Delta }_{ij}^f(f_{}(𝐫_i,\tau )f_{}(𝐫_j,\tau )f_{}(𝐫_i,\tau )f_{}(𝐫_j,\tau ))+c.c.].`$ (22)
Here $`b`$ is the holon field and $`f`$, the spinon field. $`\chi _{ij}`$, $`\mathrm{\Delta }_{ij}^f`$ and $`\mathrm{\Delta }_{ij}^b`$ are the Hubbard Stratonovich fields corresponding to the exchange, the spinon pairing and the holon pairing channels respectively. After integration over the holon and spinon fields, we obtain the total free energy,
$`F(𝐀)={\displaystyle \frac{1}{\beta }}\mathrm{ln}{\displaystyle D\chi D\mathrm{\Delta }^fD\mathrm{\Delta }^be^{(F^b(𝐀,\chi ,\mathrm{\Delta }^f,\mathrm{\Delta }^b)+F^f(\chi ,\mathrm{\Delta }^f))}},`$ (23)
where $`F^b(𝐀,\chi ,\mathrm{\Delta }^f,\mathrm{\Delta }^b)=\frac{1}{\beta }\mathrm{ln}Dbe^{S^b(𝐀,\chi ,\mathrm{\Delta }^f,\mathrm{\Delta }^b,b)}`$ is the holon free energy and $`F^f(\chi ,\mathrm{\Delta }^f)=\frac{1}{\beta }\mathrm{ln}Dfe^{S^f(\chi ,\mathrm{\Delta }^f,f)}`$, the spinon free energy.
The linear response of current to weak applied electromagnetic(EM) field is, in the energy-momentum space,
$`j_l(\omega ,𝐪)=\mathrm{\Pi }_{lm}(\omega ,𝐪)A_m(\omega ,𝐪),`$ (24)
with $`𝐣`$, the current and $`𝐀`$, the EM vector potential. Here the current response function $`\mathrm{\Pi }_{lm}(\omega ,𝐪)`$ is obtained from
$`\mathrm{\Pi }_{lm}(\omega ,𝐪)=\beta {\displaystyle \frac{^2}{A_l(\omega ,𝐪)A_m(\omega ,𝐪)}}F(𝐀)|_{𝐀=0},`$ (25)
with $`\beta =\frac{1}{k_BT}`$, where $`F(𝐀)`$ is the free energy of Eq.(23) or Eq.(LABEL:eq:free\_energy\_su2\_2). The superfluid weight $`\frac{n_s}{m^{}}`$ is defined as the transverse EM current response function in the static, long wave length limit,
$`{\displaystyle \frac{n_s}{m^{}}}={\displaystyle \frac{1}{e^2}}\underset{q0}{lim}\mathrm{\Pi }_{xx}(\omega =0,𝐪),`$ (26)
for the isotropic system.
As is shown in the second term of Eq.(18), the EM vector potential field $`A`$ is seen to modulate(is coupled with) the hopping order parameter field $`\chi _{ij}`$ in the U(1) slave-boson representation. The phase fluctuations of the hopping order parameter, i.e. $`\chi _{ij}=\chi e^{ia_{ij}}`$ are introduced. For the hole pairing of d-wave symmetry we allow the d-wave spinon pairing order parameter, $`\mathrm{\Delta }_{ji}^f=\pm \mathrm{\Delta }_f`$(the sign $`+()`$ is for the $`\mathrm{𝐢𝐣}`$ link parallel to $`\widehat{x}`$ ($`\widehat{y}`$)) and the s-wave holon pairing order parameter, $`\mathrm{\Delta }_{ji}^b=\mathrm{\Delta }_b`$. The EM current response function of the total system is obtained from the Ioffe-Larkin composition rule,
$`\mathrm{\Pi }(\omega ,𝐪)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Pi }^b(\omega ,𝐪)\mathrm{\Pi }^f(\omega ,𝐪)}{\mathrm{\Pi }^b(\omega ,𝐪)+\mathrm{\Pi }^f(\omega ,𝐪)}},`$ (27)
where $`\mathrm{\Pi }^b(\omega ,𝐪)`$ is the holon current response function to the the EM field $`A`$ and $`\mathrm{\Pi }^f(\omega ,𝐪)`$, the spinon current response function to the gauge field $`a`$ corresponding to the phase fluctuations of hopping order parameter. The holon and spinon response functions are calculated from the linear response theory(see appendices 1,2). The holon response function is evaluated to be
$`\mathrm{\Pi }_{lm}^b(\omega ,𝐪)={\displaystyle \frac{1}{N}}{\displaystyle \underset{𝐪_1}{}}[2t\chi \mathrm{cos}q_{1l}(u^b(𝐪_\mathrm{𝟏})^2n^b(E^b(𝐪_\mathrm{𝟏}))v^b(𝐪_\mathrm{𝟏})^2n^b(E^b(𝐪_\mathrm{𝟏})))\delta _{lm}`$ (33)
$`+4t^2\chi ^2e^{i\frac{q_lq_m}{2}}\mathrm{sin}(q_{1m}+{\displaystyle \frac{q_m}{2}})\mathrm{sin}(q_{1l}+{\displaystyle \frac{q_l}{2}})\times `$
$`\{(u^b(𝐪+𝐪_1)^2u^b(𝐪_1)^2u^b(𝐪+𝐪_1)u^b(𝐪_1)v^b(𝐪+𝐪_1)v^b(𝐪_1)){\displaystyle \frac{n^b(E^b(𝐪+𝐪_1))n^b(E^b(𝐪_1))}{i\omega (E^b(𝐪+𝐪_1)E^b(𝐪_1))}}`$
$`+(u^b(𝐪+𝐪_1)^2v^b(𝐪_1)^2+u^b(𝐪+𝐪_1)u^b(𝐪_1)v^b(𝐪+𝐪_1)v^b(𝐪_1)){\displaystyle \frac{n^b(E^b(𝐪+𝐪_1))n^b(E^b(𝐪_1))}{i\omega (E^b(𝐪+𝐪_1)+E^b(𝐪_1))}}`$
$`+(v^b(𝐪+𝐪_1)^2u^b(𝐪_1)^2+u^b(𝐪+𝐪_1)u^b(𝐪_1)v^b(𝐪+𝐪_1)v^b(𝐪_1)){\displaystyle \frac{n^b(E^b(𝐪+𝐪_1))n^b(E^b(𝐪_1))}{i\omega +(E^b(𝐪+𝐪_1)+E^b(𝐪_1))}}`$
$`+(v^b(𝐪+𝐪_1)^2v^b(𝐪_1)^2u^b(𝐪+𝐪_1)u^b(𝐪_1)v^b(𝐪+𝐪_1)v^b(𝐪_1)){\displaystyle \frac{n^b(E^b(𝐪+𝐪_1))n^b(E^b(𝐪_1))}{i\omega +(E^b(𝐪+𝐪_1)E^b(𝐪_1))}}\}],`$
where $`E^b(𝐪)=\sqrt{(ϵ^b(𝐪)\mu ^b)^2(J\mathrm{\Delta }_f^2\mathrm{\Delta }_b\gamma _𝐪)^2}`$ is the holon quasiparticle energy and $`ϵ^b(𝐪)=2t\chi \gamma _𝐪`$, the single holon quasiparticle energy with $`\gamma _𝐪=(\mathrm{cos}q_x+\mathrm{cos}q_y)`$. $`n^b(E)=\frac{1}{e^{\beta E}1}`$ is the boson distribution function, $`u^b(𝐪)=\frac{1}{\sqrt{2}}\sqrt{\frac{ϵ^b(𝐪)\mu ^b}{E^b(𝐪)}+1}`$ and $`v^b(𝐪)=\frac{1}{\sqrt{2}}\sqrt{\frac{ϵ^b(𝐪)\mu ^b}{E^b(𝐪)}1}`$. The spinon response function is
$`\mathrm{\Pi }_{lm}^f(\omega ,𝐪)={\displaystyle \frac{1}{N}}{\displaystyle \underset{𝐪_1}{}}[{\displaystyle \frac{J}{2}}(1p)^2\chi \mathrm{cos}q_{1l}{\displaystyle \frac{ϵ^f(𝐪_1)\mu ^f}{E^f(𝐪_1)}}(n^f(E^f(𝐪_1))n^f(E^f(𝐪_1)))\delta _{lm}`$ (39)
$`+({\displaystyle \frac{J}{2}}(1p)^2\chi )^2e^{i\frac{q_lq_m}{2}}\mathrm{sin}(q_{1m}+{\displaystyle \frac{q_m}{2}})\mathrm{sin}(q_{1l}+{\displaystyle \frac{q_l}{2}})\times `$
$`\{(n^f(E^f(𝐪+𝐪_1)){\displaystyle \frac{E^f(𝐪+𝐪_1)(i\omega +E^f(𝐪+𝐪_1))+(ϵ^f(𝐪+𝐪_1)\mu ^f)(ϵ^f(𝐪_1)\mu ^f)+\mathrm{\Delta }_f^{^{}}(𝐪+𝐪_1)\mathrm{\Delta }_f^{^{}}(𝐪_1)}{E^f(𝐪+𝐪_1)[(i\omega +E^f(𝐪+𝐪_1))^2E^f(𝐪_1)^2]}}`$
$`+(n^f(E^f(𝐪_1)){\displaystyle \frac{E^f(𝐪_1)(i\omega +E^f(𝐪_1))+(ϵ^f(𝐪+𝐪_1)\mu ^f)(ϵ^f(𝐪_1)\mu ^f)+\mathrm{\Delta }_f^{^{}}(𝐪+𝐪_1)\mathrm{\Delta }_f^{^{}}(𝐪_1)}{E^f(𝐪_1)[(i\omega +E^f(𝐪_1))^2E^f(𝐪+𝐪_1)^2]}}`$
$`(n^f(E^f(𝐪+𝐪_1)){\displaystyle \frac{E^f(𝐪+𝐪_1)(i\omega +E^f(𝐪+𝐪_1))+(ϵ^f(𝐪+𝐪_1)\mu ^f)(ϵ^f(𝐪_1)\mu ^f)+\mathrm{\Delta }_f^{^{}}(𝐪+𝐪_1)\mathrm{\Delta }_f^{^{}}(𝐪_1)}{E^f(𝐪+𝐪_1)[(i\omega +E^f(𝐪+𝐪_1))^2E^f(𝐪_1)^2]}}`$
$`(n^f(E^f(𝐪_1)){\displaystyle \frac{E^f(𝐪_1)(i\omega +E^f(𝐪_1))+(ϵ^f(𝐪+𝐪_1)\mu ^f)(ϵ^f(𝐪_1)\mu ^f)+\mathrm{\Delta }_f^{^{}}(𝐪+𝐪_1)\mathrm{\Delta }_f^{^{}}(𝐪_1)}{E^f(𝐪_1)[(i\omega +E^f(𝐪_1))^2E^f(𝐪+𝐪_1)^2]}}\}],`$
where $`E^f(𝐪)=\sqrt{(ϵ^f(𝐪)\mu ^f)^2+(\mathrm{\Delta }_f^{^{}}(𝐪))^2}`$ is the spinon quasiparticle energy, $`ϵ^f(𝐪)=\frac{J}{2}(1p)^2\chi \gamma _𝐪`$, the single spinon quasiparticle energy, and $`\mathrm{\Delta }_f^{^{}}(𝐪)=J(1p)^2\mathrm{\Delta }_f\phi _𝐪`$, the spinon pairing gap with $`\phi _𝐪=(\mathrm{cos}q_x\mathrm{cos}q_y)`$. $`n^f(E)=\frac{1}{e^{\beta E}+1}`$ is the fermion distribution function. For the U(1) theory the superfluid weight is obtained from Eq.(26) with the use of (27). In the present calculations of the superfluid weight $`J/t=0.2`$ is chosen. For other choice of $`J/t`$ we find that there exist no qualitative differences in the behavior of $`n_s/m^{}`$ vs. $`T_c`$.
In Fig.1 the computed superfluid weight displays a linear increase in doped hole concentration(rate) $`p`$ in the underdoped region(the predicted optimal doping rate is $`p_o=0.07`$). As hole doping concentration further increases, the superfluid weight first saturates and rapidly drops to $`0`$ in the heavily overdoped region. This behavior in the overdoped region is attributed to a decrease in holon pairing interaction caused by the predicted diminishing trend of the spinon pairing probability amplitude(order parameter) $`\mathrm{\Delta }_{ij}^f`$, as can be readily understood from the last term of Eq.(18). The predicted decrease of pseudogap(spin gap) temperature with hole doping rate is in excellent agreement with observation. In Fig.2, a “boomerang” shaped locus in $`\frac{n_s(T0)}{m^{}}`$ and $`T_c`$ is displayed by showing a linear relationship between the two in the underdoped region and the ‘reflex’ behavior in the overdoped region. Although not numerically agreeable, this trend is completely consistent with muon-spin-relaxation measurements. As mentioned above, the spinon pairing amplitude $`\mathrm{\Delta }_{ij}^f`$ is further reduced in the overdoped region and this, in turn, causes a decrease in holon pairing interaction and thus the superconducting charge carrier density $`n_s`$. This will allow the reflex behavior, that is, a decrease in both the superfluid weight and superconducting temperature.
For the sake of introducing the low energy phase fluctuations of the order parameters $`\chi _{ij}`$ and $`\mathrm{\Delta }_{ij}^f`$, which were not considered in the above U(1) slave-boson study, we now introduce the SU(2) slave-boson theory. In the SU(2) theory, the EM field $`A`$ is found to modulate both the hopping order parameter $`\chi _{ij}`$ and the spinon pairing order parameter $`\mathrm{\Delta }_{ij}^f`$. The phase fluctuations of both $`\chi _{ij}`$ and $`\mathrm{\Delta }_{ij}^f`$ are thus taken into account to assess the EM current response of the system. We introduce the spinon pairing order parameter of d-wave symmetry, $`\mathrm{\Delta }_{ji}^f=\pm \mathrm{\Delta }_f`$(the sign $`+()`$ is for the $`\mathrm{𝐢𝐣}`$ link parallel to $`\widehat{x}`$ ($`\widehat{y}`$)) and the holon pairing order parameter of s-wave symmetry, $`\mathrm{\Delta }_{ij;\alpha \beta }^B=\mathrm{\Delta }_b(\delta _{\alpha ,1}\delta _{\beta ,1}\delta _{\alpha ,2}\delta _{\beta ,2})`$.
In the SU(2) slave-boson theory, the response function $`\mathrm{\Pi }_{ll^{^{}}}^{b(A,a^1)}`$ of holon isospin current to the EM field $`A`$ and the gauge field $`a`$ vanishes owing to the contribution of the $`b_2`$-boson in the static and long-wavelength limit(see the appendices 3,4). Therefore the superfluid weight of the total system is given by the holon current response function only,
$`{\displaystyle \frac{n_s}{m^{}}}={\displaystyle \frac{1}{e^2}}\mathrm{\Pi }_{lm}^{b(A,A)}(\omega =0,𝐪=0).`$ (40)
Here $`\mathrm{\Pi }_{lm}^{b(A,A)}(\omega ,𝐪)`$ is computed from the use of the usual linear response theory for the holon action Eq.(LABEL:eq:su2\_holon\_action)(see Appendix 6). In Fig.3, we show the doping dependence of the superfluid weight with the choice of $`J/t=0.2`$. The superfluid weight predicted from the SU(2) theory increases faster than the U(1) case, as doped hole concentration. This is because the gauge fields(or the phase fluctuations of the order parameters) do not screen the EM field in the SU(2) theory. In Fig.4, the relation between $`T_c`$ and $`\frac{n_s(T0)}{m^{}}`$ is displayed; at low doping, $`T_c`$ increases linearly with $`\frac{n_s(T0)}{m^{}}`$ as in the U(1) theory, a plateau near the optimal doping $`p=0.13`$) and a reflex in the overdoped region is also observed. There exists only a quantitative (but not qualitative) difference between the two theories. This is because the spinon pairing amplitude $`\mathrm{\Delta }_{ij}^f`$ predicted by the SU(2) theory slowly diminishes in the overdoped region compared to the U(1) case.
In summary, we investigated a relationship between the superfluid weight and the superconducting temperature, that is, $`\frac{n_s}{m^{}}`$ vs. $`T_c`$ based on both the U(1) and SU(2) slave-boson approaches to the t-J Hamiltonian. From both theories we find a qualitatively similar boomerang shape in the path of $`\frac{n_s(T0)}{m^{}}`$ vs. $`T_c`$ by showing a linear increase in the underdoped region and a reflex behavior in the overdoped region. Such trend of boomerang behavior in high $`T_c`$ superconductors is completely consistent with $`\mu `$-SR experiments. The reflex behavior predicted by both theories is attributed to the diminishing attractive hole pairing interaction caused by the markedly reduced spin pairing order $`\mathrm{\Delta }^f`$, particularly in the heavily overdoped region, and thus to the reduction of both the superconducting charge carrier density $`n_s`$ and the superconducting temperature $`T_c`$ in the overdoped region.
|
warning/0005/quant-ph0005102.html
|
ar5iv
|
text
|
# The non-local content of quantum operations
## I Introduction
In the past, most of the research on quantum non-locality has been devoted to the issue of non-locality of quantum states. However we feel that an equally important issue is that of non-locality of quantum evolutions. That is, in parallel with the understanding of non-locality of quantum kinematics one should also develop an understanding of the non-locality of quantum dynamics.
Let us start with a simple example. Consider two qubits situated far from each other, one held by Alice and the other one by Bob. Suppose they would like to implement a two qubit quantum evolution described by the unitary operator $`U`$. (We wish to be able to apply $`U`$ on any initial state of the two qubits). With exception of the case when $`U`$ is a product of two local unitary operators, $`U=U_AU_B`$, no other quantum evolution can be accomplished by local means only. Thus almost all quantum evolutions are non-local. The main question we address in this paper is how to describe, qualitatively and quantitatively, the non-locality of quantum evolutions.
In order to be able to describe the amount of non-locality contained by the unitary operator $`U`$ we suggest the following approach. We consider that Alice and Bob, in addition of being able to perform any local operations, they also have additional resources, namely they share entangled states, and they are able to communicate classically. The question then reduces to finding out how much of these resources is needed to implement $`U`$.
The above framework has also been put-forward by Chefles, Gilson and Barnett .
We emphasise that although we have largely discussed the role of quantum entanglement above, the role of the classical communication is equally important. Understanding the character of a quantum evolution requires knowing both the amount of entanglement and the amount of classical communication needed.
## II General sufficiency conditions
First of all, it is important to note that any unitary evolution can be implemented given enough shared entanglement and classical communication. Indeed, consider the case of two qubits, one held by Alice and one by Bob. Any unitary transformation $`U`$ on these two qubits can be accomplished by having Alice teleport her qubit to Bob, Bob performing $`U`$ locally and finally Bob teleporting Alice’s qubit back to Alice. The resources needed for the two teleportation actions are: (1 e-bit plus two classical bits transmitted from Alice to Bob for the Alice to Bob teleportation) plus (1 e-bit plus two classical bits transmitted from Bob to Alice for the Bob to Alice teleportation). It is obvious now that any unitary operation involving any number of parties and any number of qubits can be accomplished by a similar procedure (teleporting all states to a single location, performing $`U`$ locally and teleporting back the qubits to their original locations).
The “double teleportation” procedure shown above is sufficient to implement any quantum evolution. The question is however whether so much resources are actually needed. We will discuss a couple of specific example below.
## III The SWAP operation on two qubits
The SWAP operation defined by:
$`U_{\mathrm{SWAP}}|\psi |\varphi =|\varphi |\psi `$ (1)
is a particularly intriguing case, since although it takes product states to product states, it is, as we now show, the most non-local operation possible in the sense described above. That is, we will prove that in order to implement a SWAP on two qubits it is not only sufficient but also necessary to use 2 e-bits plus 2 bits of classical communication from Alice to Bob plus 2 bits of classical communication from Bob to Alice.
Proof: To prove that the SWAP operation needs as non-local resources 2 e-bits, we will show that if we have an apparatus able to implement the SWAP operation we can use it in order to create 2 e-bits. Thus, since entanglement cannot be created ex nihilo, the apparatus which implements the SWAP must use 2 e-bits as an internal non-local resource.
Let us show how to generate two singlets using the SWAP operation. Firstly Alice and Bob prepare singlets locally
$`_A_a+_A_a\text{and}_B_b+_B_b,`$ (2)
Alice’s spins are labelled $`A`$ and $`a`$ and Bob’s $`B`$ and $`b`$ (here and in what follows we will leave out normalisation factors for states). Now perform the SWAP operation on spins $`A`$ and $`B`$:
$$(_A_a+_A_a)(_B_b+_B_b)$$
$`(_B_a+_B_a)(_A_b+_A_b).`$ (3)
This state contains two singlets held between Alice and Bob.
To find the classical communication resources needed to implement the SWAP operation we will adapt an argument first given in . We show that if we have an apparatus able to implement the SWAP operation we can use it in order to communicate 2 bits from Alice to Bob plus 2 bits from Bob to Alice. From this follows that it must be the case that the SWAP apparatus uses 2 bits of classical communication from Alice to Bob plus 2 bits of classical communication from Bob to Alice as an internal resource, otherwise Alice could receive information from Bob transmitted faster than light.
For suppose that the SWAP operation requires less than four bits of classical communication (two bits each way). Alice and Bob can produce an instantaneous SWAP operation which works correctly with probability greater than one sixteenth in the following way. Alice and Bob run the usual SWAP protocol, but instead of waiting for classical communication from each other, they simply guess the bits that they would have received. Since we have assumed that the SWAP operation requires less than 4 bits, the probability that Alice and Bob guess correctly is greater than one sixteenth and hence the SWAP operation also succeeds with probability greater than one sixteenth.
Thus using the protocol described previously can now use this imperfect, but instantaneous SWAP to communicate 4 bits instantaneously. The bits arrive correctly when the SWAP is implemented correctly. Hence the probability that 4 bits arrive correctly is larger than one sixteenth; 4 bits communicated correctly with probability greater than one sixteenth represents a non-zero amount of information. Thus Alice and Bob have managed to convey some information to each other instantaneously. We conclude therefore that the SWAP operation cannot be done with less that 4 bits of classical communication; otherwise it allows communication faster than the speed of light.
Earlier in this section we showed that the SWAP operation can be used to generate two singlets. We now show that the SWAP operation can be also be used to perform four bits of classical communication (two bits each way): the main idea is that of “super-dense coding” . Suppose that initially Alice and Bob share two singlets:
$`_A_B+_A_B\text{and}_a_b+_a_b.`$ (4)
Now Alice chooses one of four local unitary operations 1 (identity), $`\sigma _x`$, $`\sigma _y`$, $`\sigma _z`$ and performs it on her spin $`A`$. This causes the first singlet to be in one of the four Bell states. Bob also, independently chooses one of these four locally unitaries and performs it on his spin $`b`$, putting the second singlet into one of the Bell states. Then the SWAP operation is performed on spins $`A`$ and $`b`$. Now both Bob and Alice have one of the Bell states locally; which one they have depends on which operation the other performed. By measurement, they can work out which of the four unitaries the other performed. Thus the SWAP operation has enabled two bits of classical communication to be performed each way.
## IV The CNOT operation on two qubits
Another important quantum operation is CNOT, defined as
$$$$
(5)
$$$$
(6)
$$$$
(7)
$$.$$
(8)
As we prove below, the necessary and sufficient resources for CNOT are 1 e-bit plus 1 bit of classical communication from Alice to Bob plus 1 bit of classical communication from Bob to Alice.
Proof: Constructing a CNOT We now show how to construct the CNOT operation using one singlet and two bits of classical communication. We then show how to generate one singlet or perform two bits of classical communication using the CNOT.
Firstly we will show how, using one singlet and one bit of classical communication each way, we can perform a CNOT on the state
$`(\alpha _A+\beta _A)(\gamma _B+\delta _B)`$ (9)
i.e. transform it to
$`\alpha _A(\gamma _B+\delta _B)+\beta _A(\gamma _B+\delta _B).`$ (10)
Since the operation behaves linearly, the protocol performs the CNOT on any input state (i.e. even if the qubits are entangled with each other or with other systems).
Step 1 The first step is to append a singlet held between Alice and Bob to the state (9):
$`(\alpha _A+\beta _A)(_a_b+_a_b)(\gamma _B+\delta _B),`$ (11)
then for Alice to measure the total spin of her spins $`A`$ and $`a`$.
If the total spin is one, then the state becomes
$`(\alpha _A_a_b+\beta _A_a_b)(\gamma _B+\delta _B).`$ (12)
Now Alice disentangles the singlet spin by performing the following (local) operation:
$`_A_a_A_a;_A_a_A_a,`$ (13)
and the state becomes
$`(\alpha _A_b+\beta _A_b)(\gamma _B+\delta _B)_a.`$ (14)
If the total spin had been zero, then rather than (12) the state becomes
$`(\alpha _A_a_b+\beta _A_a_b)(\gamma _B+\delta _B).`$ (15)
In this case Alice can disentangle the $`a`$ spin by
$`_A_a_A_a;_A_a_A_a,`$ (16)
leading to
$`(\alpha _A_b+\beta _A_b)(\gamma _B+\delta _B)_a.`$ (17)
In order to get this state in the correct form, Bob needs to invert his $`b`$ spin. Thus Alice must communicate one bit to Bob to tell him whether she found total spin one or zero, and thus whether he needs to invert his spin or not.
After these operations, the state is
$`(\alpha _A_b+\beta _A_b)(\gamma _B+\delta _B)_a.`$ (18)
Step 2 Now Bob performs a CNOT on the $`b`$ and $`B`$ spins, thus the total state is
$`[\alpha _A_b(\gamma _B+\delta _B)+\beta _A_b(\gamma _B+\delta _B)]_a.`$ (19)
Step 3 Bob now measures $`\sigma _x`$ on his part of the singlet $`b`$. Either the state becomes
$`[\alpha _A(\gamma _B+\delta _B)+\beta _A(\gamma _B+\delta _B)]`$ (20)
$`_a(_b+_b),`$ (21)
or
$`[\alpha _A(\gamma _B+\delta _B)\beta _A(\gamma _B+\delta _B)]`$ (22)
$`_a(_b_b),`$ (23)
In the former case (i.e. the $`x`$ component of spin was $`+`$) we have performed the protocol as desired. In the latter, Alice needs to perform a $`\sigma _z`$ rotation by $`\pi `$. Thus Bob needs to communicate one bit to Alice to tell her whether or not to perform the rotation.
We have thus shown how to perform a CNOT using one singlet and one bit of classical communication each way.
Creating entanglement by CNOT We show now that a CNOT apparatus can be used to create 1 e-bit between Alice and Bob; thus (since entanglement cannot be increased by local operations) 1 e-bit is a necessary resource for constructing a CNOT.
Creating 1 e-bit by a CNOT is straightforward:
$$(_A+_A)_B_A_B+_A_B.$$
(24)
Classical communication by CNOT
Suppose that Alice and Bob have an apparatus which implements a CNOT and also share 1 e-bit. They can use these resources to communicate at the same time 1 classical bit from Alice to Bob and 1 classical bit from Bob to Alice. This proves (see preceding section) that communicating 1 classical bit each way is a necessary resource for constructing a CNOT.
Suppose the initial state is
$$_a_b+_a_b.$$
(25)
Alice can encode a “0” by not doing anything to the state and a “1” by flipping her qubit. Bob can encode a “0” by not doing anything to the state and a “1” by changing the phase as follows: $``$ and $``$.
The four states corresponding to the different bit combinations are thus
$$_a_b+_a_bcorrespondsto0_A0_B.$$
(26)
$$_a_b+_a_bcorrespondsto1_A0_B.$$
(27)
$$_a_b_a_bcorrespondsto0_A1_B.$$
(28)
$$_a_b_a_bcorrespondsto1_A1_B.$$
(29)
After encoding their bits, Alice and Bob apply the CNOT operation. This results in the corresponding four states
$$_a_b+_a_b=(_a+_a)_bcorrespondsto0_A0_B$$
(30)
$$_a_b+_a_b=(_a+_a)_bcorrespondsto1_A0_B$$
(31)
$$_a_b_a_b=(_a_a)_bcorrespondsto0_A1_B$$
(32)
$$_a_b_a_b=(_a_a)_bcorrespondsto1_A1_B.$$
(33)
Bob can now find out Alice’s bit by measuring his qubit in the {$`_b`$, $`_b`$} basis while Alice can find out Bob’s bit by measuring her qubit in the {$`_a+_a`$, $`_a_a`$} basis.
## V The Double CNOT operation on two qubits
One might have thought that the SWAP operation was the unique maximally non-local operation, at least in the terms used in this paper. We here demonstrate that there is another maximally non-local operator, which is the “Double CNOT”, or “DCNOT” gate, formed by acting a CNOT from particle 1 onto particle 2, and then a second CNOT from particle 2 onto particle 1. It is defined by
$$$$
(34)
$$$$
(35)
$$$$
(36)
$$.$$
(37)
To show that DCNOT is maximally non-local, we shall first demonstrate that it can be used to create 2 e-bits. We shall then show that it can be used to communicate 2 bits of information from Alice to Bob, and simultaneously to send 2 bits from Bob to Alice. The argument used for the SWAP operation then proves that to build a DCNOT we need 2 e-bits plus 2 bits of classical communication from Alice to Bob plus 2 bits of classical communication from Bob to Alice. Since any transformation on two qubits can be performed using these resources via teleportation, we will then have shown that the DCNOT is maximally non-local, in terms of resources.
Creating 2 e-bits is easy. Alice and Bob prepare singlets locally, and then perform the DCNOT on spins $`A`$ and $`B`$:
$$(_A_a+_A_a)(_B_b+_B_b)$$
$`_A_a_B_b+_A_a_B_b+_A_a_B_b+_A_a_B_b.`$ (38)
We now have a Schmidt decomposition of rank 4, ie. a 2 party state which is locally equivalent to 2 e-bits between Alice and Bob.
Transmitting 2 bits of information in both directions at the same time is a little more tricky. Alice and Bob need to have 2 e-bits in addition to the DCNOT operation. They first transform their e-bits (locally) into the state
$$_A_a_B_b+_A_a_B_b+_A_a_B_b+_A_a_B_b.$$
(39)
Alice now encodes 1 bit of information in the state by either applying, or not applying $`\sigma _z\sigma _z`$ to her 2 spins. She encodes a second bit of information by applying, or not applying $`\sigma _x`$ to her first spin, $`A`$. Bob similarly encodes two bits of information, using the transformation $`\sigma _z`$ on spin $`B`$ to encode his first bit, and $`\sigma _x\sigma _x`$ to encode his second bit.
Having encoded the information, they make it locally accessible by applying the DCNOT to spins $`A`$ and $`B`$. It is not obvious, but simple to check, that Alice and Bob now each have one of the 4 Bell states locally, and that Alice’s particular state corresponds to Bob’s encoded bits, and vice-versa.
## VI Multi-partite operations
In the previous sections we studied different bi-partite operations. What about multi-partite operations, such as the Toffoli or the Fredkin gates on three qubits? As we showed in section II, they can all be implemented by using the “double teleportation” method. On the other hand, finding the necessary resources is far more difficult than in the bi-partite case; indeed it is not possible at present. The reason is that there exist different inequivalent types of multi-partite entanglement . For example, it is known that singlets and GHZ states are inequivalent in the sense that they cannot be reversibly transformed into each other, not even in the asymptotic limit. Although GHZs (as all other entangled states) can be built out of singlets, such a procedure is wasteful. Hence, when investigating the minimal entanglement resources needed to implement multi-partite quantum operations, we have to use the different inequivalent types of entanglement. Unfortunately, at present multi-partite entanglement is far from being fully understood.
## VII “Conservation” relations
In studying the non-locality of quantum states a most important issue is that of “manipulating” entanglement, i.e. of transforming some states into others . Similarly we can ask: Given a unitary evolution, can we use it to implement some other unitary evolution?
In particular, for pure quantum states we have conservation relations . For example, when Alice and Bob share a large number $`n`$ of pairs of particles, each pair in the same state $`\mathrm{\Psi }`$, they could use these pairs to generate some other number, $`k`$, of pairs in some other state $`\mathrm{\Phi }`$. In the limit of large $`n`$, this transformation can be performed reversibly, meaning that the total amount of non-locality contained in the $`n`$ copies of the state $`\mathrm{\Psi }`$ is the same as the total amount of non-locality contained in the $`k`$ copies of the state $`\mathrm{\Phi }`$. Is something similar taking place for unitary transformations?
For unitary transformations we have not yet studied the case of the asymptotic limit, i.e. performing the same transformation $`U`$ on many pairs of particles. However, an interesting pattern emerges even at the level of a single copy.
Consider first the case of SWAP. We know what the minimal resources needed to implement a SWAP are. But suppose now that we are given a device which implements a SWAP. Could we could use it to get back the original resources needed to create the SWAP?
The balance of resources needed to implement a SWAP can be written as
$$2\text{e-bits}+2\text{bits}_{AB}+2\text{bits}_{BA}=>\text{SWAP}.$$
(40)
The question is whether
$$\text{SWAP}=>2\text{e-bits}+2\text{bits}_{AB}+2\text{bits}_{BA}\mathrm{?}$$
(41)
Though we do not have yet a complete proof, it appears that the answer to the above question is “No”. That is, combining entanglement and classical communication resources to yield a SWAP is an irreversible process \- we cannot use the SWAP to get the resources back.
On the other hand, looking back to the proof of the resources needed for SWAP, we see that we can write the following tight “implications”:
$$2\text{e-bits}+2\text{bits}_{AB}+2\text{bits}_{BA}=>1\text{SWAP}.$$
(42)
$$2\text{e-bits}+1\text{SWAP}=>2\text{bits}_{AB}+2\text{bits}_{BA}.$$
(43)
$$1\text{SWAP}=>2\text{e-bits}.$$
(44)
The first of these three implications is to be read as “given $`2\text{e-bits}`$ and $`2\text{bits}_{AB}`$ and $`2\text{bits}_{AB}`$ we can produce the SWAP operation; also if we wish to produce the SWAP operation with e-bits, and bits communicated from Alice to Bob and vice-versa, we cannot do so with fewer than $`2\text{e-bits}`$ and $`2\text{bits}_{AB}`$ and $`2\text{bits}_{AB}`$.”
The second and third implications have a slightly different meaning. For example we read the second implication as “given 1 SWAP and 2 e-bits, we can communicate 4 classical bits (two each way); also we cannot communicate more than 4 classical bits (two each way) ”. On the other hand, it does not mean that “1 SWAP and 2 e-bits are necessary for communicating 4 classical bits (two each way) ” - for example we can implement this classical communication with 2 SWAPs.
Exactly the same implications apply for the DCNOT.
$$2\text{e-bits}+2\text{bits}_{AB}+2\text{bits}_{BA}=>1\text{DCNOT}.$$
(45)
$$2\text{e-bits}+1\text{DCNOT}=>2\text{bits}_{AB}+2\text{bits}_{BA}.$$
(46)
$$1\text{DCNOT}=>2\text{e-bits}.$$
(47)
Furthermore, very similar implications can be written for the CNOT:
$$1\text{e-bit}+1\text{bit}_{AB}+1\text{bit}_{BA}=>1\text{CNOT}.$$
(48)
$$1\text{e-bit}+1\text{CNOT}=>1\text{bit}_{AB}+1\text{bit}_{BA}.$$
(49)
$$1\text{CNOT}=>1\text{e-bit}.$$
(50)
In fact these implications are very similar to the implications which describe teleportation and super-dense coding which appear, together with many other similar implications on Bennett’s famous transparency presented at almost all early quantum information conferences:
$$1\text{e-bit}+2\text{bits}_{AB}=>1\text{qubit}$$
(51)
$$1\text{e-bit}+1\text{qubit}=>2\text{bits}_{AB}$$
(52)
$$1\text{qubit}=>1\text{e-bit}$$
(53)
The above three implications (51,52,53) are generally thought to describe relations between classical information, quantum information and entanglement. However, we would like to argue that their true meaning is may be more closely related to dynamics, and that a more illuminating form is probably
$$1\text{e-bit}+2\text{bits}_{AB}=>1\text{teleportation}_{AB}$$
(54)
$$1\text{e-bit}+1\text{teleportation}_{AB}=>2\text{bits}_{AB}$$
(55)
$$1\text{teleportation}_{AB}=>1\text{e-bit}$$
(56)
We conjecture that similar relations hold between any quantum action and the resources needed to implement it, that is
$$Entanglement+ClassicalCommunication=>Action$$
(57)
$$Entanglement+Action=>ClassicalCommunication$$
(58)
$$Action=>Entanglement$$
(59)
It may be that these relations hold, in general, only in the asymptotic limit of many copies of the quantum action.
## VIII Different ways of achieving the same task
It is interesting to note that although the transformation from resources to unitary actions is irreversible, sometimes the same end product can be achieved in two different ways. For example, there are two alternative ways to implement
$$2\text{CNOTs}=>1\text{bit}_{AB}+1\text{bit}_{BA}.$$
(60)
The first way is to use one CNOT to transmit 1 classical bit from Alice to Bob and the other CNOT to transmit 1 classical bit from Bob to Alice, i.e.
$$1\text{CNOT}=>1\text{bit}_{AB}$$
(61)
and
$$1\text{CNOT}=>1\text{bit}_{BA}.$$
(62)
Another possibility is to use first one CNOT to create 1 e-bit (50) then the other CNOT plus the e-bit to transmit the 2 classical bits (49), i.e.
$$2\text{CNOTs}=>1\text{e-bit}+1\text{CNOT}=>1\text{bit}_{AB}+1\text{bit}_{BA}.$$
(63)
## IX Catalysing classical communication
A very interesting phenomenon is that of “catalysing” classical communication. This phenomenon is similar in its spirit to that of “catalysing entanglement manipulation” . An example is the following.
On its own, the SWAP can only send one bit in each direction at the same time, and cannot be used for Alice to send 2 bits to Bob, even if Bob sends no information whatsoever. That is,
$$1\text{SWAP}>2\text{bits}_{AB}.$$
(64)
However, if Alice and Bob share 1 e-bit, Alice can send 2 bits to Bob without destroying the e-bit, i.e.
$$1\text{SWAP}+1\text{e-bit}=>2\text{bits}_{AB}+1\text{e-bit}.$$
(65)
This may be done as follows. Initially Alice and Bob share a non-local singlet; Bob also prepares a second singlet locally. Alice encodes the two bits she wishes to send to Bob by performing one of the four rotations 1, $`\sigma _x`$, $`\sigma _y`$, $`\sigma _z`$ on her half of the non-local singlet. By performing the SWAP operation on Alice’s particle from the non-local singlet and one particle of the singlet that Bob has prepared locally, Alice and Bob end up with a non-local singlet held between them; also Bob can find out the two bits by measurements on the local singlet he now holds. Specifically, we begin with the state:
$$(_A_{b1}+_A_{b1})(_B_{b2}+_B_{b2}),$$
(66)
where $`A`$ is Alice’s particle, and $`B`$, $`b1`$ and $`b2`$ are Bob’s particles. Alice performs one of the rotations 1, $`\sigma _x`$, $`\sigma _y`$, $`\sigma _z`$ on her particle. They then perform the SWAP on particles $`A`$ and $`B`$, and get (if Alice performed 1):
$$(_B_{b1}+_B_{b1})(_A_{b2}+_A_{b2})$$
(67)
If Alice performed one of the other rotations, Bob will get one of the other Bell states in system ($`B`$, $`b1`$). Bob now measures that system in the Bell basis to extract the information, and Alice and Bob are left with a singlet between systems $`A`$ and $`b2`$.
In effect the SWAP acts as a double teleportation; one from Alice to Bob and one from Bob to Alice. Teleporting Alice’s qubit, in conjunction with the e-bit, implements a transmission of two bits from Alice to Bob using super-dense coding; it destroys the e-bit in the process. Simultaneously, the Bob to Alice teleportation restores the e-bit.
## X Trading one type of actions for another
An interesting question is the following. There are cases in which two different actions require the same resources. For example the resources needed for 1 SWAP are the same as for 2 CNOTs, i.e., $`2\text{e-bits}+2\text{bits}_{AB}+2\text{bits}_{BA}`$. Now, suppose we had already used the resources to build 2 CNOTs, but we wanted to change our mind and we wanted to do 1 SWAP instead. Due to the irreversibility discussed above, we cannot simply get back the original resources and use them to construct the SWAP. Is it however possible to go directly from 2 CNOTs to 1 SWAP, without going back to the original resources? As far as we are aware, the answer is “No”.
It turns out however that if we have many CNOTs it is nevertheless useful to build a SWAP from CNOTs directly rather than going back to the original resources. Indeed, to obtain the entanglement and classical communication resources needed for 1 SWAP, i.e. $`2\text{e-bits}+2\text{bits}_{AB}+2\text{bits}_{BA}`$ we need 4 CNOTs. However, it is well-known that one can construct 1 SWAP directly from 3 CNOTs. Indeed, we don’t even need 3 CNOTs, but can realize a SWAP by
$$2\text{CNOTs}+1\text{bit}_{AB}+1\text{bit}_{BA}=>1\text{SWAP}$$
(68)
which uses less non-local resources than 3 CNOTs. To see this, it suffices to note that
$$1\text{CNOT}+1\text{bit}_{AB}=>1\text{teleportation}_{AB}$$
(69)
and similarly
$$1\text{CNOT}+1\text{bit}_{BA}=>1\text{teleportation}_{BA}$$
(70)
To implement (69) Alice starts with her qubit in the state $`\mathrm{\Psi }=\alpha +\beta `$ which has to be teleported and Bob with his qubit in the state $``$. After CNOT the state becomes:
$$\mathrm{\Psi }=(\alpha +\beta )\alpha +\beta $$
(71)
Alice then measures her qubit in the $`|+>=\frac{1}{\sqrt{2}}(+)`$ and $`|>=\frac{1}{\sqrt{2}}()`$ basis and communicates the result to Bob. If $`(+)`$ then Bob’s qubit is already in the required state $`\mathrm{\Psi }=\alpha +\beta `$; if $`()`$ then Bob’s qubit is in the state $`\mathrm{\Psi }^{}=\alpha \beta `$ and Bob can obtain $`\mathrm{\Psi }`$ by changing the relative phase between $``$ and $``$ by $`\pi `$.
Note added. While completing this work we became aware of closely related work by J. Eisert, K. Jacobs, P. Papadopoulos and M. Plenio .
|
warning/0005/physics0005017.html
|
ar5iv
|
text
|
# Untitled Document
On the Indication from Pioneer 10/11 Data of an Anomalous Acceleration
Yong Gwan Yi
## Abstract
Hubble’s law, which states a linear increase in velocities with distances, can physically be understood in terms of an acceleration $`cH`$. This work proposes a connection between this “universal” acceleration seen in the solar system and the anomalous acceleration acting on the Pioneer 10/11 spacecraft, in which the Hubble constant inferred from Pioneer 10/11 data is $`87`$ km/s/Mpc. Its physical implication is discussed in relation with Mach’s principle.
By 1998, when Pioneer 10 was 71 AU away from the Sun, one team of researchers at the tracking stations published that radio metric data from Pioneer 10/11 had indicated an apparent anomalous acceleration acting on the spacecraft with a magnitude $`8.5\times 10^8`$ cm/s<sup>2</sup>, directed towards the Sun . When Pioneer 10 ventured beyond the realm of the planetary system, Anderson et al. began monitoring its orbit for evidence of the long-hypothesized Planet X. They found no such planet, but they did notice some extra tiny slowing of its outward motion. Beginning in 1980, when at 20 AU the solar radiation pressure acceleration had decreased to $`<5\times 10^8`$ cm/s<sup>2</sup>, Jet Propulsion Laboratory’s orbit determination program analysis of Pioneer 10/11 data found the biggest systematic error in the acceleration residuals. Even after all known sources of gravity and other forces were taken into account, the apparent acceleration seemed to be present in the residuals.
Ever since the effect was reported, there has been intense debate over its origin. Murphy proposes that the anomalous acceleration can be explained, at least in part, by nonisotropic radiative cooling of the spacecraft electronics. Katz argues that the anomalous acceleration may be due to anisotropic heat reflection off of the back of the antenna dish. Anderson et al. respond, these explanations fall short of accounting for the anomalous Pioneer 10 acceleration. But a few of them suppose a gas leak from thruster to be its origin. Scheffer asserts that the proposed mechanism much more likely explains the anomaly. Meanwhile, it is noted that the size of the anomaly is of the order of $`cH`$, where $`H`$ is the Hubble constant. Rosales and Östvang make attempts to develop a space-time metric which incorporates the effect of cosmic expansion. Nottale tries to tie this to the cosmological constant at the scale of the solar system. Independent of the note, I have come to see an acceleration $`cH`$ in connection with the Pioneer effect. I should like to show a possible account of the anomalous acceleration on physical considerations.
In attempts to explain the effect, my attention focused on the fact that the solar system rotating with the Galactic rotation has a centrifugal acceleration of $`1.8\times 10^8`$ cm/s<sup>2</sup>, the same order of magnitude. Moreover, the centrifugal acceleration was consistent with observation that no magnitude variation of the acceleration with distance was found, within a sensitivity of $`2\times 10^8`$ cm/s<sup>2</sup> over a range of 40 to 60 AU. The points led me to put the weight of its possible explanation in the motion of the solar system.
Non-uniform rotation of our Galaxy gives a hint on its internal motions such as local expansion or contraction while rotating, making an additional contribution to the centrifugal acceleration. It can be estimated using the experimental curve of the rotating velocity versus the distance from the axis . In the curve the gradient of velocity at the position of the solar system is seen to be about $`10`$ km/s/kpc, by which non-uniform rotation makes one order of magnitude small contribution to the centrifugal acceleration . The Coriolis effect on the moving Pioneer at 12.5 km/s is about 11% in magnitude of the centrifugal acceleration.
As no further explanation could be deduced from the Galactic rotation, I turned my attention to the motion of our Galaxy as a whole. Continuing my search for acceleration, I considered with reluctance the possibility of an acceleration in a general recession of distant galaxies. It came out clearly, how the recessional velocities could have been understood in terms of an acceleration.
The announcement by E. Hubble in 1929 of a “roughly linear relation between velocities and distances” established in most astronomers’ minds a sort of bird’s-eye view of a general recession of distant galaxies. But there is a physics to be found in the linear relation. Our information about the frequency shifts comes to us through the observation of light emitted by distant sources. The velocity of a source at distance $`r`$ is a result of velocity difference between the source at an earlier or retarded time $`tr/c`$ and the observation point at time $`t`$. Physically, Hubble’s relation states a roughly linear increase in relative velocity change due to the time of propagation $`\mathrm{}t=r/c`$: $`v=cH\mathrm{}t`$. It becomes evident that the linear increase in recessional velocities with distances is a result of longer light travel times from further distant galaxies. Hubble’s law finds a natural explanation in terms of an acceleration $`cH`$.
The times of propagation permit only the evaluation of galaxies in terms of the retarded positions and velocities. As we look further and further out into space, we see galaxies that are presumably younger and younger, the furthest naturally being those in the remotest past. The linear increase in recessional velocities with distances can therefore be put in the form of a linear decrease in relative velocities with times up to the time of observation. The relation between velocities and times up to the time of observation manifests the direction of acceleration against the recession. The general recession in deep space of distant galaxies must be slowing down at a uniform rate.
It would be of gravitational character occurring on a scale of the universe that the general recession of distant galaxies has been decelerating. From this perspective the value $`cH`$ is identified with the gravitational field of the universe as observed in the solar system. This assumption seems tenable, seeing that the spherically symmetric distribution of matter produces a constant acceleration inside the distribution. But when we identify $`cH`$ as the gravitational field of the universe, we conceive the ultimate interpretation of Slipher’s red shifts as a “universal” gravitational effect. This is because the red shifts can then be understood in terms of a “universal” gravitational potential $`cHr`$. In fact, the red shift effect is an effect of only the relative distances between sources and observation point. From the redshift-distance relation one can only infer that distant galaxies are in free fall; their states of motion remain unaccounted for. In principle, there is no objection to identifying the red shifts ultimately as gravitational red shifts caused by the gravitational field of the universe. In appreciating cosmological relevance of red shifts a change in the orientation of our thought is desirable.
On the basis of the argument we see that there is a “universal” acceleration towards the Sun of $`cH`$. We must adopt an active view—A general recession of distant galaxies is the Sun-based astronomical observation. The solar system would respond to the external gravitational field with the same magnitude, directed away from the Sun. From the general recession of distant galaxies, that is, we can realize an acceleration existing in the relative recession of our own. Pioneer 10/11 moving away from the solar system at the approximately constant velocity make themselves ideal instruments to probe for an additional acceleration existing in the solar system. To the spacecraft the equation of motion would appear as if they are moving under the influence of its inertial force. The anomalous acceleration that has appeared in Pioneer 10/11 tracking would be an inertial reaction to the solar system accelerated relative to distant galaxies. In magnitude and direction their assessment is in substantial agreement with what we should expect from Hubble’s law. Considerations lead to the conclusion that the apparent acceleration acting on the spacecraft is a reflection of the “universal” acceleration as seen in the solar system, in which the Hubble constant inferred from Pioneer 10/11 data is $`87`$ km/s/Mpc.
Of great interest is that the acceleration $`cH`$ has already been assumed in a new law of motion devised by Milgrom . He has imputed the mass discrepancy, observed in galactic systems, not to the presence of dark matter, but to a departure from Newtonian dynamics below the scale of acceleration. A success of the modified dynamics in explaining astronomical data may be interpreted as implying a need to change the law of inertia in the limit of small accelerations. In the previous consideration we have identified the acceleration ultimately as the gravitational field of the universe seen in the solar system. The consideration of the anomalous acceleration naturally leads to speculation about the inertial reference frame defined by the solar system. The issue of inertia piques curiosity.
One may inquire about the modification the anomalous acceleration would assume in the solar system of Newtonian dynamics. Apparently we are guided by a modified dynamics that imputes $`cH`$ to a departure from Newtonian dynamics:
$$\frac{GM_{}}{r^2}\frac{GM_{}}{r^2}+cH.$$
(1)
It represents an attempt to render justice to the fact that Pioneer 10/11 have been slowing down faster than predicted by Newtonian dynamics. The modification makes it obvious that inertia is due not only to the solar gravitational field but also to the gravitational field of the universe. Evidently it indicates that inertial forces do not exactly cancel solar gravitational forces for freely falling planetary systems. The paradigm is obvious. Mach’s principle happens to be true!
Mach’s principle has been the subject of some lively discussion regarding anisotropy of inertia. Cocconi and Salpeter pointed out that there is a large mass near us, the Milky Way Galaxy, and that Mach’s principle would suggest slight differences in inertial mass when a particle is accelerated toward or away from the Galactic center. In the experiments it was shown that with a precision of 1 part in $`10^{20}`$ there is no anisotropy of inertia associated with effects of mass in our Galaxy. Dicke came to defence, arguing that as Mach’s principle associates the inertial reaction with the matter distribution in the universe, an anisotropy in the inertial mass should be universal, the same for all particles. I should like to add defence: The gravitational field of the universe as observed in the solar system is the sum of the gravitational field acting on the Milky Way and the centrifugal acceleration due to rotation about the Milky Way, in which the gravitational field dominates strangely somewhat. Phenomenologically, the gravitational field of the universe as seen in the solar system directs toward the Sun. Thus, an anisotropy of inertia should be expected toward the Sun, and at present I am discussing such possibility from the anomalous acceleration seen in the Pioneer 10/11 spacecraft.
Let us consider the motion of a small body in an orbit around the Sun. The modification (1) is a phenomenological scheme which modifies the solar system into the Newtonian frame of reference which is compatible with Mach’s principle. The added inertia to the solar gravitational field leads to a differential equation for the orbit of the form
$$\frac{d^2u}{d\theta ^2}+u=\frac{mk}{l^2}\left(1+\frac{mcH}{ku^2}\right),$$
(2)
where $`m`$ is the mass of the small body, $`l`$ is the angular momentum, and $`u`$ and $`k`$ denote $`1/r`$ and $`GM_{}m`$. The second term in the round bracket is the one which distinguishes the solar system from the inertial frame of reference. Mach’s principle can be formulated in this way if the anomalous acceleration is assumed to be the inertial reaction to the matter distribution in the universe.
We may solve the inertial system equation approximately. We expand the periodic solution of the equation into a series
$$u=\alpha +\lambda \beta _1+\alpha ϵ\mathrm{cos}(\rho \theta )+\lambda \underset{n=2}{\overset{\mathrm{}}{}}\beta _n\mathrm{cos}(n\rho \theta ),$$
(3)
where $`\alpha =mk/l^2`$, $`\lambda =mcH/k`$, and $`ϵ`$ is the eccentricity of the ellipse . We substitute the series solution into the equation. For $`\lambda /u^2`$, we expand
$`{\displaystyle \frac{\lambda }{u^2}}`$ $``$ $`{\displaystyle \frac{\lambda }{\alpha ^2(1+ϵ\mathrm{cos}(\rho \theta ))^2}}`$ (4)
$``$ $`{\displaystyle \frac{\lambda }{\alpha ^2}}(12ϵ\mathrm{cos}(\rho \theta )+3ϵ^2\mathrm{cos}^2(\rho \theta )4ϵ^3\mathrm{cos}^3(\rho \theta )+\mathrm{}).`$
By comparing the $`\mathrm{cos}(\rho \theta )`$ terms we obtain the equation which determines $`\rho `$ to a first approximation. According to this calculation, the elliptical orbit of a planet referred to the Newtonian frame of reference rotates in the opposite direction as the planet moves, with a rate that is given by
$$\frac{2\pi cHa^2(1ϵ^2)^2}{GM_{}}(1+\frac{3}{2}ϵ^2+\frac{15}{8}ϵ^4+\mathrm{}),$$
(5)
where $`a`$ is the planetary semimajor axis.
Equation (5) describes the rate of precession at which the perihelion will have retarded per revolution. The precession expected from Mach’s principle increases rapidly as we move away from the Sun. For Mercury it gives the value of $`10^{\prime \prime }`$ per century and for Earth the value of $`16.34^{\prime \prime }`$. They destroy the current agreement between the general theory of relativity and the observed anomalous precessions. Strongly it casts doubt on the validity of calculation. Is my calculation erroneous? Or is there some unrecognized effect in observations?
We need to look back at the Pioneer effect. The effect could only be seen beyond 20 AU. The anomalous acceleration acting on Pioneer 10/11 could not be found until the solar radiation pressure had decreased to less than a critical value. The solar radiation pressure decreases as $`r^2`$. As indicated for the Pioneers, at distances $`>1015`$ AU it produces an acceleration that is much less than $`8\times 10^8`$ cm/s<sup>2</sup>, directed away from the Sun. Hence, even granting that the rate (5) is in principle expected, we should be aware that the inertial effect may possibly be contributing to the motion of distant planets such as Uranus, Neptune, and Pluto. On the motion of near planets would the inertial effect be entirely masked by the solar radiation pressure, and there is no prospect of its being measured.
Brans and Dicke have attempted to incorporate Mach’s principle into general relativity. They suggest field equations with a long-range scalar field produced by the total mass in the visible universe. In line with the interpretation of Mach’s principle, the long-range scalar field matches the “universal” acceleration $`cH`$ seen in the solar system. The modification (1) replaces the Schwarzschild solution by its generalization
$$1\frac{2}{c^2}\left(\frac{GM_{}}{r}\right)1\frac{2}{c^2}\left(\frac{GM_{}}{r}cHr\right).$$
(6)
We are thus led to an alternative approach by assuming that Einstein’s field equations still apply, but that the metric differs from the Schwarzschild solution by the gravitational field of the universe seen in the solar system. Just like an approximate expression $`gh`$ for gravitational potential at height $`h`$ on the Earth’s surface, so will be an expression $`cHr`$ for gravitational effects having their origin in the universe surrounding the solar system. The generalization (6) introduces a new term $`H/c`$ in addition to the relativistic term in the right hand side of (2). But it is extremely small compared to the other terms. In their theory, Brans and Dicke make mention of the gravitational red shift and the deflection of light in the context of Mach’s principle. In my view, however, these phenomena seem to be of optic nature in relation to property of the medium of propagation .
J D Anderson, P A Laing, E L Lau, A S Liu, M M Nieto, S G Turyshev,
Phys. Rev. Lett. 81 (1998) 2858.
E M Murphy, Phys. Rev. Lett. 83 (1999) 1890.
J I Katz, Phys. Rev. Lett. 83 (1999) 1892.
J D Anderson, P A Laing, E L Lau, A S Liu, M M Nieto, S G Turyshev,
Phys. Rev. Lett. 83 (1999) 1891; 1893; Phys. Rev. D65 (2002) 082004.
L K Scheffer, Phys. Rev. D67 (2003) 084021.
J L Rosales, e-print gr-qc/0212019.
D Östvang, Class. Quantum Grav. 19 (2002) 4131.
L Nottale, e-print gr-qc/0307042.
D P Clemens, Astrophys. J. 295 (1985) 422.
H Lamb, Hydrodynamics (Dover, New Yok, 1945), 6th ed. p.28.
M Milgrom, Astrophys. J. 270 (1983) 365; 371; 384;
J Bekenstein, M Milgrom, Astrophys. J. 286 (1984) 7.
G Cocconi, E E Salpeter, Nuovo Cimento 10 (1958) 3608;
Phys. Rev. Lett. 4 (1960) 176.
V W Hughes, H G Robinson, V Beltran-Lopez, Phys. Rev. Lett. 4 (1960) 342;
R W P Drever, Phil. Mag. 6 (1961) 683.
R H Dicke, Phys. Rev. Lett. 7 (1961) 359.
P G Bergmann, Introduction to the Theory of Relativity
(Prentice-Hall, New Delhi, 1977), p.215.
C Brans, R H Dicke, Phys. Rev. 124 (1961) 925.
Y G Yi, e-print physics/0006006.
|
warning/0005/hep-ph0005068.html
|
ar5iv
|
text
|
# Can LSND be included in a 3-Neutrino framework?
## Acknowledgments
One of the authors (JDV) is happy to acknowledge partial support from the TMR contract ERBFMRXCT96-0990.
Fig. 1
Fig. 2
|
warning/0005/astro-ph0005553.html
|
ar5iv
|
text
|
# A Measurement of the Temperature-Density Relation in the Intergalactic Medium Using a New Ly𝛼 Absorption Line Fitting Method The observations were made at the W.M. Keck Observatory, which is operated as a scientific partnership between the California Institute of Technology and the University of California; it was made possible by the generous support of the W.M. Keck Foundation.
## 1 INTRODUCTION
The Lyman-$`\alpha `$ forest absorption in the spectra of quasars provides a wealth of information about the properties of the intergalactic medium (hereafter, IGM). There has recently been a lot of interest in using the distribution of Doppler parameters of fitted absorption lines, measuring the total velocity dispersion of the gas, to constrain the temperature of the IGM (Schaye et al. 1999; Ricotti, Gnedin, & Shull 2000; Bryan & Machacek 2000). In this paper we develop a method to identify and to fit absorption lines, and to obtain the gas temperature from the distribution of the Doppler parameters of the lines. Our algorithm is intended to be simple to implement and be applied identically on simulations and observations.
The temperature of the IGM as a function of density is primarily determined by the balance between adiabatic cooling and photoionization heating, once ionization equilibrium with the background radiation has been established. However, during the epoch of reionization, the heating rate is higher because every atom needs to be ionized once (and the ionization can occur on a short time-scale compared to the recombination rate), and the high opacity of the low-density IGM implies that high-frequency photons are absorbed, delivering a much greater amount of heat for each ionization (e.g., Miralda-Escudé & Rees 1994; Hui & Gnedin 1997; Haehnelt & Steinmetz 1998; Abel & Haehnelt 1999; Gnedin 2000). Other sources of heating may also contribute, such as Compton heating by the X-ray background (Madau & Efstathiou 1999), or photoelectric heating by dust grains (Nath, Sethi, & Shchekinov 1999). Constraining these sources of heating is one of the two primary reasons why we are interested in measuring the temperature. The other reason is the need to make accurate predictions for the statistics of the Ly$`\alpha `$ forest flux in order to constrain cosmological parameters (e.g., Rauch et al. 1997; Weinberg et al. 1997; Croft et al. 1998; Hui 1999; McDonald & Miralda-Escudé 1999; Hui, Stebbins, & Burles 1999; Croft et al. 1999b; Weinberg et al. 1999; Croft, Hu, & Davé 1999a; Nusser & Haehnelt 2000; McDonald et al. 2000). The temperature-density relation affects the predicted relationship between the power spectrum of the transmitted flux and the power spectrum of the initial mass density perturbations (Nusser & Haehnelt 2000), as well as the predicted mean transmitted flux (Rauch et al. 1997; McDonald et al. 2000), which can be used to constrain the baryon density of the universe.
Recent Ly$`\alpha `$ forest simulations have shown that, when the structure of the absorption systems is adequately resolved, the predicted absorption line widths are smaller than observed if the temperature of the IGM is determined from photoionization equilibrium, well after reionization has ended (Theuns et al. 1998, 1999a; Bryan et al. 1999). To solve the discrepancy the temperature apparently needs to be higher. Several authors have presented measurements of the IGM temperature using different methods, generally finding values moderately higher than expected from photoionization equilibrium (Theuns et al. 1999b; Ricotti et al. 2000; Bryan & Machacek 2000; Schaye et al. 2000).
Our aim in this paper is to provide a new unambiguous measurement of the temperature, making a more exhaustive analysis than in previous work of the model uncertainties that result from comparing the observational results with a simulation. We develop a new line-fitting method as an alternative to the standard Voigt-profile fitting with line deblending, which is much faster, unambiguous, and easy to implement. Our method works by essentially assigning one line to each sufficiently deep minimum in the transmitted flux, and measuring the line width and central optical depth for each line. The gas temperature at each density is then derived from the distribution of line widths at each central optical depth. The systematic uncertainties and model dependence of the method used to derive the temperature are carefully analyzed, in a more extensive way than it was done in previous work. The new method is applied to observational data and to a simulation in exactly the same way, computing error bars due to the variance in our observed sample.
The main idea of the method to measure the temperature of the IGM was suggested by Schaye et al. (1999), Ricotti et al. (2000), and Bryan & Machacek (2000). The probability distribution of Doppler parameters, $`P(B)`$, is characterized by a lower cutoff, $`B_C`$, where $`P(B)`$ rises sharply, with very few lines having narrower Doppler parameters than this cutoff. The idea is that this cutoff is a measure of the gas temperature. In general, absorption lines have both a thermal and a hydrodynamical contribution to their breadth; however, for any set of lines with similar gas temperature, the narrowest ones will be those where the velocity field along the line of sight through the absorber is close to a caustic, so that the variation in the fluid velocity is minimized and thermal broadening dominates the observed line width. In fact, it was found by Theuns et al. (1999b) that the narrowest absorption lines are primarily thermally broadened.
A tight relationship between density and temperature in the IGM for gas at low densities, where shock heating is not very important, is expected theoretically and is found in numerical simulations (Hui & Gnedin 1997; Theuns et al. 1998). This implies that the narrowest lines at a given gas density (corresponding approximately to the optical depth at the line center) are not selected to have low gas temperature, but low fluid velocity dispersion. This justifies estimating the temperature from the lower cutoff of the Doppler parameter distribution.
In §2 we briefly describe the observational data and the simulation that we use. In §3 we describe our line fitting algorithm. In §4 we demonstrate how the line fitter works by running it on spectra from the numerical simulation. In §5 we describe our method for estimating the temperature from Doppler parameter distribution, testing the conditions under which the temperature can be recovered in a model-independent way. In §6 we use the line fitter on the observational data and give results for the measured temperatures. The results are discussed in §7. The Appendix describes further details of our line-fitting method.
## 2 THE OBSERVATIONAL DATA AND THE SIMULATION
### 2.1 Observations
We use the same set of eight quasar spectra as in Paper I . These spectra have sufficiently high resolution and signal-to-noise ratio to measure the shape of each absorption feature. The pixel noise is typically less than 5% of the continuum flux level, and frequently as low as 1%. The velocity resolution is 6.6 $`\mathrm{km}\mathrm{s}^1`$ (FWHM) and the spectra are binned in 0.04 Å pixels. More details and statistics of this data set are given in Paper I and references therein.
The seven quasars from the Rauch et al. (1997) data set have previously constructed lists of regions that are suspected of containing metal lines. Our main results include these regions in the spectra because they are not positively identified as containing metal lines.
In Paper I we defined three redshift bins: $`3.39<z<4.43`$, $`2.67<z<3.39`$, and $`2.09<z<2.67`$, with mean redshifts $`\overline{z}=3.9`$, $`\overline{z}=3.0`$, and $`\overline{z}=2.4`$. We use these same three bins in this paper, which contain approximately the same amount of data.
### 2.2 Simulation
We test the profile fitting code and the procedure to measure the temperature on the output of the Eulerian hydrodynamical simulation described in Miralda-Escudé et al. (1996) (referred to as L10 in that paper). The cosmological model used has $`\mathrm{\Omega }_0=0.4`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.6`$, $`h=0.65`$, $`\sigma _8=0.79`$, and large scale primordial power spectrum slope $`n=0.95`$. The box size of the simulation is $`10h^1`$ Mpc, and it contains $`288^3`$ cells. We use outputs from the simulation at $`z=`$4, 3, and 2.
#### 2.2.1 Generation of Simulated Spectra
Ly$`\alpha `$ spectra are computed for a large number of lines of sight along the box axes. There is one free parameter that we can vary when computing the spectra, the normalization of the optical depth, which we adjust to reproduce the mean transmitted flux of the observations that we are comparing to (the values of the mean transmitted flux are taken from Paper I ). Renormalizing the optical depth is equivalent to modifying the intensity of the ionizing background, as long as the effect of collisional ionization and the change in the gas temperature caused by the different heating rate can be neglected (see Theuns et al. 1998 for a test that these effects are in fact negligible). The optical depth is then mapped to transmitted flux using $`F=\mathrm{exp}(\tau )`$.
For each line of sight through the simulation (parallel to one of the three axes), we estimate the effects of continuum fitting by defining the maximum transmitted flux along the line of sight to be the continuum flux, $`F_c`$, and dividing the flux in all other pixels in the line by $`F_c`$. We map the 288 cells along a line of sight onto smaller cells, with their size chosen to match the observations that we want to compare to. Finally we convolve the spectra with the instrumental resolution of $`6.6\mathrm{km}\mathrm{s}^1`$, and we add Gaussian noise to each cell with a flux dependent dispersion $`n(F)`$, which is taken from Paper I .
#### 2.2.2 The $`T\mathrm{\Delta }_g`$ Relation in the Simulation
Before we describe the method we shall use to measure the temperature, it will be useful to examine the temperature-density relation in the simulation. In this paper we parameterize the mean temperature-density relation (hereafter referred to as the $`T\mathrm{\Delta }_g`$ relation ) as a power-law, $`T=T_0\mathrm{\Delta }_g^{\gamma 1}`$, where the gas over-density is $`\mathrm{\Delta }_g\rho _g/\overline{\rho }_g`$, for the purpose of measuring this relation from the data. Although the $`T\mathrm{\Delta }_g`$ relation naturally approaches a power-law form with $`\gamma 10.6`$ when the thermal evolution is determined by photoionization heating and adiabatic expansion alone (Hui & Gnedin 1997), in general it deviates significantly from a power-law.
Figure 1 shows scatter plots of T versus $`\mathrm{\Delta }_g`$ for the three redshift outputs of the simulation. The solid lines are power-law fits to the range $`1<\mathrm{\Delta }_g<2`$, to show that the $`T\mathrm{\Delta }_g`$ relation in the simulation is only roughly consistent with a power law. There is a substantial dispersion of the temperature at a given density, and the relation between the mean temperature and the density deviates from a power-law. For example, using the $`z=3`$ simulation output, the best power-law fits in the restricted ranges of density $`\mathrm{\Delta }_g=`$(0-1, 1-2, 2-3) yield $`\gamma 1=`$(0.15, 0.30, 0.39), with a mean fractional temperature deviation around the fits equal to (4%, 11%, 18%). The dispersion is due to shock-heating and to the variable expansion or contraction histories of the gas at a fixed density. The reionization of He II occurs near $`z3`$ in this simulation, heating the low-density gas to a temperature that is nearly constant with density. The parameters of the power-law fits to all three simulation outputs, in the range $`\mathrm{\Delta }_g=`$(1-2), are given in Table 1.
Because of this deviation from a power-law form of the $`T\mathrm{\Delta }_g`$ relation , a measurement of the temperature $`T_0`$ and power-law index $`\gamma `$ should be understood only as an approximation to the true mean $`T\mathrm{\Delta }_g`$ relation , near the effective density at which the measurement is made, and this effective density needs to specified. We use the following form to present our results in §6:
$$T=T_{}(\mathrm{\Delta }_g/\mathrm{\Delta }_{})^{\gamma 1},$$
(1)
where $`\mathrm{\Delta }_{}`$ is chosen so that the error bars on $`T_{}`$ and $`\gamma 1`$ are uncorrelated.
#### 2.2.3 Definition of Temperature and Density at Points in Spectra
Each pixel in a spectrum receives optical depth contributions from an extended stretch of real space along the line of sight, so no unique temperature or density can be associated with the pixel. However, we can define the temperature and the gas density at pixels in spectra to be the optical depth-weighted average over the temperature and the density of all the gas that contributes to absorption in the pixel. This definition will be used later to assign a gas temperature to any identified absorption line, which we will define to be equal to the temperature of the central pixel. Table 1 shows fits to the $`T\mathrm{\Delta }_g`$ relation for pixels, restricted over the same density range $`1<\mathrm{\Delta }_g<2`$: we see that the mean relation is almost the same as for random points in space. This is true in spite of a larger difference in the density distribution; for example, the median $`\mathrm{\Delta }_g`$ is $`(0.47,0.39,0.31)`$ for random points at $`z=(4,3,2)`$, and is $`(0.52,0.47,0.41)`$ for spectral pixels. This difference in the median density is caused by the thermal broadening and velocity dispersion in absorbers, which spread high density regions out into low density regions.
## 3 FITTING METHOD
In this section we describe the procedure that we use to identify and fit absorption lines. The temperature measurement is based on the identification of absorption lines that can be adequately fit by a single Gaussian in optical depth, over a certain interval around a point where the optical depth is maximum; the narrowest widths among the lines that can be fitted in this way will give us the gas temperature. Absorption systems that cannot be fitted by a single Gaussian must be broadened by a non-Gaussian distribution of the fluid velocity, and can therefore be discarded for the purpose of measuring the gas temperature. In contrast to the standard Voigt profile fitting approach, we make no attempt to fit the entire spectrum by superposing many absorption lines. Instead, we fit only small regions around minima of the transmitted flux, each one with a single Gaussian absorber. We therefore have a constant number of parameters to fit for each absorption line, making the algorithm simple, unambiguous, and fast.
Before we explain the procedure for identifying and fitting lines, it is useful to understand qualitatively what the results will look like.
Figure 2 shows a section of the spectrum of Q1422, with the transmitted flux indicated by the dotted line, and the fitting solutions indicated by the solid lines. Our method has selected all statistically significant maxima in optical depth that are well fitted by a single Gaussian and identified them as absorption lines and has discarded the rest. The fits are done over the regions indicated by the solid lines. We will return to this figure once we have described the process used to fit the lines.
Our fitting method consists of taking each pixel in the spectrum as a candidate for containing the center of an absorption line. After requiring several conditions and eliminating most of the pixels as candidates, a final list of absorption lines is obtained, each one having a fitting window where a fit to the three parameters of the line (line center, central optical depth, and line width) is performed. None of the fitting windows from adjacent lines can overlap in the final list, as seen in the example of Figure 2.
The first operation is to determine an integer window width, $`W`$, which sets the fitting region around each pixel $`P`$, going from $`PW`$ to $`P+W`$. The width $`W`$ is the smallest one for which the following condition is obeyed, which ensures that there is a significant decline of the flux from the average value at the edges of the fitting window to the line center:
$$\frac{1}{2}\left[F\left(P+W\right)+F\left(PW\right)\right]F\left(P\right)>E_d\sigma (P,P\pm W),$$
(2)
where
$$\sigma (P,P\pm W)\left[\sigma ^2\left(P\right)+\frac{1}{4}\left(\sigma ^2\left(P+W\right)+\sigma ^2\left(PW\right)\right)\right]^{1/2},$$
(3)
$`\sigma (P)`$ is the noise at pixel $`P`$, and $`E_d`$ is the first parameter of the fitting algorithm. $`E_d`$ is the number of “$`\sigma `$” significance required of the flux decrease.
There are of course some pixels where the condition in equation (2) is never obeyed for any width. In practice, the width $`W`$ is increased only up to some maximum value $`W_{max}`$ before the pixel is discarded as a candidate for a line center. This maximum width is chosen to be large enough so that it does not affect the results of the algorithm. In addition, the minimum value of $`W`$ is set to 2 pixels (so that the fitting region has at least 5 pixels). Some additional parameters are used in the algorithm to expedite the elimination of pixels as candidates for line fits, increasing the speed of the code without affecting the final result; these details are described in the Appendix.
For each pixel $`P`$ where a fitting width $`W`$ has been determined in this way, the line fit is performed by $`\chi ^2`$ minimization. We fit the following profile to the flux within the window:
$$F(v)=\mathrm{exp}\left[\tau _c\mathrm{exp}\left(\frac{1}{2}\frac{\left(vv_c\right)^2}{\sigma _b^2}\right)\right].$$
(4)
Here, $`v`$ is the distance from the central pixel $`P`$. The three parameters of the fit are $`\tau _c`$ (the optical depth at the line center), $`v_c`$ (the location of the line center), and $`\sigma _b`$ (the width of the line). The parameter $`v_c`$ is constrained to lie within the pixel $`P`$, so that when $`P`$ is not close to the center of the line, a good fit will not be obtained. To account for the instrumental resolution, the function in equation (4) is convolved with a Gaussian filter of width matching the resolution of the data. After the fit is performed, we impose a goodness-of-fit requirement: the probability of exceeding by chance the value of $`\chi ^2`$ for the best fit should be larger than a certain value $`P_0`$, which is a second parameter of our method. If the requirement is not satisfied, the pixel $`P`$ is discarded as a candidate for including the center of an absorption line. The two parameters $`E_d`$ and $`P_0`$ can be adjusted to optimize the temperature measurement, based on tests using numerical simulations that we will present in §4.
Because we require $`v_c`$ to be within pixel $`P`$, acceptable fits to absorption lines are only found in pixels that are indeed close to a minimum of the flux, in an absorption feature that can be adequately fitted to a single Gaussian within a certain window. Typically, the list of pixels where acceptable fits are found will include groups of a few adjacent pixels around such minima. The final step of our algorithm is to select among any group of pixels with accepted line fits that are within their own fitting windows the one that yielded the best fit. This produces the final list of absorption lines.
We can now understand the fitting example shown in Figure 2. The widths of the fitting windows, shown by the solid lines, are set by the value of $`E_d`$ (here, $`E_d=12`$) and the noise level (which in this case varies from $`0.004`$ in saturated pixels to $`0.01`$ at the continuum). The apparent maxima of absorption that have no corresponding fitted line do not increase in flux enough at their edges to satisfy the requirement in equation (2). The cluster of fitted lines near $`\lambda =4830`$Å demonstrates that our procedure does not automatically eliminate lines in blends, as long as they are clearly distinct maxima. Apparently, the requirement that the optical depth be consistent with a Gaussian curve is easily fulfilled by the peaks of all of the significant absorption lines.
Throughout this paper, we shall be expressing all results concerning the line widths in terms of the equivalent temperature, denoted as $`B`$, when the line width is assumed to be due to thermal broadening: $`B10000(\sigma _b/9.09\mathrm{km}\mathrm{s}^1)^2`$ K. This allows for an easier comparison of the results of numerical simulations and observations. Note that line widths have usually been presented in the literature in terms of the Doppler parameter, $`b=2^{1/2}\sigma _b`$.
An important property of this algorithm is that the distribution of line widths it measures should converge to a fixed answer as the signal-to-noise ratio of the observations is increased. In the limit of negligible noise and pixel size, and perfect resolution, every true minimum of the transmitted flux in the spectrum should be identified, and the fitted line width should reflect the second derivative around the minimum, because the size of the fitting window around each minimum should be very small. Moreover, the second derivative around minima is a physically well motivated quantity to obtain the gas temperature. In contrast, the Voigt profile fitting method does not converge at high signal-to-noise ratio because the number of blends assumed in an absorption system will change. Whereas the Voigt profile method attempts to fit the entire spectrum by superposing lines, our new method fits only small regions around the minima of transmitted flux as arising from a single absorber.
## 4 APPLICATION OF THE PROFILE FITTER TO THE NUMERICAL SIMULATION
### 4.1 Detailed Example of Fitting the Simulated Spectra
We now apply our method to 1500 randomly selected lines of sight through the simulation output at $`z=3`$, with the mean flux decrement, noise level and pixel size set to match the observations at $`\overline{z}=3`$ (see Paper I ). We first set the two parameters of the line-fitting algorithm to $`E_d=12`$ and $`P_0=0.01`$ (we shall analyze the optimal values of these parameters in §4.2). A total of 6378 lines are identified and successfully fitted. For each absorption line we obtain four quantities: the optical depth at the line center, $`\tau _c`$; the Doppler parameter converted to temperature units, $`B`$; the optical-depth-weighted temperature at the central pixel, $`T`$; and the optical-depth-weighted gas density at the central pixel, $`\mathrm{\Delta }_g`$.
Figure 3(a) shows $`BT`$ vs. $`\tau _c`$ for all 6378 fitted lines. The solid lines at $`BT=0`$ K and $`BT=\pm 3000`$ K are to guide the eye in evaluating the contribution of non-thermal broadening to $`B`$. The noise that has been added to the simulated spectra (see §2.2.1) is responsible for most of the absorption lines with $`B<T`$. In Figure 3(b) we show the fitted lines from the same set of spectra as Figure 3(a), but without adding noise (although the noise level that each pixel should have was still used to weight the $`\chi ^2`$ fits); the small number of lines that had $`B<T`$ in the presence of noise have been almost entirely eliminated (the few remaining lines with $`B<T`$ arise because the optical-depth-averaged temperature at the central pixel of a line can be skewed by a contribution from very hot gas).
Figure 3 demonstrates that the contribution to $`B`$ from fluid motions is generally large, and therefore the Doppler parameter of an individual line will usually overestimate the gas temperature. For $`\tau _c1`$, the distribution of $`BT`$ extends to values as low as $`1000`$ K, implying that the lower cutoff in the $`B`$ distribution should provide a good measure of the median gas temperature at a given $`\tau _c`$ \[assuming that the lowest gas temperatures do not extend very much below the median value, as is true in the simulation (Fig. 1)\]. However, for weak lines ($`\tau _c<1`$), thermal broadening is never clearly dominating, preventing a model-independent estimate of the gas temperature at the correspondingly low gas densities. The reason why the breadth of weak lines is always dominated by motions is that all the low density gas has not turned around from the Hubble expansion.
The eight histograms in Figure 4, constructed from the same fitted lines as Figure 3(a), show more quantitatively the cutoff on the $`BT`$ distribution near $`BT=0`$, for different optical depths. In the $`\tau 0.3`$ and $`\tau 0.6`$ panels of the Figure, the problem of estimating the temperature of the lower density gas is clearly seen. Over the range $`1\tau 20`$, the distribution of $`BT`$ is the desirable one for our temperature measurement: there are many lines near $`BT=0`$, but very few with $`B<T`$.
The sharpness of the cutoff in Figure 4 should of course be degraded in the observable $`B`$ distribution, owing to the scatter in the temperatures of the lines.
Figure 5 shows histograms of $`B`$ along with histograms of $`T`$. By observing the histogram of $`B`$ we would like to determine the median of $`T`$. For $`\tau _c>1`$, the cutoff on the distribution of $`B`$ appears to coincide well with the peak of the distribution of $`T`$. As we saw clearly in Figure 3, the number of lines with $`BT`$ decreases quickly for $`\tau _c<1`$.
It is interesting to note in Figure 5 the change in the distribution of the true temperature in the simulation with $`\tau _c`$. At low optical depth, the $`T`$ distribution is very narrow as a result of the simplicity of the evolution history of gas at low density, which generally expands peacefully in the voids, heated by photoionization and cooled adiabatically by its expansion. As the density increases, the heating history of the gas becomes more heterogeneous. The gas at higher densities is shock-heated more frequently and to a greater degree, and the evolution of the density itself since reionization (when the initial temperature is set) is more highly variable. For $`\tau 10`$ the range of gas temperatures increases, making the interpretation of the cutoff on the $`B`$ distribution more ambiguous.
### 4.2 Optimizing the Fitting Control Parameters
The two parameters of the line fitter, $`E_d`$ and $`P_0`$, should be optimized to give the best statistical error bars on the location of the lower cutoff of the distribution of $`B`$, which we will use to estimate the gas temperature. In order to do this optimization, we define a quality measure, $`Q=(N_gN_b)(N_g+N_b)^{1/2}`$, where $`N_g`$ and $`N_b`$ are the number of fitted lines with $`0\mathrm{K}<BT<3000`$ K, and with $`3000\mathrm{K}<BT<0`$ K, respectively. Basically, $`Q`$ is a measure of the statistical significance of the increase in the number of lines as the $`B=T`$ cutoff is crossed, computed by comparing the number of lines in bins of width 3000K on each side of the cutoff. The larger the value of $`Q`$, the more accurately we should be able to determine the temperature. For example, the set of fitted lines shown in Figures 3-5 (fitted using $`E_d=12`$ and $`P_0=0.01`$) gives $`N_g=281`$ and $`N_b=40`$, resulting in $`Q=13.5`$. When we remove the noise in Figure 3(b) we find $`N_g=236`$ and $`N_b=1`$, yielding $`Q=15.3`$. We use a bin width of $`3000`$ K because the errors in the measured temperature from our data will be of about this magnitude (see §6).
Table 2 shows the $`Q`$ values for a broad range of values of $`E_d`$, and $`P_0=0.01`$. We also list the $`Q`$ value using two different values of $`P_0`$, removing the noise from the spectra, removing the continuum fitting approximation, and using different random seeds for the added noise. We find that the changes in $`Q`$ are usually smaller than the changes that can result from simply using a different set of random numbers for the noise that is added to the spectra. We conclude that the precise values taken by the parameters are not actually very important to the results. We use $`E_d=12`$ and $`P_0=0.01`$ when we analyze spectra matching the $`\overline{z}=3`$ observational properties in the rest of this paper. From the range of $`E_d`$ with $`Q`$ close to its maximum, we chose the smallest value of $`E_d`$, $`E_d=12`$, because we expect that increasing the number of accepted lines will make the analysis procedure more robust, particularly the computation of the error bars.
We use two other redshift bins for the observational data (see §2.1), $`\overline{z}=2.4`$ and $`\overline{z}=3.9`$. Because the mean flux decrement, pixel size, and noise level are different in each bin, we determine a best value of $`E_d`$ separately for each. Tests similar to the one in Table 2, using the mean flux decrement, noise level, and pixel size matching the observations in the high and low $`\overline{z}`$ bins, show that $`E_d=9`$ and $`E_d=8`$ are the best values to use when analyzing data in the $`\overline{z}=2.4`$ and $`\overline{z}=3.9`$ bins, respectively, although the values of $`Q`$ obtained are only weakly sensitive to $`E_d`$ in each case. We fix $`P_0=0.01`$ at all three redshifts, because changing it does not significantly increase $`Q`$.
## 5 TESTS OF THE DERIVATION OF THE IGM TEMPERATURE USING THE SIMULATION
In this section we combine our line fitting method with elements of the method of Bryan & Machacek (2000) for constraining the temperature of the IGM by measuring the lower cutoff on the distribution of $`B`$. First we define the cutoff and how we associate it with an estimated temperature. Then we use the simulation to translate the observed temperature-optical depth relation into the desired $`T\mathrm{\Delta }_g`$ relation . All the results in this Section are obtained from the simulation. In §6 we present results for the temperature measured from the observational data, using the method described in this Section.
### 5.1 Locating the Cutoff on the Distribution of $`B`$
Bryan & Machacek (2000) presented a useful technique for quantifying the lower cutoff on a $`B`$ histogram like those in Figure 5. They smooth the histogram with a Gaussian filter and define the cutoff to be the location of the maximum of the first derivative of the smoothed histogram. This derivative is given by
$$\frac{dP}{dB}_i\underset{j=1}{\overset{N}{}}(B_jB_i)\mathrm{exp}\left[\frac{1}{2}\frac{\left(B_iB_j\right)^2}{\sigma _B^2}\right],$$
(5)
where the $`i`$th bin has temperature $`B_i`$, $`j`$ is the label for $`N`$ individual fitted absorption lines with fitted temperature $`B_j`$, and $`\sigma _B`$ is the smoothing length.
We plot in Figure 6 the smooth histogram of $`B`$ and its derivative, as well as the smoothed histogram of the gas temperature, for $`\sigma _B=5000`$ K, using spectra from the $`z=3`$ simulation output with the $`\overline{z}=3`$ observational noise, pixel size and mean flux decrement.
We define the Doppler parameter cutoff, $`B_C(\tau _c)`$, to be the value of $`B`$ where $`dP/dB`$, given by equation (5), is maximum, for the optical depth bin labeled by $`\tau _c`$. Our estimate of the gas temperature at optical depth $`\tau _c`$ is $`B_C(\tau _c)`$, after applying a small correction that we describe in detail in the remainder of this section.
#### 5.1.1 Error Bars on $`B_C(\tau _c)`$
We compute error bars on the location of the cutoff on the $`B`$ distribution ($`B_C`$) by bootstrap resampling (Press et al. 1992). We generate a bootstrap realization of the $`B`$ histogram, for an optical depth bin containing $`N`$ fitted lines, by selecting $`N`$ lines at random from those in the bin (with replacement) and recomputing the histogram from the new set of lines. The error on $`B_C`$ is given by the dispersion in the $`B_C`$ values measured from many bootstrap realizations of the histogram.
### 5.2 Comparison Between $`B_C(\tau _c)`$ and the Temperatures in the Simulation
In this subsection we investigate the relationship between $`B_C(\tau _c)`$ and the physical temperature, $`T`$, of the absorbers with central optical depth $`\tau _c`$. We use spectra from the simulation, where we know the optical-depth-weighted temperature and density, $`T`$ and $`\mathrm{\Delta }`$, of the absorption lines. We match the simulated spectra to the mean flux decrement, noise level and pixel size of the observations at $`\overline{z}=3`$ (see §2.2.1). In §5.3 we show how the comparisons change when the simulated spectra are matched to the properties of the $`\overline{z}=3.9`$ and $`\overline{z}=2.4`$ observations.
Figure 7 explores the meaning of $`B_C(\tau _c)`$ in detail, using the $`z=3`$ simulation results. We have separated all the fitted lines into 10 bins of the central optical depth, $`\tau _c`$, choosing the bins so that each contains an equal number of absorption lines. We smooth the $`B`$ histogram with $`\sigma _B=5000\mathrm{km}\mathrm{s}^1`$ to compute the cutoff temperature in every optical depth bin. Then, we compute the median central optical depth and gas temperature of the set of absorption lines that satisfy $`|BB_C|<5000`$K, in each optical depth bin. These sets of lines are the ones that actually determine the location of the cutoff in the $`B`$ histogram, so we examine first the relation between $`B_C`$ and their median temperature, which we denote by $`T_P`$. The values of $`B_C`$ are shown as crosses with error bars. <sup>1</sup><sup>1</sup>1The error bars are from bootstrap realizations on 1500 lines of sight. In reality, these error bars may be slightly underestimated because the mean separation between 1500 lines of sight in the simulation we use is only $`60\mathrm{km}\mathrm{s}^1`$, comparable to the flux correlation length in the spectra (Paper I ). Obtaining better statistics of the theoretical prediction for $`B_C`$ would require a larger simulation. However, for the analysis in the present paper the errors in the determination of the temperature from the observations are much larger. The temperature $`T_P`$ is shown as the filled squares (the errors on the temperature are much smaller than those on $`B_C`$).
The triangles in Figure 7 show the effect on the derived cutoff, $`B_C`$, of reducing the histogram smoothing to $`\sigma _B=3000`$ K. There is no important difference with the crosses (the error bars for this smaller smoothing are similar to those on the crosses). By experimenting with different values of $`\sigma _B`$, we have found that $`B_C`$ is not affected if $`\sigma _B`$ is reduced, although the error bars obtained are significantly larger when reducing it below $`\sigma _B=3000\mathrm{km}\mathrm{s}^1`$. Increasing $`\sigma _B`$ beyond $`5000\mathrm{km}\mathrm{s}^1`$ leads to an increase of $`B_C`$, because the smoothing is then larger than the intrinsic sharpness of the cutoff in the $`B`$ distribution. We therefore adopt $`\sigma _B=5000`$K from this point forward in the paper.
The open squares show the effect of using all the absorption lines in every optical depth bin to compute the median temperature and optical depth, rather than using lines with $`|BB_C|<5000`$K only. There is a negligible difference in the derived temperature of the absorbers as a function of $`\tau _c`$; the reason is that, as we have previously seen, broader lines are by and large systems with higher fluid velocity dispersions, but their gas temperatures are not significantly greater, except at the highest optical depths where there is a slight difference (the systematic shift to the right of the open squares relative to the filled ones is due to a larger median optical depth of the broad lines within each optical depth bin). In the rest of the paper, we always compute the medians of any properties of the absorption lines using only lines with $`|BB_C|<5000`$K.
So far, we have seen that the Doppler parameter cutoff $`B_C`$ provides a good estimator for the gas temperature of absorption systems at a given optical depth. Our ultimate goal, however, is to measure the median gas temperature at a given gas density, for randomly selected points in the IGM, which we shall henceforth refer to as $`T_R`$. The peaks in absorption are at special locations, so their median temperature, $`T_P`$, will generally not be exactly the same as $`T_R`$. To examine this question, we first compute the median density of the absorption lines in each optical depth bin, and then we calculate the median temperature of randomly selected points at this gas density. The result is shown as the dashed line in Figure 7; the dotted lines give the 25 and 75 percentiles of the temperature distribution at random points with the same gas density. Comparing the dashed line to the filled squares, we see that the fitted absorption lines with relatively large optical depth typically have $`T_P<T_R`$, i.e., they are colder than random fluid elements at the same density. We believe the reason for this effect is the characteristic double-shock structure around the absorbers (Cen et al. 1994): the gas in the highest density tube along a filament (or the highest density surface along a sheet) is located between shocks, so it has been subject to less shock-heating than the surrounding gas.
Our method of analysis of the observational data in §6 will automatically correct for this difference between the temperature of the absorption lines and the temperature at random points. This systematic difference may introduce a potential uncertainty in the derivation of the gas temperature if it depends on quantities like the resolution of the simulation, the cosmological model that is assumed, or the heating at the reionization epoch. However, the temperature difference between lines and random points is negligible compared to the observational errors we will compute for the temperature in §6, at least in the simulation we analyze here.
#### 5.2.1 Testing With Other $`T\mathrm{\Delta }_g`$ Relations
We need to test the robustness of the finding that $`B_C(\tau _c)T_P(\tau _c)`$ for $`\tau _c1`$, which we have established so far in tests on the $`z=3`$ simulation output. We can change the $`T\mathrm{\Delta }_g`$ relation that we are measuring by simply using the $`z=2`$ or $`z=4`$ outputs from the simulation (see Table 1), still creating spectra with mean flux, noise level, and pixel size matching the $`\overline{z}=3`$ observations, as described in §2.2.1. In addition to the varying $`T\mathrm{\Delta }_g`$ relations, these spectra have different amplitudes of fluctuations, and different Hubble constants.
Figure 8 shows the comparison between $`B_C(\tau _c)`$ and $`T_P(\tau _c)`$, using the $`z=2`$ and $`z=4`$ outputs of the simulation in addition to $`z=3`$. The (pentagons, squares, triangles) show $`B_C`$ for the $`z=`$(4, 3, 2) simulation output, and the (solid, long-dashed, short-dashed) line show $`T_P(\tau _c)`$. We see that $`B_C`$ traces the temperature changes at the different redshifts extremely well, tracking the different slopes at $`z=2`$ and $`z=4`$ perfectly, and matching the increased overall temperature at $`z=3`$. All three redshift outputs show the same strong increase in $`B_C`$ above the actual temperature for optical depths below unity, corresponding to $`\mathrm{\Delta }_g1`$ in each case.
### 5.3 The $`B_CT`$ Relation at Different Redshifts
We have seen that $`B_C`$ traces $`T_P`$ remarkably well for the mean flux decrement and noise level of our observational data at $`\overline{z}=3`$. We shall now verify that this is also true when the flux decrement and noise is set to the values appropriate for the other two redshift bins in which we separate the data, at $`\overline{z}=2.41`$ and $`\overline{z}=3.89`$.
We first introduce a new type of figure that shows more clearly how accurately $`B_C(\tau _c)`$ traces the gas temperature.
Figure 9(a) shows $`B_CT_R`$, where $`T_R`$ is the median temperature of random points at the density of each optical depth bin (recall that this density is defined as the median density of absorption lines in each optical depth bin that have $`|BB_C|<5000`$K). The (pentagons, squares, triangles) are obtained from the simulation outputs at $`z=`$(4, 3, 2), and are the same points as in Figure 8 except that we have subtracted $`T_R`$. The temperature is correctly traced by $`B_C`$ at $`\tau _c1`$, and drops below $`B_C`$ at lower optical depths in the same way, independently of which simulation output we use. At $`\tau _c>10`$, $`B_C`$ falls significantly below the temperature in the $`z=2`$ output (this is seen clearly using more optical depth bins). The reason is that the temperature dispersion is higher at $`z=2`$ in the simulation, causing $`B_C`$ to reflect the lower cutoff of the true temperature distribution.
The results when we fix the mean flux decrement and noise to the $`\overline{z}=2.41`$ redshift bin of the observational data (taken from Paper I ) are shown in Figure 9(b). Here, we use the parameters $`E_d=9`$ and $`P_0=0.01`$ for the line fitting algorithm (see §3). The results are shown also using all three simulation outputs. Again, we find the temperature is well traced by $`B_C`$, but over a range of $`\tau _c`$ that has shifted to lower values. This is a result of the decreased optical depth at fixed gas density. Absorption systems that are primarily thermally broadened exist above a gas density that is approximately constant at each redshift, but the corresponding optical depth varies rapidly with redshift. At high optical depths, an additional effect is important in changing the degree to which $`B_C`$ traces the gas temperature: the increased shock-heating at low redshifts (with increasing velocities of collapse) implies a higher temperature dispersion, even at a fixed gas density. Therefore, $`B_C`$ drops further below the median gas temperature as the redshift decreases.
The result for $`\overline{z}=3.9`$ is shown in Figure 9(c) (we use $`E_d=8`$ at this redshift). The Doppler parameter cutoff ($`B_C`$) now traces the temperature only at high optical depth, $`\tau _c5`$, for exactly the same reason: a fixed gas density has shifted to a significantly higher optical depth due to the increase in the mean flux decrement.
### 5.4 Determining $`T_R(\mathrm{\Delta }_g)`$ Using $`B_C(\tau _c)`$
So far, we have seen that the Doppler parameter cutoff traces the gas temperature over a reasonable range of optical depth. We have shown this to be a consequence of the presence of some absorption lines that are primarily thermally broadened, and of the small dispersion of the gas temperature. We therefore expect that this relation between $`B_C`$ and $`T_R`$ will not be significantly changed depending on the model adopted in the simulation, or when the numerical resolution is increased.
However, even over a restricted range of optical depth where $`B_C`$ and $`T_R`$ are best matched, the simulation predicts a difference between them, which we want to correct for when analyzing the observational data (although this correction could be model dependent, and will need to be compared with other simulations in future work). In addition, we need to relate the central optical depth $`\tau _c`$ to the optical-depth-weighted gas density of the absorber, $`\mathrm{\Delta }_b`$, in order to derive the $`T\mathrm{\Delta }_g`$ relation of the gas from the observed $`B_C`$ as a function of $`\tau _c`$. Given the limited amount of data that we will analyze in this paper, it will be sufficient to parameterize the $`T\mathrm{\Delta }_g`$ relation by a power-law,
$$T=T_{}(\mathrm{\Delta }_g/\mathrm{\Delta }_{})^{\gamma 1}.$$
(6)
where $`\mathrm{\Delta }_{}`$ is chosen to make the error on $`T_{}`$ and $`\gamma 1`$ uncorrelated. As discussed in §2.2.2, this power-law should be understood only as a fit to the true relation, which should be more complex. The power-law form will be adequate here given our error bars, but larger sets of data might be used to detect deviations from a power-law.
In this subsection we develop the method to derive the parameters $`T_{}`$ and $`\gamma `$, given the determination of $`B_C`$ at different optical depths.
#### 5.4.1 Accounting for the Systematic Offsets $`B_CT_R`$
To obtain an accurate estimate of the gas temperature from observations, the systematic differences between $`B_C`$ and $`T_R`$ shown in Figure 9 can be used to correct the observed $`B_C`$. However, this correction may be dependent on the model, and this dependence will not be known until a wide variety of additional simulations are analyzed. We therefore use only absorption lines over the range of $`\tau _c`$ where the offset between $`B_C`$ and $`T_R`$ is small ($`|B_CT_R|<5000`$K for all three simulation outputs). The following optical depth ranges will be used: $`0.41<\tau _c<5.4`$ for $`\overline{z}=2.4`$, $`1.0<\tau _c<19`$ for $`\overline{z}=3.0`$, and $`3.8<\tau _c<47`$ for $`\overline{z}=3.9`$.
For each fitted line obtained from the observations (within the accepted optical depth range), we determine a temperature correction at the optical depth of the line by linearly interpolating from the two adjacent points in Figure 9. The set of points used depends on the redshift bin of the observations and the simulation output we choose for the analysis. The corrected line width is $`B^{}=B\mathrm{\Delta }T(\tau _c)`$, where $`B`$ is the observed line width, $`\mathrm{\Delta }T(\tau _c)=B_{C,S}(\tau _c)T_{R,S}(\tau _c)`$, $`B_{C,S}(\tau _c)`$ is the cutoff of the $`B`$ distribution in the simulation, and $`T_{R,S}(\tau _c)`$ is the median temperature in the simulation at random points with gas density equal to the median density of the absorption lines satisfying $`|BB_C(\tau _c)|<5000`$K.
#### 5.4.2 Translating $`\tau _c`$ into $`\mathrm{\Delta }_g`$
After the corrections just described, we have an estimate of the temperature as a function of $`\tau _c`$. What we need, however, is the temperature as a function of $`\mathrm{\Delta }_g`$. This will be derived by using a transformation from $`\tau _c`$ to $`\mathrm{\Delta }_g`$ that we obtain from the simulation. This introduces an inevitable model dependence in our measurement: the relation between $`\tau _c`$ and $`\mathrm{\Delta }_g`$, given a fixed flux decrement, should essentially be subject to the same uncertainties that appear in deriving the parameter $`\mu (\mathrm{\Omega }_bh^2)^2/\mathrm{\Gamma }/H(z)`$ (where $`\mathrm{\Gamma }`$ is the photoionization rate due to the cosmic background) from the observed mean flux decrement (see Rauch et al. 1997 , Paper I ).
Figure 10 is a scatter plot of the density and optical depth of the lines fitted from the $`z=3`$ simulation output (with mean flux decrement matching the $`z=3`$ observations). The crosses show $`\mathrm{\Delta }_g`$ vs. $`\tau _c`$ for all of the lines, while the filled squares show only lines satisfying $`|BB_C(\tau _c)|<5000`$K, where $`B_C(\tau _c)`$ was determined for 10 optical depth bins as described earlier. The lines that determine the $`B`$ cutoff (shown as squares) tend to have a higher optical depth than other lines at the same density, because of their lower velocity dispersion.
To obtain the $`\mathrm{\Delta }_g\tau _c`$ relation, we assign the median density of the lines in each optical depth bin that we use to measure the temperature (i.e., those with $`|BB_C(\tau _c)|<5000`$K) with the median $`\tau _c`$ of the same lines. We use interpolation to calculate the density corresponding to any value of $`\tau _c`$ for the fitted lines in the observed spectra.
One of the quantities affecting the $`\tau _c\mathrm{\Delta }_b`$ relationship is the density-temperature relation itself, essentially because the temperature affects the recombination coefficient, which then changes the neutral fraction at a given density. In order for our determination of the $`T\mathrm{\Delta }_g`$ relation to be self-consistent, we need to change the temperatures in the simulation so that they agree with the same $`T\mathrm{\Delta }_g`$ relation . We do this in the following way: after determining a preliminary $`T\mathrm{\Delta }_g`$ relation using the simulation with the true temperatures, we modify the temperature in every cell of the simulation using the formula
$$T_i^{}=T_iT_{}(\mathrm{\Delta }_i/\mathrm{\Delta }_{})^{\gamma 1}+T_{}^{}(\mathrm{\Delta }_i/\mathrm{\Delta }_{})^{(\gamma 1)^{}},$$
(7)
where $`T_i`$ is the original temperature at cell $`i`$, $`T_i^{}`$ is the modified temperature, $`T_{}`$ and $`\gamma 1`$ are the parameters of the original $`T\mathrm{\Delta }_g`$ relation of the simulation, and $`T_{}^{}`$ and $`(\gamma 1)^{}`$ are the parameters of the new $`T\mathrm{\Delta }_g`$ relation that we wish to impose. We then iterate the application of this formula until the modified $`T\mathrm{\Delta }_g`$ relation of the simulation matches the one from the observations. This modification of the temperatures in the simulation causes only a small change in the derived $`T\mathrm{\Delta }_g`$ relation (the value of $`T_0`$ is modified by only $`5`$%).
#### 5.4.3 Fitting for $`\gamma 1`$ Without Binning
We now describe the method we use to fit the parameters $`T_{}`$ and $`\gamma 1`$ to the values of $`\mathrm{\Delta }_g`$ and $`B^{}`$ of a set of fitted lines. The simplest method would be to separate the lines in density bins, measure the cutoff $`B_C`$ in every bin, and then fit the power-law relation to the values of $`B_C`$ obtained at every bin. However, the binning could result in a degradation of the measurement errors: at least 50 lines are needed to obtain a reasonable estimate of the cutoff $`B_C`$, and since our data yield only a few hundred lines for each redshift bin, the binning in line density would need to be very coarse.
There is a simple solution to this binning problem if the cutoff on the distribution of fitted lines, in the $`B^{}\mathrm{\Delta }_g`$ plane, is described by a power-law (recall that $`B^{}`$ is the width of each fitted profile minus the expected non-thermal broadening at its optical depth, as described in §5.4.1). After we have associated a gas density $`\mathrm{\Delta }_g`$ with each fitted line and corrected their temperatures using the simulation predictions, we rotate the absorption lines in the $`B^{}\mathrm{\Delta }_g`$ plane for many different assumed values of $`\gamma 1`$, using the formula
$$B^{\prime \prime }=B^{}T_{}(\mathrm{\Delta }_g/\mathrm{\Delta }_{})^{\gamma 1}.$$
(8)
For each assumed $`\gamma 1`$, we apply the cutoff determination technique to the $`B^{\prime \prime }`$ distribution (without density binning) to find a value for the temperature and a maximum value for $`dP/dB`$ \[see eq. (5)\]. As $`\gamma 1`$ is varied, the best fit value of $`\gamma 1`$ is the one that results in the maximum value of $`dP/dB`$, i.e., the sharpest cutoff on the $`B^{\prime \prime }`$ distribution.
Note that the value of $`T_{}`$ used in equation (8) affects the size of the temperature changes in the rotation. We therefore also iterate in the determination of $`T_{}`$ and $`\gamma `$ by this procedure. In practice, the measurement of $`T_{}`$ is barely affected by the rotation in equation (8), so a single iteration is sufficient.
Before presenting the results of applying our method to the observations, it will be useful at this point to summarize the full procedure we have described for measuring $`T_{}`$ and $`\gamma 1`$ from the $`B`$ distribution of the fitted absorption lines at every redshift bin. This consists of the following steps:
1. Eliminate the absorption lines with $`\tau _c`$ outside the range where the temperature measurement is expected to be effective.
2. Correct the values of $`B`$ for all of the remaining lines using the systematic offset predicted by the simulation (Figure 9).
3. Use the $`\tau _c\mathrm{\Delta }_g`$ relation in the simulation to associate a value of $`\mathrm{\Delta }_g`$ with each fitted line.
4. Determine an initial estimate for $`T_{}`$ and $`\gamma 1`$ with the method of fitting the power-law cutoff in the $`B`$ distribution as a function of $`\tau _c`$ that was just described.
5. Modify the $`T\mathrm{\Delta }_g`$ relation in the simulation to more closely approximate the relation measured from the observations.
6. Repeat steps 2-5 until the $`T\mathrm{\Delta }_g`$ relation measured from the observations matches the relation in the simulation.
## 6 ANALYSIS OF THE OBSERVED SPECTRA
We now apply our line-fitting method to the observed spectra. Table 3 lists, for each redshift bin, the redshift range ($`z_{min}`$ to $`z_{max}`$) and mean redshift $`\overline{z}`$, the mean flux decrement $`\overline{F}`$, the mean noise, the total number of pixels, the total path length, the value of $`E_d`$ we use to fit the lines, the total number of lines fitted, and the number of lines that we actually use to measure the gas temperature in the range of central optical depth from $`\tau _{min}`$ to $`\tau _{max}`$.
The parameters $`(B,\tau _c)`$ of all the fitted lines are shown in Figure 11(a,b,c) for the redshift bins $`\overline{z}=`$(3.9, 3.0, 2.4), as outlined crosses when the lines are not in any of the regions suspected to include metal lines, and as simple crosses when they are. For the quasar KP 77 (included in the redshift bin $`\overline{z}=2.4`$), the analysis to identify potential metal lines was not done, so all lines from this quasar are shown as outlined crosses. The lower cutoff on the $`B`$ distribution is clearly visible to the eye at all three redshifts, especially when the metal lines are ignored.
In Figure 12(a,b,c) we compare the $`B`$ histograms of the observed lines, within the ranges of optical depth that we will use for the temperature measurements, to the $`B`$ histograms of fitted lines from the simulation outputs with redshifts closest to the means of the observations, in the same optical depth ranges. The observed absorption lines obviously have higher temperatures than the simulated ones in all three cases.
It is interesting to compare directly $`B_C(\tau _c)`$ from the observations and from the simulation, before we determine $`T=T_{}(\mathrm{\Delta }_g/\mathrm{\Delta }_{})^{\gamma 1}`$ using the more involved method described in §5.4.
Figure 13(a,b,c) shows the values and errors of $`B_C`$ measured from the observations and the simulation output that is nearest in redshift, using lines over the optical depth bins indicated by the horizontal error bars (the optical depth bins contain equal numbers of observed lines). The vertical dotted lines indicate the optical depth range that we use for the final temperature measurement (for the $`\overline{z}=3.9`$ analysis, the upper limit on $`\tau _c`$ is outside of the figure, and eliminates a negligible number of lines). Recall (Figure 9) that we do not expect the points with lower $`\tau _c`$ to accurately reflect the real temperature, except at the lowest redshift. The observed lines again appear to be hotter than the simulated lines. These results are listed in Table 4, along with the temperature offset $`\mathrm{\Delta }T`$ used as a correction to the temperature (see §5.4.1) and the estimated gas densities, $`\mathrm{\Delta }_g`$, that we find once the $`T\mathrm{\Delta }_g`$ relation in the simulation has been adjusted to match the observed one (determined below). If more observed spectra were available, this method of binning in optical depth would be preferable to the method in §5.4 because it does not require the assumption of a power-law $`T\mathrm{\Delta }_g`$ relation . Each bin in optical depth would be associated with the density listed in Table 4, and the value of $`B_C`$ would be corrected by the listed $`\mathrm{\Delta }T`$.
We now determine $`T_{}`$ and $`\gamma 1`$ by the method described in §5.4. In order to obtain the $`\tau \mathrm{\Delta }_g`$ relation and the temperature offsets $`\mathrm{\Delta }T`$, we can use any of the three simulation outputs at $`z=2,3,4`$ for any of the three redshift bins in which we have divided the data (although we change the mean flux decrement of the simulated spectra to the observed one at each redshift bin). The different redshift outputs of the simulation are approximately equivalent to assuming different models with a different amplitude of the power spectrum (Paper I ), so we can use the different outputs to check that our measurement of the temperature is not greatly sensitive to the model that is assumed. Paper I found that the amplitude of the initial density perturbations in our simulation needs to be reduced by about 15% to agree with the observed power spectrum of the transmitted flux, meaning that the simulation output at $`z=4`$ has fluctuations that most closely match the observed ones at $`\overline{z}=3`$, and are slightly higher than the observed fluctuations at $`\overline{z}=3.9`$. The $`\overline{z}=2.4`$ observational bin is closest in amplitude to the simulation output at $`z=3`$. The most reliable temperature results should therefore be obtained by using the $`z=`$(4, 4, 3) simulation outputs to analyze the $`\overline{z}=`$(3.9, 3.0, 2.4) observations, but we shall also give results for $`\overline{z}=`$(3.0, 2.4) analyzed using the $`z=`$(3, 2) simulation outputs.
The result at $`\overline{z}=3`$, using the $`z=4`$ simulation output to predict the $`\tau \mathrm{\Delta }_g`$ relation and the necessary correction for non-thermal broadening, is:
$$T=[20200\pm 1300]\left(\frac{\mathrm{\Delta }_g}{1.37\pm 0.11}\right)^{[0.29\pm 0.30]}\mathrm{K},$$
(9)
where the error bars on $`T_{}`$ and $`\gamma 1`$ are uncorrelated (which defines $`\mathrm{\Delta }_{}`$). The error bar on $`\mathrm{\Delta }_{}`$ reflects only the uncertainty in the Paper I determination of the mean flux decrement, which affects the relation between density and optical depth. This result implies $`T_0=18400\pm 2100`$K, where $`T_0`$ is the temperature at the mean density. The pivot density, $`\mathrm{\Delta }_{}=1.37`$, corresponds to an optical depth $`\tau _{}=1.83`$ (from the relation obtained as described in §5.4.2) When we repeat the fitting using the $`z=3`$ simulation output we find $`T=[19600\pm 1500](\mathrm{\Delta }_g/[1.24\pm 0.10])^{[0.33\pm 0.28]}`$K, or $`T_0=18300\pm 1800`$K ($`\tau _{}=1.74`$). The difference between the two values for $`\mathrm{\Delta }_{}`$, 1.37 using the $`z=4`$ simulation, and 1.24 using the $`z=3`$ simulation, are mostly a reflection of the different optical depth normalizing factors (i.e., rescalings of the baryon density or the strength of the ionizing background) needed to match the observed mean flux decrement. The normalizing factor is smaller for the $`z=4`$ simulation, giving a larger $`\mathrm{\Delta }_{}`$ for the same optical depth, because the density fluctuations are of lower amplitude, leading to less saturated absorption and more absorption in voids (see Paper I for a more detailed discussion of the optical depth normalizing factor).
All the results obtained at the three redshift bins are listed in Table 5. The data analysis at $`\overline{z}=2.4`$ and $`\overline{z}=3.9`$ is similar to the analysis at $`\overline{z}=3`$, except for the differences that we mention below.
The temperature results at $`\overline{z}=2.4`$ differ by $`2000`$ K when the $`z=2`$ simulation output is used instead of the $`z=3`$ output, once the measured values of $`T_{}`$ for the two are extrapolated to the same density. Most of this difference results from the difference in $`\mathrm{\Delta }_{}`$, and the relatively high value of $`\gamma 1`$ that is obtained at $`\overline{z}=2.4`$. The difference in $`\mathrm{\Delta }_{}`$ is caused by the different amplitudes of density fluctuations in the two simulation outputs, so we expect that the temperature derived using the $`z=3`$ simulation output, which has the correct amplitude of fluctuations, is more reliable (see Paper I ).
In order to see the evolution of the temperature with redshift, we need to obtain the temperature at a fixed over-density at each redshift. It is useful to obtain the temperature at the mean density, $`T_0`$, to compare our results to other work. However, the values of $`\mathrm{\Delta }_{}`$ are close to $`\mathrm{\Delta }_g=1.4`$ at all three redshift bins, and we can therefore have a more robust result for the temperature evolution if we examine the temperature at $`\mathrm{\Delta }_g=1.4`$, which we denote as $`T_{1.4}`$ in Table 5.
We note here that, because of the small range of optical depth used at $`\overline{z}=3.9`$, at this redshift we were forced to smooth the $`B^{\prime \prime }`$ histograms with a Gaussian filter of width $`\sigma _B=7000`$K, instead of our standard $`\sigma _B=5000`$K, in order to avoid problems with multiple, approximately equivalent, maxima of equation (5) as $`\gamma 1`$ is varied.
The primary results of this paper, the measurements of $`T_{1.4}`$ and $`\gamma 1`$, are summarized in Figures 14 and 15, respectively. These show two important conclusions: first, the temperatures are higher than the value expected if photoionization heating in equilibrium is the only heating source. Second, we find no evidence for a rapid change of the temperature with redshift.
## 7 DISCUSSION
This paper presents a measurement of the temperature-density relation of the intergalactic gas in the redshift range $`2.4<z<4`$. The new method we have developed to perform this measurement is based on the same general idea as the previous work by Schaye et al. (1999), Ricotti et al. (2000), and Bryan & Machacek (2000): provided that there is a tight relation between the temperature and density of the gas, absorption lines of similar central optical depth should have little dispersion in their thermal broadening, and the varying line widths should correspond to variable amounts of hydrodynamic broadening. Occasionally, some absorption lines will be subject to only a small degree of hydrodynamic broadening; this will typically happen when most of the atomic hydrogen occurs near a velocity caustic along the line of sight. We therefore expect the histogram of line widths to show a rapid increase near the value of the Doppler parameter corresponding to the gas temperature. In the absence of noise, every line should be wider than the thermal broadening width, so at least an upper limit to the temperature can be obtained unambiguously.
The tests we have performed using a numerical simulation of the Ly$`\alpha `$ forest, based on a CDM model that successfully reproduces the observations of large-scale structure at present, confirm this general idea. However, they show that this method to recover the gas temperature works efficiently only over a limited range of line optical depths, which corresponds approximately to a range of gas over-density $`1\mathrm{\Delta }_g3`$. At lower densities, the gas is generally in Hubble expansion and this effect dominates the contribution to the line widths in essentially all the lines. The minimum line widths can therefore only provide an upper limit to the temperature of this low-density gas. Of course, the simulation can in principle be used to correct for the effect of line broadening due to Hubble expansion, and to obtain the temperature by subtracting the hydrodynamic contribution to the minimum line widths. However, the results can then strongly depend on the assumed model and the numerical resolution of the simulation, especially as the thermal broadening becomes a small effect compared to the expansion.
At very high densities, the increasing dispersion of the temperature at a given gas density can result in a large difference between the typical gas temperature and the “lower cutoff” in the line width histogram. This implies again that the recovery of the median gas temperature from the distribution of line widths is highly sensitive to the model assumptions that affect the temperature dispersion and the turbulent motions in the gas.
There are two main differences between the method we use to measure the gas temperature, and that used by previous authors. First, our line detection algorithm avoids the necessity of the Voigt-profile fitting method to “deblend” lines, by simply throwing out any “absorption lines” that do not correspond to a clearly identified minimum in the transmitted flux, or that do not have a large enough region around that minimum that is adequately fitted by a simple Gaussian in optical depth. In the method we use here, every line width is essentially a measurement of the second derivative around a minimum in the flux. This is one important reason why the total number of lines we identify is much lower than Schaye et al. (2000), even though we use nearly the same data set. The second difference is that we restrict the absorption lines we use to lie within the range of central optical depth where the correction that needs to be applied to the temperature measured from the histogram of line widths (as described in §5.4.1) is small.
These differences explain our substantially increased error bars in measuring $`T_0`$ and $`\gamma `$, relative to those of Schaye et al. (2000). However, we believe our error bars are more reliable and model-independent, for the reasons we have discussed. We also note here that, even though we cannot rule out a substantial change of the temperature we have measured depending on the numerical simulation of the Ly$`\alpha `$ forest that is used for comparing to the observational data, this possibility appears unlikely for several reasons, in addition to the arguments explained before about the small size of the correction $`\mathrm{\Delta }T`$ that needs to be applied to the line width cutoff $`B_c`$ to obtain the gas temperature (see §5.4.1). The gas temperature of the simulation we use is lower than that observed only by a small amount: in the simulation, the temperature at the mean density is $`T_0`$=(14000, 15800, 12800)K at $`z=`$(4, 3, 2), while our measurement from the observational data is $`T_0=(17400\pm 3900`$, $`18400\pm 2100`$, $`17400\pm 1900`$)K for $`\overline{z}=`$(3.9, 3.0, 2.4). The amplitude of the flux power spectrum of the simulation is also very close to the observationally determined one (Paper I ). We have also shown that we obtain nearly identical results for the temperature measurement when using different simulation outputs to analyze the same observations (see Figures 14 and 15) meaning that the dependence on the amplitude of the power spectrum is weak. Our result for the gas temperature might also be affected by the limited resolution of the simulation we use (with a comoving cell size of $`35`$ Kpc). We have not yet performed a convergence test for the effects of resolution on the Ly$`\alpha `$ forest; however, Schaye et al. (2000) find that a mean particle spacing of $`45`$ Kpc in their SPH simulations is sufficient for convergence.
We now compare our results for the evolution of the gas temperature with previous measurements. Ricotti et al. (2000) find a temperature at the mean density (from their Figure 12b) $`T_0`$($`18600_{6900}^{+10900}`$, $`23400_{5200}^{+10400}`$, $`17000_{9600}^{+22800}`$)K, at $`z`$ (3.6, 2.75, 1.9). Considering their large statistical error bars, our results appear to agree well with theirs, although the true temperature at $`z2.75`$ must be at the low end of their error bar. Schaye et al. (2000) give their temperature results as 16 separate points, in their Figure 5. We read off their points in each of our redshift bins and create error weighted averages for comparison with our results, finding $`T_0`$($`12200\pm 1700`$, $`17300\pm 1400`$, $`14000\pm 1300`$) at $`z`$(3.75, 3.2, 2.4). In order to compare these with our temperatures, which have lower errors at $`\mathrm{\Delta }_g=1.4`$, we extrapolate the Schaye et al. (2000) results to $`\mathrm{\Delta }_g=1.4`$ using their measured value of $`\gamma 1`$, obtaining $`T_{1.4}13600\pm 2000`$K at $`\overline{z}3.75`$. Our result at $`\overline{z}=3.9`$, $`T_{1.4}=20100\pm 2800`$K, is higher than theirs by $`1.9\sigma `$. Because of this, we do not find evidence for the increase of the temperature with time between $`z=4`$ and $`z=3`$ that Schaye et al. (2000) reported. Our results are consistent with a constant temperature.
Actually, the data set analyzed by Schaye et al. (2000) is almost identical to the one we analyze in this paper, with 7 of the 10 quasars used in the two papers being identical. However, our methods of analysis are very different, so that even the statistical error bars from the two analysis are largely independent. The main difference, as mentioned before, is in the number of lines that are identified. For example, we use 98 lines to measure the temperature in our $`z>3.4`$ bin, while Schaye et al. (2000) use about 550 lines in a comparable bin, even though $``$ 80% of the data that they use is identical to ours.
### 7.1 HeII Reionization as a Heating Mechanism
Is the value of the temperature we have measured in agreement with the known sources of heating and cooling in the intergalactic gas? The evolution of the temperature is determined by the equation:
$$\frac{d\mathrm{log}T}{Hdt}=2\left(1\frac{1}{3}\frac{d\mathrm{log}\mathrm{\Delta }}{Hdt}\right)+\frac{2}{3kHT}(L_{He}+L_HL_{CMB}L_RL_{ff}L_a),$$
(10)
where the cooling and heating rates per particle are denoted as follows: $`L_{He}`$ is the heating by He II photoionization, $`L_H`$ is the heating by H I photoionization, $`L_{CMB}`$ is the cooling off the microwave background, $`L_R`$ is the cooling by recombination, $`L_{ff}`$ is the cooling by free-free emission, and $`L_a`$ is the cooling due to line excitation and collisional ionization. We have separated the terms for He II and H I photoionization because He II plays a dominant role for heating, but the other cooling terms include both hydrogen and helium.
To see what the important terms for the thermal balance of the IGM are, we now evaluate all the heating and cooling terms at the conditions where we have measured the temperature most accurately: a temperature $`T=20000`$ K at a gas density $`\mathrm{\Delta }=1.4`$, at $`z=3`$. It is convenient to define the quantities $`T^{}2/(3kHT)L`$, for each subscript corresponding to every heating and cooling term. We assume the model $`\mathrm{\Omega }_bh^2=0.019`$, $`\mathrm{\Omega }_0=0.3`$, $`\mathrm{\Lambda }_0=0.7`$, and $`H_0=65\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$. We evaluate first the total cooling: the dominant term is adiabatic cooling, which is equal to $`2`$ on the right-hand-side of equation (10) if we assume a rate of expansion equal to the Hubble rate (i.e., a constant $`\mathrm{\Delta }`$). This assumption is of course not exact; every gas element expands at a different rate, causing a dispersion in the temperature-density relation. However, at a density $`\mathrm{\Delta }=1.4`$, the average rate of expansion is in fact not very different from the Hubble rate. In addition, when we consider the evolution of the temperature at a fixed $`\mathrm{\Delta }`$, the $`T/\mathrm{\Delta }`$ term in the total temperature derivative of the left-hand-side of equation (10) should partly compensate for the effect of expansion (canceling it exactly if $`\gamma 1=2/3`$).
The next most important contribution is cooling off the microwave background, which is given by $`T_{CMB}^{}=(8\sigma _TaT_{CMB}^4)/(3Hm_ec)(n_e/n)=0.58`$ (where $`\sigma _T`$ is the Thompson cross section, $`m_e`$ the electron mass, $`n_e`$ the electron density, and $`n`$ the total particle density). Notice that this term becomes more important at higher redshift, growing as $`(1+z)^{5/2}`$. We compute the other cooling rates using the formulae given in Black (1981). Recombination yields $`T_R^{}=0.22`$, and free-free emission $`T_{ff}^{}=0.08`$. The atomic processes of line cooling and collisional ionization are completely negligible at this low density, for the H I photoionization rate $`\mathrm{\Gamma }10^{12}\mathrm{s}^1`$ that is obtained from the observed abundance of quasars. The total cooling rate is therefore $`T^{}=2.88`$. If this temperature is being kept roughly constant, as indicated by the measurement we have presented here, then the heating terms should approximately balance the total cooling.
To evaluate the heating rate, we first assume ionization equilibrium; we will discuss later how the heating from He II can be increased if the He II reionization is still in progress. The heating term due to ionization can then be expressed in terms of the recombination rate: $`L_H=<E_H>(\alpha _Hn_e)(n_H/n)`$, and the analogous expression for He II , where $`\alpha _H`$ and $`n_H`$ are the recombination coefficient and the number density of hydrogen, and $`E_H`$ is the mean energy of the absorbed photons minus the ionization potential. The mean energy $`E_H`$ depends on the spectrum of the ionizing background, and we evaluate it as follows: assuming a background intensity per unit frequency $`J_\nu \nu ^\beta `$ from the ionization edge at $`\nu _0`$ to some maximum frequency $`\nu _m=q_m\nu _0`$, and approximating also the photoionization cross section as $`\sigma (\nu )\nu ^3`$, we find
$$<E_H>=I_H\left[\frac{1q_m^{\beta 2}}{1q_m^{\beta 3}}\frac{\beta +3}{\beta +2}1\right],$$
(11)
where $`I_H`$ is the ionization potential. We use $`\beta =0`$ and $`q_m=4`$ for both H I and He II , which adequately approximates the shape of the spectrum found in numerical calculations when the emitted spectrum is a quasar power-law (with $`\beta =1.5`$), and the effect of absorption by Lyman limit systems is taken into account (Miralda-Escudé & Ostriker 1990; Haardt & Madau 1996). This yields $`<E_H>=0.43I_H`$ and $`<E_{He}>=0.43I_{He}`$. We then obtain: $`T_H^{}=0.67`$, and $`T_{He}^{}=1.11`$. The total heating therefore falls short to compensate for cooling by a factor $`1.6`$.
An obvious way to increase the heating rate is to assume that the He II reionization is not yet complete; in other words, that there are patches of low-density gas in the IGM where all the helium is in the form of He II . As discussed in Miralda-Escudé & Rees (1994), there are two reasons why the heating rate is higher during the reionization, relative to the case of photoionization equilibrium. The first is that the ionization rate needs to be higher simply because every He II ion needs to be ionized once during the course of the He II reionization. The second is that all the hard photons will now be absorbed by a random He II ion in the IGM, up to the frequency $`\nu _m=q_m\nu _0`$ where the mean-free-path through the He II IGM reaches the horizon length. For the baryon density we use and at $`z=3`$, and assuming also that about 50% of all the helium is in the form of He II in the diffuse IGM (and not in dense clouds having a small covering factor over the Hubble length), this maximum frequency is given by $`q_m=13`$, i.e., a frequency 13 times higher than the He II ionization edge. The mean energy of the absorbed photons is therefore equal to the mean energy of the emitted photons up to this maximum frequency, without weighting them with the photoionization cross section:
$$<E_{He,r}>=I_{He}\left[\frac{1q_m^{\beta ^{}+1}}{1q_m^\beta ^{}}\frac{\beta ^{}}{\beta ^{}1}1\right],$$
(12)
where the subscript $`r`$ in $`<E_{He,r}>`$ indicates the mean energy per absorption during reionization, and the emitted spectrum from sources is $`J_\nu \nu ^\beta ^{}`$. For $`q_m=13`$ and $`\beta ^{}=1.5`$, we obtain $`<E_{He,r}>=0.71I_{He}`$. Since the recombination rate for He II at the mean density is equal to $`1.6`$ times the Hubble rate $`H^1`$ at $`z=3`$, and if the reionization is occurring over $``$ a Hubble time near $`z3`$, it is reasonable to expect that the additional heating rate due to He II reionization is comparable to the heating rate due to balancing recombinations of He II . We therefore conclude that the heating from He II reionization can reasonably account for the IGM temperature we have determined here.
### 7.2 Usefulness for Measuring the Baryon Density from the Lyman Alpha Forest
One of the applications that the development of the new theory of the Ly$`\alpha `$ forest based on structure formation has had is to provide a measurement of the baryon density through its effect on the mean transmitted flux. For a fixed distribution of temperature, over-density, and peculiar velocities in the IGM, the Ly$`\alpha `$ optical depth at any point in the spectrum is proportional to $`\left(\mathrm{\Omega }_Bh^2\right)^2H(z)^1\mathrm{\Gamma }_{12}^1`$, where $`\mathrm{\Gamma }=10^{12}\mathrm{\Gamma }_{12}\mathrm{s}^1`$ is the photoionization rate due to the cosmic ionizing background. To be specific, we define the parameter
$$\omega _B\mathrm{\Omega }_Bh^2\left[\left(\frac{\mathrm{\Omega }_0h^2}{0.3\times 0.65^2}\right)^{1/2}\mathrm{\Gamma }_{12}\right]^{1/2},$$
(13)
\[where we have used $`H(z)H_0(1+z)^{3/2}\mathrm{\Omega }_0^{1/2}`$, which is highly accurate at the relevant redshifts and in a flat universe\]. As discussed by Hernquist et al. (1996), Miralda-Escudé et al. (1996), Rauch et al. (1997), Weinberg et al. (1997), and Paper I , a measurement of $`\omega _B`$ can be translated into a lower bound on $`\mathrm{\Omega }_bh^2`$ by using the contribution to the ionizing background from known quasars as a lower bound on $`\mathrm{\Gamma }_{12}`$. This lower limit is on the high side of the range of $`\mathrm{\Omega }_b`$ that is allowed by primordial nucleosynthesis: $`\mathrm{\Omega }_bh^20.02`$ (Rauch et al. 1997 ; Paper I ).
One of the main model uncertainties in deriving the relationship between the parameter $`\omega _B`$ and the predicted mean transmitted flux is the mean IGM temperature. A higher temperature implies a lower recombination coefficient, and therefore a lower neutral hydrogen density. This implies that in order to reproduce a given observed mean transmitted flux, the mean density $`\mathrm{\Omega }_B`$ needs to be increased further to compensate the reduced recombination coefficient. The measurements of the IGM temperature reported here and in Schaye et al. (2000), and Ricotti et al. (2000), all coincide in finding temperatures that are high compared to what is expected if the IGM is heated by photoionization and is in ionization equilibrium. As we have discussed, these higher temperatures can probably be understood as a result of the He II reionization. independently of its cause, the higher temperature implies an even higher value of $`\omega _B`$ than was obtained previously, which we can easily determine by modifying the temperature in the simulation to match the observations, as described earlier in §5.3. We find that, when the temperature in the simulation is increased to match the observed one, the derived value of $`\omega _B`$ is increased slightly to $`\omega _B=`$($`0.0336\pm 0.0020`$, $`0.0288\pm 0.0023`$, $`0.0248\pm 0.0017`$) at $`\overline{z}=`$(3.9, 3.0, 2.4), from the previous values (0.0329, 0.0274, 0.0245) when the temperatures in the simulations are not modified (Paper I ). The errors are derived from the observational errors in the determination of the mean flux decrement from Paper I .
Changing the temperature of the simulation affects the value of $`\omega _B`$ not only by modifying the recombination coefficient, but by increasing the amount of thermal broadening, which can spread the absorption in saturated regions to the outskirts of absorption lines, increasing the mean absorption for a given $`\omega _B`$. We find that the effect of thermal broadening is less important than the effect of the reduced recombination coefficient. As an example, if we replace $`\alpha (T)\alpha (T+3000K)`$ in every pixel in the simulation, the inferred $`\omega _B`$ increases by 0.0017 (where $`\alpha (T)`$ is the recombination coefficient), while the replacement $`\sigma _b(T)\sigma _b(T+3000K)`$ changes $`\omega _B`$ by -0.00046 (where $`\sigma _b(T)`$ is the dispersion of the Gaussian thermal broadening). Dynamical effects caused by the increased pressure of the gas would probably go in the same direction as the thermal broadening, since they would tend to spread the gas in absorption systems over wider regions. However, it seems unlikely that any such dynamical effects (which can only be investigated by running the same simulation with different temperatures) can be more important than the thermal broadening effect. The examples shown in Theuns et al. (1999b) (see their Figure 6) appear to confirm that the dynamical effects of increased pressure are not more important than the increased thermal broadening when the gas temperature is raised.
With the statistical error bars we have obtained on the $`T\mathrm{\Delta }_g`$ relation , we can place more conservative lower bounds on $`\omega _B`$ than the ones obtained in Paper I . The lowest allowed value of $`\omega _B`$ needed to account for a given observed mean transmitted flux is obtained when $`T_{}`$ is minimum and $`\gamma 1`$ is maximum in equation (1), because that yields the minimum temperature for the low-density gas that determines the optical depth in unsaturated regions of the Ly$`\alpha `$ spectrum. We set $`T_{}`$ equal to the measured value minus twice the statistical error bar given in Table 5, and $`\gamma 1=0.6`$, which is the value valid when the IGM has been in photoionization equilibrium for a long time (Hui & Gnedin 1997). Any uniform heating of the IGM, such as that caused by reionization, should give rise to a lower $`\gamma 1`$, although shock-heating can increase $`\gamma 1`$ above $`0.6`$, the simulations show that this happens only at high enough gas densities that the Ly$`\alpha `$ absorption is already saturated. The error in our observational determination of $`\gamma 1`$ is too large to give us a better constraint than $`\gamma <0.6`$ (see Table 5).
The results of this exercise are $`\omega _B>`$(0.0270, 0.0192, 0.0209 ) for $`\overline{z}=`$(3.9, 3.0, 2.4), at 95% confidence, including the error from the mean flux decrement and the temperature measurements (added in quadrature). Using the lower bound obtained in Rauch et al. (1997) of $`\mathrm{\Gamma }12>0.7`$ in the range $`2<z<3`$, obtained by counting only radiation from the observed quasars, and not including the power-law extrapolation of the quasar luminosity function that has been observed only at redshifts $`z<2`$, we obtain $`\mathrm{\Omega }_Bh^2>0.017`$. This result is still consistent with the determinations of the deuterium abundance (Burles & Tytler 1998). However, if the quasar luminosity function extends to low luminosities with a similar power-law slope as observed at $`z<2`$, or if emission from galaxies increases significantly the intensity of the ionizing background, then the higher baryon density implied would come into conflict with the primordial nucleosynthesis predictions and the observed deuterium abundance.
In summary, we have reached the following conclusions:
1. The temperature of the IGM is $`20000\pm 2000`$K at density 1.4 times the mean, independent of redshift, although an increase of $`3500`$K from $`z=3.9`$ to $`z=3.0`$ cannot be ruled out.
2. The high temperature cannot be explained by heating in ionization equilibrium, and probably indicates on-going He II reionization.
3. The contribution of temperature uncertainty to the uncertainty in the baryon density required by the observed mean flux decrement in the Ly$`\alpha `$ forest is now well constrained.
We thank Adam Steed and David Weinberg for helpful comments on the manuscript.
## Appendix A THE PROFILE FITTER
This algorithm has three input parameters that control how the fitting proceeds: $`E_d`$ sets the amount that the flux must increase from the center point to the edges of the window before a fit will be attempted, $`W_{min}`$ sets the minimum size of the window within which a fit is performed, and $`P_0`$ sets the quality of fit that will be accepted. In this paper we set $`W_{min}=2`$.
Three more input parameters effect the speed of the code but are not important to the results: $`E_s`$ controls the degree of symmetry around a central pixel that is required for a fit to be attempted, $`E_c`$ sets the level of flux decrease, from the center pixel to the window edges, at which a point will be eliminated from consideration for fitting, and $`W_{max}`$ sets the maximum allowed window size. These parameters are set to values large enough that they do not actually eliminate any profiles that would otherwise be accepted.
Before we describe the algorithm in detail, a few more terms must be introduced: We are going to fit pieces of the spectrum that have center point $`P`$ and extend $`\pm W`$ pixels to either side of $`P`$. The width of the fitting window, $`W`$, will be adjustable but constrained to $`W_{min}WW_{max}`$. The transmitted flux at a point $`P`$ is $`F(P)`$. The error in the sum or difference of the flux at two points $`P_1`$ and $`P_2`$ is $`\sigma (P_1,P_2)=[\sigma (P_1)^2+\sigma (P_2)^2]^{1/2}`$, where $`\sigma (P)`$ is the observational error in the flux at point $`P`$. The minimum acceptable probability for $`\chi ^2`$ is $`P_0`$ (we need to define $`P_0`$ by the probability because there will be varying numbers of degrees of freedom in the fits).
For a given spectrum the algorithm that we use is the following (the reader should keep in mind that, except for the added complication of setting the window position and width, this procedure just fits a single Voigt profile to each absorption maximum by $`\chi ^2`$ minimization):
1. Scan along the spectrum pixel by pixel searching for places where $`|F(PW_{min})F(P+W_{min})|<E_s\sigma (P+W_{min},PW_{min})`$. Also require that $`F(P)F(P\pm W_{min})<E_c\sigma (P,P\pm W_{min})`$, where the flux at the $`\pm W_{min}`$ points is averaged. These places are candidates for a symmetric, non-concave profile.
2. If there is a significant increase in flux at the edges of the window, so that $`F(P)F(P\pm W_{min})>E_d\sigma (P,P\pm W_{min})`$, go ahead and fit Equation (4) to the absorption. If there isn’t a significant increase try to expand the window.
3. To expand the window require that symmetry is maintained when $`W`$ is increased, i.e., $`|F(PW)F(P+W)|<E_s\sigma (P+W,PW)`$. If the window can be expanded return to step 2 to check if a fit can be done with the enlarged window, i.e., if $`F(P)F(P\pm W)>E_d\sigma (P,P\pm W)`$.
4. If the region can’t be fit, but also can’t be expanded, eliminate the candidate point. Also eliminate the point if the window size has been increased to $`W_{max}`$ without meeting the requirement for fitting.
5. Set initial parameters for the fit using $`F(P)`$ to set $`\tau _c`$ and $`[F(P+W)+F(PW)]/2`$ to set $`\sigma _b`$. Set $`v_c=0`$. If $`F(P)<0`$ set $`\tau _c=10`$.
6. Minimize $`\chi ^2`$ using the flux values and their error bars in the range of points between $`P+W`$ and $`PW`$. Require that $`|v_c|<0.5`$ pixels (outside this range is covered by other candidate points).
7. Eliminate candidate if $`P(>\chi ^2,\nu )<P_0`$.
8. If $`P`$ falls within $`W`$ of a previously accepted candidate, eliminate the candidate with a smaller value of $`P(>\chi ^2,\nu )`$. This does not eliminate any independent profiles because a single Gaussian would not fit if the window contained multiple lines.
|
warning/0005/hep-th0005093.html
|
ar5iv
|
text
|
# References
String-Network type BPS states have been analyzed in the past few years in string theory (with or without branes), as well as in other quantum field theories \- . On the other hand, superconducting strings have been studied in field theories as they are expected to play important role in the evolution of our universe at early time. Construction of such configurations in string theory requires coupling of macroscopic strings to electromagnetic fields. Electromagnetically charged macroscopic strings are known to exist in heterotic string theory for some time . In a recent paper , these were extended to the type IIB string theories. Moreover, it was also shown that BPS networks of such objects with 1/4 supersymmetry can also be constructed.
In this paper we analyze various aspects of such electrically charged string networks. We first show that in the above construction, the electric-currents flowing through the strings are conserved on the 3-string junctions. In the absence of complete supergravity solutions for such networks, these are done by examining the current flowing through them far away from any of the 3-string junction. Thus, although a redistribution of the current flow as well those of electrostatic charges are expected to take place in the complete solution near the junction, we observe that the electric-currents emerging far away from these junctions are conserved. We then argue the consistency of charged string from the point of view of world volume theories as well. We discuss this aspect for charged-string networks which are constructed by applying a Lorentz transformation involving time and an internal direction. It has earlier been pointed out that configurations obtained by applying above transformations represent genuinely charged string network. In particular, by comparing the world-volume and supergravity expression for the energy-densities, we show that indeed a fundamental (charged) string ends on a D-string to form a 3-string junction.
Finally, we present a new charged string solution in 8-dimensional type II theories. These solutions are obtained by observing a mapping between the heterotic and type II supergravity actions, after suitable truncations. Since type II classical solutions also map on to the corresponding solutions in the heterotic theory, the BPS nature of type II solutions is guaranteed. However, it is possible that they preserve different set of supersymmetries than the ones obtained by another mapping between the truncated type IIB and heterotic theory, which was used in . To show the difference between the solutions presented in this paper and that of , we notice that in the charged macroscopic string solutions of and , the charges are acquired by fields which are identified as Kaluza-Klein (KK) gauge fields: $`G_{\mu i}\pm B_{\mu i}`$, for $`i=1,2`$, representing the internal directions. In the present case however, charges are assigned within an $`SL(3)`$ multiplet of gauge fields. Since two such multiplets are formed by combinations $`(G_{\mu 1},B_{\mu 2}^{NS},B_{\mu 2}^{RR})`$ and $`(G_{\mu 2},B_{\mu 1}^{NS},B_{\mu 1}^{RR})`$, the solutions presented in this paper are necessarily different from the ones in . Furthermore, we discuss charge and current conservations around these junctions.
We now start with the discussion of the electric current conservations. Let us consider a charged string solution in D-dimensions which is given by a supergravity configuration with non-zero 2-form $`B_{\mu \nu }^a`$ and 1-form $`A_\mu ^I`$. Only nonzero components of $`B_{\mu \nu }^a(r)`$ are $`B_{0,(D1)}^a(r)`$ and that of gauge field $`A_\mu ^I`$ are $`A_0^I(r)`$ and $`A_{D1}^I(r)`$. Here $`r`$ denotes the radial coordinate in $`(D2)`$ dimensions transverse to the string and the superscripts $`a,I`$ on $`B_{\mu \nu }`$ and $`A_\mu `$ distinguish between various two-form and one-form fields respectively. Let us first focus on the 2-form fields. The charges associated to the fields are defined by
$$Z^a=𝑑\mathrm{\Omega }{}_{}{}^{}H_{\mathrm{\Omega }}^{a}.$$
(1)
Here $`H^a=dB^a`$ and $``$ denotes the Hodge dual. The integral above is taken over the $`(D2)`$ transverse directions of the string.
Now let us think of a 3-string junction with its three legs coupling to different values of $`B^a`$’s, denoted as $`B_1^a`$, $`B_2^a`$ and $`B_3^a`$. The corresponding charges are denoted as $`Z_1^a`$, $`Z_2^a`$ and $`Z_3^a`$. Then, using above identifications, the charge conservation: $`Z_{(1)}^a+Z_{(2)}^a=Z_{(3)}^a`$ implies
$$_{\mathrm{}}𝑑\mathrm{\Omega }_{(1)}{}_{}{}^{}H_{\mathrm{\Omega }_{(1)}}^{a}+_{\mathrm{}}𝑑\mathrm{\Omega }_{(2)}{}_{}{}^{}H_{\mathrm{\Omega }_{(2)}}^{a}=_{\mathrm{}}𝑑\mathrm{\Omega }_{(3)}{}_{}{}^{}H_{\mathrm{\Omega }_{(3)}}^{a},$$
(2)
where we have now introduced subscripts $`\mathrm{\Omega }_{(i)}`$’s for the angular variable in the transverse space. Also, the subscripts of the integrals imply that they are evaluated far away from the junction. In other words, the charge-conservation condition of the above type follows by evaluating the expression (1) along any one of the string of a 3-string junction (far away from it) and then by sliding the large spherical surface through the junction to the other side while deforming it in a manner to surround the remaining two strings.
Same statements are true for the case of 1-form charges as well, when one considers the charges corresponding to non-zero $`A_{(D1)}^I`$’s. The field-strengths corresponding to $`A_{(D1)}^I`$ are given by $`F_{(D1)r}^I`$. The corresponding charges are given as:
$$J_0^I=𝑑\mathrm{\Omega }{}_{}{}^{}F_{0\mathrm{\Omega }}^{I}$$
(3)
where we have now kept the ‘time’-index on the charge to differentiate them with other changes defined below. The charge conservation (which finally amounts to the electric-current conservation) then implies the condition:
$$J_{0}^{I}{}_{(1)}{}^{}+J_{0}^{I}{}_{(2)}{}^{}=J_{0}^{I}{}_{(3)}{}^{}$$
(4)
As in the case of 2-form charge conservation condition (2), eqn. (4) now follows from the following equation:
$$_{\mathrm{}}𝑑\mathrm{\Omega }_{(1)}{}_{}{}^{}F_{0\mathrm{\Omega }_{(1)}}^{I}+_{\mathrm{}}𝑑\mathrm{\Omega }_{(2)}{}_{}{}^{}F_{0\mathrm{\Omega }_{(2)}}^{I}=_{\mathrm{}}𝑑\mathrm{\Omega }_{(3)}{}_{}{}^{}F_{0\mathrm{\Omega }_{(3)}}^{I}.$$
(5)
Finally we discuss the case of gauge charges associated with 1-form components $`A_0^I`$’s. Now the nonzero field-strengths are: $`F_{0r}^I`$ and their Hodge-duals are given as: $`{}_{}{}^{}F_{(D1)\mathrm{\Omega }}^{I}`$. The corresponding charges are now given as:
$$q_{(D1)}^I=𝑑\mathrm{\Omega }{}_{}{}^{}F_{(D1)\mathrm{\Omega }}^{I}.$$
(6)
An important difference between two 1-form charges defined in eqns. (3) and (6) is that the later one depends on the direction along which the string lies as this charge is being measured by the value of a field strength at a large orthogonal distance from the string. The corresponding consistency condition of the charged 3-string junction is the force balance condition, depending on their orientations. We now discuss these aspects for examples presented earlier in .
We start by analyzing the 3-string junctions of charged macroscopic strings in $`D=9`$, discussed in section-(3.1) of . These are parameterized by a single solution-generating parameter $`\alpha `$. Moreover, its action can be identified in ten-dimensions simply as a Lorentz-transformation involving time-coordinate $`x^0`$ and and an internal direction: $`x^9`$. The consistency of the network of such charged-string solutions is already known. Explicit solution for the electrically charged fundamental string ( henceforth $`(1,0)`$ string ), appearing in the networks, is presented in section-3.1 of . We only write down the 1-form potentials:
$$\widehat{A}_t^1=\frac{C\mathrm{sinh}\alpha \mathrm{cosh}\alpha }{2(r^5+C\mathrm{cosh}^2\alpha )},\widehat{A}_8^1=0,$$
(7)
$$\widehat{A}_t^2=0,\widehat{A}_8^2=\frac{C\mathrm{sinh}\alpha }{2(r^5+C)}.$$
(8)
As we notice, the above $`(1,0)`$-string solution is characterized by two gauge fields for this $`D=9`$ example. They come from KK reduction of the ten-dimensional metric and antisymmetric tensor fields. Then the $`SL(2,Z)`$ duality in ten-dimensions generates, one more gauge field identified as the one coming from the KK reduction of the ten dimensional RR sector antisymmetric tensor. Nonzero gauge field components in the final $`(p,q)`$-string solution can then be written as:
$$\widehat{A}_t^1=\frac{C\mathrm{\Delta }_q^{1/2}\mathrm{sinh}\alpha \mathrm{cosh}\alpha }{2(r^5+C\mathrm{\Delta }_q^{1/2}\mathrm{cosh}^2\alpha )},\widehat{A}_8^1=0,$$
(9)
$$\widehat{A}_t^2=0,\widehat{A}_{i8}^2=\frac{C\mathrm{sinh}\alpha }{2(r^5+C\mathrm{\Delta }_q^{1/2})}(M_0^1)_{ij}q_j,i=1,2,$$
(10)
where $`i=1,2`$ stand for the $`B_{\mu \nu }`$’s in NS-NS and R-R sectors of type IIB in ten dimension and $`M`$ is a $`2\times 2`$ matrix parameterizing $`SL(2)/SO(2)`$ moduli. In the above equations, $`\mathrm{\Delta }_q=q_i(M_0^1)_{ij}q_j`$. By denoting the currents originating from the metric and the antisymmetric tensors respectively as $`J^1`$ and $`J^2`$, the electric current in a $`(p,q)`$ string is:
$$J^1=0,J_i^2=C\frac{\mathrm{sinh}\alpha }{2}(M_0^1)_{ij}q_j,i=1,2.$$
(11)
We now observe that, apart from an $`O(d,d)`$ factor $`\mathrm{sinh}\alpha /2`$, the above electric currents are proportional to the 2-form charges of the strings. As a result, the electric-current conservations hold at the junctions directly due to the conservation of 2-form charges.
We now discuss the electric-current conservation in string networks constructed by starting with the $`(1,0)`$ charged-macroscopic strings which are T-dual to the ones we mentioned in the last paragraph. An analysis of the supersymmetry condition in this case implies that such solutions are possible in $`D=8`$. One then has six gauge fields, two each from the KK components of $`G_{\mu \nu }`$, $`B_{\mu \nu }^{NS}`$ and $`B_{\mu \nu }^R`$. We denote the corresponding electric-currents as $`J_a^1,J_a^{2NS},J_a^{2R}`$, $`a=1,2`$. Then the results of section-(3.2) in imply the following values of the electric-currents for the $`(1,0)`$ string solution (in a complex notation):
$$J_1^1+iJ_2^1=\frac{\mathrm{\Delta }_q^{1/2}}{2}\mathrm{sinh}\alpha e^{i\theta },J_1^{2i}+iJ_2^{2i}=0$$
(12)
where an extra superscript $`i`$ in $`J^2`$, in the second equation above, stands for $`NSNS`$ and $`RR`$ components and the parameter $`\theta `$ is an angular $`O(2,2)`$ parameter, identified as a spatial rotation among the compactified coordinates $`x^8`$ and $`x^9`$. However the network construction requires this parameter to be identical to the one associated with an $`SL(2,Z)`$ transformation, namely eqn.(16) of , which generates a $`(p,q)`$-string solution from a $`(1,0)`$ one.
Then for 3-string junctions, nonzero electric-current along a $`(p,q)`$-string-prong is given by:
$$J_1^1+iJ_2^1=\frac{1}{2}\mathrm{sinh}\alpha e^{\varphi /2}(pq\tau )$$
(13)
We then once again see the electric-current conservations in the networks of such strings, following directly from the 2-form charge conservations.
Unlike electric-currents, there is no direct way to examine the status of the consistency of (static) electric-charge densities coming from the gauge field components $`A_t^1`$ in (9), as these are expected to be redistributed in the full supergravity solution near a 3-string junction. One way to settle the issue is to analyze supergravity solutions of such 3-string junctions. In the absence of these solutions, for the moment, we note that the string tension of the charged macroscopic $`(p,q)`$-strings discussed above in eqns.(7-11) and (12-13) are given by
$$T_{p,q}=\mathrm{\Delta }_q^{1/2}C(\mathrm{cosh}\alpha +1)/2.$$
(14)
Since the $`O(d,d)`$ parameter is only an overall factor, the tension balance continues to hold for both kinds of string-networks mentioned above. This is because the spatial orientations of these strings in a network are identical to the one for the neutral ones.
In view of applications in our later analysis , we now write down the induced world-sheet energy momentum tensors corresponding to general charged F-string solution. They are given in terms of $`O(d,d)`$ parameters $`\alpha ,\beta `$:
$`T_{00}`$ $`=`$ $`C\mathrm{cosh}\alpha \mathrm{cosh}\beta ,`$
$`T_{11}`$ $`=`$ $`C,`$
$`T_{01}`$ $`=`$ $`{\displaystyle \frac{C}{2}}(\mathrm{cosh}\alpha \mathrm{cosh}\beta ).`$ (15)
To get the corresponding answer for the examples considered in (7-11) and (12-13), we have to set $`\beta =\pm \alpha `$. Therefore, for these configurations, the $`(00)`$ component of world-sheet stress tensor reduces to
$$(T_{00})_{p,q}=\frac{1}{2}C\mathrm{\Delta }_{q}^{}{}_{}{}^{\frac{1}{2}}(1+\mathrm{cosh2}\alpha ).$$
(16)
This should not, however, be compared with the tension calculated in (14) as $`(T_{00})_{p,q}`$ receives contributions from string tension as well as from the gauge fields associated with electrically charged strings. However we note that for $`\alpha =0`$, the expression (16) matches with the string tension, as expected in the neutral case.
We now discuss the consistency of the charged string solutions from the point of view of the D-string world-sheet theory. In this context, we consider the example of the $`D=9`$ charged string networks formed by a Lorentz-boost, parameterized by the parameter $`\alpha `$ mentioned above. Then, a classical solution representing a 3-string junction of charged macroscopic strings in the world-sheet theory, is given by the application of the Lorentz-boost on the solution of and can be written as in eqns.(17) and (19) below:
$`A_0`$ $`=`$ $`gx^1\mathrm{cosh}\alpha ,\mathrm{\Phi }A_9=gx^1\mathrm{sinh}\alpha ,x^1>0`$ (17)
$`=`$ $`0=0,x^1<0,`$
where $`\mathrm{\Phi }`$ is the scalar coming from the dimensional reduction of the world-volume gauge field from $`D=10`$ to $`D=9`$. The choice of the Lorentz-transformation parameter ‘$`\alpha `$’ is fixed through the results in section-(3.1), in particular (3.13) of .
To maintain supersymmetry one has to excite one more world-volume field, identified as the coordinate representing the F-string. Following ,, in this case we have
$`X^8`$ $`=gx^1,x^1>0,`$ (18)
$`=0,x^1<0.`$ (19)
This is a $`1/2`$ supersymmetric solution in the world-volume theory. The supersymmetry condition for this solution is obtained from that of the neutral string by the above Lorentz-transformation.
We now evaluate the energy of this configuration to identify it with the expression for $`T_{00}`$ of the F-string given in (16). In order to proceed, following , we first write down the expression of the Hamiltonian associated with the above configuration. After evaluating the expressions, one gets:
$$H=\frac{1}{2}(1+\mathrm{cosh2}\alpha )H_0,$$
(20)
where $`H_0`$ is the Hamiltonian associated with the neutral 3-string junction of . The first term inside the bracket in the right hand side corresponds to the contribution of $`X^8`$ to the classical action whereas the second term is the combined contribution from $`A_0`$ and $`X^9`$. As a result, the energy expression is modified by a factor $`(1+\mathrm{cosh2}\alpha )/2`$, which precisely coincides with the energy density of the charged string given in (16), when retricted to $`(1,0)`$-string. We have therefore given a world-volume argument in favor of the existence of the 3-string junction solutions of charged macroscopic strings by identifying the relevant variables in the two approaches, namely $`T_{00}`$ in (16) and $`H`$ in (20). In the limit $`\alpha =0`$, one also reproduces: $`H_0=T_fx^8`$, a result following from the analysis of , with $`T_f`$ being the F-string tension. In other words, in , the world-volume energy was associated with the string tension of a ‘spike’-configuration interpreted as an F-string. We observe that similar interpretation holds in the case of charged strings as well, provided one takes into account the contribution of the charges in the supergravity solution.
We now give explicit construction for some new charged string configuration in eight dimensions. We further discuss how various conservation laws are satisfied around the junction for such strings.
Following , we start with a truncated version of eight dimensional type IIB action. In Einstein frame the action can be written as
$`S={\displaystyle }d^8x[`$ $`R{\displaystyle \frac{1}{2}}\{(\sigma )^2+(\varphi _1)^2+(\varphi _2)^2\}`$ (21)
$`{\displaystyle \frac{1}{12}}\{e^{\varphi _1+\frac{1}{\sqrt{3}}\varphi _2}H_{3}^{}{}_{}{}^{(1)}{}_{}{}^{2}+e^{\varphi _1+\frac{1}{\sqrt{3}}\varphi _2}H_{3}^{}{}_{}{}^{(2)}{}_{}{}^{2}+e^{\frac{2}{\sqrt{3}}\varphi _2}H_{3}^{}{}_{}{}^{(3)}{}_{}{}^{2}\}`$
$`{\displaystyle \frac{1}{4}}e^\sigma \{e^{\varphi _1\frac{1}{\sqrt{3}}\varphi _2}F_{2}^{}{}_{}{}^{(1)}{}_{}{}^{2}+e^{\varphi _1\frac{1}{\sqrt{3}}\varphi _2}F_{2}^{}{}_{}{}^{(2)}{}_{}{}^{2}+e^{\frac{2}{\sqrt{3}}\varphi _2}F_{2}^{}{}_{}{}^{(3)}{}_{}{}^{2}\}`$
$`{\displaystyle \frac{1}{4}}e^\sigma \{e^{\varphi _1\frac{1}{\sqrt{3}}\varphi _2}_{2}^{}{}_{}{}^{(1)}{}_{}{}^{2}+e^{\varphi _1\frac{1}{\sqrt{3}}\varphi _2}_{2}^{}{}_{}{}^{(2)}{}_{}{}^{2}+e^{\frac{2}{\sqrt{3}}\varphi _2}_{2}^{}{}_{}{}^{(3)}{}_{}{}^{2}\}].`$
We would like to make few comments about the origin of different fields in this action. The action contains three scalars $`\sigma `$, $`\varphi _1`$ and $`\varphi _2`$. They are certain linear combinations of ten dimensional dilaton, and the two scalars that originate due to compactification from ten to eight dimensions. Three different three-form field strengths are denoted above as $`H_3^{(i)}`$. Furthermore, there are six two-form field strengths. Out of them $`F_2^{(i)}`$ come from reduction of various antisymmetric tensors in ten dimension. The other set $`_2^{(i)}`$ have their KK origin. In order to keep our discussion simple, we have set all the other fields to zero including the zero forms (axions) that appear in the eight dimensional action. Various details of eight dimensional type IIB supergravity action can be found in . As discussed previously, type IIB string in eight dimensions has $`SL(3,R)`$ symmetry. Defining $`H_3=dB_2,F_2=dA_1`$ and $`=d𝒜_1`$, it is easy to see that (21) is invariant under
$`g_{\mu \nu }g_{\mu \nu },\sigma \sigma ,`$
$`𝚲𝚲^T,𝒜_\mathrm{𝟏}𝚲𝒜_\mathrm{𝟏},`$
$`𝐀_1𝚲𝐀_1,𝐁_2(𝚲^1)^T𝐁_2,`$ (22)
where $`𝚲`$ is a global $`SL(3,R)`$ matrix. $``$ is a matrix with diagonal entries $`(e^{\varphi _1+\frac{1}{\sqrt{3}}\varphi _2},`$ $`e^{\varphi _1+\frac{1}{\sqrt{3}}\varphi _2},e^{\frac{2}{\sqrt{3}}\varphi _2})`$. In (22), $`𝒜_\mathrm{𝟏}`$ is defined as three dimensional column matrix with entries $`𝒜_1^{(1)}`$, $`𝒜_1^{(2)}`$ and $`𝒜_1^{(3)}`$. We have also defined $`𝐀_1`$ and $`𝐁_2`$ in a similar manner. In the following, we will be using only $`SO(3)`$ subgroup $`𝚲`$ of $`SL(3,R)`$. This can be represented by Euler angles $`\theta ,\varphi `$ and $`\psi `$.
The macroscopic string solution of this theory in Einstein frame can be written down as
$`ds^2={\displaystyle \frac{1}{[1+NG(r)]^{\frac{2}{3}}}}`$ $`[dt^2+(dx^7)^2]+{\displaystyle \frac{q^2G(r)}{4N[1+NG(r)]^{\frac{5}{3}}}}[dt+dx^7]^2`$ (23)
$`+[1+NG(r)]^{\frac{1}{3}}(dr^2+r^2d\mathrm{\Omega }_5^2),`$
with
$`\varphi _1={\displaystyle \frac{1}{2}}\mathrm{log}[1+NG(r)],\varphi _2={\displaystyle \frac{1}{2\sqrt{3}}}\mathrm{log}[1+NG(r)]`$
$`B_{t7}^{}{}_{}{}^{(1)}={\displaystyle \frac{NG(r)}{[1+NG(r)]}},𝒜_t^{(2)}=𝒜_7^{(2)}={\displaystyle \frac{qG(r)}{2[1+NG(r)]}},`$ (24)
All other fields are set to zero. In the above expressions, $`G(r)=\frac{1}{4\omega _5r^4},N=M\mathrm{cosh}^2\frac{\delta }{2},q=M\mathrm{sinh}\delta `$. Here $`\omega _5`$ is the unit volume of the $`5`$-sphere. The tension of the string is given by $`T=N`$. Using the asymptotic behavior of various fields, we find that the NS-NS two form charge $`(Z)`$, the electric charge $`(Q)`$ and the electric current $`(J)`$ for the solution are given respectively by:
$$Z=N,Q=q,\mathrm{and}J=q.$$
(25)
Notice that the electric charge and current are same for the solution. One way to obtain this solution is to first embed eight dimensional heterotic string theory in type IIB string theory. Then we can translate the charged heterotic string solutions of in terms of type IIB variables.
Now, using the symmetry of eight dimensional type IIB strings, one can construct an $`SL(3,Z)`$ multiplet of above string solution. We do not give this explicitly, since it is straightforward to write them down. We directly write down the charges $`𝐙^{}`$ and and current $`𝐉^{}`$ that follow from the above configuration:
$`𝐙^{}`$ $`=`$ $`\left(\begin{array}{c}Z_{}^{}{}_{}{}^{(1)}\\ Z_{}^{}{}_{}{}^{(2)}\\ Z_{}^{}{}_{}{}^{(3)}\end{array}\right)=N_{(z_1,z_2,z_3)}\left(\begin{array}{c}\mathrm{cos}\theta \mathrm{cos}\varphi \\ \mathrm{sin}\theta \mathrm{cos}\varphi \\ \mathrm{sin}\varphi \end{array}\right),`$
$`𝐉^{}`$ $`=`$ $`\left(\begin{array}{c}J_{}^{}{}_{}{}^{(1)}\\ J_{}^{}{}_{}{}^{(2)}\\ J_{}^{}{}_{}{}^{(3)}\end{array}\right)=q_{(z_1,z_2,z_3)}\left(\begin{array}{c}\mathrm{sin}\theta \mathrm{cos}\psi \mathrm{cos}\theta \mathrm{sin}\varphi \mathrm{sin}\psi \\ \mathrm{cos}\theta \mathrm{cos}\psi \mathrm{sin}\theta \mathrm{sin}\varphi \mathrm{sin}\psi \\ \mathrm{cos}\varphi \mathrm{sin}\psi \end{array}\right),`$ (26)
where $`N_{(z_1,z_2,z_3)}=\sqrt{z_1^2+z_2^2+z_3^2}N`$ and $`q_{(z_1,z_2,z_3)}=\sqrt{z_1^2+z_2^2+z_3^2}q`$. Parameters $`\theta ,\varphi ,\psi `$ in equation (26) are the Euler angles, as mentioned earlier. We would now like to identify $`𝐙^{}=(z_1,z_2,z_3)^TM\mathrm{cosh}^2\delta /2`$. This, in turn, fixes part of the $`SO(3)`$ group parameters $`\theta `$ and $`\varphi `$. Namely,
$`\mathrm{cos}\theta \mathrm{cos}\varphi ={\displaystyle \frac{z_1}{_{i=1}^3\sqrt{z_i^2}}},\mathrm{sin}\theta \mathrm{cos}\varphi ={\displaystyle \frac{z_2}{_{i=1}^3\sqrt{z_i^2}}},\mathrm{sin}\varphi ={\displaystyle \frac{z_3}{_{i=1}^3\sqrt{z_i^2}}}.`$ (27)
Notice that in this way, $`Z_{(1,0,0)}`$ string corresponds to electrically charged F-string, $`Z_{(0,1,0)}`$ is an electrically charged D-string, and, $`Z_{(0,0,1)}`$ is a ten dimensional D-3 brane wrapped on two internal circles. The subscript on $`Z`$ in the last line denotes their $`z`$ quantum numbers. In a similar manner, we can define the electric charge of the configuration in (26) as $`𝐉^{}=(q_1,q_2,q_3)^TM\mathrm{sinh}\delta `$. However, $`q`$’s are not independent quantities. Rather, they are determined by $`z`$’s and one of the $`SO(3)`$ group parameter $`\psi `$. Explicitly, using (27) in (26) for currents, we get
$`\left(\begin{array}{c}J_1^{}\\ J_2^{}\\ J_3^{}\end{array}\right)=\left(\begin{array}{c}\frac{z_2\sqrt{z_1^2+z_2^2+z_3^2}\mathrm{cos}\psi z_1z_3\mathrm{sin}\psi }{\sqrt{z_1^2+z_2^2}}\\ \frac{z_1\sqrt{z_1^2+z_2^2+z_3^2}\mathrm{cos}\psi z_2z_3\mathrm{sin}\psi }{\sqrt{z_1^2+z_2^2}}\\ \sqrt{z_1^2+z_2^2}\mathrm{sin}\psi \end{array}\right).`$ (28)
Here we note that one can get different $`𝐉^{}`$’s for different value of $`\psi `$.
In order to construct a junction configuration, one can consider a special class of solutions, namely where a $`(z_1,0,0)`$ and $`(0,z_2,0)`$ strings meet. From the $`𝐙`$ charge conservation, we see that resulting string must be a $`(z_1,z_2,0)`$ string. Furthermore, in order to analyze the stability of a junction of three such strings we notice that their string tensions are given by expressions:
$$T_{(z_1,z_2,z_3)}=\sqrt{z_1^2+z_2^2+z_3^2}T_{(1,0,0)}\mathrm{cosh}^2\frac{\delta }{2},$$
(29)
where $`T_{(1,0,0)}`$ is the tension of electrically neutral $`(1,0,0)`$ string, and, from (25), we see that it is given by $`T_{(1,0,0)}=M`$. Once again, since $`\delta `$ does not mix with $`z_1`$ and $`z_2`$, various angles between the strings in a network would be same as their electrically neutral counterparts. Beside $`Z`$ charge conservation and tension balance, our string junction have to satisfy other constraints as discussed before. One of them comes from electric current conservation. We notice that in general the electric charge is not conserved unless $`\psi =0`$. Thus the only allowed charged string junction in this class is for $`\psi =0`$, when
$`𝐉_{}^{}{}_{(z_1,z_2,0)}{}^{}=𝐉_{}^{}{}_{(z_1,0,0)}{}^{}+𝐉_{}^{}{}_{(0,z_2,0)}{}^{}.`$ (30)
This is the ‘Kirchoff’s law’ for the junction. It simply says that the algebraic sum of the currents around the junction must be zero. At this stage, the restriction on $`\psi `$ for charge conservation might seem unnatural. However, we should notice that we started with a very special class of solutions (23). We thus believe that the restriction on $`\psi `$ is an artifact of restricting ourselves within this special class of configuration. We expect that such restrictions can be avoided if we look for more general class of string junctions. Now, turning back to (24), we see that the seed solution that we started with has $`𝒜_t^{(2)}=𝒜_7^{(2)}`$. This in turn, leads to a condition on the charge densities similar to the one in (30) for the currents.
Now, for charge string junctions satisfying the above conservation and stability criteria, we can put them together to construct a string network as in . However, unlike in the previous cases, in our case these conditions only guarantee their classical stability properties. In addition, as in , one has to examine supersymmetry property as well to find out if these are BPS string networks or not. In the later case, they will decay into other BPS states. It is of interest to examine if the corresponding final states are again string-networks and whether they are built out of charged or neutral strings. These statements can be made more precise by thinking of string networks on tori. Then the final mass formula for a ‘particle-like’ object will carry the overall factor $`\mathrm{cosh}^2\frac{\delta }{2}`$ appearing in the string tension in (29) which has a minimum at $`\delta =0`$. One however needs a more careful study, including quantum corrections, to clarify this further.
We conclude by stating that one of our main motivation for studying charged, current carrying junction configurations and their networks is to set a framework for understanding entropy associated with the network when compactified on two-torus. Electrically charged networks that we discuss in this paper can be viewed as excitations over neutral networks. These excitations, in some examples, preserve a fraction of original supersymmetry. We believe (as was in the case of identifying string states associated with black hole entropy; see for example , ) that identification of the degrees of freedom for such excitations will play important role in understanding entropy associated with network on torus. We hope to return back to this issue in the future.
Acknowledgements: We have benefitted from discussions with Sunil Mukhi and Ashoke Sen.
|
warning/0005/math0005247.html
|
ar5iv
|
text
|
# 1 Introduction
## 1 Introduction
Even in the setting of infinite dimensional dynamics many of the dynamical objects of interest are low dimensional, e.g. equilibria, periodic orbits, connecting orbits, horseshoes, etc. In this paper we introduce techniques which, in principle, allow for the rigorous verification of such solutions for a wide variety of partial differential equations. Our approach is to combine rigorous computer calculations with topological invariants to obtain accurate existence statements. To demonstrate these techniques we have chosen to study the the Kuramoto-Sivashinsky equation
$$u_t=\nu u_{xxxx}u_{xx}+2uu_x(t,x)[0,\mathrm{})\times (\pi ,\pi )$$
(1)
subject to periodic and odd boundary conditions
$$u(t,\pi )=u(t,\pi )\mathrm{and}u(t,x)=u(t,x).$$
(2)
The following theorem is a prototype for the results which can be obtained.
###### Theorem 1.1
Let $`u(x)=_{k=1}^{28}a_k\mathrm{sin}(kx)`$ where the $`a_k`$ are given in Table 1. Then, for $`\nu =0.1`$ there exists an equilibrium $`u^{}(x)`$ for (1) such that
$$u^{}u_{L^2}<2.71547\times 10^{13}\mathrm{and}u^{}u_{C^0}<2.06706\times 10^{13}.$$
Having stated this theorem we now try to put the result into the context of the goals of our methods. To begin with it needs to be emphasized that the computations which lead to this result are rigorous in the sense that we have employed interval arithmetic to overcome all errors introduced by the fact that we are using floating point arithmetic in our calculations.
As it will become clear in the later sections, this result is obtained by studying the full partial differential equation rather than attempting to solve a boundary value problem. While from the point of view of traditional numerical analysis this approach may appear inefficient, it is an important point. To be more precise, our method does not attempt to directly approximate any particular solution to the partial differential equation. Rather we essentially compute the Conley index of a compact region, called an isolating neighborhood, of phase space. The diameter of this region provides the error bounds stated in the theorem. The index guarantees the existence of the equilibrium solution.
The Conley index theory is a far reaching topological generalization of Morse theory. In particular, this index can be used to prove the existence of periodic orbits, connecting orbits, and chaotic dynamics . It has been numerically observed that for various parameter values (1) contains these types of dynamical objects. In principle, combining earlier rigorous numerical methods with the techniques described in this paper and the above mentioned index theorems will lead to rigorous proofs of the existence of periodic orbits and even chaotic dynamics. However, we do not pursue these more complicated structures in this paper for two reasons. First, finding the appropriate isolating neighborhoods is more complicated in these cases and our goal here is to emphasize the fundamental ideas associated with the methods. The second, and more important point, is that a straightforward application of the earlier numerical methods would lead to large computations - which we believe can be avoided by alternative methods (see for example ). This latter point is currently being investigated.
Returning to our discussion of Theorem 1.1, an obvious question concerns the stability of $`u^{}`$. For this we have no definitive answer. As was indicated before our method does not directly approximate $`u^{}`$ and therefore we do not obtain uniqueness results or hyperbolicity results. On the other hand, being a generalization of the Morse index the Conley index does contain some information about the stability or instability of the dynamics in the isolating neighborhood. Thus, what can be asserted is the following. Assume that $`u^{}`$ is a hyperbolic fixed point, i.e. all eigenvalues have nonzero real part, and that $`u^{}`$ is the only solution which remains within either the $`L^2`$ or $`C^0`$ bounds of $`u`$ for all time, then $`u^{}`$ has exactly two unstable eigenvalues, i.e. its unstable manifold is two dimensional. We hope to treat this problem in a subsequent paper.
It should be mentioned that even though we are doing the computations via an approximation of the full partial differential equation, we never integrate the equations. Rather, as will be made clear in Section 2 the computations are reduced to solving a set of inequalities. It is for this reason that we are able to get such sharp bounds on the equilibria. As the following theorem demonstrates we can, in fact obtain bounds on the level of the floating point accuracy.
###### Theorem 1.2
One can compute a sequence $`a_1,a_2,\mathrm{},a_{30}`$ and the function $`u(x)=_{k=1}^{30}a_k\mathrm{sin}(kx)`$, such that for $`\nu =0.75`$ there exists an equilibrium $`u^{}(x)`$ for (1) such that $`u^{}u_{L^2}<1.26281\times 10^{15}`$ and $`u^{}u_{C^0}<9.57396\times 10^{16}`$.
By now it is a well demonstrated principle that the asymptotic behavior of a wide variety of infinite dimensional dynamical systems is finite dimensional . The Kuramoto-Sivashinsky equation (1) is a particularly well studied example of such a system . In fact, it is known that (1) posses an inertial manifold and therefore, that there exists a family of ordinary differential equations that exactly describes the asymptotic dynamics. Unfortunately, the estimates for the dimension of these manifolds make them impractical for our purposes .
We mention these methods to emphasize that our approach does not directly make use of any of these results. What appears to be essential for our techniques is that the spectrum of the linear operator for the evolution equation is not clustered near the imaginary axis. This is in contrast to the inertial manifold techniques which strongly rely on gap conditions of the spectrum or cone conditions from the flow. Our approach is to use the computer to restrict our attention to that portion of phase space in which the desired dynamics (for this paper the fixed points) occur. Obviously, by restricting the phase space one can get much better estimates. This sets up a loop by which one can continuously improve the estimates until the desired bounds are reached.
Our analysis of the fixed points for (1) was motivated in part by the work of Jolly, Keverkidis, and Titi . In particular, using a 12 mode traditional Galerkin approximation, they produced a bifurcation diagram for $`\nu (0.057,\mathrm{})`$. We used their reported solutions to test our methods. In particular, as is indicated below we were able to find and prove the existence of an equilibrium point on each of their stable branches. Unfortunately, we used a fairly primitive search procedure and therefore missed a few unstable branches. Our expectation is that by combining our methods with a continuation package, one could produce a rigorous bifurcation diagram with fairly precise bounds in a computationally inexpensive manner.
Below we include some of the steady states we found.
* $`\nu =0.5`$. Two stable unimodal fixed points
* $`\nu =0.3`$. Two stable unimodal fixed points
* $`\nu =0.127`$, $`\nu =0.125`$. A stable and unstable bimodal fixed point. Negative branch is stable, positive one is unstable with apparently two-dimensional unstable manifold.
Our primitive search procedure did not find a solution on the bi-tri branch.
* $`\nu =0.1`$ An bimodal stable and unstable (2 unstable directions) and two unstable trimodal fixed points (both with 1-dimensional unstable manifold)
We did not find an unstable branch connecting bi-tri branch with quadrimodal branch.
* $`\nu =0.08`$ A bimodal stable (neg. branch) and unstable (2 unstable directions) fixed points. A pair of stable fixed points close to $`R_3t_2`$ (see ). A pair of unstable trimodal fixed points (1 unstable direction).
We did not find a branch connecting bi-tri and quadimodal branches.
* $`\nu =0.0666..`$, $`\nu =0.063`$ Two unstable bimodal points, two stable trimodal points and two stable solutions apparently belonging to the giant branch.
We are lacking two unstable branches which are present in .
* $`\nu =0.062`$, Two stable trimodal points and two stable points from giant branch
* $`\nu =0.045`$, Two stable points from giant branch and pairs of unstable tri- and quadrimodal fixed points
* $`\nu =0.04`$, Two stable giant fixed points. Two stable quadrimodal fixed points. Two unstable trimodal points.
* $`\nu =0.029`$ Two unstable quadrimodal points.
## 2 The Method
Our method begins with the reduction of the full dynamical system to a lower dimensional system which can be studied numerical. In particular, we begin with a nonlinear evolution equation in a Hilbert space $`H`$ ($`L^2`$ in our treatment of Kuramoto-Sivashinsky) of the form
$$\frac{du}{dt}=F(u)$$
(3)
where domain of $`F`$ is dense in $`H`$. Furthermore, we assume that $`\{\varphi _ii=0,1,\mathrm{}\}`$ forms a complete orthogonal basis for $`H`$.
In the case of the Kuramoto-Sivashinsky equation $`F(u)=Lu+B(u,u)`$, where $`L`$ is a linear part and $`B`$ is a nonlinear part, the functions $`\{\phi _i\}`$ are chosen to be eigenvalues of $`L`$.
Fix $`m\text{}`$ and let
$$P=P_m:HX_m=X$$
be the orthogonal projection from $`H`$ onto the finite dimensional subspace spanned by $`\{\varphi _1,\varphi _2,\mathrm{},\varphi _m\}`$. Let
$$Q=Q_m:=IP:HY=Y_m$$
be the complementary orthogonal projection. Finally, let
$$A_k:H\text{}$$
be the orthogonal projection onto the subspace generated by $`\varphi _k`$.
Given $`uH`$, let $`Pu=p`$ and $`Qu=q`$. Equation (3) can be rewritten as
$`{\displaystyle \frac{dp}{dt}}`$ $`=`$ $`PF(p,q)`$ (4)
$`{\displaystyle \frac{dq}{dt}}`$ $`=`$ $`QF(p,q)`$ (5)
The strategy adopted is fairly simple: study the dynamics of the low dimensional Galerkin projection (4) to draw conclusions about the dynamics of (3). Before turning to the precise conditions, consider the following heuristic description of our approach.
Let $`WX=X_m`$. For $`j>m`$, let $`W_jX_j`$ such that $`P(P_j^1(W_j))=W`$, (i.e. $`W_j=W(IP)W_j`$). Similarly, let $`VY`$ and set $`V_j=Q_j(V)`$. Furthermore, given $`q_jV_j`$ assume that $`lim_j\mathrm{}q_j=0`$. Our only knowledge concerning the higher order modes or “tails” of the solutions is that they project into $`V`$. This implies that our knowledge of the vector field is reduced to the following differential inclusion
$$\frac{dp}{dt}PF(p,V)$$
where $`pW`$. Numerical calculations on this equation are used to find topological invariants (the Conley index, the fixed point index) which guarantee the existence of specific dynamics, e.g. fixed points, periodic orbits, symbolic dynamics, positive entropy, etc. It is simultaneously argued that the topological invariant is the same for any Galerkin system of the form
$$\frac{dp_j}{dt}PF(p_j,V_j)$$
where $`p_jW_j`$. Thus, the same dynamical object exists for each sufficiently high Galerkin approximation. Finally, it is shown that the limit of these objects leads to the desired dynamics for the full system (3).
### 2.1 Self-consistent Bounds
As one might expect the orthomormal basis $`\{\varphi _i\}`$ and the sets $`W`$ and $`V`$ must be chosen with care. The first issue that needs to be dealt with is analytic in nature - solutions to the ordinary differential equations must approximate solutions of the partial differential equation. This leads to the following definition.
###### Definition 2.1
Let $`m,M\text{}`$ with $`mM`$. A compact set $`WX_m`$ and a sequence of pairs $`\{a_k^\pm \text{}a_k^{}<a_k^+,k\text{}\}`$ form self-consistent apriori bounds for (3) if the following conditions are satisfied:
For $`k>M`$, $`a_k^{}<0<a_k^+`$.
Let $`\widehat{a}_k:=\mathrm{max}|a_k^\pm |`$ and set $`\widehat{u}=_{k=0}^{\mathrm{}}\widehat{a}_k\varphi _k`$. Then, $`\widehat{u}H`$. In particular, $`\widehat{u}<\mathrm{}`$.
The function $`uF(u)`$ is continuous on
$$W\underset{k=m+1}{\overset{\mathrm{}}{}}[a_k^{},a_k^+]H.$$
In practice $`W_{k=1}^m[a_k^{},a_k^+]`$. Given self-consistent apriori bounds $`W`$ and $`\{a_k^\pm \}`$, let
$$V:=\underset{k=m+1}{\overset{\mathrm{}}{}}[a_k^{},a_k^+]Y_m.$$
Our goal is to numerically solve (4) on $`W`$ and draw conclusions about the dynamics of (3) on the set $`WVH`$. To do this we will make use of the following results, the first two of them are obvious.
###### Lemma 2.2
Given self-consistent apriori bounds $`W`$ and $`\{a_k^\pm \}`$, $`WV`$ is a compact subset of $`H`$.
###### Lemma 2.3
Given self-consistent apriori bounds $`W`$ and $`\{a_k^\pm \}`$, $`WV`$, then
$$\underset{n\mathrm{}}{lim}Q_n(F(u))=0,\text{uniformly for }uWV$$
###### Proposition 2.4
Let $`W`$ and $`\{a_k^\pm \}`$ be self-consistent bounds for (3). A function $`a:[0,T]WV`$ given by
$$a(t):=\underset{k=0}{\overset{\mathrm{}}{}}a_k(t)\varphi _k$$
is a solution to (3), if and only if, for each $`k\text{}`$ and all $`t[0,T]`$
$$\frac{da_k}{dt}=A_kF(a).$$
(6)
Proof. ($``$) This direction follows directly from the projection of (3) onto each of the basis elements.
($``$) Assume that (6) is satisfied for each $`k\text{}`$ and all $`t[0,T]`$. Let
$$a(t):=\underset{k=0}{\overset{\mathrm{}}{}}a_k(t)\varphi _kH$$
First observe that from C3 it follows immediately that $`_{k=1}^{\mathrm{}}\frac{da_k}{dt}\varphi _k=F(a)H`$.
It needs to be shown that
$$\frac{da}{dt}=\underset{h0}{lim}\frac{a(t+h)a(t)}{h}=\underset{k=0}{\overset{\mathrm{}}{}}\frac{da_k}{dt}\varphi _k.$$
This is equivalent to showing that
$$\underset{h0}{lim}\left|\frac{1}{h}\underset{k=1}{\overset{\mathrm{}}{}}(a_k(t+h)a_k(t))\varphi _k\underset{k=1}{\overset{\mathrm{}}{}}\frac{da_k}{dt}\varphi _k\right|=0$$
for all $`t[0,T]`$.
Fix $`h>0`$, then for any $`n\text{}`$
$`\left|{\displaystyle \frac{1}{h}}{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}(a_k(t+h)a_k(t))\varphi _k{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{da_k}{dt}}\varphi _k\right|`$
$`\left|{\displaystyle \frac{1}{h}}{\displaystyle \underset{k=1}{\overset{n}{}}}(a_k(t+h)a_k(t))\varphi _k{\displaystyle \underset{k=1}{\overset{n}{}}}{\displaystyle \frac{da_k}{dt}}\varphi _k\right|`$
$`+\left|{\displaystyle \frac{1}{h}}{\displaystyle \underset{k=n+1}{\overset{\mathrm{}}{}}}(a_k(t+h)a_k(t))\varphi _k\right|+\left|{\displaystyle \underset{k=n+1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{da_k}{dt}}\varphi _k\right|`$
We will estimate the three terms on the right hand side separately. From lemma 2.3 it follows for a given $`ϵ>0`$ there exists $`n_0`$ such that $`n>n_0`$ implies
$$\left|\underset{k=n+1}{\overset{\mathrm{}}{}}\frac{da_k}{dt}\varphi _k\right|=\left|Q_n(F(a))\right|<ϵ/3.$$
From now on fix $`n>n_0`$. Again lemma 2.3 and the mean value theorem implies
$`\left|{\displaystyle \underset{k=n+1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{h}}(a_k(t+h)a_k(t))\varphi _k\right|`$ $`=`$ $`\left|{\displaystyle \underset{k=n+1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{da_k}{dt}}(t+\theta _kh)\varphi _k\right|`$
$`=`$ $`|Q_n(F(a(t+\theta _kh))|<ϵ/3.`$
Finally, for $`h`$ sufficiently small,
$$\left|\frac{1}{h}\underset{k=1}{\overset{n}{}}(a_k(t+h)a_k(t))\varphi _k\underset{k=1}{\overset{n}{}}\frac{da_k}{dt}\varphi _k\right|<ϵ/3$$
and hence the desired limit is obtained.
### 2.2 Conley Index
Proposition 2.4 indicates that given self-consistent apriori bounds $`W`$ and $`\{a_k^\pm \}`$, finite time solutions to (6) are solutions to the full partial differential equation. Thus, the goal of this paper is to find solutions to (6). Of course, numerically one can only study (4) restricted to $`W`$ and then argue that the resulting numerical solution is an approximation to a solution to (3). Hence, rather than attempting to approximate specific trajectories in $`W`$ directly, the objective is to compute a Conley index for (4) and then show that this index information is sufficient to guarantee a solution for (3).
In order to describe this index the following definitions are needed. Let $`\phi :\text{}\times \text{}^m\text{}^m`$ be a continuous flow generated by a differential equation $`\dot{z}=f(z)`$.
###### Definition 2.5
A compact set $`N\text{}^n`$ is an isolating neighborhood if
$$\mathrm{Inv}(N,\phi ):=\{zN\phi (\text{},z)N\}\mathrm{int}N.$$
If, in addition, for any $`zN`$, there exists $`t_z>0`$ such that
$$\phi ((0,t_z),z)N=\mathrm{}\mathrm{or}\phi ((t_z,0),z)N=\mathrm{},$$
(7)
then $`N`$ is an isolating block. Given an isolating neighborhood $`N`$, the associated maximal invariant set $`\mathrm{Inv}(N,\phi )`$ is an isolated invariant set.
The easiest way to verify the existence of an isolating block is through local sections.
###### Definition 2.6
$`\mathrm{\Xi }\text{}^n`$ is a local section for $`\phi `$ if for some $`ϵ>0`$
$$\phi :(ϵ,ϵ)\times \mathrm{\Xi }\phi ((ϵ,ϵ),\mathrm{\Xi })$$
(8)
is a homeomorphism and $`\phi ((ϵ,ϵ),\mathrm{\Xi })`$ is an open subset of $`\text{}^n`$.
A special form of local section is a hypersurface which is transverse to the flow. More formally, let $`\mathrm{\Xi }\text{}^n`$ be an $`n1`$ dimensional manifold with normal vector $`\mu (z)`$ at $`z\mathrm{\Xi }`$. $`\mathrm{\Xi }`$ is a local section if for each $`z\mathrm{\Xi }`$,
$$\mu (z)f(z)0.$$
(9)
It is straightforward to check that $`N`$ is an isolating block if $`N`$ can be written as the union of the closure of local sections with the property that (7) is satisfied at every point in the intersection of the closure of the sections.
In this paper the focus is both on proving the existence of equilibria and providing tight bounds on the location of the equilibria. To do this requires have good isolating blocks. With this in mind consider the linear ordinary differential equation
$$\dot{z}=Bz,z\text{}^n.$$
(10)
Assume that the origin is a hyperbolic fixed point. Without loss of generality it can be assumed that $`B`$ is in Jordan normal form. Generically, to each real eigenvalue there is associated a 1-dimensional eigenspace and to each pair of complex conjugate eigenvalues there is an associated 2-dimensional eigenspace. Thus, $`\text{}^n`$ can be decomposed into the product of eigenspaces, i.e.
$$\text{}^n=V_1\times V_2\times \mathrm{}\times V_k$$
where $`V_i`$ is either or $`\text{}^2`$. In what follows we will use the following notation, $`z_iV_i`$, $`i=1,\mathrm{},k`$, and if $`V_i\text{}^2`$, then $`z_i=(x_i,y_i)`$.
Our interest is not on the dynamics of (10) on the entire phase space, but rather on a prescribed compact subset. Since our goal is to understand the equilibria of (10) consider a neighborhood of the origin,
$$N=I_1\times I_2\times \mathrm{}\times I_k$$
where
$$I_i:=\{\begin{array}{cc}[b_i^{},b_i^+],b_i^{}<0<b_i^+\hfill & \text{if }V_i\text{}\text{,}\hfill \\ \{(x_i,y_i)V_i\sqrt{x_i^2+y_i^2}b_i,b_i>0\}\hfill & \text{if }V_i\text{}^2\text{.}\hfill \end{array}$$
The following result is obvious, but to make a point crucial to the results of this paper we will provide the proof.
###### Lemma 2.7
The compact set $`N`$ is an isolating block for (10).
Proof. Since $`B`$ is in Jordan normal form the system decouples according to the decomposition $`\text{}^n=V_1\times V_2\times \mathrm{}\times V_k`$.
If $`V_i=\text{}`$, then (10) reduces to $`\dot{z}_i=\lambda _iz_i`$. Since $`B`$ is hyperbolic $`\lambda _i0`$, and hence at $`z_i=b_i^\pm `$ (9) becomes
$$\lambda _ib_i^\pm 0.$$
If $`V_i=\text{}^2`$, then (10) reduces to
$`\dot{x}_i`$ $`=`$ $`\alpha _ix_i+\beta _iy_i`$
$`\dot{y}_i`$ $`=`$ $`\beta _ix_i+\alpha _iy_i`$
where by hyperbolicity $`\alpha _i0`$. So again for $`\sqrt{x_i^2+y_i^2}=b_i`$ (9) becomes
$$(x_i,y_i)(\alpha _ix_i+\beta _iy_i,\beta _ix_i+\alpha _iy_i)^t=\alpha _ib_i^20.$$
To see why this trivial argument is of importance, consider the more interesting example of
$$\dot{z}=Bz+f(z)+E(z)$$
(11)
where $`f:\text{}^n\text{}^n`$ is $`o(z^2)`$ at $`0`$ and $`E`$ represents a known bounded error. In our situation $`E`$ arises from numerical errors and approximations. More precisely, we assume that there are known constants $`c_i`$ such that
$$\underset{zN}{sup}E_i(z)c_i.$$
Observe that a sufficient condition for $`N`$ to be an isolating block for (11) is the following: for each $`i`$ such that $`V_i=\text{}`$,
$$\lambda _ib_i^\pm +f_i(z)+E_i(z)$$
(12)
has the same sign as $`\lambda _ib_i^\pm `$ over the set $`\{zNz_i=b_i^\pm \}`$, and for each $`i`$ such that $`V_i=\text{}^2`$,
$$(x_i,y_i)(\alpha x_i+\beta y_i+f_{i_1}(x)+E_{i_1}(x),\beta x_i+\alpha y_i+f_{i_2}(x)+E_{i_2}(x))^t$$
(13)
has the same sign as $`\alpha _i`$ over the set $`\{zN\sqrt{x_i^2+y_i^2}=b_i\}`$.
For the linear case the eigenvalues of $`B`$ are assumed to be known and the $`b_i^\pm `$ can be chosen arbitrarily. Therefore, one can also interpret (12) and (13) as providing a set of inequalities that if simultaneously solved for $`b_i^\pm `$ provide an isolating block even in the context of numerical errors and approximations. In particular, finding isolating blocks need not involve numerically solving the ordinary differential equation.
In itself the knowledge that $`N`$ is an isolating block does not imply anything about $`\mathrm{Inv}(N,\phi )`$. To gain information concerning the isolated invariant set we will make use of the Conley index. For our purposes we need only a very small portion of the index theory and so we give a minimal operational definition (see for further information).
###### Definition 2.8
Let $`N`$ be an isolating block and let $`N=L^+L^{}`$ where $`L^\pm `$ are closed sets. Furthermore, assume that $`zL^{}`$ implies that
$$\phi ((0,ϵ),z)N=\mathrm{}$$
for a sufficiently small $`ϵ>0`$. Similarly, assume that if $`zL^+`$, then
$$\phi ((ϵ,0),z)N=\mathrm{}$$
for a sufficiently small $`ϵ>0`$. The Conley index of $`S=\mathrm{Inv}(N,\phi )`$ is
$$CH_{}(S):=H_{}(N,L).$$
No knowledge of relative homology groups is required for the applications described in this paper. The following theorem gives a formula for the index of a hyperbolic fixed point.
###### Proposition 2.9
Let $`q`$ be the number of eigenvalues of $`B`$ with positive real part. Assume that for all $`i=1,\mathrm{},k`$ either the condition associated with (12) or the condition associated with (13) are satisfied . Then,
$$CH_j(\mathrm{Inv}(N,\phi ))\{\begin{array}{cc}\text{}\hfill & \text{if }j=q\hfill \\ 0\hfill & \text{otherwise.}\hfill \end{array}$$
The following theorem due to McCord \[13, Corollary 5.9\], indicates that if the Conley index is a that of Proposition 2.9, then there exists a fixed point in $`N`$.
###### Theorem 2.10
If the Conley index has the form
$$CH_j(\mathrm{Inv}(N,\phi ))\{\begin{array}{cc}\text{}\hfill & \text{if }j=q\hfill \\ 0\hfill & \text{otherwise,}\hfill \end{array}$$
for some $`q`$, then $`N`$ contains a fixed point.
To indicate how these index ideas will be used in this paper let us return to the system (11). Observe that the only assumption on the error term $`E`$ was that it is bounded, therefore, it is no longer apriori true that the origin is a fixed point or that there even exists a fixed point to (11). On the other hand the sets $`N`$ and $`L`$ remain unchanged. Therefore, the Conley index implies the existence of the fixed point.
### 2.3 A Singular Perturbation Result
As was indicated in the previous section, it is possible to find an isolating block for a finite dimensional ordinary differential equation about a fixed point by solving an appropriate set of inequalities. However, to do this requires a good estimate of the location of the fixed point, knowledge of the eigenvalues, the ability to evaluate the nonlinear terms, and estimates of associated errors. Therefore, the dimension to which one can hope to apply this procedure is obviously limited. In this section we will describe a singular perturbation result which allows one to “lift” the index computations of the previous sections to arbitrary dimensions.
The definition of self-consistent bounds related individual solutions of the infinite family of ordinary differential equations to solutions in partial differential equation. We now need to extend this definition in order to know that the index computations we perform for the finite dimensional approximation have implications for the partial differential equation.
###### Definition 2.11
Let $`m,M\text{}`$ with $`mM`$. A pair of compact sets $`NWX_m`$ and a sequence of pairs $`\{a_k^\pm \text{}a_k^{}<a_k^+,k\text{}\}`$ are topologically self-consistent if $`W`$ and $`\{a_k^\pm \}`$ are self-consistent apriori bounds and the following conditions are satisfied.
Let $`uW_{k=m+1}^{\mathrm{}}[a_k^{},a_k^+]`$. Then, for $`k>m`$
$$A_ku=a_k^\pm A_kF(u)0.$$
(14)
$`N`$ is an isolating block for (4) for all $`q_{k>m}[a_k^{},a_k^+]`$.
For Kuramoto-Sivashinsky we will make use of the following stricter form of C4.
Let $`uW_{k=m+1}^{\mathrm{}}[a_k^{},a_k^+]`$. Then, for $`k>m`$
$`A_ku=a_k^+`$ $``$ $`A_kF(u))<0`$ (15)
$`A_ku=a_k^{}`$ $``$ $`A_kF(u))>0.`$ (16)
Using the line of reasoning that was described in the analysis of (11), condition C5 can be replaced by the following assumption.
Let $`N`$ be an isolating block for (4). Let $`\nu ^\pm (p)`$ be the outward normal at $`pL^\pm `$. If $`uW_{k=m+1}^{\mathrm{}}[a_k^{},a_k^+]`$ such that $`PuL^\pm `$, then
$$PF(u)\nu ^+(Pu)<0PF(u)\nu ^{}(Pu)>0.$$
We shall now discuss two singular perturbation results. The first allows one to lift isolating blocks.
###### Theorem 2.12
Let $`m,M\text{}`$ with $`mM`$. Assume $`NWX_m`$ and the sequence of pairs $`\{a_k^\pm \text{}a_k^{}<a_k^+,k\text{}\}`$ are topologically self-consistent. Fix an integer $`r>m`$. Then for any $`q=_{k=r+1}^{\mathrm{}}q_k\varphi _k`$, such that $`q\mathrm{\Pi }_{k=r+1}^{\mathrm{}}[a_k^{},a_k^+]`$ and $`q_k=0`$ for $`k>M`$ the set
$$\stackrel{~}{N}:=N\times [a_{m+1}^{},a_{m+1}^+]\times [a_{m+2}^{},a_{m+2}^+]\times \mathrm{}\times [a_r^{},a_r^+]$$
is an isolating block for the system of equations
$$\dot{x}_k=A_kF(\underset{i=1}{\overset{k}{}}x_i\varphi _i+q)k=1,\mathrm{},r$$
(17)
where $`x𝐑^r`$.
Proof. Let $`u=(w,v)W_{k=m+1}^r[a_k^{},a_k^+]`$. From C1 it follows that $`u+qW\mathrm{\Pi }_{k=m+1}^{\mathrm{}}[a_k^{},a_k^+]`$. If $`u\stackrel{~}{N}`$, then either $`w`$ is in $`N`$ or $`v`$ is in $`_{k=m+1}^r[a_k^{},a_k^+]`$. In the first case C5 forces the vector field to be transverse at the boundary. In the second case transversality follows from C4.
###### Remark 2.13
For $`r>M`$ equations (17) are the Galerkin projection of $`\dot{u}=F(u)`$.
The direction of the vector field influences the index computation. With this in mind define
$$\mathrm{dir}(k):=\{\begin{array}{cc}1\hfill & \text{if }A_ku=a_k^+A_kF(u))<0\text{ and}\hfill \\ & \text{ }A_ku=a_k^{}A_kF(u))>0\hfill \\ 0\hfill & \text{if }A_ku=a_k^+A_kF(u))<0\text{ and}\hfill \\ & \text{ }A_ku=a_k^{}A_kF(u))<0\hfill \\ 0\hfill & \text{if }A_ku=a_k^+A_kF(u))>0\text{ and}\hfill \\ & \text{ }A_ku=a_k^{}A_kF(u))>0\hfill \\ 1\hfill & \text{if }A_ku=a_k^+A_kF(u))>0\text{ and}\hfill \\ & \text{ }A_ku=a_k^{}A_kF(u))<0\hfill \end{array}$$
###### Theorem 2.14
Let $`m,M\text{}`$ with $`mM`$. Assume $`NWX_m`$ and the sequence of pairs $`\{a_k^\pm \text{}a_k^{}<a_k^+,k\text{}\}`$ are topologically self-consistent. Fix an integer $`r>m`$. Let $`q=_{k=r+1}^{\mathrm{}}q_k\varphi _k`$, such that $`q\mathrm{\Pi }_{k=r+1}^{\mathrm{}}[a_k^{},a_k^+]`$ and $`q_k=0`$ for $`k>M`$. Let
$$\stackrel{~}{N}:=N\times [a_{m+1}^{},a_{m+1}^+]\times [a_{m+2}^{},a_{m+2}^+]\times \mathrm{}\times [a_r^{},a_r^+]$$
Consider the dynamical system induced by (17).
If for some $`j\{m+1,\mathrm{},r\}`$, $`\mathrm{dir}(j)=0`$, then
$$CH_{}(\mathrm{Inv}(\stackrel{~}{N}))=0.$$
Assume that for all $`j\{m+1,\mathrm{},r\}`$, $`\mathrm{dir}(j)0`$, and let $`d`$ be the number of $`j\{m+1,\mathrm{},r\}`$ such that $`\mathrm{dir}(j)=1`$, then
$$CH_{s+d}(\mathrm{Inv}(\stackrel{~}{N}))CH_s(\mathrm{Inv}(N)).$$
(18)
Proof. By Theorem 2.12, $`\stackrel{~}{N}`$ is an isolating block.
We will present the proof of the second part of the theorem, only.
Assume that for all $`j\{m+1,\mathrm{},r\}`$, $`\mathrm{dir}(j)0`$. Let $`𝒥:=\{jm<jr,\mathrm{dir}(j)=1\}`$. Set
$`\stackrel{~}{L}^{}:=`$ $`\left(L^{}\times {\displaystyle \underset{k=m+1}{\overset{r}{}}}[a_k^{},a_k^+]\right)`$
$`{\displaystyle \underset{j𝒥}{}}\left(N\times {\displaystyle \underset{k=m+1}{\overset{j1}{}}}[a_k^{},a_k^+]\times \{a_j^\pm \}\times {\displaystyle \underset{k=j+1}{\overset{r}{}}}[a_k^{},a_k^+]\right)`$
and
$`\stackrel{~}{L}^+:=`$ $`\left(L^+\times {\displaystyle \underset{k=m+1}{\overset{r}{}}}[a_k^{},a_k^+]\right)`$
$`{\displaystyle \underset{j𝒥}{}}\left(N\times {\displaystyle \underset{k=m+1}{\overset{j1}{}}}[a_k^{},a_k^+]\times \{a_j^\pm \}\times {\displaystyle \underset{k=j+1}{\overset{r}{}}}[a_k^{},a_k^+]\right)`$
Let $`\phi :\text{}\times \text{}^r\text{}^r`$, be any flow generated by
$$\dot{a}_k=A_kF(u+q)k=1,\mathrm{},r$$
where $`a_k=A_ku`$ and $`uW_{k=m+1}^r[a_k^{},a_k^+]`$.
Clearly, if $`P_ruL^{}`$, then $`\phi ((0,ϵ),P_ru)\stackrel{~}{N}`$ for small $`ϵ>0`$. Similarly, if $`P_ruL^+`$, then $`\phi ((ϵ,0),P_ru)\stackrel{~}{N}`$ for small $`ϵ>0`$.
Let $`u=(w,v)W_{k=m+1}^r[a_k^{},a_k^+]`$. If $`u\stackrel{~}{N}`$, then either $`w`$ is in $`N`$ or $`v`$ is in $`_{k=m+1}^r[a_k^{},a_k^+]`$. Therefore, $`N=\stackrel{~}{L}^+\stackrel{~}{L}^{}`$.
Thus, the Conley index of $`\mathrm{Inv}(\stackrel{~}{N})`$ is given by
$$CH_{}(\mathrm{Inv}(\stackrel{~}{N})H_{}(\stackrel{~}{N},\stackrel{~}{L}^{}).$$
A simple argument using the Mayer-Vietoris sequence gives the desired homology groups.
Observe that C4a implies that $`\mathrm{dir}(k)=1`$ for all $`k>m`$. Therefore, one has the following result.
###### Corollary 2.15
Assume $`NWX_m`$ and the sequence of pairs $`\{a_k^\pm \text{}a_k^{}<a_k^+,k\text{}\}`$ are topologically self-consistent and satisfy C4a. Fix an integer $`r>m`$ and let
$$\stackrel{~}{N}:=N\times [a_{m+1}^{},a_{m+1}^+]\times [a_{m+2}^{},a_{m+2}^+]\times \mathrm{}\times [a_r^{},a_r^+].$$
Then
$$CH_{}(\mathrm{Inv}(\stackrel{~}{N}))CH_{}(\mathrm{Inv}(N)).$$
The following theorem is used for all the results described in the Introduction.
###### Theorem 2.16
Assume $`NWX_m`$ and the sequence of pairs $`\{a_k^\pm \text{}a_k^{}<a_k^+,k\text{}\}`$ are topologically self-consistent and satisfy C4a. Assume
$$CH_j(\mathrm{Inv}(N,\phi ))\{\begin{array}{cc}\text{}\hfill & \text{if }j=l\text{,}\hfill \\ 0\hfill & \text{otherwise,}\hfill \end{array}$$
for some $`l`$, then there exists
$$u^{}N\times \underset{k=m+1}{\overset{\mathrm{}}{}}[a_k^{},a_k^+],$$
a fixed point for the partial differential equation (3).
Proof. Combining Theorems 2.14 and 2.10, immediately gives that for each $`r>M`$ there exists a fixed point
$$z_rN\times \underset{k=m+1}{\overset{r}{}}[a_k^{},a_k^+]$$
for the Galerkin projection onto the first $`r`$ coordinates.
Since $`N\times _{k=m+1}^{\mathrm{}}[a_k^{},a_k^+]`$ is compact the collection $`\{z_rr=m+1,m+2,\mathrm{}\}`$ contains a limit point $`u^{}`$. From the continuity of $`P_nF`$ on $`W_{k=m+1}^{\mathrm{}}[a_k^{},a_k^+]`$ it follows that $`P_nF(u^{})=0`$ for each $`n\text{}`$. By Proposition 2.4 $`u^{}`$ is a fixed point for (3).
### 2.4 Remarks on Related Work
We are aware of at least two other results that are closely related to the methods described in this Section. The first is work of L. Cesari from the early 60’s which in spirit is very similar to ours. His method can be characterized as follows . Let $`B`$ be a Banach space. Let $`X`$ be a finite dimensional subspace of $`B`$ and let $`P:BX`$ be a projection. Let $`\stackrel{~}{N}B`$ be closed with the property that $`P\stackrel{~}{N}=NX`$ is compact and for every $`xN`$, $`P^1(x)\stackrel{~}{N}`$ is closed. Consider a continuous map $`f:\stackrel{~}{N}B`$. The goal is to find fixed points for $`f`$ by studying the behavior of the projection of the map onto $`X`$.
It is obvious that $`u^{}`$ is a fixed point of $`f`$, if and only if $`Pu^{}=Pfu^{}`$, and $`u^{}=Pu^{}+(IP)f(u^{})`$.
Cesari’s method applies if and only if the following three conditions are satisfied:
1. For each $`xN`$,
$$P+(IP)f:P^1(x)\stackrel{~}{N}P^1(x)\stackrel{~}{N}$$
is a contraction.
2. Given condition (i), for each $`xN`$, there exists a unique $`u(x)\stackrel{~}{N}`$ such that $`Pu(x)+(IP)f(u(x))=u(x)`$. The function $`u:N\stackrel{~}{N}`$ is continuous.
3. There are no fixed points of $`Pfu:NN`$ on the boundary of $`N`$.
Relating this back to the context of this paper, observe that a fixed point for the partial differential equation is a fixed point for any nonzero constant time map of the corresponding semi-flow. (iii) is closely related to the condition C5. As stated (ii) is not well defined unless (i) holds. C4 is the analogous assumption to (i) and differs in two significant ways. A necessary condition to have a contraction, is for the stronger assumption of C4a to hold. However, C4a is not sufficient. An important point is that we do not make any assumptions on the direction of the vector field within $`\stackrel{~}{N}`$. Thus, condition C4a is in principle easier to verify than (i). On the other hand, this makes it clear that we cannot guarantee uniqueness of the fixed point given our assumptions.
The other work is due to C. Conley and P. Fife and is closely related to Theorem 2.14. Formulas of the form (18) are classical in the context of product systems (see ). In one finds a similar formula, but in that context at the parameter value for which one computes the index in the lower dimensional system, there is no higher dimensional dynamics defined. However, the higher dimensional system is defined for an arbitrarily small perturbation. The key idea is that in the proper context the lower dimensional system is normally hyperbolic. In this paper we circumvent this type of assumption using isolating blocks, C5, and imposing C4.
## 3 Estimates for Kuramoto-Sivashinsky equation
As the Hilbert space $`H`$ for the Kuramoto-Sivashinsky equation (1) we choose the subspace of $`L^2(\pi ,\pi )`$ consisting of $`2\pi `$-periodic and odd functions.
Since $`u(t,x)`$ is odd its Fourier expansion takes the form
$$u(t,x)=\underset{k=\mathrm{}}{\overset{k=\mathrm{}}{}}b_k(t)\mathrm{exp}(ikx)$$
(19)
Since $`u(t,x)`$ is real, $`b_k=\overline{b}_k`$. Substituting (19) into (1) gives the following equations
$$\dot{b}_k=(k^2\nu k^4)b_k+\text{i}k\underset{m=\mathrm{}}{\overset{m=\mathrm{}}{}}b_mb_{km}$$
(20)
Since we are interested in solutions with odd symmetry it follows that $`b_k`$ are pure imaginary. Let
$$a_k:=\sqrt{1}b_k.$$
Then, $`a_k=a_k`$ and $`a_0=0`$ which results in the following infinite system of ordinary differential equations
$$\dot{a}_k=k^2(1\nu k^2)a_kk\underset{n=1}{\overset{k1}{}}a_na_{kn}+2k\underset{n=1}{\overset{\mathrm{}}{}}a_na_{n+k}k=1,2,3,\mathrm{}$$
(21)
We will use these equations to draw rigorous conclusions about Kuramoto-Sivashinsky by finding self-consistent apriori bounds for (1) that satisfy the stronger condition C4a and then applying Theorem 2.16. To do this, however, we will need to understand the errors contributed by ignoring the higher modes and the errors introduced by the use of floating point arithmetic.
Let $`m,M\text{}`$ be fixed with $`mM`$. Let $`W\text{}^m`$ and $`\{a_k^\pm \text{}k\text{}\}`$ satisfy conditions C1 \- C3 with the added constraints that
$$W=\underset{k=1}{\overset{m}{}}[a_k^{},a_k^+]$$
and
$$a_k^\pm =\pm \frac{C_s}{k^s},k>M$$
(22)
for some constant $`C_s>0`$ and integer $`s>1`$.
Though technically incorrect, it is perhaps useful for the reader to think of the numerical approximation of the dynamics being computed with respect to the finite dimensional system
$$\dot{a}_k=k^2(1\nu k^2)a_kk\underset{n=1}{\overset{k1}{}}a_na_{kn}+2k\underset{n=1}{\overset{Mk}{}}a_na_{n+k}k=1,\mathrm{},m$$
(23)
where $`a_k=(a_k^{}+a_k^+)/2`$ for $`k=m+1,\mathrm{}M`$. In doing so it becomes clear that there are essentially three levels of approximation that need to be dealt with. The first involves the terms in the infinite tail $`\{a_kk>M\}`$. These are completely absent from (23) and therefore must be absorbed as a fixed error term (think of the term $`E(x)`$ in (11)). The power decay rule (22) will be used to determine this quantity. The second, involves the terms $`\{a_km<kM\}`$. In principle, one could set $`m=M`$, however our strategy is to try to obtain better estimates for these terms than can be expected by the general decay of (22). However, these terms act as constants and hence can be viewed as parameters for the system (23). Finally, the terms $`\{a_kk=1,\mathrm{},m\}`$ are the actual variables for the dynamical system being studied. It should also be kept in mind that we need to lift the index information, and therefore need to be able to verify C4a for all $`k`$.
Of course, our goal is that of rigorous computations. Therefore each of the above mentioned $`a_i`$ is actually an interval. The intervals associated with $`\{a_kk=1,\mathrm{},m\}`$ are essentially determined by the floating point approximations. For $`k>m`$, the intervals are
$$a_k=[a_k^{},a_k^+].$$
We will let
$$|a_k|:=\mathrm{max}\{|a_k^{}|,|a_k^+|\}.$$
To compute the above mentioned errors we return to (21) and observe that in addition to the linear part there is a finite sum of terms
$$FS(k)=\underset{n=1}{\overset{k1}{}}a_na_{kn}$$
(24)
and an infinite sum of terms
$$IS(k)=\underset{n=1}{\overset{\mathrm{}}{}}a_na_{n+k}.$$
(25)
Obviously bounds on these terms are necessary.
### 3.1 $`1kM`$
Since $`FS(k)`$ is a finite sum and we have already chosen the interval values for $`a_n`$, we can explicitly compute $`FS(k)`$. Perhaps it is worth noting that to evaluate $`FS(k)`$ only involves the intervals $`\{a_nn=1,\mathrm{},M1\}`$ which are chosen in such a way that we expect them to be reasonably good approximations of the actual terms.
###### Lemma 3.1
Assume $`1kM`$. Then,
$`IS(k)`$ $``$ $`{\displaystyle \underset{n=1}{\overset{Mk}{}}}a_na_{k+n}+C_s{\displaystyle \underset{n=Mk+1}{\overset{M}{}}}{\displaystyle \frac{|a_n|}{(k+n)^s}}[1,1]+`$
$`{\displaystyle \frac{C_s^2}{(k+M+1)^s(s1)M^{s1}}}[1,1]`$
Proof. By definition,
$$IS(k)=\underset{n=1}{\overset{Mk}{}}a_na_{k+n}+\underset{n=Mk+1}{\overset{M}{}}a_na_{k+n}+\underset{n=M+1}{\overset{\mathrm{}}{}}a_na_{k+n}.$$
With regard to the second sum
$`{\displaystyle \underset{n=Mk+1}{\overset{M}{}}}a_na_{k+n}`$ $``$ $`{\displaystyle \underset{n=Mk+1}{\overset{M}{}}}|a_n|{\displaystyle \frac{C_s}{(k+n)^s}}[1,1].`$
Finally, the third sum produces
$`{\displaystyle \underset{n=M+1}{\overset{\mathrm{}}{}}}a_na_{k+n}`$ $``$ $`{\displaystyle \underset{n=M+1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{C_s}{n^s}}{\displaystyle \frac{C_s}{(n+k)^s}}[1,1]`$
$``$ $`{\displaystyle \frac{C_s^2}{(k+M+1)^s}}[1,1]{\displaystyle \underset{n=M+1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n^s}}`$
$``$ $`{\displaystyle \frac{C_s^2}{(k+M+1)^s(s1)M^{s1}}}[1,1].`$
In above derivation we used the following estimate
$$\underset{n=M+1}{\overset{\mathrm{}}{}}\frac{1}{n^s}<_M^{\mathrm{}}\frac{dx}{n^s}=\frac{1}{(s1)M^{s1}}$$
###### Remark 3.2
This estimate and some of those that follow can be improved by noting that
$$\underset{n=M+1}{\overset{\mathrm{}}{}}\frac{1}{n^s(n+k)^s}<_M^{\mathrm{}}\frac{dx}{x^s(x+k)^s}.$$
Of course, the right hand side has an explicit rational expression, but it is rather complicated for large $`s`$ and so was not utilized here.
A simple extension of Lemma 3.1 leads to the following corollary.
###### Corollary 3.3
Let $`1km`$. Then,
$`{\displaystyle \underset{n=mk+1}{\overset{\mathrm{}}{}}}a_na_{n+k}`$ $``$ $`{\displaystyle \underset{n=mk+1}{\overset{Mk}{}}}a_na_{n+k}+C_s{\displaystyle \underset{n=Mk+1}{\overset{M}{}}}{\displaystyle \frac{|a_n|}{(k+n)^s}}[1,1]+`$
$`{\displaystyle \frac{C_s^2}{(k+M+1)^s(s1)M^{s1}}}[1,1]`$
Observe that collorary 3.3 estimates the error in the vector field due to the Galerkin projection, namely
$$A_k(p+q)A_kF(p)=2k\underset{n=mk+1}{\overset{\mathrm{}}{}}a_na_{n+k}$$
(26)
### 3.2 $`k>M`$
Throughout this section it is assumed that $`k>M`$. Let
$$e(k):=\{\begin{array}{cc}1\hfill & \text{if }k\text{ is even,}\hfill \\ 0\hfill & \text{if }k\text{ is odd.}\hfill \end{array}$$
###### Lemma 3.4
Let $`M<k2M`$. Then,
$$FS(k)2\underset{n=kM}{\overset{k/2}{}}a_na_{kn}+e(k)a_{k/2}^2+2C_s\underset{n=1}{\overset{kM1}{}}\frac{|a_n|}{(kn)^s}[1,1].$$
Proof. Expanding (24) gives
$`FS(k)`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{kM1}{}}}a_na_{kn}+{\displaystyle \underset{n=kM}{\overset{M}{}}}a_na_{kn}+{\displaystyle \underset{n=M+1}{\overset{k1}{}}}a_na_{kn}`$
$``$ $`2C_s{\displaystyle \underset{n=1}{\overset{kM1}{}}}{\displaystyle \frac{|a_n|}{(kn)^s}}[1,1]+{\displaystyle \underset{n=kM}{\overset{M}{}}}a_na_{kn}`$
$``$ $`{\displaystyle \underset{n=kM}{\overset{k/2}{}}}a_na_{kn}+e(k)a_{k/2}^2+2C_s{\displaystyle \underset{n=1}{\overset{kM1}{}}}{\displaystyle \frac{|a_n|}{(kn)^s}}[1,1].`$
###### Lemma 3.5
Let $`k>2M`$. Then,
$$FS(k)\frac{C_s}{k^{s1}}\left(\frac{2^{s+1}}{2M+1}\underset{n=1}{\overset{M}{}}|a_n|+\frac{C_s4^s}{(2M+1)^{s+1}}+\frac{C_s2^s}{(s1)M^s}\right)[1,1].$$
Proof. From (24) it follows that
$`FS(k)`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{k1}{}}}a_na_{kn}`$
$`=`$ $`2{\displaystyle \underset{n=1}{\overset{k/2}{}}}a_na_{kn}+e(k)a_{k/2}^2`$
$`=`$ $`2{\displaystyle \underset{n=1}{\overset{M}{}}}a_na_{kn}+2{\displaystyle \underset{n=M+1}{\overset{k/2}{}}}a_na_{kn}+e(k)a_{k/2}^2.`$
Each of these terms will be estimated separately. The first one results in:
$`{\displaystyle \underset{n=1}{\overset{M}{}}}a_na_{kn}`$ $``$ $`{\displaystyle \underset{n=1}{\overset{M}{}}}{\displaystyle \frac{|a_n|C_s}{(kM)^s}}[1,1]`$
$``$ $`{\displaystyle \frac{C_s}{k^s(1M/k)^s}}[1,1]{\displaystyle \underset{n=1}{\overset{M}{}}}|a_n|`$
$``$ $`{\displaystyle \frac{2^sC_s}{k^s}}[1,1]{\displaystyle \underset{n=1}{\overset{M}{}}}|a_n|`$
$``$ $`{\displaystyle \frac{2^sC_s}{k^{s1}(2M+1)}}[1,1]{\displaystyle \underset{n=1}{\overset{M}{}}}|a_n|.`$
The second term leads to
$`{\displaystyle \underset{n=M+1}{\overset{k/2}{}}}a_na_{kn}`$ $``$ $`C_s^2{\displaystyle \underset{n=M+1}{\overset{k/2}{}}}{\displaystyle \frac{1}{n^s(kn)^s}}[1,1]`$
$`=`$ $`{\displaystyle \frac{C_s^2}{k^s}}{\displaystyle \underset{n=M+1}{\overset{k/2}{}}}{\displaystyle \frac{1}{n^s(1n/k)^s}}[1,1]`$
$``$ $`{\displaystyle \frac{C_s^22^s}{k^s}}{\displaystyle \underset{n=M+1}{\overset{k/2}{}}}{\displaystyle \frac{1}{n^s}}[1,1]`$
$``$ $`{\displaystyle \frac{C_s^22^s}{k^s}}{\displaystyle _M^{\mathrm{}}}{\displaystyle \frac{dx}{x}}[1,1]`$
$`=`$ $`{\displaystyle \frac{C_s^22^s}{k^s(s1)M^{s1}}}[1,1]`$
$``$ $`{\displaystyle \frac{C_s^22^{s1}}{k^s(s1)M^s}}[1,1].`$
Finally, the third term gives rise to
$$e(k)a_{k/2}^2\frac{C_s^22^{2s}}{k^{2s}}[1,1]\frac{1}{k^{s1}}\frac{C_s^24^s}{(2M+1)^{s+1}}[1,1].$$
Turning now to the infinite sum we can obtain the following estimate.
###### Lemma 3.6
Let $`k>M`$. Then,
$$IS(k)\frac{C_s}{k^{s1}(M+1)}\left(\frac{C_s}{(M+1)^{s1}(s1)}+\underset{n=1}{\overset{M}{}}|a_n|\right)[1,1].$$
Proof. From (25) it follows that
$$IS(k)=\underset{n=1}{\overset{M}{}}a_na_{k+n}+\underset{n=M+1}{\overset{\mathrm{}}{}}a_na_{k+n}.$$
As in the previous case, each term is treated separately.
$`{\displaystyle \underset{n=1}{\overset{M}{}}}a_na_{k+n}`$ $``$ $`C_s{\displaystyle \underset{n=1}{\overset{M}{}}}{\displaystyle \frac{|a_n|}{(k+n)^s}}[1,1]`$
$`=`$ $`{\displaystyle \frac{C_s}{k^s}}{\displaystyle \underset{n=1}{\overset{M}{}}}{\displaystyle \frac{|a_n|}{(1+n/k)^s}}[1,1]`$
$``$ $`{\displaystyle \frac{C_s}{k^{s1}(M+1)}}{\displaystyle \underset{n=1}{\overset{M}{}}}|a_n|[1,1].`$
The remaining term leads to
$`{\displaystyle \underset{n=M+1}{\overset{\mathrm{}}{}}}a_na_{k+n}`$ $`=`$ $`C_s^2{\displaystyle \underset{n=M+1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{m^s(k+m)^s}}[1,1]`$
$``$ $`{\displaystyle \frac{C_s^2}{(M+1)^s}}{\displaystyle _M^{\mathrm{}}}{\displaystyle \frac{dx}{(k+x)^s}}[1,1]`$
$`=`$ $`{\displaystyle \frac{C_s^2}{(M+1)^s(s1)(k+M)^{s1}}}[1,1]`$
$``$ $`{\displaystyle \frac{C_s^2}{(M+1)^s(s1)k^{s1}}}[1,1].`$
### 3.3 Refining the Self-Consistent Bounds
The proof of our results obviously depends on having good self-consistent bounds and the precision of the final result is determined directly by these bounds. For this reason it is important to have a process by which these bounds can be improved. With this in mind consider an initial sequence of bounds $`\{a_k^\pm \}`$ which defines the sets
$$W=\underset{k=1}{\overset{m}{}}[a_k^{},a_k^+]\mathrm{and}V=\underset{k=m+1}{\overset{\mathrm{}}{}}[a_k^{},a_k^+].$$
We will also assume that
$$1<\nu m^2.$$
This condition means that the Fourier modes for $`km`$ are linearly stable.
We shall describe the refinement procedure under the assumption of C4a. In particular, we need that our sequence $`\{a_k^\pm \}`$ satisfy $`\mathrm{dir}(k)=1`$ for all $`k>m`$. We also assume that since we can numerically solve (4), that the estimates for $`W`$ are reasonably good.
We will inductively adjust $`a_k\pm `$, for $`k=m+1,\mathrm{},M`$, beginning with $`a_{m+1}^\pm `$, as follows. Let $`aWV`$ such that $`a_k=a_k^+`$. To satisfy C4a requires that $`\dot{a}_k<0`$, i.e.
$$k^2(1\nu k^2)a_k^+kFS(k)+2kIS(k)<0.$$
This is equivalent to requiring
$$a_k^+>\frac{2IS(k)FS(k)}{k^3(\nu k^2)}.$$
(27)
Of course, our goal is to make $`a_k^+`$ as small as possible within the constraints imposed by the approximations. Since we are iteratively improving our bounds, it is reasonable to assume that a worst case equality is the best guess at this stage in the procedure. Note, we are not claiming a proof at this point, we just are seeking good bounds which later will be verified to be self-consistent bounds. So using Lemma 3.1, define $`f_k^\pm `$ to be bounds for $`2IS(k)FS(k)`$,
$`[f_k^{},f_k^+]:=2{\displaystyle \underset{n=1}{\overset{Mk}{}}}a_na_{k+n}+2C_s{\displaystyle \underset{n=Mk+1}{\overset{M}{}}}{\displaystyle \frac{|a_n|}{(k+n)^s}}[1,1]+`$
$`{\displaystyle \frac{2C_s^2}{(k+M+1)^s(s1)M^{s1}}}[1,1]{\displaystyle \underset{n=1}{\overset{k1}{}}}a_na_{kn}.`$
The new value of $`a_k^+`$ is given by
$$a_k^+:=\frac{f_k^+}{k^3(\nu k^2)}.$$
A similar argument suggests setting
$$a_k^{}:=\frac{f_k^{}}{k^3(\nu k^2)}.$$
This approach works up to $`k=M`$. Recall that for $`k>M`$, we set $`a_k^\pm =\pm C_s/k^s`$. Here our goal is to improve the power of convergence, i.e. we want to increase $`s`$. Again, since we are trying to satisfy C4a the basic inequality which needs to be satisfied is (27). The estimates for $`FS(k)`$ and $`IS(k)`$ for $`k>M`$ obviously are crucial here. However, we had two sets of estimates one for $`M<k2M`$ and the other for $`k>2M`$. Thus, we need to choose the worst of both estimates. This is done as follows.
Given an interval $`I\text{}`$ let
$$|I|:=\underset{xI}{sup}|x|.$$
With the estimate from Lemma 3.4 in mind define
$$D_1(k):=\left|2\underset{n=kM}{\overset{k/2}{}}a_na_{kn}+e(k)a_{k/2}^2\right|+2C_s\underset{n=1}{\overset{kM1}{}}\frac{|a_n|}{(kn)^s}.$$
Combining this with the estimate on $`IS(k)`$ given by Lemma 3.6 and multiplying by $`k^{s1}`$ leads to the following definition
$$D_1:=\frac{2C_s}{(M+1)}\left(\frac{2C_s}{(M+1)^{s1}(s1)}+\underset{n=1}{\overset{M}{}}|a_n|\right)+\underset{M<k2M}{\mathrm{max}}k^{s1}D_1(k).$$
Turning now to the bounds for $`k>2M`$, Lemmas 3.5, 3.6, and again multiplying by $`k^{s1}`$ suggests setting
$`D_2`$ $`:=`$ $`{\displaystyle \frac{2C_s}{(M+1)}}\left({\displaystyle \frac{2C_s}{(M+1)^{s1}(s1)}}+{\displaystyle \underset{n=1}{\overset{M}{}}}|a_n|\right)+`$
$`{\displaystyle \frac{C_s}{k^{s1}}}({\displaystyle \frac{2^{s+1}}{2M+1}}{\displaystyle \underset{n=1}{\overset{M}{}}}|a_n|+{\displaystyle \frac{C_s4^s}{(2M+1)^{s+1}}}+{\displaystyle \frac{C_s2^s}{(s1)M^s}}.)`$
From Lemmas 3.4, 3.5, and 3.6 we obtain the following result.
###### Corollary 3.7
For $`k>M`$ and $`D_s:=\mathrm{max}\{D_1,D_2\}`$
$$\left|FS(k)+2IS(k)\right|<\frac{D_s}{k^{s1}}.$$
We will use this corollary to update the decay rate for the tail terms. Again, we want (27) (which gave us C4a) to hold for all $`k>M`$. It is sufficient that
$$a_k^+>\frac{\left|FS(k)+2IS(k)\right|}{k^3(\nu (M+1)^2)}$$
and therefore it is sufficient that
$$a_k^+>\frac{D_s}{k^{s1}}\frac{1}{k^3}\frac{1}{\nu (M+1)^2}=\frac{1}{k^{s+2}}\frac{D_s}{\nu (M+1)^2}.$$
There is a similar inequality for $`a_k^{}`$.
Setting this to an equality we can define
$$a_k^\pm :=\pm \frac{C_{s+2}}{k^{s+2}}\mathrm{where}C_{s+2}:=\frac{D_s}{\nu (M+1)^2}$$
(28)
for $`k>M`$.
## 4 A Typical Proof
This section describes the proof of the following result.
###### Theorem 4.1
Let
$$u(x)=\frac{1}{\sqrt{2}}\mathrm{sin}x\frac{1}{8}\mathrm{sin}2x.$$
For $`\nu =0.75`$ there exists an equilibrium solution $`u^{}(x)`$ to (1) such that
$$u^{}u_{L^2}<0.052\mathrm{and}u^{}u_{C^0}<0.05$$
The reader should observe that this is a weaker version of Theorem 1.2. However, we present its proof since it contains all the essential features, but with a very low dimensional approximation. The rest of the results described in the introduction were proved in a similar manner.
The first step is to choose $`m`$ the dimension of the Galerkin approximation and $`M`$ the level to which we make specific choices for the $`\{a_k^\pm \}`$. For $`m=2`$, (23) reduces to the system of equations
$`\dot{a}_1`$ $`=`$ $`{\displaystyle \frac{1}{4}}a_1+2a_1a_2`$ (29)
$`\dot{a}_2`$ $`=`$ $`8a_22a_1^2.`$
A simple algebraic computation shows that this system has exactly three fixed points: $`(0,0)`$-unstable, with one-dimensional unstable manifold and two attracting fixed points $`u_\pm =(\pm \frac{1}{\sqrt{2}},\frac{1}{8})`$. Theorem 4.1 is obtained by studying the dynamics of (29) in a neighborhood of $`(\frac{1}{\sqrt{2}},\frac{1}{8})`$.
The next step is to obtain self-consistent apriori bounds for (21). It is unrealistic to expect that goods bounds can be obtained immediately. Thus, we make a reasonable guess for bounds and then try to improve them. Let
$$W=[\frac{1}{\sqrt{2}}0.1,\frac{1}{\sqrt{2}}+0.1]\times [\frac{1}{8}0.1,\frac{1}{8}+0.1].$$
The initial estimates for $`[a_k^{},a_k^+]`$ are given in Table 2. The formula used to derive these initial estimates will be presented in Section 5. The reason for delaying the presentation is to emphasize the fact that the initial estimates are only estimates. Obviously, choosing good estimates allows for faster convergence and choosing terrible estimates will probably result in a failure of convergence.
Beginning with the data in Table 2, the refinement procedure described in Section 3.3 is used to update $`a_k^\pm `$ for $`k>2`$. After three iterations one obtains the estimates given in Table 3. It can now be checked that $`W`$ and $`\{a_k^\pm \}`$ from Table 3 form self-consistent bounds.
What should be clear at this point is that the uncertainty contributed by the terms not in the Galerkin projection are extremely small. Obviously, at this point most of the uncertainty is due to the size of $`W`$.
Having controlled the errors from the Galerkin truncation, the next step is to obtain an isolating block $`N`$ which is topologically self consistent with $`W`$ and $`\{a_k\}`$. In determing $`N`$ we use the vector field (29). Of course, the correct equations are given by (21):
$`\dot{a}_1`$ $`=`$ $`{\displaystyle \frac{1}{4}}a_1+2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}a_na_{n+1}`$
$`\dot{a}_2`$ $`=`$ $`8a_k2a_1^2+4{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}a_na_{n+2}.`$
Using Corollary 3.3 one obtains the following bounds on the errors, $`ϵ_i`$, $`i=1,2`$,
$$ϵ_1=[0.00955626,7.0600510^{10}]\mathrm{and}ϵ_2=[1.54410^8,0.0697171].$$
Thus, the equations for which the isolating neighborhood should be found are
$`\dot{a}_1`$ $`=`$ $`{\displaystyle \frac{1}{4}}a_1+2a_1a_2+ϵ_1`$ (30)
$`\dot{a}_2`$ $`=`$ $`8a_22a_1^2+ϵ_2.`$
The construction of the isolating block around $`(1/\sqrt{2},1/8)`$ is easier if one works in coordinates determined by the eigenfunctions of the linearized equations at the fixed point. The eigenvalues are
$$\lambda _1=4+2\sqrt{3},\lambda _2=42\sqrt{3}$$
with corresponding unit eigenvectors
$$v_1=\frac{1}{\sqrt{158\sqrt{3}}}(1,\frac{\lambda _1}{\sqrt{2}})^t,v_2=\frac{1}{\sqrt{15+8\sqrt{3}}}(1,\frac{\lambda _2}{\sqrt{2}})^t.$$
Let $`T`$ be the affine change of variables from $`(x,y)^t`$ in this new basis to the original variables $`(a_1,a_2)^t`$. Then, on the set $`T^1(W)`$ (30) becomes
$`\dot{x}=(4+2\sqrt{3})x+f_1(x,y)+\stackrel{~}{ϵ}_1`$ (31)
$`\dot{y}=(42\sqrt{3})y+f_2(x,y)+\stackrel{~}{ϵ}_2`$
where $`f_1`$ and $`f_2`$ are polynomials containing only terms of degree two and $`\stackrel{~}{ϵ}_1`$, $`\stackrel{~}{ϵ}_2`$ are obtained from $`ϵ_1`$, $`ϵ_2`$ by the transformation $`T`$. In particular,
$$\stackrel{~}{ϵ}_1=[0.0110098,0.0152184],\stackrel{~}{ϵ}_2=[0.0764461,0.00397075]$$
Set $`\stackrel{~}{W}=[0.0748016,0.0748016]^2`$, then $`T(\stackrel{~}{W})W`$. Thus, an isolating block $`\stackrel{~}{N}\stackrel{~}{W}`$ (which satisfies the error constraints) will give rise to an isolating block $`N=T(\stackrel{~}{N})W`$ such that $`N`$, $`W`$ and $`\{a_k\}`$ are topologically consistent. Since for each $`k`$, $`\mathrm{dir}(k)=1`$,
$$CH_j(\mathrm{Inv}(N))\{\begin{array}{cc}\text{}\hfill & \text{if }j=0\text{,}\hfill \\ 0\hfill & \text{otherwise.}\hfill \end{array}$$
Thus, by Theorem 2.16, there exists the desired fixed point in
$$u^{}N\times \underset{k=2}{\overset{\mathrm{}}{}}[a_k^{},a_k^+].$$
Thus, all that remains is to construct $`\stackrel{~}{N}`$. In the $`(x,y)`$ coordinates the error constraints become
$`f_1(x,y)+\stackrel{~}{ϵ}_1`$ $``$ $`(b_x,B_x):=(0.0334071,0.0376157)\text{for }(x,y)\stackrel{~}{W}`$
$`f_2(x,y)+\stackrel{~}{ϵ}_2`$ $``$ $`(b_y,B_y):=(0.0764461,0.00397075)\text{for }(x,y)\stackrel{~}{W}`$
This implies that bounds on the derivative are given by
$$\lambda _1(x+\frac{b_x}{\lambda _1})<\dot{x}<\lambda _1(x+\frac{B_x}{\lambda _1})\mathrm{and}\lambda _2(y+\frac{b_y}{\lambda _2})<\dot{y}<\lambda _2(y+\frac{B_y}{\lambda _2}).$$
Observe that $`\lambda _i<0`$, hence the box
$`\stackrel{~}{N}`$ $`=`$ $`[{\displaystyle \frac{b_x}{\lambda _1}},{\displaystyle \frac{B_x}{\lambda _1}}]\times [{\displaystyle \frac{b_y}{\lambda _2}},{\displaystyle \frac{B_y}{\lambda _2}}]`$
$`=`$ $`[0.0623385,0.0701918]\times [0.0132425,0.00353264]\stackrel{~}{W}`$
is an isolating block.
## 5 Obtaining the Initial Estimates
Before beginning this section we want to once again emphasize that the proofs of the theorems in this paper are in principle independent of this section. On the other hand, good initial guesses greatly improve the speed of convergence. The estimates described in what follows apparently provide excellent initial values for the self-consistent bounds.
We will follow the arguments from to produce estimates for errors in the Galerkin projection. Kuramoto-Sivashinsky can be written in the form
$$u_t+\nu AuA^{1/2}u+2B(u,u)=0$$
where
$$A=\frac{^4}{x^4},A^{1/2}=\frac{^2}{x^2},B(u,v)=u\frac{v}{x}.$$
While $`A^{1/4}\frac{}{x}`$, it is still the case that
$$|A^{1/4}u|_2=|\frac{u}{x}|_2.$$
To simplify the notation, let
$$|u|=|u|_2,u=|A^{1/4}u|_2.$$
Since we are interested in bounded invariant sets we can without loss of generality assume the following apriori bounds for the invariant set under consideration:
$$|u(t)|\rho _0,u(t)\rho _1,\text{for }t>T(u(0))$$
(32)
We will make use of the following inequality \[10, Lemma 1.4\]
$$|(B(u,v),w)|\sqrt{2}|u|^{1/2}u^{1/2}v|w|$$
(33)
The eigenvalues of $`A`$ are $`\lambda _n=n^4`$, $`n\text{}`$ and the corresponding complete family of orthonormal eigenfunctions are $`\{\frac{1}{\sqrt{\pi }}\mathrm{sin}(nx)\}`$. Let $`P=P_m`$ be an orthogonal projection on first $`m`$ eigenfunctions and set $`Q=IP`$.
Using the decomposition
$$p=Pu,q=Qu$$
and the same abuse of notation the equation for $`q`$ is
$$\dot{q}=\nu Aq+A^{1/2}q2QB(u,u).$$
(34)
###### Theorem 5.1
Under the assumptions stated above if $`m`$ is large enough such that $`\lambda _{m+1}>\frac{1}{\nu ^2}`$, then
$$\underset{t\mathrm{}}{lim\; sup}|q(t)|\frac{2\sqrt{2}\rho _0^{1/2}\rho _1^{3/2}}{\lambda _{m+1}(\nu \lambda _{m+1}^{\frac{1}{2}})}.$$
Proof: Beginning with (34) and taking a scalar product with $`q`$ gives
$$(\frac{dq}{dt}|q)=\nu (Aq|q)+(A^{1/2}q|q)2(B(u,u)|q).$$
Therefore,
$$\frac{1}{2}\frac{d}{dt}|q|^2\nu (Aq|q)+(A^{1/2}q|q)+2|(B(u,u)|q)|$$
(35)
Observe that
$$(Aq|q)=(A^{1/2}q|A^{1/2}q)=|A^{1/2}q|^2$$
(36)
and
$`(A^{1/2}q|q)`$ $`=`$ $`{\displaystyle \underset{n=m+1}{\overset{\mathrm{}}{}}}\lambda _n^{1/2}|q_n|^2`$ (37)
$`=`$ $`\lambda _{m+1}^{1/2}{\displaystyle \underset{n=m+1}{\overset{\mathrm{}}{}}}\lambda _{m+1}^{1/2}\lambda _n^{1/2}|q_n|^2`$
$``$ $`\lambda _{m+1}^{1/2}{\displaystyle \underset{n=m+1}{\overset{\mathrm{}}{}}}\lambda _n|q_n|^2`$
$`=`$ $`\lambda _{m+1}^{1/2}|A^{1/2}q|^2`$
Thus, (35) (36) and (37) imply that
$$\frac{d}{dt}|q|^22\nu |A^{1/2}q|^2+2\lambda _{m+1}^{1/2}|A^{1/2}q|^2+4|(B(u,u)|q)|.$$
(38)
From (33) it follows that
$$|(B(u,u),q)|\sqrt{2}|u|^{1/2}u^{3/2}|q|.$$
While, (32) implies that for $`t>T(u)`$
$$|(B(u,u),q)|\sqrt{2}\rho _0^{1/2}\rho _1^{3/2}|q|$$
Therefore, (38) becomes
$$\frac{d}{dt}|q|^22(\nu \lambda _{m+1}^{1/2})|A^{1/2}q|^2+4\sqrt{2}\rho _0^{1/2}\rho _1^{3/2}|q|$$
Observe that $`|A^{1/2}q|^2\lambda _{m+1}|q|^2`$ and by assumption $`\nu \lambda _{m+1}^{1/2}>0`$. Thus,
$`{\displaystyle \frac{d}{dt}}|q|^2`$ $``$ $`2(\nu \lambda _{m+1}^{1/2})\lambda _{m+1}|q|^2+4\sqrt{2}\rho _0^{1/2}\rho _1^{3/2}|q|`$
$`=`$ $`(4\sqrt{2}\rho _0^{1/2}\rho _1^{3/2}2(\nu \lambda _{m+1}^{1/2})\lambda _{m+1}|q|)|q|`$
Thus, for
$$|q|>\frac{4\sqrt{2}\rho _0^{1/2}\rho _1^{3/2}}{2(\nu \lambda _{m+1}^{1/2})\lambda _{m+1}}$$
$`\frac{d}{dt}|q|^2<0`$ and hence,
$$\underset{t>\mathrm{}}{lim\; sup}|q(t)|\frac{4\sqrt{2}\rho _0^{1/2}\rho _1^{3/2}}{2(\nu \lambda _{m+1}^{1/2})\lambda _{m+1}}.$$
###### Remark 5.2
Because an orthonormal collection eigenvectors were used for the calculations in this section, the coefficients $`q_k`$ and $`a_k`$ differ by a scaling, i.e.
$$q_k=2\sqrt{\pi }a_n.$$
Theorem 5.1 can be used as follows. For a fixed $`\nu `$, numerical experiments can suggest values for $`\rho _0`$ and $`\rho _1`$. For $`m<kM`$ one can use the formula
$$a_k^\pm :=\mathrm{min}\{\pm \rho _0,\pm \frac{\rho _1}{k},\pm \frac{4\sqrt{2\pi \rho _0\rho _1^3}}{k^4(\nu k^{}2)}\}.$$
For $`k>M`$, one defines
$$C_4:=\frac{4\sqrt{2\pi \rho _0\rho _1^3}}{\nu (M+1)^2}.$$
|
warning/0005/math0005208.html
|
ar5iv
|
text
|
# On the consistency of the definable tree property on ℵ₁
## 1. Introduction
A well known result of Aronszajn is the existence of an Aronszajn tree on $`\mathrm{}_1`$, i.e. there exist an $`\omega _1`$tree with no uncountable branch. The construction uses the axiom of choice and therefore does not give a definable such tree.
Sierpinski and Kurepa proved that if Ramsey theorem holds for $`\kappa `$ then $`\kappa `$ is a strong limit cardinal. Erdős proved that such a $`\kappa `$ is inaccessible. They also provided counterexamples to Ramsey theorem on small cardinals. These counterexamples explicitly used a well-ordering of $`P(\kappa )`$. The proof raises the question whether we can find a definable counterexample.
In a straightforward generalization of Aronszajn’s proof, Specker proved the existence of Aronszajn trees for every succesor $`\kappa ^+`$ s.t $`\kappa ^{<\kappa }=\kappa `$. This raised the question whether the GCH was needed for this result. Mitchell and Silver have proved that the tree property on $`\mathrm{}_2`$ is equiconsistent with the existence of a weakly compact cardinal. Magidor and Shelah have proved the consistency of the tree property on $`\mathrm{}_{\omega +1}`$ from the existence of very large cardinals.
Mitchell’s forcing for the tree property works well if one tries to obtain the tree property on two non-consecutive cardinals (e.g. $`\mathrm{}_2`$ and $`\mathrm{}_4`$). However, his methods fail to prove the consistency of the tree property on $`\mathrm{}_2`$ and $`\mathrm{}_3`$ together. By a result of Magidor this is indeed a substential difficulty since the consistency strength of the tree property on $`\mathrm{}_2`$ and $`\mathrm{}_3`$ together is much higher then a weakly compact (e.g. it implies the existence of $`0^{\mathrm{}}`$). Abraham has proved the consistency of the tree property on both $`\mathrm{}_2`$ and $`\mathrm{}_3`$ from a supercompact cardinal and a weakly compact cardinal above it. This result was extended by Cummings and Foremann which proved the consistency of the tree property on all $`\mathrm{}_n`$’s from many supercompacts.
Kunen, and Shelah and Harrington considered the consistency strength of obtaining Lebesgue measurability of projective sets together with Martin’s axiom. They proved the equiconsistency of this theory with the existence of weakly compact cardinals. This shows that adding Martin’s axiom to projective measurability increases the consistency strength.
In this paper we consider the consistency strength of the definable tree property on $`\mathrm{}_1`$, i.e. the existence of a model in which every $`\omega _1`$tree which is first order definable (with parameters) over $`(H_{\omega _1},\epsilon )`$, has a cofinal branch. Our proof method answers also a related question regarding definable counterexamples to Ramsey theorem on $`\mathrm{}_1`$.
In section (2) we define the exact meaning of a definable $`\kappa `$tree, and study some variations used in this paper. In section (3) we define the $`\mathrm{\Pi }_1^1`$ reflecting cardinals and derive an extension property which is similar to the extension property of weakly compact cardinal. We also bound the consistency strength of the existence of a $`\mathrm{\Pi }_1^1`$ reflecting cardinal by the existence of a Mahlo cardinal. In section (4) we show that by forcing with the well known Levy collapse of a $`\mathrm{\Pi }_1^1`$ reflecting cardinal $`\kappa `$ to $`\mathrm{}_1`$, we obtain a model of the definable tree property. We also prove that in this model a definable Ramsey theorem on $`\mathrm{}_1`$ holds. In section (5) we show that our assumptions on $`\kappa `$ are necessary by proving that if the definable tree property on $`\mathrm{}_1`$ holds, then
$$L\mathrm{}_1\text{ is a }\mathrm{\Pi }_1^1\text{ reflecting cardinal}.$$
In section (6) we prove the equiconsistency of the definable tree property on $`\mathrm{}_1+`$ MA with the existence of a weakly compact cardinal. This is done by exploiting the methods of Kunen and Shelah and Harrington . We also comment that adding any reasonable failure of GCH does not add to the consistency strength. Finally in section (7) we comment on the consistency of the definable tree property on higher cardinals with GCH. We describe a forcing to get the definable tree property on all $`\mathrm{}_n`$’s together with GCH by using $`\omega `$ many $`\mathrm{\Pi }_1^1`$ reflecting cardinals.
## 2. Definable $`\kappa `$trees
In this section we define various notions of definable $`\kappa `$trees and study the relationship between these notions. First define the usual notion of a $`\kappa `$tree.
###### Definition 2.1.
* A tree is a partially ordered set$`(T,<_T)`$ such that for any $`tT`$ the set $`\{sT|s<_Tt\}`$ of predecessors of $`t`$ is well-ordered under $`<_T`$, and there is a root $`rT`$ such that for any $`tT`$, such that $`tr`$, $`r<_Tt`$.
* The $`\alpha `$’th level of $`T`$ denoted by $`T_\alpha `$ is the set of elements of T whose set of $`T`$predecessors has order type $`\alpha `$.
* A tree $`(T,<_T)`$ is a $`\kappa `$tree if $`|T|=\kappa `$ and for every $`\alpha `$ $`|T_\alpha |<\kappa `$.
* A branch is a $`<_T`$ linearly ordered subset of $`T`$.
* A cofinal branch is a branch which intersects every level of $`T`$.
Next we would like to give a definition of a definable $`\kappa `$tree. Three different notions naturally arise
###### Definition 2.2.
* A $`\kappa `$tree is definable in the strict sense if its underlying set is $`\kappa `$, and $`<_T`$ is $`𝚺_\omega \left((H_\kappa ,)\right)`$.
* A $`\kappa `$tree is definable in the wide sense if its underlying set $`T`$ and $`<_T`$ are both $`𝚺_\omega \left((H_\kappa ,)\right)`$ and $`T`$ has definable cardinality $`\kappa `$, i.e. there is a bijection $`f:\kappa T`$ which is $`𝚺_\omega \left((H_\kappa ,)\right)`$.
* A $`\kappa `$tree is definable in the very wide sense if its underlying set $`T`$ and $`<_T`$ are both $`𝚺_\omega \left((H_\kappa ,)\right)`$.
Obviously being definable in the wide sense implies being definable in the very wide sense, and being definable in the strict sense implies being definable in the wide sense. The following proposition states that being definable in the wide sense is almost equivalent to being definable in the strict sense.
###### Proposition 2.1.
If $`(T,<_T)`$ is a $`\kappa `$tree definable in the wide sense then there is a tree $`(\kappa ,<_T^{})`$ isomorphic to $`(T,<_T)`$ which is definable in the strict sense.
Proof :Let $`f`$ be the definable bijection $`f:\kappa T`$, and let $`\psi (x,y,z)`$ be the definition of $`<_T`$ from the parameter $`z`$. Now define $`<_T^{}`$ by $`\alpha <_T^{}\beta `$ iff $`\psi (f(\alpha ),f(\beta ),z)`$. The rest trivially follows.
Using this proposition we will not distinguish between the strict and the wide notions of definability. Also notice that if there is a well-ordering of $`H_\kappa `$ which is $`𝚺_\omega \left((H_\kappa ,)\right)`$ then the last two notions coincide. However in general this is not the case, and we will distinguish between definability in the very wide sense and definability in the strict sense which we will adopt as the definition of definability. So from now on, a definable $`\kappa `$tree is a tree definable in the strict sense.
## 3. $`\mathrm{\Pi }_1^1`$ reflecting cardinals
Let $`\kappa `$ be a cardinal. We say that $`\kappa `$ is $`\mathrm{\Pi }_n^m`$ reflecting, if $`\kappa `$ is inaccessible and for every $`AV_\kappa `$ definable over $`V_\kappa `$ (with parameters) and for every $`\mathrm{\Pi }_n^m`$ sentence $`\mathrm{\Phi }`$, such that
$$(V_\kappa ,\epsilon ,A)\mathrm{\Phi }$$
there is an $`\alpha <\kappa `$ such that
$$(V_\alpha ,\epsilon ,AV_\alpha )\mathrm{\Phi }.$$
The $`\mathrm{\Pi }_n^m`$ reflecting cardinals are a lightface analog of the $`\mathrm{\Pi }_n^m`$ indescribable cardinals. However they have a much weaker consistency strength. If $`\kappa `$ is $`\mathrm{\Pi }_1^1`$ reflecting, then an easy consequence is the fact that $`\kappa `$ is the $`\kappa ^{}th`$ inaccessible cardinal, since inaccessibility is a $`\mathrm{\Pi }_1^1`$ property, and for each $`\alpha <\kappa `$ being the $`\alpha ^{}th`$ inaccessible is also a $`\mathrm{\Pi }_1^1`$ property. Next we prove the following lemma:
###### Lemma 3.1.
Suppose $`\kappa `$ is a Mahlo cardinal then for every $`m,n`$,
$`\{\alpha <\kappa |\alpha \text{ is a }\mathrm{\Pi }_n^m\text{ reflecting cardinal}\}`$ is stationary.
Proof :Let $`\kappa `$ be Mahlo, and let $`C`$ be a club. We shall find a $`\mu C`$ such that $`\mu `$ is a $`\mathrm{\Pi }_n^m`$ reflecting cardinal. Let $`e:(\mathrm{\Phi },\psi ,a)\kappa `$ be an enumeration of triples $`(\mathrm{\Phi },\psi ,a)`$ s.t. $`\mathrm{\Phi }`$ is a $`\mathrm{\Pi }_n^m`$ sentence, $`\psi (x,a)`$ is a first order formula in the free variable $`x`$, $`aV_\kappa `$. Also assume without loss of generality that $`aV_{e(\mathrm{\Phi },\psi ,a)}`$ and $`e(\mathrm{\Phi },\psi ,a)<rank(a)^{\mathrm{}}`$, where $`\alpha ^{\mathrm{}}`$ is the least inaccessible above $`\alpha `$, and that $`e`$ is $`11`$ and onto $`\kappa `$. For each triple $`(\mathrm{\Phi },\psi ,a)`$ define $`g(e(\mathrm{\Phi },\psi ,a))`$ by:
(3.1)
$$g(e(\mathrm{\Phi },\psi ,a))=\{\begin{array}{c}\text{ The least }\rho C\text{ such that }e(\mathrm{\Phi },\psi ,a)<\rho \text{ and }\hfill \\ (V_\rho ,\epsilon ,S)(xS\text{ iff }\psi (x,a))\mathrm{\Phi },\hfill \\ \text{ if there is such a }\rho .\hfill \\ \\ e(\mathrm{\Phi },\psi ,a)+1\text{ otherwise}.\hfill \end{array}$$
$`g:\kappa \kappa `$ satisfies for every $`\alpha `$, $`\alpha <g(\alpha )`$ by its definition. Since $`\kappa `$ is Mahlo the set of inaccessibles is stationary so there is an inaccessible $`\mu C`$ s.t. $`g^{\prime \prime }\mu \mu `$. We claim that $`\mu `$ is $`\mathrm{\Pi }_n^m`$ reflecting. By definition $`\mu `$ is inaccessible. Suppose that $`SV_\mu `$ is defined by $`xS\text{ iff }V_\kappa \psi (x,a)`$. Assume that $`(V_\mu ,\epsilon ,S)\mathrm{\Phi }(S)`$. Since $`\mu C`$ there is a $`\rho C`$ s.t. $`(V_\rho ,\epsilon ,S\rho )\mathrm{\Phi }`$ and $`S\rho `$ is defined by $`\psi (x,a)`$. Hence $`g(e(\mathrm{\Phi },\psi ,a))`$ is defined by the first case, since $`g^{\prime \prime }\mu \mu `$ and $`e(\mathrm{\Phi },\psi ,a)<\mu `$. Now by the definition of $`e`$, we obtain that $`g(e(\mathrm{\Phi },\psi ,a))<\mu `$. Hence $`\mathrm{\Phi }`$ reflects to $`(V_{g(e(\mathrm{\Phi },\psi ,a))},\epsilon ,Sg(e(\mathrm{\Phi },\psi ,a)))`$.
The following theorem is the analog of the Hanf-Scott theorem on the indescribability of weakly compact cardinals.
###### Definition 3.1.
$`\kappa `$ has the extension property iff for every $`n`$ and for every $`AV_\kappa `$ first order definable over $`V_\kappa `$ with parameters from $`V_\kappa `$, there is a transitive ser $`X`$, and $`A^XX`$ such that $`\kappa X`$ and $`(V_\kappa ,\epsilon ,A)_n(X,\epsilon ,A^X)`$.
Remark. One can show that $`\kappa `$ has the extension property if and only if for every $`n`$ there is a transitive $`\mathrm{\Sigma }_n`$-elementary end extension of $`V_\kappa `$. containing $`\kappa `$, that is, $`\kappa X`$ and $`(V_\kappa ,\epsilon )_n(X,\epsilon )`$.
Note that the fact that $`\kappa `$ has the extension property is $`\mathrm{\Sigma }_1^1`$.
###### Theorem 3.2.
A cardinal $`\kappa `$ is $`\mathrm{\Pi }_1^1`$ reflecting iff $`\kappa `$ is inaccessible and has the extension property.
The proof is identical to the proof of the Hanf-Scott theorem (see pp. 59-60).
Proof :
Assume that $`\kappa `$ is inaccessible and has the extension property. Let $`AV_\kappa `$ be definable in $`V_\kappa `$. Let $`\mathrm{\Phi }`$ be a $`\mathrm{\Pi }_1^1`$ formula, say of the form $`Y\varphi (Y)`$, where $`\varphi (Y)`$ is first order formula in the predicates $`Y`$ and $`A`$, and assume that $`(V_\kappa ,\epsilon ,A)\mathrm{\Phi }`$. Let $`n`$ be large enough so that the sentence $`\alpha \left(YV_{\alpha +1}\left((V_\alpha ,\epsilon ,AV_\alpha )\varphi (Y)\right)\right)`$ is $`\mathrm{\Sigma }_n`$ in $`L_\epsilon (A)`$, the ambient language of $`(V_\kappa ,\epsilon ,A)`$. Now let $`(X,\epsilon ,A^X)`$ be given by the extension property, that is,
(3.2)
$$(V_\kappa ,\epsilon ,A)_n(X,\epsilon ,A^X)$$
where $`X`$ is transitive and $`\kappa X`$. Therefore $`A^XV_\kappa =A`$. Also note that $`V_{\kappa +1}^XV_{\kappa +1}`$ and $`V_\kappa ^X=V_\kappa `$ (for $`n`$ large enough). Now since $`(V_\kappa ,\epsilon ,A)\mathrm{\Phi }`$ it follows that
$$(X,\epsilon ,A^X)YV_{\kappa +1}\left((V_\kappa ^X,\epsilon ,A^XV_\kappa ^X)\varphi (Y)\right)$$
and thus
$$(X,\epsilon ,A^X)\alpha \left(YV_{\alpha +1}\left((V_\alpha ,\epsilon ,A^XV_\alpha )\varphi (Y)\right)\right).$$
Hence by (3.2)
$$(V_\kappa ,\epsilon ,A)\alpha ((V_\alpha ,\epsilon ,AV_\alpha )\mathrm{\Phi }).$$
Hence
$$(V_\alpha ,\epsilon ,AV_\alpha )\mathrm{\Phi }.$$
Thus $`\kappa `$ is $`\mathrm{\Pi }_1^1`$ reflecting.
For the other direction let $`\kappa `$ be a $`\mathrm{\Pi }_1^1`$ reflecting cardinal, and fix some $`n<\omega `$. Let $`\sigma `$ be a $`\mathrm{\Pi }_1^1`$ formula expressing the failure of the extension property relative to $`n`$, that is there is no transitive $`\mathrm{\Sigma }_n`$-elementary extension $`X`$ of $`V_\kappa `$ containing $`\kappa `$ (see the above remark). Let $`\tau `$ be a $`\mathrm{\Pi }_1^1`$ formula expressing the inaccessibility of $`\kappa `$. Let
$$C=\{\alpha <\kappa |(V_\alpha ,\epsilon )_n(V_\kappa ,\epsilon )\}.$$
$`C`$ is a club subset of $`\kappa `$ which is $`\mathrm{\Sigma }_{n+1}`$ definable in $`V_\kappa `$ . Let $`\psi `$ be the formula expressing the fact that $`C`$ is a club. Therefore,
$$(V_\kappa ,\epsilon ,C)\sigma \tau \psi .$$
By $`\mathrm{\Pi }_1^1`$ reflection there exists an $`\alpha `$ such that
$$(V_\alpha ,\epsilon ,C\alpha )\sigma \tau \psi .$$
Hence $`\alpha `$ is an inaccessible cardinal below $`\kappa `$ which is a limit point of $`C`$. Hence $`\alpha C`$, so $`(V_\alpha ,\epsilon )_n(V_\kappa ,\epsilon )`$. Let $`X_0=V_\alpha \{V_\alpha \}`$ and construct an elementary submodel $`(X^{},E)(V_\kappa ,\epsilon )`$, containing $`X_0`$ of cardinality $`\alpha `$. Let $`(X^{},E)`$ be the Skolem hull of $`X_0`$ inside $`(V_\kappa ,\epsilon )`$. Let $`(X,\epsilon )`$ be the transitive collapse of $`(X^{},E)`$. It follows that $`(V_\alpha ,\epsilon )_n(X,\epsilon )`$, and $`V_\alpha X`$ hence
$$(V_\alpha ,\epsilon )\neg \sigma $$
and this is a contradiction.
We finish this section with the following observation on $`\mathrm{\Pi }_1^1`$ reflecting cardinals:
###### Lemma 3.3.
Let $`\kappa `$ be a $`\mathrm{\Pi }_1^1`$ reflecting cardinal. Let $`T`$ be a $`\kappa `$tree definable in the very wide sense, then $`T`$ has a cofinal branch.
Proof :Note that since $`\kappa `$ is inaccessible $`V_\kappa =H_\kappa `$. Let $`T`$ be a definable $`\kappa `$-tree definable in the very wide sense. Let $`n`$ be large enough such that the assertion “$`\alpha T\text{ has a cofinal branch of length }\alpha \text{.}`$” is $`𝚺_n`$ over $`(V_\kappa ,\epsilon ,T)`$. By theorem 3.2 there is a transitive structure $`(X,\epsilon ,T)`$ such that
(3.3)
$$(V_\kappa ,\epsilon ,T)_n(X,\epsilon ,T^X).$$
Since (for $`n`$ large enough) $`V_\kappa ^X=V_\kappa X`$, it follows that $`T^XV_\kappa ^X=T`$. Now, since $`T`$ is a $`\kappa `$-tree it follows that
$$(V_\kappa ,\epsilon ,T)\alpha T\text{ has a cofinal branch of length }\alpha \text{.}$$
Therefore,
$$(X,\epsilon ,T^X)\alpha T^X\text{ has a branch of length }\alpha \text{.}$$
Since $`\kappa X`$ we see that $`(X,\epsilon ,T^X)T^X\text{ has a branch }b\text{ of length }\kappa `$. Because $`T^XV_\kappa =T`$, this branch $`b`$ is really a cofinal branch through $`T`$.
## 4. The forcing construction
In this section we describe the forcing construction and prove that the extended model satisfies the tree property for $`\omega _1`$trees first order definable over $`(H_{\omega _1},\epsilon )`$. Let $`\kappa `$ be a $`\mathrm{\Pi }_1^1`$ reflecting cardinal in $`V`$. Let
$$\text{P}=\text{ Coll }(\omega ,<\kappa )$$
be the Levy collapse of $`\kappa `$ to $`\omega _1`$ and for every $`\alpha <\kappa `$ let
(4.1)
$$\text{P}_\alpha =\text{ Coll }(\omega ,<\alpha )$$
be an initial segments of the forcing. Let $`G`$ be a P generic filter. Our main theorem is
###### Theorem 4.1.
$`V[G]`$ “every definable $`\omega _1`$tree T has a cofinal branch”. Moreover if $`V=L`$ then $`V[G]`$ “every $`\omega _1`$tree T definable in the wide sense over $`(H_{\omega _1},\epsilon )`$ has a cofinal branch”.
Define the following definable partition relation :
###### Definition 4.1.
$`\mathrm{}_1\stackrel{def}{}(\mathrm{}_1)_\alpha ^m`$ iff every partition of $`[\mathrm{}_1]^m`$ into $`\alpha `$ sets which is first order definable over $`(H_{\omega _1},\epsilon )`$ (with parameters in $`H_{\omega _1}`$), has a homogeneous set of size $`\mathrm{}_1`$.
In order to prove theorem 4.1 we shall first prove a definable Ramsey theorem in $`V[G]`$.
###### Lemma 4.2.
$`V[G]\mathrm{}_1\stackrel{def}{}(\mathrm{}_1)_2^2`$
Proof :Let $`F:[\mathrm{}_1]^22`$ be a definable function in $`V[G]`$, defined by
(4.2)
$$F(\{\alpha ,\beta \})=i\mathrm{\Phi }(x,\alpha ,\beta ,i)$$
where $`\mathrm{\Phi }`$ is a $`𝚺_n`$ formula relativized to $`(H_{\omega _1})^{V[G]}`$, and the parameter $`x`$ can be taken as a real, that is, a function $`x:\omega \omega `$. By the $`\kappa `$c.c. of the Levy collapse for every real $`xV[G]`$ there is an $`\epsilon <\kappa `$ such that $`xV[G_\epsilon ]`$, where $`G_\epsilon `$ is an initial segment of the generic, which is generic for $`\text{P}_\epsilon =\text{ Coll }(\omega ,<\epsilon )`$. Moreover by a result of Solovay (see proposition 10.21) $`V[G]=V[G_\epsilon ][H]`$, and $`H`$ is generic for the Levy collapse $`\text{ Coll }(\omega ,<\kappa )`$. By the homogeneity of the Levy collapse every set of ordinals definable in $`V[G]`$ with parameters in $`V[G_\epsilon ]`$ is definable in $`V[G_\epsilon ]`$, and for every formula $`\mathrm{\Psi }`$ we can compute another formula $`\mathrm{\Phi }`$ such that:
(4.3)
$$V[G]\mathrm{\Psi }(x,\alpha ,\beta ,i)V[G_\epsilon ]\mathrm{\Phi }(x,\alpha ,\beta ,i)$$
for all $`\alpha ,\beta <\kappa `$ and $`i\{0,1\}`$. Fix such an $`\epsilon `$. Let $`\stackrel{~}{\sigma }`$ be a $`𝐏_\epsilon `$ name for $`x`$. Now we define a $`\kappa `$tree $`T`$ which is definable in $`V_\kappa `$. For each $`\alpha <\kappa `$ let
(4.4)
$$\begin{array}{cc}\stackrel{~}{h}T_\alpha & \\ & \stackrel{~}{h}\text{ is a }𝐏_\epsilon \text{ name for a function from }\alpha \text{ to }\{0,1\}\text{ and }\\ & \mu _0\mu \delta >\mu \beta <\alpha p𝐏_\epsilon \\ & p_{𝐏_\epsilon }\mathrm{\Phi }(\stackrel{~}{\sigma },\beta ,\mu _0,\stackrel{~}{h}(\beta ))p_{𝐏_\epsilon }\mathrm{\Phi }(\stackrel{~}{\sigma },\beta ,\delta ,\stackrel{~}{h}(\beta ))\end{array}$$
and define $`T=_{\alpha <\kappa }T_\alpha `$. We shall write $`\stackrel{~}{h}_\alpha `$ to denote that $`\stackrel{~}{h}_\alpha T_\alpha `$.
The ordering of $`T`$ is
(4.5)
$$\stackrel{~}{h}_\alpha \stackrel{~}{h}_\beta _{𝐏_\epsilon }\stackrel{~}{h}_\alpha \stackrel{~}{h}_\beta $$
The tree $`T`$ is a $`\kappa `$tree. To prove this first observe that for every $`\alpha <\kappa `$ $`T_\alpha \mathrm{}`$ since for every $`\alpha `$ there are at most $`2^{2^{\alpha +|𝐏_\epsilon |}}`$ many $`𝐏_\epsilon `$ names for such functions, and therefore we are partitioning $`\kappa `$ into less than $`\kappa `$ many subsets according to the possible values of $`F(\{\beta ,\delta \}):\beta <\alpha `$. Secondly $`|T_\alpha |2^\alpha <\kappa `$. $`T`$ is definable in the very wide sense in $`H_\kappa `$ by (4.4). Therefore by lemma (3.3) $`T`$ has a cofinal branch $`\stackrel{~}{h}_\alpha :\alpha <\omega _1`$. Work now in $`V[G]`$. Let $`\stackrel{~}{h}_\alpha (G_\epsilon )`$ denote the realization of the name $`\stackrel{~}{h}_\alpha `$ in $`V[G_\epsilon ]`$. Let $`h=_{\alpha <\kappa }\stackrel{~}{h}_\alpha (G_\epsilon )`$, h is a function from $`\kappa =\mathrm{}_1^{V[G]}`$ to $`\{0,1\}`$. Define
(4.6)
$$A_\alpha =\{\alpha <\gamma <\mathrm{}_1|\beta <\alpha F(\{\beta ,\gamma \})=h(\beta )\}$$
then for every $`\alpha `$, $`|A_\alpha |=\mathrm{}_1`$ by the definition of $`\stackrel{~}{h}_\alpha `$. Moreover $`A_\alpha :\alpha <\mathrm{}_1`$ is a decreasing sequence of sets. We construct $`H_0`$ by induction on $`\alpha <\mathrm{}_1`$. Let $`\beta _0=0`$. For each $`\alpha `$ let $`\gamma _\alpha =sup\left\{\beta _i\right|i<\alpha \}`$. Let $`\beta _\alpha =\mathrm{min}A_{\gamma _\alpha }`$. Let $`H_0=\left\{\beta _\alpha \right|\alpha <\mathrm{}_1\}`$. By the definition of $`H_0`$, for every $`\alpha <\beta H_0`$ we have $`F(\{\alpha ,\beta \})=h(\alpha )`$. Let $`l`$ be minimal such that $`|h^1(l)H_0|=\mathrm{}_1`$. Now $`H=\left\{\alpha H_0\right|h(\alpha )=l\}`$ is the homogeneous set.
Note that the proof really gives the following consequence
$$V[G]\mathrm{}_1\stackrel{def}{}(\mathrm{}_1)_\mathrm{}_0^2.$$
If $`V=L`$, or there is a definable well-ordering of $`H_\kappa `$ in $`V`$, then the same proof yields the following:
###### Lemma 4.3.
Assume $`V=L`$. Let $`h,AV[G]`$ be such that $`|A|=\mathrm{}_1`$, and $`h:[A]^22`$ is a partition of $`[A]^2`$ into two parts, where both $`A`$ and $`h`$ are first order definable (with parameters) over $`(H_{\omega _1},\epsilon )^{V[G]}`$. Then there is a $`BA`$ homogeneous for $`h`$ and $`|B|=\mathrm{}_1`$.
To derive theorem 4.1 we follow the proof that a weakly compact cardinal has the tree property (see lemma 29.6 of ), replacing the partition property, of a weakly compact cardinal with the definable partition lemma (4.2).
proof of theorem 4.1
Let $`T=(\mathrm{}_1,<_T)`$ be a definable tree on $`\mathrm{}_1`$, i.e. there is a formula $`\mathrm{\Psi }(\alpha ,\beta ,z)`$ such that
(4.7)
$$\alpha <_T\beta (H_{\omega _1},\epsilon )\mathrm{\Psi }(\alpha ,\beta ,z).$$
We extend the partial tree ordering $`<_T`$ into a linear ordering as follows:
$`\underset{¯}{\alpha \beta }`$ iff
* $`\alpha <_T\beta `$ or
* $`\alpha ,\beta `$ are $`<_T`$ incomparable and if $`\zeta `$ is the first level where the predecessors of $`\alpha ,\beta `$, $`\alpha _\zeta ,\beta _\zeta `$ are distinct then $`\alpha _\zeta <\beta _\zeta `$.
$``$ is first order definable in $`(H_{\omega _1},\epsilon )`$ using the definition of $`<_T`$. Now define a partition of $`[\mathrm{}_1]^2`$ by
(4.8)
$$F(\{\alpha ,\beta \})=1\alpha <\beta \text{ agrees with }\alpha \beta .$$
Since both $``$ and $`<_T`$ are definable $`F`$ is definable as well. Hence there is $`H\mathrm{}_1`$ which is homogenous for $`F`$, with $`|H|=\mathrm{}_1`$. Let
(4.9)
$$B=\left\{x\mathrm{}_1\right||\{\alpha H|x<_T\alpha \}|=\mathrm{}_1\}.$$
Since every level is countable there are members of $`B`$ of every level. If we’ll prove that any two members of $`B`$ are $`<_T`$ comparable, then $`B`$ will be the $`\mathrm{}_1`$-branch. Let $`x,yB`$ be $`<_T`$ incomparable elements. Assume, without loss of generality, that $`xy`$. Since both $`x,y`$ have $`\mathrm{}_1`$ many $`<_T`$ succesors in $`H`$ we can find $`\alpha ,\beta ,\mu H`$ such that $`\alpha <_T\beta <_T\mu `$, $`x<_T\alpha ,\mu `$ and $`y<_T\beta `$. By the definition of $``$ we get $`\alpha \beta `$ and $`\mu \beta `$. Thus $`F(\{\alpha ,\beta \})=1`$ and $`F(\{\beta ,\mu \})=0`$ , contradicting the fact that $`H`$ is homogeneous for $`F`$. Finally note that if we force over $`L`$ theorem 4.1 can be strengthened to trees definable in the very wide sense. The proof is identical using lemma 4.3 instead of 4.2.
## 5. The lower bound
In this section we prove that the definable tree property implies the consistency of a $`\mathrm{\Pi }_1^1`$ reflecting cardinal. In this section let $`\mathrm{}_1`$ denote $`\mathrm{}_1^V`$.
###### Theorem 5.1.
If $`\mathrm{}_1`$ has the definable tree property then
$$L\mathrm{}_1\text{ is a }\mathrm{\Pi }_1^1\text{ reflecting cardinal}.$$
First we prove that
(5.1)
$$L\mathrm{}_1\text{ is inaccessible}.$$
Assume that $`\mathrm{}_1`$ is not inaccessible in $`L`$ then there is an $`x\omega ^\omega `$ such that $`\mathrm{}_1=\mathrm{}_1^{L[x]}`$ (see proposition 11.5). However inside $`L[x]`$ there is a special Aronszajn tree $`T`$ which is definable from the well ordering of $`(H_{\omega _1},\epsilon )^{L[x]}`$, which is itself $`𝚺_1`$ definable over $`(H_{\omega _1},\epsilon )^{L[x]}`$ . However $`T`$ cannot have a cofinal branch in $`V`$ since this implies $`\mathrm{}_1^{L[x]}<\mathrm{}_1`$. Similarly, by relativization, for every real $`x`$ $`\mathrm{}_1`$ is inaccessible in $`L[x]`$.
Next we prove that $`\mathrm{}_1`$ is a $`\mathrm{\Pi }_1^1`$reflecting cardinal in $`L`$. The proof is based on an idea from . We define a tree using the $`𝚺_n`$ definable power set of $`\mathrm{}_1`$. Note that since $`\mathrm{}_1`$ is inaccessible in $`L`$, by (5.1), we have $`(H_{\mathrm{}_1^V})^L=L_{\mathrm{}_1^V}=\left(V_{\mathrm{}_1^V}\right)^L`$. From a cofinal branch in the tree we define an ultrafilter on the $`𝚺_n`$ definable subsets of $`\mathrm{}_1`$, and construct an “ultrapower” of $`L_\mathrm{}_1`$, using only functions $`𝚺_n`$ definable over $`(H_{\omega _1},\epsilon )`$. Note that this is a definable tree in the strong sense since there is a definable well-oredering of the underlying set. Let $`A_\alpha |\alpha <\mathrm{}_1`$ be a definable enumeration of $`P^L(\mathrm{}_1)𝚺_n`$ of order type $`\mathrm{}_1`$. Define a tree $`T`$ of functions by
$$fTf:\tau \{0,1\},\tau <\mathrm{}_1\text{ and }|_{\alpha <\tau }A_\alpha ^{f(\alpha )}|=\mathrm{}_1$$
where $`A^0=A`$, and $`A^1=\mathrm{}_1\backslash A`$. The ordering on $`T`$ is $`f<_Tg`$ iff $`fg`$. Since $`\mathrm{}_1`$ is inaccessible in $`L`$ the tree is an $`\omega _1`$tree. Since the truth of $`\mathrm{\Sigma }_n`$ formulas is a $`\mathrm{\Sigma }_{n+1}`$ definable, the tree $`T`$ is $`𝚺_k`$ definable over $`(H_{\omega _1},\epsilon )`$, for some $`k`$. Hence by the definable tree property it has a cofinal branch producing a function $`b:\mathrm{}_1\{0,1\}`$. $`b`$ defines an ultrafilter $`U`$ on $`P^L(\mathrm{}_1)𝚺_n`$ by
$$A_\alpha Ub(\alpha )=0$$
and this ultrafilter is countably complete on $`P^L(\mathrm{}_1)𝚺_n`$. The “ultrapower” is now defined by
(5.2)
$$\begin{array}{cc}fult(L_\mathrm{}_1,U)\hfill & f:\mathrm{}_1L_\mathrm{}_1,fL\hfill \\ & \text{ and }f\text{ is }𝚺_n\text{ definable over }(H_{\omega _1},\epsilon )\hfill \end{array}$$
(5.3)
$$fg\{\alpha |f(\alpha )=g(\alpha )\}U$$
and
(5.4)
$$fEg\{\alpha |f(\alpha )g(\alpha )\}U.$$
The “ultrapower” is wellfounded by the completeness of the ultrafilter and there is a $`\mathrm{\Sigma }_n`$ embedding $`j:L_\mathrm{}_1_nult(L_\mathrm{}_1,U)`$, by the proof of $`Ł`$os theorem. However since $`L_\mathrm{}_1V=L`$ if $`n`$ is large enough $`ult(L_\mathrm{}_1,U)V=L`$, and hence its transitive collapse is really $`L_\alpha `$ for some $`\alpha >\mathrm{}_1`$. Therfore $`L_\mathrm{}_1_nL_\alpha `$. This gives the desired extension property, and by the equivalence of the extension property and $`\mathrm{\Pi }_1^1`$ reflection $`\mathrm{}_1`$ is $`\mathrm{\Pi }_1^1`$ reflecting in $`L`$. Finally by relativizing we obtain that for every real $`x`$, $`\mathrm{}_1`$ is $`\mathrm{\Pi }_1^1`$ reflecting in $`L[x]`$.
## 6. The definable tree property and Martin’s axiom
In this section we investigate the consistency strengh of the Definable Tree Property on $`\mathrm{}_1`$ together with Martin’s Axiom and large continuum. The exposition is based on the Shelah-Harrington paper . The main theorem is the following :
###### Theorem 6.1.
Let $`\mathrm{}_0<\lambda `$ be a cardinal satisfying $`\lambda ^{<\lambda }=\lambda `$. The following are equiconsistent :
* The definable tree property on $`\mathrm{}_1+MA+2^\mathrm{}_0=\lambda `$
* $`\mathrm{}_1`$ is weakly compact cardinal in $`L`$.
For the $`21`$ direction, we use lemma 5.1 which proves that the definable tree property implies that for every real $`x`$ $`\mathrm{}_1^{L[x]}<\mathrm{}_1`$. Now we finish by the following result from :
###### Theorem 6.2.
Assume MA then either there exists a real $`x`$ such that $`\mathrm{}_1^{L[x]}=\mathrm{}_1`$, or $`\mathrm{}_1`$ is weakly compact in $`L`$.
The proof of the other direction follows closely Kunen’s forcing for the consistency of “$`MA+`$ Every set in L(R) is Lebesgue measureable $`+2^\mathrm{}_0=\lambda `$”. The model we use is Kunen’s model, and we only prove that the definable tree property holds in that model. For completeness we present the full proofs of Kunen’s basic observations, as presented in . Let $`\kappa `$ be a weakly compact cardinal.
###### Lemma 6.3.
If $`B`$ is a complete Boolean Algebra with the $`\kappa `$c.c., and if $`XB`$ satisfies $`|X|<\kappa `$, then there is a complete subalgebra $`\overline{B}`$, s.t. $`X\overline{B}`$ and $`|\overline{B}|<\kappa `$.
Proof :Since $`B`$ satisfies the $`\kappa `$c.c., and $`\kappa ^{<\kappa }=\kappa `$, there is $`B^{}`$ a complete subalgebra of $`B`$ such that $`XB^{}`$ and $`|B^{}|\kappa `$. Without loss of generality assume that $`B^{}\kappa `$. Let $`D`$ be the set of maximal antichains of $`B^{}`$. $`D[\kappa ]^{<\kappa }`$. By $`\mathrm{\Pi }_1^1`$ reflection there is an $`\alpha <\kappa `$ s.t. $`B^{}\alpha `$ is $`<\alpha `$ complete and $`D[\alpha ]^{<\alpha }`$ is the set of maximal antichains of $`B^{}\alpha `$. Therefore $`B^{}\alpha `$ is a complete subalgebra of $`B`$. If we choose $`\alpha >sup(X)`$ then $`XB^{}\alpha `$.
###### Lemma 6.4.
If $`P_0,P_1`$ are two complete Boolean algebras with $`\kappa `$c.c., then $`P_0\times P_1`$ has the $`\kappa `$c.c.
Proof :Let $`<p_\alpha ^0,p_\alpha ^1>|\alpha <\kappa `$ be a sequence of elements of $`P_0\times P_1`$. Define $`F:[\kappa ]^22\times 2`$. $`F(\{\alpha ,\beta \})(i)=0\text{ iff }p_\alpha ^i,p_\beta ^i`$ are compatible. A size $`\kappa `$ homogeneous set for $`F`$ gives a size $`\kappa `$ set of pairwise compatible elements, since by the $`<\kappa `$c.c of $`P_i`$ the homogeneous color is $`<0,0>`$.
Assume that $`\nu <\lambda =\lambda ^{<\lambda }`$. To obtain the model we iterate $`\lambda `$ many times $`\kappa `$c.c. forcings of size $`<\lambda `$ using finite support, the same way this is done for $`MA`$ (). By lemma (6.4) we can assume, without loss of generality, that each forcing appears $`\lambda `$ many times in the iteration. Denote the iteration by $`B`$. To prove that every $`\mathrm{}_1`$ tree which is ordinal definable from a real has a branch in $`V^B`$, we first prove that Ramsey theorem for $`\mathrm{}_1`$ holds for $`L(R)`$ partitions. Then we finish the same way as in the proof of theorem (4.1).
Let $`F:[\mathrm{}_1]^22`$ be definable in $`V^B`$ from a real $`x`$, and ordinal parameters. By lemma (6.3) $`x`$ is generic for a countable subalgebra $`\overline{B}`$. Moreover since each forcing in the iteration appears unboundedly many times we can assume , without loss of generality, that $`B/\overline{B}`$ is an homogeneous iteration over $`V^{\overline{B}}`$ which is built the same way as $`B`$ is built over $`V`$. Therefore every value of $`F`$ is decided inside $`V^{\overline{B}}`$. Since $`|\overline{B}|<\kappa `$ we can construct (in $`V`$) a $`\kappa `$ tree of $`\overline{B}`$ names of possible values of $`F`$, the same way as we have built it in lemma (4.2). Now by the weak compactness of $`\kappa `$ we can find a branch through that tree. Finally we use the branch to obtain the homogeneous set the same as we have done in lemma (4.2).
Note that since $`\kappa `$ is weakly compact we are able to find homogeneous sets for every $`L(R)`$ coloring, and not just first order definable over $`H_{\omega _1}`$.
Finally we remark that like in Solovay’s proof of Lebesgue measurability , just adding the negation of the Continuum Hypothesis does not add to the consistency strength. We have to notice that the product of collapsing $`\kappa `$ to $`\omega _1`$ and then adding $`\lambda `$ Cohen reals satisfies the $`\kappa `$c.c. The product is also homogeneous enough. Moreover every real belongs to a small generic extension, which is complemented by a homogeneous forcing. Hence the same situation as in theorem 4.1 still holds, and the argument there can be carried out.
## 7. The definable tree property on higher cardinals
This section contains several remarks regarding the definable tree property on cardinals above $`\mathrm{}_1`$. By Mitchel’s and Silver’s results , the tree property on $`\mathrm{}_2`$ is equiconsistent with the existence of a weakly compact cardinal. However it is well known that assuming GCH, or even $`\lambda ^{<\lambda }=\lambda `$, we have a $`\lambda ^+`$ special Aronszajn tree. This is proved by the same construction as Aronszajn’s construction on $`\mathrm{}_1`$, using the fact that under this hypothesis there is a universal linear order of cardinality $`\lambda `$.
We will prove that assuming a $`\mathrm{\Pi }_1^1`$ reflecting cardinal the definable tree property on $`\mathrm{}_2`$ is consistent with GCH. Then we consider the property of having the definable tree property on succesive cardinals. We will generalize our forcing argument to prove that if it is consistent that there are $`\omega `$ $`\mathrm{\Pi }_1^1`$ reflecting cardinals then it is consistent to have the definable tree property on $`\mathrm{}_i:1i<\omega `$.
###### Theorem 7.1.
Assume $`\kappa `$ is a $`\mathrm{\Pi }_1^1`$ reflecting cardinal in L. Let $`P=Coll(\mathrm{}_1,<\kappa )`$. Let $`G`$ be $`P`$ generic over L. Then
$$L[G]\text{GCH and }\mathrm{}_2\text{ has the definable tree property.}$$
Proof :The proof is identical to the $`\mathrm{}_1`$ case. Using the homogeneity of the Levy collapse, and the fact that $`\kappa `$ is $`\mathrm{\Pi }_1^1`$ reflecting. We just build a tree of possible values for the definable function. Then using the definabiltiy of the function we obtain definability of the tree and thus we can use the branch to prove the definable version of Ramsey theorem for $`\mathrm{}_2`$. The fact that GCH holds in $`V[G]`$, is proved in the usual way, by assuming $`GCH`$ in the ground model.
To obtain the definable tree property on all $`\mathrm{}_n`$’s from a sequence of $`\omega `$ $`\mathrm{\Pi }_1^1`$ reflecting cardinals, $`\kappa _0<\kappa _1<\mathrm{}`$ just iterate the forcings $`Coll(\mathrm{}_n,<\kappa _n)`$ with finite support. The proof that the definable tree property holds for every n, is a straightforward generalization and will not be given here.
## 8. Acknowledgement
I would like to thank Menachem Magidor for helpful discussions, and the anonymous referee for his comments which greatly improved the presentation of the results.
|
warning/0005/math0005051.html
|
ar5iv
|
text
|
# 1 Introduction. Basic definitions
## 1 Introduction. Basic definitions
In this paper, notions of compatible and almost compatible Riemannian and pseudo-Riemannian metrics, which are motivated by the theory of compatible (local and nonlocal) Poisson structures of hydrodynamic type and generalize the notion of flat pencil of metrics (this notion plays an important role in the theory of integrable systems of hydrodynamic type and the Dubrovin theory of Frobenius manifolds , see also ), are introduced and studied. Compatible metrics generate compatible Poisson structures of hydrodynamic type (these structures are local for flat metrics and they are nonlocal if the metrics are not flat ) and their investigation is necessary for the theory of integrable systems of hydrodynamic type. In “nonsingular” case, when eigenvalues of pair of metrics are distinct, in this paper the complete explicit description of compatible and almost compatible metrics is obtained. The “singular” case of coinciding eigenvalues of pair of metrics is considerably more complicated for the complete analysis and has still not been completely studied. Nevertheless, the problem on two-component compatible flat metrics is completely investigated. All such pairs, both “nonsingular” and “singular”, are classified by ours. In this paper we present the complete description of nonsingular pairs of two-component flat metrics. The problems of classification of compatible flat metrics and compatible metrics of constant Riemannian curvature are of particular interest, in particular, from the viewpoint of the theory of Hamiltonian systems of hydrodynamic type. More detailed classification results for these problems will be published in another paper. In the given paper we prove that the approach proposed by Ferapontov in for the study of flat pencils of metrics can be also applied (with the corresponding modifications and corrections) to pencils of metrics of constant Riemannian curvature and to the general compatible Riemannian and pseudo-Riemannian metrics. We also correct a mistake which is in in the criterion of compatibility of local nondegenerate Poisson structures of hydrodynamic type (or, in other words, compatibility of flat metrics).
We shall use both contravariant metrics $`g^{ij}(u)`$ with upper indices, where $`u=(u^1,\mathrm{},u^N)`$ are local coordinates, $`1i,jN`$, and covariant metrics $`g_{ij}(u)`$ with lower indices, $`g^{is}(u)g_{sj}(u)=\delta _j^i.`$ The indices of coefficients of the Levi–Civita connections $`\mathrm{\Gamma }_{jk}^i(u)`$ (the Riemannian connections generated by the corresponding metrics) and tensors of Riemannian curvature $`R_{jkl}^i(u)`$ are raised and lowered by the metrics corresponding to them:
$$\mathrm{\Gamma }_k^{ij}(u)=g^{is}(u)\mathrm{\Gamma }_{sk}^j(u),\mathrm{\Gamma }_{jk}^i(u)=\frac{1}{2}g^{is}(u)\left(\frac{g_{sk}}{u^j}+\frac{g_{js}}{u^k}\frac{g_{jk}}{u^s}\right),$$
$$R_{kl}^{ij}(u)=g^{is}(u)R_{skl}^j(u),R_{jkl}^i(u)=\frac{\mathrm{\Gamma }_{jl}^i}{u^k}+\frac{\mathrm{\Gamma }_{jk}^i}{u^l}\mathrm{\Gamma }_{pk}^i(u)\mathrm{\Gamma }_{jl}^p(u)+\mathrm{\Gamma }_{pl}^i(u)\mathrm{\Gamma }_{jk}^p(u).$$
###### Definition 1.1
Two contravariant metrics $`g_1^{ij}(u)`$ and $`g_2^{ij}(u)`$ of constant Riemannian curvature $`K_1`$ and $`K_2`$, respectively, are called compatible if any linear combination of these metrics
$$g^{ij}(u)=\lambda _1g_1^{ij}(u)+\lambda _2g_2^{ij}(u),$$
(1.1)
where $`\lambda _1`$ and $`\lambda _2`$ are arbitrary constants such that $`det(g^{ij}(u))0`$, is a metric of constant Riemannian curvature $`\lambda _1K_1+\lambda _2K_2`$ and the coefficients of the corresponding Levi–Civita connections are related by the same linear formula:
$$\mathrm{\Gamma }_k^{ij}(u)=\lambda _1\mathrm{\Gamma }_{1,k}^{ij}(u)+\lambda _2\mathrm{\Gamma }_{2,k}^{ij}(u).$$
(1.2)
We shall also say in this case that the metrics $`g_1^{ij}(u)`$ and $`g_2^{ij}(u)`$ form a pencil of metrics of constant Riemannian curvature.
###### Definition 1.2
Two Riemannian or pseudo-Riemannian contravariant metrics $`g_1^{ij}(u)`$ and $`g_2^{ij}(u)`$ are called compatible if for any linear combination of these metrics
$$g^{ij}(u)=\lambda _1g_1^{ij}(u)+\lambda _2g_2^{ij}(u),$$
(1.3)
where $`\lambda _1`$ and $`\lambda _2`$ are arbitrary constants such that $`det(g^{ij}(u))0`$, the coefficients of the corresponding Levi–Civita connections and the components of the corresponding tensors of Riemannian curvature are related by the same linear formula:
$$\mathrm{\Gamma }_k^{ij}(u)=\lambda _1\mathrm{\Gamma }_{1,k}^{ij}(u)+\lambda _2\mathrm{\Gamma }_{2,k}^{ij}(u),$$
(1.4)
$$R_{kl}^{ij}(u)=\lambda _1R_{1,kl}^{ij}(u)+\lambda _2R_{2,kl}^{ij}(u).$$
(1.5)
We shall also say in this case that the metrics $`g_1^{ij}(u)`$ and $`g_2^{ij}(u)`$ form a pencil of metrics.
###### Definition 1.3
Two Riemannian or pseudo-Riemannian contravariant metrics $`g_1^{ij}(u)`$ and $`g_2^{ij}(u)`$ are called almost compatible if for any linear combination of these metrics (1.3) relation (1.4) is fulfilled.
###### Definition 1.4
Two Riemannian or pseudo-Riemannian metrics $`g_1^{ij}(u)`$ and $`g_2^{ij}(u)`$ are called nonsingular pair of metrics if the eigenvalues of this pair of metrics, that is, the roots of the equation
$$det(g_1^{ij}(u)\lambda g_2^{ij}(u))=0,$$
(1.6)
are distinct.
These definitions are motivated by the theory of compatible Poisson brackets of hydrodynamic type. In the case if the metrics $`g_1^{ij}(u)`$ and $`g_2^{ij}(u)`$ are flat, that is, $`R_{1,jkl}^i(u)=R_{2,jkl}^i(u)=0,`$ relation (1.5) is equivalent to the condition that an arbitrary linear combination of the flat metrics $`g_1^{ij}(u)`$ and $`g_2^{ij}(u)`$ is also a flat metric and Definition 1.2 is equivalent to the well-known definition of a flat pencil of metrics or, in other words, a compatible pair of local nondegenerate Poisson structures of hydrodynamic type (see also ). In the case if the metrics $`g_1^{ij}(u)`$ and $`g_2^{ij}(u)`$ are metrics of constant Riemannian curvature $`K_1`$ and $`K_2`$, respectively, that is,
$$R_{1,kl}^{ij}(u)=K_1(\delta _k^i\delta _l^j\delta _l^i\delta _k^j),R_{2,kl}^{ij}(u)=K_2(\delta _k^i\delta _l^j\delta _l^i\delta _k^j),$$
relation (1.5) gives the condition that an arbitrary linear combination of the metrics $`g_1^{ij}(u)`$ and $`g_2^{ij}(u)`$ (1.3) is a metric of constant Riemannian curvature $`\lambda _1K_1+\lambda _2K_2`$ and Definition 1.2 is equivalent to Definition 1.1 of a pencil of metrics of constant Riemannian curvature or, in other words, a compatible pair of the corresponding nonlocal Poisson structures of hydrodynamic type which were introduced and studied by the author and Ferapontov in . Compatible metrics of more general type correspond to compatible pairs of nonlocal Poisson structures of hydrodynamic type which were introduced and studied by Ferapontov in . They arise, for example, if we shall use a recursion operator generated by a pair of compatible Poisson structures of hydrodynamic type and determining, as is well-known, an infinite sequence of corresponding Poisson structures.
As was earlier noted by the author in , condition (1.5) follows from condition (1.4) in the case of certain special reductions connected with the associativity equations (see also Theorem 3.2 below). Of course, it is not by chance. Under certain very natural and quite general assumptions on metrics (it is sufficient but not necessary, in particular, that the eigenvalues of the pair of the metrics under consideration are distinct), compatibility of the metrics follows from their almost compatibility but, generally speaking, in the general case, it is not true even for flat metrics (we shall present here below the corresponding counterexamples). Correspondingly, we would like to emphasize that condition (1.4) which is considerably more simple than condition (1.5) “almost” guarantees compatibility of metrics and deserves a separate investigation but, in the general case, it is necessary to require also the fulfillment of condition (1.5) for compatibility of the corresponding Poisson structures of hydrodynamic type. It is also interesting to find out, does condition (1.5) guarantee the fulfillment of condition (1.4) or not.
## 2 Compatible local Poisson structures of <br>hydrodynamic type
The local homogeneous Poisson bracket of the first order, that is, the Poisson bracket of the form
$$\{u^i(x),u^j(y)\}=g^{ij}(u(x))\delta _x(xy)+b_k^{ij}(u(x))u_x^k\delta (xy),$$
(2.1)
where $`u^1,\mathrm{},u^N`$ are local coordinates on a certain smooth $`N`$-dimensional manifold $`M`$, is called a local Poisson structure of hydrodynamic type or Dubrovin–Novikov structure . Here, $`u^i(x),1iN,`$ are functions (fields) of a single independent variable $`x`$, and the coefficients $`g^{ij}(u)`$ and $`b_k^{ij}(u)`$ of bracket (2.1) are smooth functions on $`M`$.
In other words, for arbitrary functionals $`I[u]`$ and $`J[u]`$ on the space of fields $`u^i(x),1iN,`$ a bracket of the form
$$\{I,J\}=\frac{\delta I}{\delta u^i(x)}\left(g^{ij}(u(x))\frac{d}{dx}+b_k^{ij}(u(x))u_x^k\right)\frac{\delta J}{\delta u^j(x)}𝑑x$$
(2.2)
is defined and it is required that this bracket is a Poisson bracket, that is, it is skew-symmetric:
$$\{I,J\}=\{J,I\},$$
(2.3)
and satisfies the Jocobi identity
$$\{\{I,J\},K\}+\{\{J,K\},I\}+\{\{K,I\},J\}=0$$
(2.4)
for arbitrary functionals $`I[u]`$, $`J[u]`$ and $`K[u]`$. The skew-symmetry (2.3) and the Jacobi identity (2.4) impose very strict conditions on the coefficients $`g^{ij}(u)`$ and $`b_k^{ij}(u)`$ of bracket (2.2) (these conditions will be considered below).
The local Poisson structures of hydrodynamic type (2.1) were introduced and studied by Dubrovin and Novikov in . In this paper, they proposed a general Hamiltonian approach to the so-called homogeneous systems of hydrodynamic type, that is, to evolutionary quasilinear systems of first-order partial differential equations
$$u_t^i=V_j^i(u)u_x^j$$
(2.5)
that corresponds to structures (2.1).
This Hamiltonian approach was motivated by the study of the equations of Euler hydrodynamics and the Whitham averaging equations that describe the evolution of slowly modulated multiphase solutions of partial differential equations .
Local bracket (2.2) is called nondegenerate if $`det(g^{ij}(u))0`$. For general nondegenerate brackets of form (2.2), Dubrovin and Novikov proved the following important theorem.
###### Theorem 2.1 (Dubrovin, Novikov )
If $`det(g^{ij}(u))0`$, then bracket (2.2) is a Poisson bracket, that is, it is skew-symmetric and satisfies the Jacobi identity if and only if
* $`g^{ij}(u)`$ is an arbitrary flat pseudo-Riemannian contravariant metric (a metric of zero Riemannian curvature),
* $`b_k^{ij}(u)=g^{is}(u)\mathrm{\Gamma }_{sk}^j(u),`$ where $`\mathrm{\Gamma }_{sk}^j(u)`$ is the Riemannian connection generated by the contravariant metric $`g^{ij}(u)`$ (the Levi–Civita connection).
Consequently, for any local nondegenerate Poisson structure of hydrodynamic type, there always exist local coordinates $`v^1,\mathrm{},v^N`$ (flat coordinates of the metric $`g^{ij}(u)`$) in which the coefficients of the brackets are constant:
$$\stackrel{~}{g}^{ij}(v)=\eta ^{ij}=\mathrm{const},\stackrel{~}{\mathrm{\Gamma }}_{jk}^i(v)=0,\stackrel{~}{b}_k^{ij}(v)=0,$$
(2.6)
that is, the bracket has the constant form
$$\{I,J\}=\frac{\delta I}{\delta v^i(x)}\eta ^{ij}\frac{d}{dx}\frac{\delta J}{\delta v^j(x)}𝑑x,$$
(2.7)
where $`(\eta ^{ij})`$ is a nondegenerate symmetric constant matrix:
$$\eta ^{ij}=\eta ^{ji},\eta ^{ij}=\mathrm{const},det(\eta ^{ij})0.$$
On the other hand, as early as 1978, Magri proposed a bi-Hamiltonian approach to the integration of nonlinear systems . This approach demonstrated that the integrability is closely related to the bi-Hamiltonian property, that is, to the property of a system to have two compatible Hamiltonian representations. As was shown by Magri in , compatible Poisson brackets generate integrable hierarchies of systems of differential equations. Therefore, the description of compatible Poisson structures is very urgent and important problem in the theory of integrable systems. In particular, for a system, the bi-Hamiltonian property generates recurrent relations for the conservation laws of this system.
Beginning from , quite extensive literature (see, for example, and the necessary references therein) has been devoted to the bi-Hamiltonian approach and to the construction of compatible Poisson structures for many specific important equations of mathematical physics and field theory. As far as the problem of description of sufficiently wide classes of compatible Poisson structures of defined special types is concerned, apparently the first such statement was considered in , (see also , ). In those papers, the present author posed and completely solved the problem of description of all compatible local scalar first-order and third-order Poisson brackets, that is, all Poisson brackets given by arbitrary scalar first-order and third-order ordinary differential operators. These brackets generalize the well-known compatible pair of the Gardner–Zakharov–Faddeev bracket , (first-order bracket) and the Magri bracket (third-order bracket) for the Korteweg–de Vries equation.
In the case of homogeneous systems of hydrodynamic type, many integrable systems possess compatible Poisson structures of hydrodynamic type. The problems of description of these structures for particular systems and numerous examples were considered in many papers (see, for example, ). In particular, in Nutku studied a special class of compatible two-component Poisson structures of hydrodynamic type and the related bi-Hamiltonian hydrodynamic systems. In Ferapontov classified all two-component homogeneous systems of hydrodynamic type possessing three compatible local Poisson structures of hydrodynamic type.
In the general form, the problem of description of flat pencil of metrics (or, in other words, compatible nondegenerate local Poisson structures of hydrodynamic type) was considered by Dubrovin in , in connection with the construction of important examples of such flat pencils of metrics, generated by natural pairs of flat metrics on the spaces of orbits of Coxeter groups and on other Frobenius manifolds and associated with the corresponding quasi-homogeneous solutions of the associativity equations. In the theory of Frobenius manifolds introduced and studied by Dubrovin , (they correspond to two-dimensional topological field theories), a key role is played by flat pencils of metrics, possessing a number of special additional (and very strict) properties (they satisfy the so-called quasi-homogeneity property). In addition, in Dubrovin proved that the theory of Frobenius manifolds is equivalent to the theory quasi-homogeneous compatible nondegenerate Poisson structures of hydrodynamic type. The general problem of compatible nondegenerate local Poisson structures was also considered by Ferapontov in .
The author’s papers are devoted to the general problem of classification of local Poisson structures of hydrodynamic type, to integrable nonlinear systems which describe such compatible Poisson structures and to special reductions connected with the associativity equations.
###### Definition 2.1 (Magri )
Two Poisson brackets $`\{,\}_1`$ and $`\{,\}_2`$ are called compatible if an arbitrary linear combination of these Poisson brackets
$$\{,\}=\lambda _1\{,\}_1+\lambda _2\{,\}_2,$$
(2.8)
where $`\lambda _1`$ and $`\lambda _2`$ are arbitrary constants, is also always a Poisson bracket. In this case, one can say also that the brackets $`\{,\}_1`$ and $`\{,\}_2`$ form a pencil of Poisson brackets.
Correspondingly, the problem of description of compatible nondegenerate local Poisson structures of hydrodynamic type is pure differential-geometric problem of description of flat pencils of metrics (see , ).
In , Dubrovin presented all the tensor relations for the general flat pencils of metrics. First, we introduce the necessary notation. Let $`_1`$ and $`_2`$ be the operators of covariant differentiation given by the Levi–Civita connections $`\mathrm{\Gamma }_{1,k}^{ij}(u)`$ and $`\mathrm{\Gamma }_{2,k}^{ij}(u)`$, generated by the metrics $`g_1^{ij}(u)`$ and $`g_2^{ij}(u)`$, respectively. The indices of the covariant differentials are raised and lowered by the corresponding metrics: $`_1^i=g_1^{is}(u)_{1,s}`$, $`_2^i=g_2^{is}(u)_{2,s}`$. Consider the tensor
$$\mathrm{\Delta }^{ijk}(u)=g_1^{is}(u)g_2^{jp}(u)\left(\mathrm{\Gamma }_{2,ps}^k(u)\mathrm{\Gamma }_{1,ps}^k(u)\right),$$
(2.9)
introduced by Dubrovin in , .
###### Theorem 2.2 (Dubrovin , )
If metrics $`g_1^{ij}(u)`$ and $`g_2^{ij}(u)`$ form a flat pencil, then there exists a vector field $`f^i(u)`$ such that the tensor $`\mathrm{\Delta }^{ijk}(u)`$ and the metric $`g_1^{ij}(u)`$ have the form
$$\mathrm{\Delta }^{ijk}(u)=_2^i_2^jf^k(u),$$
(2.10)
$$g_1^{ij}(u)=_2^if^j(u)+_2^jf^i(u)+cg_2^{ij}(u),$$
(2.11)
where $`c`$ is a certain constant, and the vector field $`f^i(u)`$ satisfies the equations
$$\mathrm{\Delta }_s^{ij}(u)\mathrm{\Delta }_l^{sk}(u)=\mathrm{\Delta }_s^{ik}(u)\mathrm{\Delta }_l^{sj}(u),$$
(2.12)
where
$$\mathrm{\Delta }_k^{ij}(u)=g_{2,ks}(u)\mathrm{\Delta }^{sij}(u)=_{2,k}_2^if^j(u),$$
(2.13)
and
$$(g_1^{is}(u)g_2^{jp}(u)g_2^{is}(u)g_1^{jp}(u))_{2,s}_{2,p}f^k(u)=0.$$
(2.14)
Conversely, for the flat metric $`g_2^{ij}(u)`$ and the vector field $`f^i(u)`$ that is a solution of the system of equations (2.12) and (2.14), the metrics $`g_2^{ij}(u)`$ and (2.11) form a flat pencil.
The proof of this theorem immediately follows from the relations that are equivalent to the fact that the metrics $`g_1^{ij}(u)`$ and $`g_2^{ij}(u)`$ form a flat pencil and are considered in flat coordinates of the metric $`g_2^{ij}(u)`$ , .
In my paper , an explicit and simple criterion of compatibility for two Poisson structures of hydrodynamic type is formulated, that is, it is shown what explicit form is sufficient and necessary for the Poisson structures of hydrodynamic type to be compatible.
For the moment, we are able to formulate such explicit general criterion only namely in terms of Poisson structures but not in terms of metrics as in Theorem 2.2.
###### Lemma 2.1 ()
(An explicit criterion of compatibility for Poisson structures of hydrodynamic type) Any local Poisson structure of hydrodynamic type $`\{I,J\}_2`$ is compatible with the constant nondegenerate Poisson bracket (2.7) if and only if it has the form
$$\{I,J\}_2=\frac{\delta I}{\delta v^i(x)}\left(\left(\eta ^{is}\frac{h^j}{v^s}+\eta ^{js}\frac{h^i}{v^s}\right)\frac{d}{dx}+\eta ^{is}\frac{^2h^j}{v^sv^k}v_x^k\right)\frac{\delta J}{\delta v^j(x)}𝑑x,$$
(2.15)
where $`h^i(v),1iN,`$ are smooth functions defined on a certain neighbourhood.
We do not require in Lemma 2.1 that the Poisson structure of hydrodynamic type $`\{I,J\}_2`$ is nondegenerate. Besides, it is important to note that this statement is local.
In 1995, in the paper , Ferapontov proposed an approach to the problem on flat pencils of metrics, which is motivated by the theory of recursion operators, and formulated the following theorem as a criterion of compatibility of nondegenerate local Poisson structures of hydrodynamic type:
###### Theorem 2.3 ()
Two local nondegenerate Poisson structures of hydrodynamic type given by flat metrics $`g_1^{ij}(u)`$ and $`g_2^{ij}(u)`$ are compatible if and only if the Nijenhuis tensor of the affinor $`v_j^i(u)=g_1^{is}(u)g_{2,sj}(u)`$ vanishes, that is,
$$N_{ij}^k(u)=v_i^s(u)\frac{v_j^k}{u^s}v_j^s(u)\frac{v_i^k}{u^s}+v_s^k(u)\frac{v_i^s}{u^j}v_s^k\frac{v_j^s}{u^i}=0.$$
(2.16)
Besides, it is noted in the remark in that if the spectrum of $`v_j^i(u)`$ is simple, the vanishing of the Nijenhuis tensor implies the existence of coordinates $`R^1,\mathrm{},R^N`$ for which all the objects $`v_j^i(u)`$, $`g_1^{ij}(u)`$, $`g_2^{ij}(u)`$ become diagonal. Moreover, in these coordinates the $`i`$th eigenvalue of $`v_j^i(u)`$ depends only on the coordinate $`R^i`$. In the case when all the eigenvalues are nonconstant, they can be introduced as new coordinates. In these new coordinates $`\stackrel{~}{v}_j^i(R)=\mathrm{diag}(R^1,\mathrm{},R^N)`$, $`\stackrel{~}{g}_2^{ij}(R)=\mathrm{diag}(g^1(R),\mathrm{},g^N(R))`$, $`\stackrel{~}{g}_1^{ij}(R)=\mathrm{diag}(R^1g^1(R),\mathrm{},R^Ng^N(R))`$.
Unfortunately, in the general case the Theorem 2.3 is not true and, correspondingly, it is not a criterion of compatibility of flat metrics. Generally speaking, compatibility of flat metrics does not follow from the vanishing of the corresponding Nijenhuis tensor. In this paper we shall present the corresponding counterexamples. In the general case, as it will be shown here, the Theorem 2.3 is actually a criterion of almost compatibility of flat metrics that does not guarantee compatibility of the corresponding nondegenerate local Poisson structures of hydrodynamic type. But if the spectrum of $`v_j^i(u)`$ is simple, that is, all the eigenvalues are distinct, then we proves here that the Theorem 2.3 is not only true but also can be essentially generalized for the case of arbitrary compatible Riemannian or pseudo-Riemannian metrics, in particular, for the especially important cases in the theory of systems of hydrodynamiic type, namely, the cases of metrics of constant Riemannian curvature or the metrics generating the general nonlocal Poisson structures of hydrodynamic type. Namely, the following theorem is true.
###### Theorem 2.4
* If for any linear combination (1.3) of two metrics $`g_1^{ij}(u)`$ and $`g_2^{ij}(u)`$ condition (1.4) is fulfilled, then the Nijenhuis tensor of the affinor
$$v_j^i(u)=g_1^{is}(u)g_{2,sj}(u)$$
vanishes. Thus, for compatible and almost compatible metrics, the corresponding Nijenhuis tensor always vanishes.
* If a pair of metrics $`g_1^{ij}(u)`$ and $`g_2^{ij}(u)`$ is nonsingular, that is, the roots of the equation
$$det(g_1^{ij}(u)\lambda g_2^{ij}(u))=0,$$
(2.17)
are distinct, then it follows from the vanishing of the Nijenhuis tensor of the affinor $`v_j^i(u)=g_1^{is}(u)g_{2,sj}(u)`$ that the metrics $`g_1^{ij}(u)`$ and $`g_2^{ij}(u)`$ are compatible. Thus, a nonsingular pair of metrics is compatible if and only if the metrics are almost compatible.
## 3 Almost compatible metrics and Nijenhuis tensor
Let us consider two contravariant Riemannian or pseudo-Riemannian metrics $`g_1^{ij}(u)`$ and $`g_2^{ij}(u)`$, and also the corresponding coefficients of the Levi–Civita connections $`\mathrm{\Gamma }_{1,k}^{ij}(u)`$ and $`\mathrm{\Gamma }_{2,k}^{ij}(u)`$.
We introduce the tensor
$$M^{ijk}(u)=g_1^{is}(u)\mathrm{\Gamma }_{2,s}^{jk}(u)g_2^{js}(u)\mathrm{\Gamma }_{1,s}^{ik}(u)g_1^{js}(u)\mathrm{\Gamma }_{2,s}^{ik}(u)+g_2^{is}(u)\mathrm{\Gamma }_{1,s}^{jk}(u).$$
(3.1)
It follows from the following representation that $`M^{ijk}(u)`$ is actually a tensor:
$`M^{ijk}(u)=g_1^{is}(u)g_2^{jp}(u)(\mathrm{\Gamma }_{2,ps}^k(u)\mathrm{\Gamma }_{1,ps}^k(u))`$
$`g_1^{js}(u)g_2^{ip}(u)(\mathrm{\Gamma }_{2,ps}^k(u)\mathrm{\Gamma }_{1,ps}^k(u))=\mathrm{\Delta }^{ijk}(u)\mathrm{\Delta }^{jik}(u).`$ (3.2)
###### Lemma 3.1
The tensor $`M^{ijk}(u)`$ vanishes if and only if the metrics $`g_1^{ij}(u)`$ and $`g_2^{ij}(u)`$ are almost compatible.
Let us introduce the affinor
$$v_j^i(u)=g_1^{is}(u)g_{2,sj}(u)$$
(3.3)
and consider the Nijenhuis tensor of this affinor
$$N_{ij}^k(u)=v_i^s(u)\frac{v_j^k}{u^s}v_j^s(u)\frac{v_i^k}{u^s}+v_s^k(u)\frac{v_i^s}{u^j}v_s^k\frac{v_j^s}{u^i},$$
(3.4)
following , where were similarly considered the affinor $`v_j^i(u)`$ and its Nijenhuis tensor for two flat metrics.
###### Theorem 3.1
The metrics $`g_1^{ij}(u)`$ and $`g_2^{ij}(u)`$ are almost compatible if and only if the corresponding Nijenhuis tensor $`N_{ij}^k(u)`$ (3.4) vanishes.
###### Lemma 3.2
The following identities are always fulfilled:
$$g_{1,sp}(u)N_{rq}^p(u)g_2^{ri}(u)g_2^{qj}(u)g_2^{sk}(u)=M^{kij}(u)+M^{jki}(u)+M^{jik}(u),$$
(3.5)
$`2(M^{jki}(u)+M^{jik}(u))=`$
$`g_{1,sp}(u)N_{rq}^p(u)g_2^{ri}(u)g_2^{qj}(u)g_2^{sk}(u)+g_{1,sp}(u)N_{rq}^p(u)g_2^{rk}(u)g_2^{qj}(u)g_2^{si}(u),`$ (3.6)
$$2M^{kij}(u)=g_{1,sp}(u)N_{rq}^p(u)g_2^{ri}(u)g_2^{qj}(u)g_2^{sk}(u)g_{1,sp}(u)N_{rq}^p(u)g_2^{rk}(u)g_2^{qj}(u)g_2^{si}(u).$$
(3.7)
###### Corollary 3.1
The tensor $`M^{ijk}(u)`$ vanishes if and only if the Nijenhuis tensor (3.4) vanishes.
In the papers , the present author studied reductions in the general problem on compatible flat metrics, connected with the associativity equations, namely, the following ansatz in formula (2.15):
$$h^i(v)=\eta ^{is}\frac{\mathrm{\Phi }}{v^s},$$
where $`\mathrm{\Phi }(v^1,\mathrm{},v^N)`$ is a function of $`N`$ variables.
Correspondingly, in this case the metrics have the form:
$$g_1^{ij}(v)=\eta ^{ij},g_2^{ij}(v)=\eta ^{is}\eta ^{jp}\frac{^2\mathrm{\Phi }}{v^sv^p}.$$
(3.8)
###### Theorem 3.2
If metrics (3.8) are almost compatible, then they are always compatible. Moreover, in this case, the metric $`g_2^{ij}(v)`$ is necessarily also flat, that is, metrics (3.8) form a flat pencil of metrics. Condition of almost compatibility of metrics (3.8) has the form
$$\eta ^{sp}\frac{^2\mathrm{\Phi }}{v^pv^i}\frac{^3\mathrm{\Phi }}{v^sv^jv^k}=\eta ^{sp}\frac{^2\mathrm{\Phi }}{v^pv^k}\frac{^3\mathrm{\Phi }}{v^sv^jv^i}$$
(3.9)
and coincides with the condition of compatible deformation of two Frobenius algebras which was derived and studied by the author in .
In particular, in the author’s papers , it is proved that in the two-component case ($`N=2`$), for $`\eta ^{ij}=\epsilon ^i\delta ^{ij},`$ $`\epsilon ^i=\pm 1,`$ condition (3.9) is equivalent to the following linear second-order partial differential equation with constant coefficients:
$$\alpha \left(\epsilon ^1\frac{^2\mathrm{\Phi }}{(v^1)^2}\epsilon ^2\frac{^2\mathrm{\Phi }}{(v^2)^2}\right)=\beta \frac{^2\mathrm{\Phi }}{v^1v^2},$$
(3.10)
where $`\alpha `$ and $`\beta `$ are arbitrary constants.
## 4 Compatible metrics and Nijenhuis tensor
Let us prove the second part of Theorem 2.4. In the previous section it is proved, in particular, that it always follows from compatibility (and even almost compatibility) of metrics that the corresponding Nijenhuis tensor vanishes.
Assume that a pair of metrics $`g_1^{ij}(u)`$ and $`g_2^{ij}(u)`$ is nonsingular, that is, the eigenvalues of this pair of the metrics are distinct, and assume also that the corresponding Nijenhuis tensor vanishes. Let us prove that, in this case, the metrics $`g_1^{ij}(u)`$ and $`g_2^{ij}(u)`$ are compatible (their almost compatibility follows from the previous section).
Obviously, that the eigenvalues of the pair of the metrics $`g_1^{ij}(u)`$ and $`g_2^{ij}(u)`$ coincide with the eigenvalues of affinor $`v_j^i(u)`$. But it is well-known that if all eigenvalues of an affinor are distinct, then it always follows from the vanishing of the Nijenhuis tensor of this affinor that there exist local coordinates such that in these coordinates the affinor reduces to a diagonal form in the corresponding neighbourhood (see also ).
So we can consider further that the affinor $`v_j^i(u)`$ is diagonal in the local coordinates $`u^1,\mathrm{},u^N`$,
$$v_j^i(u)=\lambda ^i(u)\delta _j^i,$$
(4.1)
where is no summation over the index $`i`$, and $`\lambda ^i(u),i=1,\mathrm{},N,`$ are the eigenvalues of the pair of metrics $`g_1^{ij}(u)`$ and $`g_2^{ij}(u)`$, which are assumed distinct:
$$\lambda ^i\lambda ^j,\mathrm{for}ij.$$
(4.2)
###### Lemma 4.1
If the affinor (3.3) is diagonal in a certain local coordinates and all its eigenvalues are distinct, then, in these coordinates, the metrics $`g_1^{ij}(u)`$ and $`g_2^{ij}(u)`$ are also necessarily diagonal.
Actually, we have
$$g_1^{ij}(u)=\lambda ^i(u)g_2^{ij}(u).$$
It follows from symmetry of the metrics $`g_1^{ij}(u)`$ and $`g_2^{ij}(u)`$ that for any $`i,j`$
$$(\lambda ^i(u)\lambda ^j(u))g_2^{ij}(u)=0,$$
(4.3)
where is no summation over indices, that is,
$$g_2^{ij}(u)=g_1^{ij}(u)=0\mathrm{for}ij.$$
###### Lemma 4.2
Let an affinor $`w_j^i(u)`$ be diagonal in a certain local coordinates $`u=(u^1,\mathrm{},u^N)`$, that is, $`w_j^i(u)=\mu ^i(u)\delta _j^i`$.
* If all the eigenvalues of an diagonal affinor are distinct, that is, $`\mu ^i(u)\mu ^j(u)`$ for $`ij`$, then the Nijenhuis tensor of this affinor vanishes if and only if the $`i`$th eigenvalue $`\mu ^i(u)`$ depend only on the coordinate $`u^i.`$
* If all the eigenvalues coincide, then the Nijenhuis tensor vanishes.
* In the general case of a diagonal affinor, the Nijenhuis tensor vanishes if and only if
$$\frac{\mu ^i}{u^j}=0$$
(4.4)
for all $`i,j`$ such that $`\mu ^i(u)\mu ^j(u).`$
Actually, for any diagonal affinor $`w_j^i(u)=\mu ^i(u)\delta _j^i,`$ the Nijenhuis tensor $`N_{ij}^k(u)`$ has the form:
$$N_{ij}^k(u)=(\mu ^i\mu ^k)\frac{\mu ^j}{u^i}\delta ^{kj}(\mu ^j\mu ^k)\frac{\mu ^i}{u^j}\delta ^{ki}$$
(no summation over indices). Thus, the Nijenhuis tensor vanishes if and only if for all $`i,j`$
$$(\mu ^i(u)\mu ^j(u))\frac{\mu ^i}{u^j}=0,$$
where is no summation over indices.
It follows immediately from Lemmas 4.1 and 4.2 that for any nonsingular pair of almost compatible metrics there always exist local coordinates in which the metrics have the form
$$g_2^{ij}(u)=g^i(u)\delta ^{ij},g_1^{ij}(u)=\lambda ^i(u^i)g^i(u)\delta ^{ij},\lambda ^i=\lambda ^i(u^i),i=1,\mathrm{},N.$$
Moreover, we derive immediately that any diagonal metrics of the form $`g_2^{ij}(u)=g^i(u)\delta ^{ij}`$ and $`g_1^{ij}(u)=f^i(u^i)g^i(u)\delta ^{ij}`$ for any nonzero functions $`f^i(u^i),`$ $`i=1,\mathrm{},N,`$ (they can be here, for example, coinciding nonzero constants, that is, the pair of metrics may be “singular”) are always almost compatible. We shall prove now that they are always compatible. Then Theorem 2.4 will be completely proved.
Consider diagonal metrics $`g_2^{ij}(u)=g^i(u)\delta ^{ij}`$ and $`g_1^{ij}(u)=f^i(u^i)g^i(u)\delta ^{ij},`$ where $`f^i(u^i),`$ $`i=1,\mathrm{},N,`$ are arbitrary functions of a single variable, which are not equal to zero identically, and consider their arbitrary linear combination
$$g^{ij}(u)=(\lambda _2+\lambda _1f^i(u^i))g^i(u)\delta ^{ij},$$
where $`\lambda _1`$ and $`\lambda _2`$ are arbitrary constants such that $`det(g^{ij}(u))0.`$
Let us prove that relation (1.5) is always fulfilled for the corresponding tensors of Riemannian curvature.
Recall that for any diagonal metric $`\mathrm{\Gamma }_{jk}^i(u)=0`$ if all the indices $`i,j,k`$ are distinct. Correspondingly, $`R_{kl}^{ij}(u)=0`$ if all the indices $`i,j,k,l`$ are distinct. Besides, as a result of the well-known symmetries of the tensor of Riemannian curvature we have:
$$R_{kl}^{ii}(u)=R_{kk}^{ij}(u)=0,$$
$$R_{il}^{ij}(u)=R_{li}^{ij}(u)=R_{li}^{ji}(u)=R_{il}^{ji}(u).$$
Thus, it is sufficient to prove relation (1.5) only for the following components of the tensor of Riemannian curvature: $`R_{il}^{ij}(u)`$, where $`ij,`$ $`il`$.
For any diagonal metric $`g_2^{ij}(u)=g^i(u)\delta ^{ij}`$ we have
$$\mathrm{\Gamma }_{2,ik}^i(u)=\mathrm{\Gamma }_{2,ki}^i(u)=\frac{1}{2g^i(u)}\frac{g^i}{u^k},\mathrm{for}\mathrm{any}i,k;$$
$$\mathrm{\Gamma }_{2,jj}^i(u)=\frac{1}{2}\frac{g^i(u)}{(g^j(u))^2}\frac{g^j}{u^i},ij.$$
$`R_{2,il}^{ij}(u)=g^i(u)R_{2,iil}^j(u)=`$
$`g^i(u)\left({\displaystyle \frac{\mathrm{\Gamma }_{2,il}^j}{u^i}}+{\displaystyle \frac{\mathrm{\Gamma }_{2,ii}^j}{u^l}}{\displaystyle \underset{s=1}{\overset{N}{}}}\mathrm{\Gamma }_{2,si}^j(u)\mathrm{\Gamma }_{2,il}^s(u)+{\displaystyle \underset{s=1}{\overset{N}{}}}\mathrm{\Gamma }_{2,sl}^j(u)\mathrm{\Gamma }_{2,ii}^s(u)\right).`$ (4.5)
It is necessary to consider separately two different cases.
1) $`jl`$.
$`R_{2,il}^{ij}(u)=g^i(u)\left({\displaystyle \frac{\mathrm{\Gamma }_{2,ii}^j}{u^l}}\mathrm{\Gamma }_{2,ii}^j(u)\mathrm{\Gamma }_{2,il}^i(u)+\mathrm{\Gamma }_{2,jl}^j(u)\mathrm{\Gamma }_{2,ii}^j(u)+\mathrm{\Gamma }_{2,ll}^j(u)\mathrm{\Gamma }_{2,ii}^l(u)\right)=`$
$`{\displaystyle \frac{1}{2}}g^i(u){\displaystyle \frac{}{u^l}}\left({\displaystyle \frac{g^j(u)}{(g^i(u))^2}}{\displaystyle \frac{g^i}{u^j}}\right)+{\displaystyle \frac{1}{4}}{\displaystyle \frac{g^j(u)}{(g^i(u))^2}}{\displaystyle \frac{g^i}{u^j}}{\displaystyle \frac{g^i}{u^l}}`$
$`{\displaystyle \frac{1}{4g^i(u)}}{\displaystyle \frac{g^i}{u^j}}{\displaystyle \frac{g^j}{u^l}}+{\displaystyle \frac{1}{4}}{\displaystyle \frac{g^j(u)}{g^i(u)g^l(u)}}{\displaystyle \frac{g^l}{u^j}}{\displaystyle \frac{g^i}{u^l}}.`$ (4.6)
Respectively, for the metric
$$g^{ij}(u)=(\lambda _2+\lambda _1f^i(u^i))g^i(u)\delta ^{ij},$$
we obtain (we use here that all the indices $`i,j,l`$ are distinct):
$`R_{il}^{ij}(u)=(\lambda _2+\lambda _1f^j(u^j))[{\displaystyle \frac{1}{2}}g^i(u){\displaystyle \frac{}{u^l}}\left({\displaystyle \frac{g^j(u)}{(g^i(u))^2}}{\displaystyle \frac{g^i}{u^j}}\right)+{\displaystyle \frac{1}{4}}{\displaystyle \frac{g^j(u)}{(g^i(u))^2}}{\displaystyle \frac{g^i}{u^j}}{\displaystyle \frac{g^i}{u^l}}`$
$`{\displaystyle \frac{1}{4g^i(u)}}{\displaystyle \frac{g^i}{u^j}}{\displaystyle \frac{g^j}{u^l}}+{\displaystyle \frac{1}{4}}{\displaystyle \frac{g^j(u)}{g^i(u)g^l(u)}}{\displaystyle \frac{g^l}{u^j}}{\displaystyle \frac{g^i}{u^l}}]=\lambda _1R^{ij}_{1,il}(u)+\lambda _2R^{ij}_{2,il}(u).`$ (4.7)
2) $`j=l`$.
$`R_{2,ij}^{ij}(u)=g^i(u)({\displaystyle \frac{\mathrm{\Gamma }_{2,ij}^j}{u^i}}+{\displaystyle \frac{\mathrm{\Gamma }_{2,ii}^j}{u^j}}\mathrm{\Gamma }_{2,ii}^j(u)\mathrm{\Gamma }_{2,ij}^i(u)`$
$`\mathrm{\Gamma }_{2,ji}^j(u)\mathrm{\Gamma }_{2,ij}^j(u)+{\displaystyle \underset{s=1}{\overset{N}{}}}\mathrm{\Gamma }_{2,sj}^j(u)\mathrm{\Gamma }_{2,ii}^s(u))=`$
$`{\displaystyle \frac{1}{2}}g^i(u){\displaystyle \frac{}{u^i}}\left({\displaystyle \frac{1}{g^j(u)}}{\displaystyle \frac{g^j}{u^i}}\right)+{\displaystyle \frac{1}{2}}g^i(u){\displaystyle \frac{}{u^j}}\left({\displaystyle \frac{g^j(u)}{(g^i(u))^2}}{\displaystyle \frac{g^i}{u^j}}\right)+{\displaystyle \frac{1}{4}}{\displaystyle \frac{g^j(u)}{(g^i(u))^2}}{\displaystyle \frac{g^i}{u^j}}{\displaystyle \frac{g^i}{u^j}}`$
$`{\displaystyle \frac{1}{4}}{\displaystyle \frac{g^i(u)}{(g^j(u))^2}}{\displaystyle \frac{g^j}{u^i}}{\displaystyle \frac{g^j}{u^i}}+{\displaystyle \frac{1}{4g^j(u)}}{\displaystyle \frac{g^j}{u^i}}{\displaystyle \frac{g^i}{u^i}}{\displaystyle \underset{si}{}}{\displaystyle \frac{1}{4}}{\displaystyle \frac{g^s(u)}{g^i(u)g^j(u)}}{\displaystyle \frac{g^j}{u^s}}{\displaystyle \frac{g^i}{u^s}}.`$ (4.8)
Respectively, for the metric
$$g^{ij}(u)=(\lambda _2+\lambda _1f^i(u^i))g^i(u)\delta ^{ij},$$
we obtain (we use here that the indices $`i,j`$ are distinct):
$`R_{ij}^{ij}(u)={\displaystyle \frac{1}{2}}(\lambda _2+\lambda _1f^i(u^i))g^i(u){\displaystyle \frac{}{u^i}}\left({\displaystyle \frac{1}{g^j(u)}}{\displaystyle \frac{g^j}{u^i}}\right)+`$
$`{\displaystyle \frac{1}{2}}g^i(u){\displaystyle \frac{}{u^j}}\left({\displaystyle \frac{(\lambda _2+\lambda _1f^j(u^j))g^j(u)}{(g^i(u))^2}}{\displaystyle \frac{g^i}{u^j}}\right)+{\displaystyle \frac{1}{4}}(\lambda _2+\lambda _1f^j(u^j)){\displaystyle \frac{g^j(u)}{(g^i(u))^2}}{\displaystyle \frac{g^i}{u^j}}{\displaystyle \frac{g^i}{u^j}}`$
$`{\displaystyle \frac{1}{4}}(\lambda _2+\lambda _1f^i(u^i)){\displaystyle \frac{g^i(u)}{(g^j(u))^2}}{\displaystyle \frac{g^j}{u^i}}{\displaystyle \frac{g^j}{u^i}}+{\displaystyle \frac{1}{4g^j(u)}}{\displaystyle \frac{g^j}{u^i}}{\displaystyle \frac{((\lambda _2+\lambda _1f^i(u^i))g^i)}{u^i}}`$
$`{\displaystyle \frac{1}{4g^i(u)}}{\displaystyle \frac{g^i}{u^j}}{\displaystyle \frac{((\lambda _2+\lambda _1f^j(u^j)))g^j}{u^j}}{\displaystyle \underset{si,sj}{}}{\displaystyle \frac{1}{4}}{\displaystyle \frac{(\lambda _2+\lambda _1f^s(u^s))g^s(u)}{g^i(u)g^j(u)}}{\displaystyle \frac{g^j}{u^s}}{\displaystyle \frac{g^i}{u^s}}=`$
$`\lambda _1R_{1,ij}^{ij}(u)+\lambda _2R_{2,ij}^{ij}(u).`$ (4.9)
Theorem 2.4 is proved. Moreover, the complete explicit description of nonsingular pairs of compatible and almost compatible metrics is obtained and the following theorem is proved:
###### Theorem 4.1
An arbitrary nonsingular pair of metrics is compatible if and only if there exist local coordinates $`u=(u^1,\mathrm{},u^N)`$ such that $`g_2^{ij}(u)=g^i(u)\delta ^{ij}`$ and $`g_1^{ij}(u)=f^i(u^i)g^i(u)\delta ^{ij},`$ where $`f^i(u^i),`$ $`i=1,\mathrm{},N,`$ are arbitrary functions of single variables (of course, in the case of nonsingular pair of metrics, these functions are not equal to each other if they are constants and they are not equal identically to zero).
Let us consider here the problem on nonsingular pairs of compatible flat metrics. It follows from Theorem 4.1 that it is sufficient to classify flat metrics of the form $`g_2^{ij}(u)=g^i(u)\delta ^{ij}`$ and $`g_1^{ij}(u)=f^i(u^i)g^i(u)\delta ^{ij},`$ where $`f^i(u^i),`$ $`i=1,\mathrm{},N,`$ are arbitrary functions of single variables.
The problem of description of diagonal flat metrics, that is, flat metrics $`g_2^{ij}(u)=g^i(u)\delta ^{ij},`$ is a classical problem of differential geometry. This problem is equivalent to the problem of description of curvilinear orthogonal coordinate systems in a pseudo-Euclidean space and it was studied in detail and mainly solved in the beginning of the 20th century (see ). Locally, such coordinate systems are determined by $`n(n1)/2`$ arbitrary functions of two variables (see , ). Recently, Zakharov showed that the Lamé equations describing curvilinear orthogonal coordinate systems can be integrated by the inverse scattering method (see also an algebraic-geometric approach in ). The condition that the metric $`g_1^{ij}(u)=f^i(u^i)g^i(u)\delta ^{ij}`$ is also flat gives exactly $`n(n1)/2`$ equations which are linear with respect to the functions $`f^i(u^i)`$. Note that, in this case, components (4.6) of tensor of Riemannian curvature automatically vanish as a result of formula (4.7). And the vanishing of components (4.8) gives the corresponding $`n(n1)/2`$ equations. In particular, in the case $`N=2`$ this completely solves the problem of description of nonsingular pairs of compatible flat metrics. In the next section we give their complete description. But without a doubt this problem is integrable in general for any $`N`$. We are going to devote to this question another work. In particular, it is very interesting to classify all the $`n`$-orthogonal curvilinear coordinate systems in a pseudo-Euclidean space, for which the functions $`f^i(u^i)=(u^i)^n`$ define compatible flat metrics (respectivly, separately for $`n=1`$; $`n=1,2`$; $`n=1,2,3,`$ and so on).
## 5 Two-component compatible flat metrics
We present here the complete description of nonsingular pairs of two-component compatible flat metrics (see also , , , where an integrable homogeneous system of hydrodynamic type, describing all the two-component compatible flat metrics, was derived and investigated).
It is proved above that for any nonsingular pair of two-component compatible metrics $`g_1^{ij}(u)`$ and $`g_2^{ij}(u)`$ there always exist local coordinates $`u^1,\mathrm{},u^N`$ such that
$$(g_2^{ij}(u))=\left(\begin{array}{cc}\frac{\epsilon ^1}{(b^1(u))^2}& 0\\ 0& \frac{\epsilon ^2}{(b^2(u))^2}\end{array}\right),(g_1^{ij}(u))=\left(\begin{array}{cc}\frac{\epsilon ^1f^1(u^1)}{(b^1(u))^2}& 0\\ 0& \frac{\epsilon ^2f^2(u^2)}{(b^2(u))^2}\end{array}\right),$$
(5.1)
where $`\epsilon ^i=\pm 1,i=1,2;`$ $`b^i(u)`$ and $`f^i(u^i),i=1,2,`$ are arbitrary nonzero functions of the corresponding single variables.
###### Lemma 5.1
An arbitrary diagonal metric $`g_2^{ij}(u)`$ (5.1) is flat if and only if the functions $`b^i(u),`$ $`i=1,2,`$ are solutions of the following linear system:
$$\frac{b^2}{u^1}=\epsilon ^1\frac{F}{u^2}b^1(u),\frac{b^1}{u^2}=\epsilon ^2\frac{F}{u^1}b^2(u),$$
(5.2)
where $`F(u)`$ is an arbitrary function.
###### Theorem 5.1
The metrics $`g_1^{ij}(u)`$ and $`g_2^{ij}(u)`$ (5.1) form a flat pencil of metrics if and only if the functions $`b^i(u),i=1,2,`$ are solutions of the linear system (5.2), where the function $`F(u)`$ is a solution of the following linear equation:
$$2\frac{^2F}{u^1u^2}(f^1(u^1)f^2(u^2))+\frac{F}{u^2}\frac{df^1(u^1)}{du^1}\frac{F}{u^1}\frac{df^2(u^2)}{du^2}=0.$$
(5.3)
In the case, if the eigenvalues of the pair of the metrics $`g_1^{ij}(u)`$ and $`g_2^{ij}(u)`$ are not only distinct but also are not constants, we can always choose local coordinates such that $`f^1(u^1)=u^1,`$ $`f^2(u^2)=u^2`$ (see also remark in ). In this case, equation (5.3) has the form
$$2\frac{^2F}{u^1u^2}(u^1u^2)+\frac{F}{u^2}\frac{F}{u^1}=0.$$
(5.4)
Let us continue this recurrent procedure for the metrics $`G_{n+1}^{ij}(u)=v_j^i(u)G_n^{ij}(u)`$ with the help of the affinor $`v_j^i(u)=u^i\delta _j^i.`$
###### Theorem 5.2
Three metrics
$$(G_n^{ij}(u))=\left(\begin{array}{cc}\frac{\epsilon ^1(u^1)^n}{(b^1(u))^2}& 0\\ 0& \frac{\epsilon ^2(u^2)^n}{(b^2(u))^2}\end{array}\right),n=0,1,2,$$
(5.5)
form a flat pencil of metrics (pairwise compatible) if and only if the functions $`b^i(u),i=1,2,`$ are solutions of the linear system (5.2), where
$$F(u)=c\mathrm{ln}(u^1u^2),$$
(5.6)
$`c`$ is an arbitrary constant. Already the metric $`G_3^{ij}(u)`$ is flat only in the most trivial case, when $`c=0,`$ and, respectively, $`b^1=b^1(u^1),`$ $`b^2=b^2(u^2)`$.
The metric $`G_3^{ij}(u)`$ is a metric of nonzero constant Riemannian curvature $`K0`$ if and only if
$$(b^1(u))^2=(b^2(u))^2=\frac{\epsilon ^2}{4K}(u^1u^2),\epsilon ^1=\epsilon ^2,c=\pm \frac{1}{2}.$$
(5.7)
## 6 Almost compatible metrics which are not compatible
###### Lemma 6.1
Two-component diagonal conformally Euclidean metric
$$g^{ij}(u)=\mathrm{exp}(a(u))\delta ^{ij},1i,j2,$$
is flat if and only if the function $`a(u)`$ is harmonic, that is,
$$\mathrm{\Delta }a\frac{^2a}{(u^1)^2}+\frac{^2a}{(u^2)^2}=0.$$
(6.1)
In particular, the metric $`g_1^{ij}(u)=\mathrm{exp}(u^1u^2)\delta ^{ij},1i,j2`$, is flat. Obviously, that the flat metrics $`g_1^{ij}(u)=\mathrm{exp}(u^1u^2)\delta ^{ij},1i,j2`$, and $`g_2^{ij}(u)=\delta ^{ij},1i,j2,`$ are almost compatible, for them the Nijenhuis tensor (3.4) vanishes. But it follows from Lemma 6.1 that these metrics are not compatible, their sum is not a flat metric.
Similarly it is possible to construct also other counterexamples to Theorem 2.3. Moreover, the following statement is true.
###### Proposition 6.1
Any nonconstant real harmonic function $`a(u)`$ defines a pair of almost compatible metrics $`g_1^{ij}(u)=\mathrm{exp}(a(u))\delta ^{ij},1i,j2`$, and $`g_2^{ij}(u)=\delta ^{ij},1i,j2,`$ which are not compatible. These metrics are compatible if and only if $`a=a(u^1\pm iu^2).`$
Let us construct also almost compatible metrics of constant Riemannian curvature, which are not compatible.
###### Lemma 6.2
Two-component diagonal conformally Euclidean metric
$$g^{ij}(u)=\mathrm{exp}(a(u))\delta ^{ij},1i,j2,$$
is a metric of constant Riemannian curvature $`K`$ if and only if the function $`a(u)`$ is a solution of the Liouville equation
$$\mathrm{\Delta }a\frac{^2a}{(u^1)^2}+\frac{^2a}{(u^2)^2}=2Ke^{a(u)}.$$
(6.2)
###### Proposition 6.2
For the metrics $`g_1^{ij}(u)=\mathrm{exp}(a(u))\delta ^{ij},1i,j2`$, and $`g_2^{ij}(u)=\delta ^{ij},1i,j2,`$ the corresponding Nijenhuis tensor vanishes, that is, they are always almost compatible. But they are compatible metrics of constant Riemannian curvature $`K0`$ and $`0`$, respectively, if and only if the function $`a(u)`$ is constant.
Note, that all the one-component “metrics” are always compatible, and all the one-component local Poisson structures of hydrodynamic type are also always compatible. Let us construct for any $`N>1`$ examples of almost compatible metrics which are not compatible.
###### Proposition 6.3
The metrics $`g_1^{ij}(u)=b(u)\delta ^{ij},1i,jN`$ and $`g_2^{ij}(u)=\delta ^{ij},1i,jN,`$ where $`b(u)`$ is an arbitrary function, are always almost compatible, the corresponding Nijenhuis tensor vanishes. But they are compatible real metrics only in the most trivial case, when the function $`b(u)`$ is constant. Complex metrics are compatible if and only if either the function $`b(u)`$ is constant, or $`N=2`$ and $`b(u)=b(u^1\pm iu^2)`$.
## 7 Compatible metrics and nonlocal Poisson structures of hydrodynamic type
Nonlocal Poisson structures of hydrodynamic type were introduced and studied in the work of the present author and Ferapontov (see also ). They have the following form:
$$\{I,J\}=\frac{\delta I}{\delta u^i(x)}\left(g^{ij}(u(x))\frac{d}{dx}+b_k^{ij}(u(x))u_x^k+Ku_x^i\left(\frac{d}{dx}\right)^1u_x^j\right)\frac{\delta J}{\delta u^j(x)}𝑑x,$$
(7.1)
where $`K`$ is an arbitrary constant.
The bracket of form (7.1) is called nondegenerate if $`det(g^{ij})(u)0`$.
###### Theorem 7.1 ()
If $`det(g^{ij})(u)0`$, then bracket (7.1) is a Poisson bracket, that is, it is skew-symmetric and satisfies the Jacobi identity, if and only if
* $`g^{ij}(u)`$ is an arbitrary pseudo-Riemannian contravariant metric of constant Riemannian curvature $`K`$,
* $`b_k^{ij}(u)=g^{is}(u)\mathrm{\Gamma }_{sk}^j(u),`$ where $`\mathrm{\Gamma }_{sk}^j(u)`$ is the Riemannian connection generated by the contravariant metric $`g^{ij}(u)`$ (the Levi–Civita connection).
###### Proposition 7.1
Nonlocal nondegenerate Poisson brackets of form (7.1) are compatible if and only if their metrics are compatible.
In Ferapontov introduced and studied more general nonlocal Poisson brackets of hydrodynamic type, namely, the brackets of the following form:
$`\{I,J\}={\displaystyle }{\displaystyle \frac{\delta I}{\delta u^i(x)}}(g^{ij}(u(x)){\displaystyle \frac{d}{dx}}+b_k^{ij}(u(x))u_x^k`$
$`+{\displaystyle \underset{\alpha =1}{\overset{L}{}}}(w^\alpha )_k^i(u)u_x^k\left({\displaystyle \frac{d}{dx}}\right)^1(w^\alpha )_s^j(u)u_x^s){\displaystyle \frac{\delta J}{\delta u^j(x)}}dx,det(g^{ij})(u)0.`$ (7.2)
###### Theorem 7.2 ()
Bracket (7.2) is a Poisson bracket, that is, it is skew-symmetric and satisfies the Jacobi identity, if and only if
* $`b_k^{ij}(u)=g^{is}(u)\mathrm{\Gamma }_{sk}^j(u),`$ where $`\mathrm{\Gamma }_{sk}^j(u)`$ is the Riemannian connection generated by the contravariant metric $`g^{ij}(u)`$ (the Levi–Civita connection),
* the pseudo-Riemannian metric $`g^{ij}(u)`$ and the set of affinors $`(w^\alpha )_j^i(u)`$ satisfy the relations:
$$g_{ik}(u)(w^\alpha )_j^k(u)=g_{jk}(u)(w^\alpha )_i^k(u),\alpha =1,\mathrm{},L,$$
(7.3)
$$_k(w^\alpha )_j^i(u)=_j(w^\alpha )_k^i(u),\alpha =1,\mathrm{},L,$$
(7.4)
$$R_{kl}^{ij}(u)=\underset{\alpha =1}{\overset{L}{}}\left((w^\alpha )_k^i(u)(w^\alpha )_l^j(u)(w^\alpha )_k^j(u)(w^\alpha )_l^i(u)\right).$$
(7.5)
Moreover, the family of affinors $`w^\alpha (u)`$ is commutative: $`[w^\alpha ,w^\beta ]=0.`$
###### Proposition 7.2
If nonlocal Poisson brackets of form (7.2) are compatible, then their metrics are also compatible.
Actually, if nonlocal Poisson brackets of form (7.2) are compatible, then it follows from the conditions of compatibility and from Theorem 7.2 that, first, relation (1.4) is fulfilled and, secondly, the curvature tensor for the metric $`g^{ij}(u)=\lambda _1g_1^{ij}(u)+\lambda _2g_2^{ij}(u)`$ has the form
$`R_{kl}^{ij}(u)={\displaystyle \underset{\alpha =1}{\overset{L_1}{}}}\lambda _1\left((w_1^\alpha )_k^i(u)(w_1^\alpha )_l^j(u)(w_1^\alpha )_k^j(u)(w_1^\alpha )_l^i(u)\right)+`$
$`{\displaystyle \underset{\alpha =1}{\overset{L_2}{}}}\lambda _2\left((w_2^\alpha )_k^i(u)(w_2^\alpha )_l^j(u)(w_2^\alpha )_k^j(u)(w_2^\alpha )_l^i(u)\right)=\lambda _1R_{1,kl}^{ij}(u)+\lambda _2R_{2,kl}^{ij}(u).`$
Apparently, the converse statement is also always true (this is only our conjecture which is not strictly proved in the most general case at this moment).
Relations (7.3)–(7.5) are nothing but the Gauss–Peterson–Codazzi equations for $`N`$-dimensional surfaces $`M`$ with flat normal connections in a pseudo-Euclidean space $`E^{N+L}`$. Here $`g_{ij}(u)`$ is the first fundamental form of the surface $`M`$, and $`w^\alpha (u)`$ are the Weingarten operators . We consider more in detail the case of compatible nonlocal Poisson structures of hydrodynamic type that correspond to surfaces with holonomic net of curvature lines (see ). This case is the most interesting for applications (here we do not give numerous important examples, see, for example, in , , or in the author’s survey ).
###### Proposition 7.3
Let two nonlocal Poisson brackets of form (7.2) correspond to surfaces with holonomic net of curvature lines and be given in coordinates of curvature lines. In this case, if the corresponding pair of metrics is nonsingular, then the nonlocal Poisson structures are compatible if and only if their metrics are compatible.
In this case the metrics $`g_1^{ij}(u)=g_1^i(u)\delta ^{ij}`$ and $`g_2^{ij}(u)=g_2^i(u)\delta ^{ij}`$, and also the Weingarten operators $`(w_1^\alpha )_j^i(u)=(w_1^\alpha )^i(u)\delta _j^i`$ and $`w_2^\alpha (u)=(w_2^\alpha )^i(u)\delta _j^i`$ are diagonal in the coordinates under consideration. For any such “diagonal” case, condition (7.3) is automatically fulfilled, all the Weingarten operators commute, conditions (7.4) and (7.5) have the following form, respectively:
$$2g^i(u)\frac{(w^\alpha )^i}{u^k}=((w^\alpha )^i(w^\alpha )^k)\frac{g^i}{u^k}\mathrm{for}\mathrm{all}ik,$$
(7.6)
$$R_{ij}^{ij}(u)=\underset{\alpha =1}{\overset{L}{}}(w^\alpha )^i(u)(w^\alpha )^j(u).$$
(7.7)
It follows from nonsingularity of the pair of the metrics and from compatibility of the metrics that the corresponding Nijenhuis tensor vanishes and there exist functions $`f^i(u^i),i=1,\mathrm{},N,`$ such that:
$$g_1^i(u)=f^i(u^i)g_2^i(u).$$
Using relations (7.6) and (7.7), it is easy to prove that in this case it follows from compatibility of the metrics that an arbitrary linear combination of nonlocal Poisson brackets under consideration is also a Poisson bracket.
Centre for Nonlinear Studies,
L.D.Landau Institute for Theoretical Physics,
Russian Academy of Sciences,
ul. Kosygina, 2,
Moscow, 117940 Russia
e-mail: mokhov@genesis.mi.ras.ru, mokhov@landau.ac.ru
Department of Mathematics,
University of Paderborn,
Paderborn, Germany
e-mail: mokhov@uni-paderborn.de
|
warning/0005/nucl-th0005052.html
|
ar5iv
|
text
|
# How to Renormalize the Gap Equation in High Density QCD
## Abstract
We discuss two technical issues related to the gap equation in high-density QCD: i) how to obtain the asymptotic solution with well controlled approximations, and ii) the renormalization of four-quark operators in the high-density effective field theory.
preprint: DOE/ER/40561-95-INT00 NT@UW-00-14
Recently Son obtained the leading exponential behavior of the superconducting gap in QCD at asymptotically high density using both an indirect renormalization group argument and a direct QCD calculation. Subsequently many papers confirmed Son’s result. However, to the best of our knowledge, the analytical determinations of the gap suffer from two flaws. First, the gap equation is divergent and so must be regularized and renormalized. Most treatments have taken the baryon density as a sharp cutoff. This might cause some concern since it leads to the possibility of contaminating the low energy physics of the gap with ad hoc high energy physics. Issues of cutoff sensitivity have been addressed in Ref. and Ref. . Second, the solution of the gap equation at momenta large compared to the gap has been obtained by assuming (what appears to be) a particularly unhealthy approximation which allows the integral equation for the gap to be expressed as a simple differential equation. The goals of this paper are modest. Using cutoff regularization we define a renormalized gap equation and we find the exact asymptotic solution, thus excising the flaws contained in previous determinations. Our results confirm Son’s original analysis. We also obtain the asymptotic solution using dimensional regularization with minimal subtraction.
Many degrees of freedom, including antiparticles and hard gluons, are not dynamical on the Fermi surface. Hence it is sensible to work with an effective theory of QCD appropriate to the scales in question. Explicit construction of the effective field theory appropriate for momentum scales below $`2\mu `$, can be found in papers by Hong. The fermions in this effective theory live on the two-dimensional Fermi surface and so depend only on the parallel momentum, $`q_{}`$. The gluons on the other hand propagate in directions perpendicular to the Fermi surface as well and therefore also depend on the perpedicular momentum, $`q_{}`$. Hence we should treat the effective field theory as a superposition of two-dimensional theories, one for each direction on the Fermi surface, interacting through four-dimensional gluons and contact operators. Only the graphs shown in Fig. (1a) and (1b) need be calculated if we are interested in the leading exponential behavior of the gap. The sum of these graphs is (after rotating to Euclidean space)
$`\mathrm{\Delta }(p_{})`$ $`=`$ $`{\displaystyle \frac{d^2q_{}}{(2\pi )^2}\frac{\mathrm{\Delta }(q_{})}{q_{}^2+\mathrm{\Delta }(q_{})^2}\left\{\frac{2g^2}{3}\frac{d^2q_{}}{(2\pi )^2}\left(\frac{1}{\stackrel{}{q}_{}^{\mathrm{\hspace{0.33em}2}}+\frac{\pi }{4}M_d^2|p_0q_0|/|\stackrel{}{q}_{}|}+\frac{1}{\stackrel{}{q}_{}^{\mathrm{\hspace{0.33em}2}}+M_d^2}\right)+\stackrel{~}{D}\right\}},`$ (1)
where the propagators within the parentheses represent magnetic and electric gluon exchanges, respectively, and
$`M_d^2={\displaystyle \frac{N_fg^2\mu ^2}{2\pi ^2}}.`$ (2)
The effects of hard gluon exchange, antiparticle exchange and any residual gauge dependence are represented by the coefficient, $`\stackrel{~}{D}`$, of a four-fermi operator in the two-dimensional effective theory.
It is clear that the integration over $`q_{}`$ is divergent. We will first regulate the gap equation with a sharp cutoff, $`\mathrm{\Lambda }_{}`$. In principle there is a cutoff associated with the parallel integration as well and therefore in general we have $`\stackrel{~}{D}=\stackrel{~}{D}(\mathrm{\Lambda }_{},\mathrm{\Lambda }_{})`$. It is straightforward to do the integration over $`q_{}`$ and $`\stackrel{}{q}_{}`$. We obtain
$`\mathrm{\Delta }(p_0)`$ $`=`$ $`{\displaystyle \frac{1}{2}}𝒞^2{\displaystyle _\mathrm{\Lambda }_{}^\mathrm{\Lambda }_{}}𝑑q_0{\displaystyle \frac{\mathrm{\Delta }(q_0)}{\sqrt{q_0^2+\mathrm{\Delta }(q_0)^2}}}\left(\mathrm{log}{\displaystyle \frac{_\mathrm{\Lambda }_{}}{|p_0q_0|}}+{\displaystyle \frac{1}{2}}D(\mathrm{\Lambda }_{},\mathrm{\Lambda }_{})\right)`$ (3)
where
$`_\mathrm{\Lambda }_{}={\displaystyle \frac{4(\mathrm{\Lambda }_{})^6}{\pi M_d^5}};\stackrel{~}{D}={\displaystyle \frac{g^2D}{18\pi }};𝒞={\displaystyle \frac{g}{3\sqrt{2}\pi }}.`$ (4)
Since $`\mathrm{\Delta }(p_0)`$ is independent of the cutoff, $`D`$ runs according to
$`D(\mathrm{\Lambda }_{},\mathrm{\Lambda }_{})=D(\eta ,\mathrm{\Lambda }_{})6\mathrm{log}{\displaystyle \frac{\mathrm{\Lambda }_{}^2}{\eta ^2}}.`$ (5)
Naive dimensional analysis suggests that $`D(\mathrm{\Lambda }_{}=2\mu ,\mathrm{\Lambda }_{})`$ is small and can therefore be dropped at leading order. The effect of the running of $`D`$ appears at next order . We then have
$`_{2\mu }={\displaystyle \frac{2^{10}\sqrt{2}\pi ^4\mu }{N_{f}^{}{}_{}{}^{5/2}g^5}};D(2\mu ,\mathrm{\Lambda }_{})=D(\mathrm{\Lambda }_{}).`$ (6)
Dropping all subscripts for simplicity the gap equation becomes
$`\mathrm{\Delta }(p)`$ $`=`$ $`{\displaystyle \frac{1}{2}}𝒞^2{\displaystyle _0^\mathrm{\Lambda }}𝑑q{\displaystyle \frac{\mathrm{\Delta }(q)}{\sqrt{q^2+\mathrm{\Delta }(q)^2}}}\left(\mathrm{log}{\displaystyle \frac{^2}{|p^2q^2|}}+D(\mathrm{\Lambda })\right).`$ (7)
Because of the log singularity at $`qp`$, the integral is dominated by momenta $`qp`$. Therefore in the asymptotic region defined by $`p\mathrm{\Delta }(p)`$, we can take $`q\mathrm{\Delta }(q)`$ under the integral. Hence asymptotically the gap equation becomes the homogeneous integral equation
$`\mathrm{\Delta }(p)`$ $`=`$ $`{\displaystyle \frac{1}{2}}𝒞^2{\displaystyle _\mathrm{\Delta }^\mathrm{\Lambda }}{\displaystyle \frac{dq}{q}}\mathrm{\Delta }(q)\left(\mathrm{log}{\displaystyle \frac{^2}{|p^2q^2|}}+D(\mathrm{\Lambda })\right).`$ (8)
Note that although the integral equation is homogeneous, we have by necessity introduced an infrared cutoff, $`\mathrm{\Delta }`$, which we will see is related to the overall normalization of the gap. In order to proceed<sup>*</sup><sup>*</sup>* The usual way of proceeding is to make the approximation $`\mathrm{log}|p^2q^2|=\mathrm{log}p^2\theta (p^2q^2)+\mathrm{log}q^2\theta (q^2p^2)`$. We do NOT make this approximation in this paper., consider the derivative of the gap equation:
$`\mathrm{\Delta }^{}(p)`$ $`=`$ $`𝒞^2p{\displaystyle _\mathrm{\Delta }^\mathrm{\Lambda }}{\displaystyle \frac{dq}{q}}{\displaystyle \frac{\mathrm{\Delta }(q)}{(q^2p^2)}}.`$ (9)
Note that since the counterterm is no longer present, this integral is convergent as $`\mathrm{\Lambda }\mathrm{}`$.
We make an ansatz of the form $`\mathrm{\Delta }(p)=p^z`$ with $`z`$ a complex number. Inserting this solution in eq. (9) leads to the indicial equation:
$`z`$ $`=`$ $`𝒞^2p^{2z}{\displaystyle _\mathrm{\Delta }^\mathrm{\Lambda }}𝑑q{\displaystyle \frac{q^{z1}}{(q^2p^2)}}.`$ (10)
The integral is straightforward to evaluate. It is given by
$`\mathrm{𝑅𝑒}{\displaystyle \frac{1}{1\mathrm{exp}2\pi i(z1)}}{\displaystyle 𝑑q\frac{q^{z1}}{(q^2p^2+iϵ)}},`$ (11)
where the contour in the complex q-plane is taken to enclose the poles at $`piϵ`$ and $`p+iϵ`$ while avoiding the branch point at the origin. For $`\mathrm{\Delta }p\mathrm{\Lambda }`$ we can take the limits $`\mathrm{\Delta }0`$ and $`\mathrm{\Lambda }\mathrm{}`$ since the corrections are suppressed by powers of $`p/\mathrm{\Lambda }`$ and $`\mathrm{\Delta }/p`$. The resulting indicial equation is transcendental
$`z`$ $`=`$ $`{\displaystyle \frac{\pi }{2}}𝒞^2\mathrm{cot}{\displaystyle \frac{\pi z}{2}}.`$ (12)
Therefore the exact asymptotic solution of the gap equation is $`p^z`$ with z satisfying eq. (12). For small z this has the solution
$`z=\pm i𝒞.`$ (13)
The general asymptotic solution to the gap equation in the weak coupling limit can then be written as
$`\mathrm{\Delta }(p)=A\mathrm{sin}(𝒞\mathrm{log}{\displaystyle \frac{\mathrm{\Lambda }^{}}{p}}),`$ (14)
where the constants $`A`$ and $`\mathrm{\Lambda }^{}`$ are to be determined.
We will now determine the constant $`A`$. In the region $`\mathrm{\Delta }p<2\mu `$, $`\mathrm{\Delta }(p)`$ has one maximum, which we will assume is at $`p=\overline{\mathrm{\Delta }}`$ with $`\mathrm{\Delta }\overline{\mathrm{\Delta }}<\mathrm{\Lambda }`$. This determines $`A=\mathrm{\Delta }(\overline{\mathrm{\Delta }})`$. We can find $`\overline{\mathrm{\Delta }}`$ by plugging the general solution, eq. (14), into eq. (9) which leads to the equation
$`0=\mathrm{\Delta }^{}(p)|_{p=\overline{\mathrm{\Delta }}}=𝒞^2\overline{\mathrm{\Delta }}\mathrm{\Delta }(\overline{\mathrm{\Delta }}){\displaystyle _\mathrm{\Delta }^\mathrm{\Lambda }}{\displaystyle \frac{dq}{q}}\mathrm{sin}(𝒞\mathrm{log}{\displaystyle \frac{\mathrm{\Lambda }^{}}{q}}){\displaystyle \frac{1}{(q^2\overline{\mathrm{\Delta }}^2)}}.`$ (15)
A straightforward computation then gives
$`\mathrm{cos}x_\mathrm{\Delta }+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(𝒞^2+4n^2)}}(𝒞^2[\mathrm{exp}\left({\displaystyle \frac{n}{𝒞}}(\pi 2x_\mathrm{\Delta })\right)\mathrm{cos}x_\mathrm{\Delta }+\mathrm{exp}({\displaystyle \frac{n}{𝒞}}(\pi 2x_\mathrm{\Lambda }))\mathrm{cos}x_\mathrm{\Lambda }]`$ (16)
$`+2𝒞n[\mathrm{exp}\left({\displaystyle \frac{n}{𝒞}}(\pi 2x_\mathrm{\Delta })\right)\mathrm{sin}x_\mathrm{\Delta }\mathrm{exp}({\displaystyle \frac{n}{𝒞}}(\pi 2x_\mathrm{\Lambda }))\mathrm{sin}x_\mathrm{\Lambda }])=0`$ (17)
where
$`x_\mathrm{\Delta }=𝒞\mathrm{log}{\displaystyle \frac{\mathrm{\Lambda }^{}}{\mathrm{\Delta }}};x_\mathrm{\Lambda }=𝒞\mathrm{log}{\displaystyle \frac{\mathrm{\Lambda }^{}}{\mathrm{\Lambda }}}.`$ (18)
The expression under the sum is exponentially suppressed in the QCD coupling $`g`$. Therefore, to leading order $`x_\mathrm{\Delta }=\pi /2`$ and we can identify $`\overline{\mathrm{\Delta }}`$ with $`\mathrm{\Delta }`$. This determines $`A=\mathrm{\Delta }(\mathrm{\Delta })`$ and the asymptotic solution is
$`\mathrm{\Delta }(p)=\mathrm{\Delta }(\mathrm{\Delta })\mathrm{sin}(𝒞\mathrm{log}{\displaystyle \frac{\mathrm{\Lambda }^{}}{p}});\mathrm{\Delta }(\mathrm{\Delta })=\mathrm{\Delta }=\mathrm{\Lambda }^{}\mathrm{exp}\left({\displaystyle \frac{\pi }{2𝒞}}\right).`$ (19)
We now determine $`\mathrm{\Lambda }^{}`$. We can rewrite eq. (8) as
$`\mathrm{\Delta }(p)`$ $`=`$ $`{\displaystyle \frac{1}{2}}𝒞^2{\displaystyle _\mathrm{\Delta }^\mathrm{\Lambda }}{\displaystyle \frac{dq}{q}}\mathrm{\Delta }(q)\mathrm{log}{\displaystyle \frac{\overline{}^2}{|p^2q^2|}}+{\displaystyle \frac{1}{2}}𝒞^2\left(D(\mathrm{\Lambda })D(\overline{})\right){\displaystyle _\mathrm{\Delta }^\mathrm{\Lambda }}{\displaystyle \frac{dq}{q}}\mathrm{\Delta }(q)`$ (20)
where
$`\overline{}^2=^2\mathrm{exp}(D(\overline{})).`$ (21)
The condition that $`\mathrm{\Delta }(p)`$ be independent of the choice of cutoff leads to the renormalization group equation
$`\mathrm{\Lambda }{\displaystyle \frac{d}{d\mathrm{\Lambda }}}\left(D(\mathrm{\Lambda })D(\overline{})\right)=𝒞\mathrm{tan}(𝒞\mathrm{log}{\displaystyle \frac{\mathrm{\Lambda }^{}}{\mathrm{\Lambda }}})\left[\left(D(\mathrm{\Lambda })D(\overline{})\right)+\mathrm{log}{\displaystyle \frac{\overline{}^2}{\mathrm{\Lambda }^2}}\right]`$ (22)
where we have used the asymptotic solution, eq. (19). The solution of this equation is
$`D(\mathrm{\Lambda })D(\overline{})={\displaystyle \frac{2}{𝒞}}\left[\mathrm{tan}(𝒞\mathrm{log}{\displaystyle \frac{\mathrm{\Lambda }^{}}{\mathrm{\Lambda }}}){\displaystyle \frac{\mathrm{sin}(𝒞\mathrm{log}\frac{\mathrm{\Lambda }^{}}{\overline{}})}{\mathrm{cos}(𝒞\mathrm{log}\frac{\mathrm{\Lambda }^{}}{\mathrm{\Lambda }})}}\right]\mathrm{log}{\displaystyle \frac{\overline{}^2}{\mathrm{\Lambda }^2}}.`$ (23)
Notice that the difference between counterterms evaluated at different choices of scale is of order $`g^2`$. The counterterm has a strong cutoff dependence except when $`\mathrm{\Lambda }\overline{}`$. The peculiar running of the counterterm is very similar to the running of a counterterm which arises at leading order in effective field theory treatments of the three-body problem in nuclear physics.
Choosing $`\mathrm{\Lambda }=\overline{}`$, the renormalized gap equation is
$`\mathrm{\Delta }(p)`$ $`=`$ $`{\displaystyle \frac{1}{2}}𝒞^2{\displaystyle _0^\overline{}}𝑑q{\displaystyle \frac{\mathrm{\Delta }(q)}{\sqrt{q^2+\mathrm{\Delta }(q)^2}}}\mathrm{log}{\displaystyle \frac{\overline{}^2}{|p^2q^2|}}.`$ (24)
We can now find $`\mathrm{\Lambda }^{}`$ by plugging the general solution, eq. (19), into the asymptotic form of the renormalized gap equation, eq. (24), evaluated at $`p=\mathrm{\Lambda }^{}`$. We then have
$`0=\mathrm{\Delta }(\mathrm{\Lambda }^{})={\displaystyle \frac{1}{2}}𝒞^2\mathrm{\Delta }(\mathrm{\Delta }){\displaystyle _\mathrm{\Delta }^\overline{}}{\displaystyle \frac{dq}{q}}\mathrm{sin}(𝒞\mathrm{log}{\displaystyle \frac{\mathrm{\Lambda }^{}}{q}})\mathrm{log}{\displaystyle \frac{\overline{}^2}{|\mathrm{\Lambda }_{}^{}{}_{}{}^{2}q^2|}}.`$ (25)
A straightforward computation then gives
$`0=\mathrm{sin}x_\overline{}2x_\overline{}\mathrm{cos}x_\overline{}+{\displaystyle \frac{1}{2}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n}}{\displaystyle \frac{1}{(𝒞^2+4n^2)}}(𝒞^3\mathrm{exp}\left({\displaystyle \frac{2nx_\overline{}}{𝒞}}\right)\mathrm{cos}x_\overline{}`$ (26)
$`2𝒞^2[\mathrm{exp}\left({\displaystyle \frac{2nx_\overline{}}{𝒞}}\right)\mathrm{sin}x_\overline{}+n\mathrm{exp}\left({\displaystyle \frac{n\pi }{𝒞}}\right)])`$ (27)
where
$`x_\overline{}=𝒞\mathrm{log}{\displaystyle \frac{\mathrm{\Lambda }^{}}{\overline{}}}.`$ (28)
The expression under the sum is exponentially suppressed in the QCD coupling $`g`$. Therefore, to leading order $`x_\overline{}=0`$ and we can identify $`\mathrm{\Lambda }^{}=\overline{}`$.
The final form for the asymptotic solution is thus
$`\mathrm{\Delta }(p)=\mathrm{\Delta }(\mathrm{\Delta })\mathrm{sin}(𝒞\mathrm{log}{\displaystyle \frac{\overline{}}{p}}),`$ (29)
and the gap is
$`\mathrm{\Delta }=\overline{}\mathrm{exp}\left({\displaystyle \frac{\pi }{2𝒞}}\right)={\displaystyle \frac{2^{10}\sqrt{2}\pi ^4\mu }{N_{f}^{}{}_{}{}^{5/2}g^5}}\mathrm{exp}\left({\displaystyle \frac{D(2\mu ,)}{2}}\right)\mathrm{exp}\left({\displaystyle \frac{3\pi ^2}{\sqrt{2}g}}\right).`$ (30)
This is Son’s result aside from the additional contribution to the prefactor from the counterterm, which is expected to be a number of order one.
One may wonder whether use of cutoff regularization in dense QCD violates important symmetries. Consistent implementation of dimensional regularization would erase these concerns. We will see that dimensional regularization with minimal subtraction is a quick way of obtaining the asymptotic solution directly from the gap equation. We can continue the four-dimensional measure to an $`2+n`$-dimensional measure
$`{\displaystyle \frac{d^2q_{}}{(2\pi )^2}\frac{d^nq_{}}{(2\pi )^n}}.`$ (31)
The gap equation relevant to dimensional regularization with minimal subtraction is
$`\mathrm{\Delta }(p_{})`$ $`=`$ $`{\displaystyle \frac{d^2q_{}}{(2\pi )^2}\frac{\mathrm{\Delta }(q_{})}{q_{}^2+\mathrm{\Delta }(q_{})^2}\left\{\frac{2g^2}{3}\frac{d^nq_{}}{(2\pi )^n}\left(\frac{1}{\stackrel{}{q}_{}^{\mathrm{\hspace{0.33em}2}}+\frac{\pi }{4}M_d^2|p_0q_0|/|\stackrel{}{q}_{}|}+\frac{1}{\stackrel{}{q}_{}^{\mathrm{\hspace{0.33em}2}}+M^2}\right)+\stackrel{~}{D}^{\overline{MS}}\right\}}.`$ (32)
The integrals of the gluon propagators in $`n`$-dimensions are
$`{\displaystyle \frac{d^nq_{}}{(2\pi )^n}\frac{1}{\stackrel{}{q}_{}^{\mathrm{\hspace{0.33em}2}}+A/|\stackrel{}{q}|}}={\displaystyle \frac{1}{6\pi }}\left({\displaystyle \frac{\lambda ^3}{A}}\right)^ϵ{\displaystyle \frac{\mathrm{\Gamma }\left(ϵ\right)\mathrm{\Gamma }\left(1ϵ\right)}{\mathrm{\Gamma }\left(1\frac{3ϵ}{2}\right)}};`$ (33)
$`{\displaystyle \frac{d^nq_{}}{(2\pi )^n}\frac{1}{\stackrel{}{q}_{}^{\mathrm{\hspace{0.33em}2}}+M^2}}={\displaystyle \frac{1}{6\pi }}\left({\displaystyle \frac{\lambda ^3}{M^3}}\right)^ϵ\mathrm{\Gamma }\left({\displaystyle \frac{3ϵ}{2}}\right),`$ (34)
where $`ϵ=(2n)/3`$ and $`\lambda `$ is a renormalization scale. Absorbing the $`1/ϵ`$ pole into the counterterm, we can then define the regularized gap equation
$`\mathrm{\Delta }(p)`$ $`=`$ $`{\displaystyle \frac{1}{2}}𝒞^2{\displaystyle _0^{\mathrm{}}}𝑑q{\displaystyle \frac{\mathrm{\Delta }(q)}{\sqrt{q^2+\mathrm{\Delta }(q)^2}}}\left(\mathrm{log}{\displaystyle \frac{_\lambda ^2}{|p^2q^2|}}+D^{\overline{MS}}(\lambda )\right)`$ (35)
where
$`_\lambda ={\displaystyle \frac{4(\lambda )^6}{\pi M_{d}^{}{}_{}{}^{5}}};\stackrel{~}{D}^{\overline{MS}}={\displaystyle \frac{g^2D^{\overline{MS}}}{18\pi }}.`$ (36)
This equation was obtained in Ref. . The counterterm runs according to
$`D^{\overline{MS}}(\lambda )=D^{\overline{MS}}(\eta )6\mathrm{log}{\displaystyle \frac{\lambda ^2}{\eta ^2}}.`$ (37)
Physical quantities are $`\lambda `$-independent so we choose $`\lambda =2\mu `$. We again consider the asymptotic gap equation
$`\mathrm{\Delta }(p)`$ $`=`$ $`{\displaystyle \frac{1}{2}}𝒞^2{\displaystyle _\mathrm{\Delta }^{\mathrm{}}}𝑑q{\displaystyle \frac{\mathrm{\Delta }(q)}{q}}\left(\mathrm{log}{\displaystyle \frac{_{2\mu }^2}{|p^2q^2|}}+D^{\overline{MS}}(2\mu )\right).`$ (38)
Say the asymptotic solution is of the form $`\mathrm{\Delta }(p)=p^z`$. All scales then appear multiplied by power law divergences which vanish in minimal subtraction. Hence it is appropriate to return to the unsubtracted expression for the log in eq. (34). We then obtain
$`p^z=𝒞^2\mathrm{\Gamma }(ϵ){\displaystyle _0^{\mathrm{}}}𝑑q{\displaystyle \frac{q^{z1}}{(q^2p^2)^ϵ}}.`$ (39)
Again this integral is straightforward to evaluate and leads to
$`p^z=𝒞^2\mathrm{\Gamma }(ϵ)\left[{\displaystyle \frac{p^{z2ϵ}}{2}}\mathrm{cos}\pi (ϵ{\displaystyle \frac{z}{2}}){\displaystyle \frac{\mathrm{\Gamma }\left(ϵ\frac{z}{2}\right)\mathrm{\Gamma }\left(\frac{z}{2}\right)}{\mathrm{\Gamma }\left(ϵ\right)}}\right],`$ (40)
which as expected reduces to
$`z={\displaystyle \frac{\pi }{2}}𝒞^2\mathrm{cot}{\displaystyle \frac{\pi z}{2}}`$ (41)
in the limit $`ϵ0`$. Notice that the original $`\mathrm{\Gamma }(ϵ)`$ pole from the integration over perpendicular momenta cancels a $`1/\mathrm{\Gamma }(ϵ)`$ zero from the integration over $`q`$. The asymptotic solution is again
$`\mathrm{\Delta }(p)=\mathrm{\Delta }(\mathrm{\Delta })\mathrm{sin}(𝒞\mathrm{log}{\displaystyle \frac{\mathrm{\Lambda }^{}}{p}}),`$ (42)
where the prefactor is fixed by the argument presented above. Now fixing $`\mathrm{\Lambda }^{}`$ is trivial since there is only one scale in the problem. We have
$`\mathrm{\Lambda }^{}=_{2\mu }\mathrm{exp}{\displaystyle \frac{D^{\overline{MS}}(2\mu )}{2}}`$ (43)
from which follows
$`\mathrm{\Delta }={\displaystyle \frac{2^{10}\sqrt{2}\pi ^4\mu }{N_{f}^{}{}_{}{}^{5/2}g^5}}\mathrm{exp}\left({\displaystyle \frac{D^{\overline{MS}}(2\mu )}{2}}\right)\mathrm{exp}\left({\displaystyle \frac{3\pi ^2}{\sqrt{2}g}}\right).`$ (44)
Evidently the counterterm does not depend on the renormalization scale in the longitudinal direction in minimal subtraction. As argued previously, the counterterm is of order $`g^0`$ and can be dropped at this order in the expansion.
We thank Martin Savage for valuable conversations. This work is supported in part by the U.S. Dept. of Energy under Grants No. DE-FG03-97ER4014 and DOE-ER-40561.
|
warning/0005/hep-th0005168.html
|
ar5iv
|
text
|
# 1 Introduction
## 1 Introduction
Since the conjecture of Maldacena on the duality between the large $`𝒩`$ limits of certain conformal field theories (CFT) in $`d`$ dimensions and the superstring theory ,in a certain limit, on the product of $`d+1`$ dimensional anti-de Sitter (AdS) spaces with spheres an enormous amount of research has been done on AdS/CFT dualities . The most studied example of this duality is between the $`N=4`$ super Yang-Mills in $`d=4`$ and the IIB superstring over $`AdS_5\times S^5`$ in the large $`𝒩`$ limit. In it was pointed out how the conjecture of Maldacena can be understood on the basis of some work done long time ago on Kaluza-Klein supergravity theories. Referring to for details and references to the earlier work let us recall the salient features of the earlier work that bear directly on the Maldacena conjecture. In the unitary supermultiplets of the $`d=4`$ AdS supergroups $`OSp(2N/4,R)`$ were constructed and the spectrum of the $`S^7`$ compactification of eleven dimensional supergravity was shown to fit into an infinite tower of short unitary supermultiplets of $`OSp(8/4,R)`$. The ultra-short singleton supermultiplet of $`OSp(8/4,R)`$ sits at the bottom of this infinite tower of Kaluza-Klein modes and decouple from the spectrum as local gauge degrees of freedom . However , even though it decouples from the spectrum as local gauge modes, one can generate the entire spectrum of 11-dimensional supergravity over $`S^7`$ by tensoring $`p`$ copies (“colors”) ($`p=2,3,4,\mathrm{}`$) of singleton supermultiplets and restricting oneself to ”CPT self-conjugate ” vacuum supermultiplets. <sup>2</sup><sup>2</sup>2 As will be explained below a vacuum supermultiplet corresponds to a unitary representation whose lowest weight vector is the Fock vacuum.
The compactification of 11-d supergravity over the four sphere $`S^4`$ down to seven dimensions was studied in and its spectrum was shown to fall into an infinite tower of unitary supermultiplets of $`OSp(8^{}/4)`$ with the even subgroup $`SO(6,2)\times USp(4)`$ in . Again the vacuum doubleton supermultiplet of $`OSp(8^{}/4)`$ decouples from the spectrum as local gauge degrees of freedom <sup>3</sup><sup>3</sup>3 See the next section for the distinction between singleton and doubleton supermultiplets.. It consists of five scalars, four fermions and a self-dual two form field . The entire physical spectrum of 11-dimensional supergravity over $`S^4`$ is obtained by simply tensoring an arbitrary number (colors) of the doubleton supermultiplets and restricting oneself to the vacuum supermultiplets .
The spectrum of the $`S^5`$ compactification of ten dimensional IIB supergravity was calculated in . Again the entire spectrum falls into an infinite tower of massless and massive unitary supermultiplets of $`N=8`$ $`AdS_5`$ superalgebra $`SU(2,2/4)`$ . The ”CPT self-conjugate” doubleton supermultiplet of $`N=8`$ $`AdS`$ superalgebra decouples from the physical spectrum as local gauge degrees of freedom. By tensoring it with itself repeatedly and restricting oneself to the $`CPT`$ self-conjugate vacuum supermultiplets one generates the entire spectrum of Kaluza-Klein states of ten dimensional IIB supergravity on $`S^5`$.
The authors of pointed out that the CPT self-conjugate $`N=8`$ $`AdS_5`$ doubleton supermultiplet does not have a Poincare limit in five dimensions and its field theory exists only on the boundary of $`AdS_5`$ which can be identified with the $`d=4`$ Minkowski space . Furthermore, they pointed out, for the first time, that the doubleton field theory of $`SU(2,2/4)`$ is the conformally invariant $`N=4`$ super Yang-Mills theory in $`d=4`$. Similarly, the singleton supermultiplet of $`OSp(8/4,R)`$ and the doubleton supermultiplet of $`OSp(8^{}/4)`$ do not have a Poincare limit in $`d=4`$ and $`d=7`$, respectively, and their field theories are conformally invariant theories in one lower dimension <sup>4</sup><sup>4</sup>4 see for references . Thus we see that at the level of physical states the proposal of Maldacena is perfectly consistent with the above mentioned results if we assume that the spectrum of the superconformal field theories fall into (”$`CPT`$ self-conjugate” ) vacuum supermultiplets. Remarkably, this is equivalent to assuming that the spectrum consists of “color” singlet supermultiplets. !
## 2 Massless and Massive Supermultiplets of Anti-de Sitter Supergroups
The Poincaré limit of the remarkable representations (singletons) of the $`d=4AdS`$ group $`SO(3,2)`$ discovered by Dirac are known to be singular . However, the tensor product of two singleton unitary irreducible representations of $`SO(3,2)`$ decomposes into an infinite set of massless unitary irreducible representations which do have a smooth Poincaré limit . Similarly, the tensor product of two singleton supermultiplets of $`N`$ extended $`AdS_4`$ supergroup $`OSp(N/4,R)`$ decomposes into an infinite set of massless supermultiplets which do have a Poincaré limit in five dimensions . The $`AdS`$ groups $`SO(d1,2)`$ in higher dimensions than four that do admit supersymmetric extensions have doubleton representations only. The doubleton supermultiplets of extended $`AdS`$ supergroups in $`d=5(SU(2,2/N))`$ and $`d=7(OSp(8^{}/2N))`$ share the same remarkable features of the singleton supermultiplets of $`d=4`$ $`AdS`$ supergroups i.e the tensor product of any two doubletons decompose into an infinite set of massless supermultiplets . In $`d=3`$ the $`AdS`$ group $`SO(2,2)`$ is not simple and is isomorphic to $`SO(2,1)\times SO(2,1)`$. Since each $`SO(2,1)`$ factor can be extended to a simple superalgebra with some internal symmetry group one has a rich variety of $`AdS`$ supergroups in $`d=3`$ . Since locally we have the isomorphisms $`SO(2,1)SL(2,R)SU(1,1)Sp(2,R)`$ the $`AdS`$ supergroups in $`d=3`$ (and hence in $`d=2`$) admit singleton representations .
Since the Poincaré mass operator is not an invariant ( Casimir) operator of the $`AdS`$ group the following definition of a massless representation (or supermultiplet) of an $`AdS`$ group (or supergroup) was proposed in :
A representation (or a supermultiplet) of an $`AdS`$ group (or supergroup) is massless if it occurs in the decomposition of the tensor product of two singleton or two doubleton representations (or supermultiplets).
The tensor product of more than two copies of the singleton or doubleton supermultiplets of $`AdS`$ supergroups decompose into an infinite set of massive supermultiplets in the respective dimensions as has been amply demonstrated within the Kaluza-Klein supergravity theories and more recently . A noncompact group that admits only doubleton representations can always be embedded in a larger noncompact group that admits singleton representations. In such cases the singleton representation of the larger group decomposes , in general, into an infinite tower of doubleton representations of the subgroup.
## 3 Unitary Lowest Weight Representations of Noncompact Groups
A representation of a non-compact group is said to be of the lowest weight type if the spectrum of at least one of its generators is bounded from below within the representation space. A non-compact simple group G admits unitary lowest weight representations (ULWR) if and only if its quotient space G/H with respect to its maximal compact subgroup H is an hermitian symmetric space . Thus the complete list of simple non-compact groups G that admit ULWRS’s follows from the list of irreducible hermitian symmetric spaces which we give below:
The Lie algebra $`g`$ of a non-compact group G that admits ULWR’s has a 3-grading with respect to the Lie algebra $`h`$ of its maximal compact subgroup H i.e
$`g=g^1g^0g^{+1}`$ (3 - 1)
where $`g^0=h`$ and we have the formal commutation relations
$`[g^{(m)},g^{(n)}]g^{(m+n)}m,n=1,0`$
and $`g^{(m)}0`$ for $`|m|>1.`$ In the general oscillator construction of unitary lowest weight representations (ULWR) of non-compact groups was given. Particular cases of the oscillator construction for certain representations of some special groups such as $`SU(1,1)`$ and $`SU(2)`$ had previously appeared in the physics literature.
To construct the ULWR’s one first realizes the generators of the noncompact group $`G`$ as bilinears of bosonic oscillators transforming in a certain representation of $`H`$. Then in the corresponding Fock space $``$ one chooses a set of states $`|\mathrm{\Omega }>`$, referred to as the ”lowest weight vector” (lwv), which transforms irreducibly under $`H`$ and which are annihilated by the generators belonging to the $`g^1`$ space. Then by acting on $`|\mathrm{\Omega }>`$ repeatedly with the generators belonging to the $`g^{+1}`$ space one obtains an infinite set of states
$$|\mathrm{\Omega }>,g^{+1}|\mathrm{\Omega }>,g^{+1}g^{+1}|\mathrm{\Omega }>,\mathrm{}$$
(3 - 2)
This set of states forms the basis of an irreducible unitary lowest weight representation of $`g`$. (The irreducibility of the representation of $`g`$ follows from the irreducibility of the lwv $`|\mathrm{\Omega }>`$ under $`h`$. The bosonic oscillators $`a_i(r)`$ satisfy the canonical commutation relations
$$\begin{array}{c}\\ [a_i(r),a^j(s)]=\delta _i^j\delta _{rs}i,j=1,\mathrm{},n\\ [a_i(r),a_j(s)]=0r,s=1,\mathrm{},p\end{array}$$
(3 - 3)
where the upper indices $`i,j,k,\mathrm{}`$ are the indices in the representation $`R`$ of $`h`$ under which the oscillators transform and $`r,s..=1,2,..p`$ label the different sets of oscillators. We denote the creation (annihilation) operators with upper (lower) indices $`i,j,..`$, respectively :
$`a_i(r)^{}a^i(r)`$
Generally, $`R`$ is the fundamental representation of $`h`$ and we shall refer to $`p`$ as the number of colors. The generators are color singlet bilinears but the lowest weight vector $`|\mathrm{\Omega }>`$ and hence the infinite tower of vectors belonging to the corresponding ULWR can carry color. Depending on the non-compact group the minimal number $`p`$ of colors required to realize the generators can be one or two. If $`p_{min}=1`$, we shall call the corresponding unitary irreducible representations singletons and if $`p_{min}=2`$, they will be referred to as doubletons . The non-compact groups $`Sp(2n,R)`$ admit singleton unitary irreducible representations while the groups $`SO^{}(2n)`$ and $`SU(n,m)`$ admit doubleton unitary irreducible representations . We should note that the “remarkable representations” of the four dimensional $`AdS`$ group $`SO(3,2)`$ with the covering group $`Sp(4,R)`$ discovered by Dirac are simply the singletons. While when $`p_{min}=1`$ for a given non-compact group there exist only two singletons , one finds infinitely many doubletons for $`p_{min}=2`$. The two singletons of $`Sp(4,R)`$ can be associated with spin zero and spin $`\frac{1}{2}`$ fields. On the other hand the $`d=7AdS`$ group $`SO^{}(8)=SO(6,2)`$ admits infinitely many doubletons corresponding to fields of arbitrarily large spin . However, we should note that the doubleton fields are not of the form of the most general higher spin fields in $`d=7`$. Their decomposition with respect to the little group $`SU(4)Spin(6)`$ in $`d=7`$ correspond to those representations of $`SU(4)`$ whose Young-Tableaux have only one row . Whereas the general massive higher spin fields correspond to the representations of the little group with arbitrary Young-Tableaux.
If one replaces the bosonic oscillators with fermionic ones, then the above construction leads to the unitary representations of the compact forms of the corresponding groups. One finds that the compact $`USp(2n)`$ admits doubleton (unitary irreducible) representations (finitely many) while the group $`SO(2n)`$ admits two singleton (unitary irreducible) representations . The singletons of $`SO(2n)`$ are the two spinor representations. In general the compact group $`USp(2n)`$ admits $`n`$ non-trivial doubleton representations. For $`USp(4)`$ they are the spinor representation (4) and the adjoint representation (10). The two singletons (spinors) of $`SO(2n)`$ combine to form the unique singleton (spinor representation) of $`SO(2n+1)`$.
## 4 Unitary Lowest Weight Representations of Noncompact Supergroups
The extension of the oscillator method to the construction of the ULWR’s of non-compact supergroups with a three-graded structure with respect to a maximal compact subsupergroup was given in <sup>5</sup><sup>5</sup>5 A non-compact supergroup is defined as a supergroup whose even subgroup has a non-compact subgroup. . This method was further developed and applied to space-time supergroups and Kaluza-Klein supergravity theories in the eighties . The general construction of the ULWR’s of the noncompact supergroup $`OSp(2n/2m,R)`$ with the even subgroup $`SO(2n)\times Sp(2m,R)`$ was studied in and the ULWR’s of $`OSp(2n^{}/2m)`$ with the even subgroup $`SO^{}(2n)\times USp(2m)`$ in reference . More recently a detailed study of the unitary supermultiplets of the supergroups $`SU(2,2/4)`$ and of $`OSp(8^{}/4)`$ relevant to AdS/CFT dualities in M-theory was given in and , respectively.
Consider now the Lie superalgebra $`g`$ of a non-compact supergroup $`G`$ that has a 3-graded structure with respect to a compact subsuperalgebra $`g^0`$ of maximal rank
$`g=g^1g^0g^{+1}`$
To construct the ULWR’s of $`g`$ we first realize its generators as bilinears of a set of superoscillators $`\xi _A(\xi ^A)`$ whose first $`m`$ components are bosonic and the remaining $`n`$ components are fermionic
$`\xi _A(r)=\left(\begin{array}{cc}a_i(r)& \\ \alpha _\mu (r)& \end{array}\right)\xi ^A(r)=\left(\begin{array}{cc}a^i(r)& \\ \alpha ^\mu (r)& \end{array}\right)`$ (4 - 5)
$`i,j=1,\mathrm{},m;\mu ,\nu =1,\mathrm{},n`$
$`r,s=1,\mathrm{},p.`$
which satisfy the supercommutation relations
$`[\xi _A(r),\xi ^B(s)\}=\delta _A^B\delta _{rs}`$ (4 - 6)
where \[ , } means an anti-commutator for any two fermionic oscillators and a commutator otherwise. Furthermore we have
$`[\xi _A(r),\xi _B(s)\}=0=[\xi ^A(r),\xi ^B(s)\}`$ (4 - 7)
Generally the operators belonging to the $`g^1`$ and $`g^{+1}`$ spaces are realized as super di-annihilation and di-creation operators respectively. Consider now a lowest weight vector $`|\mathrm{\Omega }>`$, that transforms irreducibly under $`g^0`$ and is annihilated by $`g^1`$ operators. Acting on $`|\mathrm{\Omega }>`$ with the $`g^{+1}`$ operators repeatedly one generates an infinite set of states that form the basis of a ULWR of $`g`$
$`g^1|\mathrm{\Omega }>=0,g^0|\mathrm{\Omega }>=|\mathrm{\Omega }^{}>`$
(4 - 8)
$`\{\mathrm{ULWR}\}\{|\mathrm{\Omega }>,g^{+1}|\mathrm{\Omega }>,g^{+1}g^{+1}|\mathrm{\Omega }>,\mathrm{}\}`$ (4 - 9)
The resulting ULWR is uniquely labelled by $`|\mathrm{\Omega }>`$. A supergroup $`g`$ admits singleton or doubleton unitary irreducible representations depending on whether $`p_{min}=1`$ or $`p_{min}=2`$, respectively. For example the non-compact supergroup $`OSp(2n/2m,R)`$ with even subgroup $`SO(2n)\times Sp(2m,R)`$ admits singleton representations. The non-compact supergroup $`OSp(2n^{}/2m)`$ with even subgroup $`SO(2n)^{}\times USp(2m)`$ admits doubleton representations, as does the supergroup $`SU(n,m/p)`$ with even subgroup $`S(U(n,m)\times U(p))`$. There exist only two irreducible singleton supermultiplets of the non-compact supergroup $`OSp(2n/2m,R)`$ . On the other hand, the supergroups $`OSp(2n^{}/2m)`$ and $`SU(n,m/p)`$ admit infinitely many irreducible doubleton supermultiplets .
In contrast to the situation with noncompact groups, not all noncompact supergroups that have ULWRs admit a three grading with respect to a compact subsupergroup of maximal rank. The method of was generalized to the case when the noncompact supergroup admits a 5-grading with respect to a compact subsupergroup of maximal rank in . For example, the superalgebra of $`OSp(2n+1/2m,R)`$ admits a 5-grading with respect to its compact subsuperalgebra $`U(n/m)`$ , but it does not admit a three grading with respect to a compact subsuperalgebra of maximal rank for general $`n`$ and $`m`$. All finite dimensional non-compact supergroups do admit a 5-grading with respect to a compact subsupergroup of maximal rank .
$`g=g^2g^1g^0g^{+1}g^{+2}`$ (4 - 10)
## 5 Generalized space-times defined by Jordan algebras
### 5.1 Generalized Rotation, Lorentz and Conformal Groups
The twistor formalism in four-dimensional space-time $`(d=4)`$ leads naturally to the representation of four vectors in terms of $`2\times 2`$ Hermitian matrices over the field of complex numbers $`𝐂`$. In particular, the coordinate four vectors $`x_\mu `$ can be represented as :
$$x=x_\mu \sigma ^\mu $$
(5 - 1)
Since the Hermitian matrices over the field of complex numbers close under the symmetric anti-commutator product we can regard the coordinate vectors as elements of a Jordan algebra denoted as $`J_2^𝐂`$ . Then the rotation, Lorentz and conformal groups in $`d=4`$ can be identified with the automorphism , reduced structure and Möbius ( linear fractional) groups of the Jordan algebra of $`2\times 2`$ complex Hermitian matrices $`J_2^𝐂`$ . The reduced structure group $`Str_0(J)`$ of a Jordan algebra $`J`$ is simply the invariance group of its norm form $`N(J)`$. (The structure group $`Str(J)=Str_0(J)\times SO(1,1)`$ ,on the other hand, is simply the invariance group of $`N(J)`$ up to an overall constant scale factor.) Furthermore, this interpretation allows one to define generalized space-times whose coordinates are parametrized by the elements of Jordan algebras <sup>6</sup><sup>6</sup>6More generally one can define spacetimes coordinatized by the elements of Jordan triple systems and study their symmetry groups . However, in this talk we restrict ourselves to spacetimes (superspaces) coordinatized by Jordan algebras ( Jordan superalgebras) . The rotation $`Rot(J)`$, Lorentz $`Lor(J)`$ and conformal $`Con(J)`$ groups of these generalized space-times are then identified with the automorphism $`Aut(J)`$, reduced structure $`Str_0(J)`$ and Möbius Mö(J) groups of the corresponding Jordan algebra <sup>7</sup><sup>7</sup>7Similar algebraic structures appear also in the study of internal U-duality groups of extended supergravity theories .. Denoting as $`J_n^𝐀`$ the Jordan algebra of $`n\times n`$ Hermitian matrices over the division algebra $`𝐀`$ and the Jordan algebra of Dirac gamma matrices in $`d`$ ( Euclidean) dimensions as $`\mathrm{\Gamma }(d)`$ one finds the following symmetry groups of generalized space-times defined by simple Jordan algebras:
The symbols $`𝐑`$, $`𝐂`$, $`𝐇`$, $`𝐎`$ represent the four division algebras. For the Jordan algebras $`J_n^𝐀`$ the norm form is the determinantal form ( or its generalization to the quaternionic and octonionic matrices). For the Jordan algebra $`\mathrm{\Gamma }(d)`$ generated by Dirac gamma matrices $`\mathrm{\Gamma }_i(i=1,2,\mathrm{}d)`$
$$\{\mathrm{\Gamma }_i,\mathrm{\Gamma }_j\}=\delta _{ij}\mathrm{𝟏};i,j,\mathrm{}=1,2,\mathrm{},d$$
(5 - 2)
the norm of a general element $`x=x_0\mathrm{𝟏}+x_i\mathrm{\Gamma }_i`$ of $`\mathrm{\Gamma }(d)`$ is quadratic and given by
$$N(x)=x\overline{x}=x_0^2x_ix_i$$
(5 - 3)
where $`\overline{x}=x_0\mathrm{𝟏}x_i\mathrm{\Gamma }_i`$. Its automorphism, reduced structure and Möbius groups are simply the rotation, Lorentz and conformal groups of $`(d+1)`$-dimensional Minkowski spacetime. One finds the following special isomorphisms between the Jordan algebras of $`2\times 2`$ Hermitian matrices over the four division algebras and the Jordan algebras of gamma matrices:
$$J_2^𝐑\mathrm{\Gamma }(2);J_2^𝐂\mathrm{\Gamma }(3);J_2^𝐇\mathrm{\Gamma }(5);J_2^𝐎\mathrm{\Gamma }(9)$$
(5 - 4)
The Minkowski spacetimes they correspond to are precisely the critical dimensions for the existence of super Yang-Mills theories as well as of the classical Green-Schwarz superstrings. These Jordan algebras are all quadratic and their norm forms are precisely the quadratic invariants constructed using the Minkowski metric.
### 5.2 Covariant Quantum Fields over Generalized Spacetimes and the ULWR’s of Their Conformal Groups
A remarkable fact about Table 2 is that the maximal compact subgroups of the generalized conformal groups of formally real Jordan algebras are simply the compact forms of their structure groups (generalized Lorentz group times dilatations). Furthermore, they all admit unitary representations (positive energy) of the lowest weight type. <sup>8</sup><sup>8</sup>8 Similarly, the generalized conformal groups defined by Hermitian Jordan triple systems all admit unitary irreducible representations of the lowest weight type . For example, the conformal group of the Jordan algebra $`J_2^𝐂`$ corresponding to the four dimensional Minkowski space is $`SU(2,2)`$ with a maximal compact subgroup $`SU(2)\times SU(2)\times U(1)`$ which is simply the compact form of the structure group $`SL(2,𝐂)\times SO(1,1)`$. In it was explicitly shown how to go from the compact $`SU(2)\times SU(2)\times U(1)`$ basis of the ULWR’s of $`SU(2,2)`$ to the manifestly covariant $`SL(2,𝐂)\times SO(1,1)`$ basis. The transition from the compact to the covariant basis corresponds simply to going from a ”particle” basis to a coherent state basis of the ULWR. The coherent states are labelled by the elements of $`J_2^𝐂`$ i.e by the coordinates of four dimensional Minkowski space. One can then establish a one-to-one correspondence between irreducible ULWR’s of $`SU(2,2)`$ and the fields transforming irreducibly under the Lorentz group $`SL(2,𝐂)`$ with a definite conformal dimension. Thus one can associate with irreducible ULWR’s of $`SU(2,2)`$ fields transforming covariantly under the Lorentz group with a definite conformal dimension.
Similarly, the conformal group $`SO^{}(8)`$ of the Jordan algebra $`J_2^𝐇`$ parametrizing the six dimensional Minkowski space has a maximal compact subgroup $`U(4)`$ which is the compact form of the structure group $`SU^{}(4)\times SO(1,1)`$. In it was shown explicitly how to go from the compact $`U(4)`$ basis of the ULWR’s of $`SO^{}(8)`$ to the non-compact basis $`SU^{}(4)\times SO(1,1)`$ which is simply the Lorentz group in six dimensions times dilatations. The coherent states of the non-compact basis are again labelled by the elements of $`J_2^𝐇`$, i.e the coordinates of 6d Minkowski space. Thus each irreducible ULWR of $`SO^{}(8)`$ can be identified with a field transforming covariantly under the Lorentz group $`SU^{}(4)`$ with a definite conformal dimension.
The results obtained explicitly for the conformal groups of $`J_2^𝐂`$ and $`J_2^𝐇`$ extend to the conformal groups of all formally real Jordan algebras and of Hermitian Jordan triple systems . The general theory can be summarized as follows: Let $`g`$ be the Lie algebra of the conformal group of a formally real Jordan algebra and $`g^0`$ the Lie algebra of its maximal compact subgroup. Then $`g`$ has a three-graded decomposition with respect to $`g^0`$:
$$g=g^{}+g^0+g^+$$
where the grading is determined by the ”conformal energy operator”. Now let $`n^0`$ be the Lie algebra of the structure group of the Jordan algebra or triple system. Then $`g`$ has a 3-graded decomposition with respect to $`n^0`$ as well:
$$g=n^{}+n^0+n^+$$
(5 - 5)
where the grading is defined by the generator of scale transformations. In the compact basis an irreducible ULWR of $`Conf(J)`$ is uniquely determined by a lowest weight vector $`|\mathrm{\Omega }`$ transforming irreducibly under the maximal compact subgroup $`K`$ that is annihilated by the operators belonging to $`g^{}`$
$$g^{}|\mathrm{\Omega }=0$$
(5 - 6)
As was done explicitly for the conformal groups in 4 and 6 dimensions one can show that there exists a complex rotation operator $`W`$ in the representation space with the property that the vector $`W|\mathrm{\Omega }`$ is annihilated by all the generators belonging to $`n^{}`$
$$n^{}W|\mathrm{\Omega }=0$$
(5 - 7)
and it transforms in a finite dimensional non-unitary representation of the non-compact structure group. Remarkably the transformation properties of $`W|\mathrm{\Omega }`$ under the structure group coincide with the transformation properties of $`|\mathrm{\Omega }`$ under the maximal compact subgroup $`K`$. In particular, the conformal dimension of the vector $`W|\mathrm{\Omega }`$ is simply the negative of the conformal energy of $`|\mathrm{\Omega }`$. If one chooses a basis $`e_\mu `$ for the Jordan algebra $`J`$ and denote the generators of generalized translations in the space $`n^+`$ corresponding to $`e_\mu `$ as $`P_\mu `$, then the coherent states defined by the action of generalized translations on $`W|\mathrm{\Omega }`$
$$|\mathrm{\Phi }(x_\mu :=e^{ix^\mu P_\mu }W|\mathrm{\Omega }$$
(5 - 8)
form the covariant basis of the ULWR of the generalized conformal group $`Con(J)`$ <sup>9</sup><sup>9</sup>9We should note that the (super) coherent states associated with ULWR’s of non-compact (super) groups introduced in are labelled by complex (super) ”coordinates” in the compact basis. These (super) coordinates parametrize the (super) hermitian symmetric space $`G/H`$. . The coherent states $`|\mathrm{\Phi }(x_\mu `$ labelled by the coordinates correspond to conformal fields transforming covariantly under the Lorentz group with a definite conformal dimension. Since the state $`W|\mathrm{\Omega }`$ is annihilated by the generators of special conformal transformations $`K_\mu `$ belonging to the space $`n^{}`$ this proves that the irreducible ULWR’s are equivalent to representations induced by finite dimensional irreps of the Lorentz group with a definite conformal dimension and trivial special conformal transformation properties. This generalizes the well-known construction of the positive energy representations of the four dimensional conformal group $`SU(2,2)`$ to all generalized conformal groups of formally real Jordan algebras and Hermitian Jordan triple systems. They are simply induced representations with respect to the maximal parabolic subgroup $`Str(J)S_J`$ where $``$ denotes semi-direct product and $`S_J`$ is the Abelian subgroup generated by generalized special conformal transformations.
We should perhaps note that the generalized Poincaré groups associated with the spacetimes defined by Jordan algebras are of the form
$$𝒫𝒢(J):=Lor(J)T_J$$
(5 - 9)
where $`T_J`$ is the Abelian subgroup generated by generalized translations $`P_\mu `$. For quadratic Jordan algebras, $`\mathrm{\Gamma }(d)`$ , $`𝒫𝒢(\mathrm{\Gamma }(d))`$ is simply the Poincaré group in $`d`$ dimensional Minkowski space. The group $`𝒫𝒢(\mathrm{\Gamma }(d))`$ has a quadratic Casimir operator $`M^2=P_\mu P^\mu `$ which is simply the mass operator. For Jordan algebras $`J`$ of degree $`n`$ the generalized Poincaré group $`𝒫𝒢(J)`$ has a Casimir invariant of order $`n`$ constructed out of the generalized translation generators $`P_\mu `$. For example for the exceptional Jordan algebra $`J_3^𝐎`$ the corresponding Casimir invariant is cubic and has the form
$$M^3=C_{\mu \nu \rho }P^\mu P^\nu P^\rho $$
(5 - 10)
where $`C_{\mu \nu \rho }`$ is the symmetric invariant tensor of the generalized Lorentz group $`E_{6(26)}`$ of $`J_3^𝐎`$.
## 6 Generalized superspaces defined by Jordan superalgebras and their symmetry supergroups
The generalized space-times defined by Jordan algebras can be extended to define generalized superspaces over Jordan superalgebras and super Jordan triple systems . A Jordan superalgebra is a $`Z_2`$ graded algebra $`J=J^0+J^1`$ with a supersymmetric product
$$\begin{array}{c}\\ ab=(1)^{\alpha \beta }ba\\ \\ aJ^\alpha ,bJ^\beta ;\alpha ,\beta =0,1\end{array}$$
(6 - 1)
which satisfies the identity
$$(1)^{\alpha \gamma }[L_{ab},L_c\}+(1)^{\beta \alpha }[L_{bc},L_a\}+(1)^{\gamma \beta }[L_{ca},L_b\}=0$$
(6 - 2)
where the mixed bracket \[ , } denotes the usual Lie superbracket and $`L_a`$ denotes left multiplication by the element $`a`$ of $`J`$. Jordan superalgebras have been classified by Kac .
One defines the generalized superspaces by multiplying the even elements of a Jordan superalgebra $`J`$ by real coordinates and their odd elements by Grassmann coordinates . The rotation, Lorentz and conformal supergroups of these generalized superspaces are then given the the automorphism, reduced structure and Möbius supergroups of $`J`$. A complete list of these supergroups was given in . We reproduce this list in Table 3.
The conformal groups of formally real Jordan algebras all admit ULWR’s and as explained above one can associate with each irreducible ULWR a covariant conformal field with a definite conformal dimension. Hence we shall restrict ourselves to those Jordan superalgebras or super Jordan triple systems whose conformal supergroups admit unitary representations of the lowest weight type. The general theory for the construction of the unitary lowest weight representations of non-compact supergroups was given in , both in a compact particle state basis as well as the compact super-coherent state basis. The coherent states defined in for non-compact groups $`G`$ are labelled by the complex variables parametrizing the hermitian symmetric space $`G/H`$ where $`H`$ is the maximal compact subgroup. On the other hand the coherent states defined in for $`SU(2,2)`$ and in for $`OSp(8^{}|4)`$ as well as their generalizations to all non-compact groups discussed in the previous section are labelled by real (generalized) coordinates of the (generalized) space-times on which $`G`$ acts as a (generalized) conformal group.
The even subgroup of (generalized) conformal supergroups $`SCon(JX)`$ are of the form $`G\times K`$ where $`G`$ is the (generalized) conformal group and $`K`$ is some compact internal symmetry group. The ULWR’s of $`SCon(JX)`$ decompose into a set of irreducible ULWR’s of $`G\times K`$. By acting on the lowest weight vectors of the irreducible ULWR’s of $`G\times K`$ with the operator
$$e^{ix^\mu P_\mu }W$$
(6 - 3)
one obtains a set of coherent states transforming covariantly under the (generalized) Lorentz group $`Lor(J)`$ with definite conformal dimension. Thus the irreducible ULWR’s of $`SCon(JX)`$ correspond simply to a supermultiplet of fields transforming irreducibly under $`Lor(J)\times K`$ with definite conformal dimension. If one starts from the compact super-coherent state basis of $`SCon(JX)`$ and goes over to the non-compact basis one obtains a ”superstate” which corresponds to a superfield built out of covariant fields multiplied by appropriate Grassmann parameters <sup>10</sup><sup>10</sup>10 Recently, a number of papers studied the supermultiplets of conformal supergroups in 3,4 and 6 dimensions using the formalism of superfields . Writing a ULWR of a non-compact conformal supergroup as a superfield corresponds to going to a super-coherent state basis for the corresponding ULWR in the oscillator formalism. . It is also possible to define covariant super-coherent states directly by acting on the lowest weight vector $`|\mathrm{\Omega }`$ of the ULWR of $`SCon(JX)`$ by the operator
$$e^{ix^\mu P_\mu +\theta ^\alpha Q_\alpha }W$$
(6 - 4)
where $`Q_\alpha `$ are the (generalized) ”Poincaré” supersymmetry generators. A detailed formulation of the covariant super-coherent state basis of the ULWR’s of generalized conformal supergroups will be given elsewhere .
Before concluding I should point out that the simple yet powerful oscillator method for the construction of the ULWR’s of non-compact superconformal groups can be given a dynamical realization in terms of twistorial or super-twistorial fields such that the (super)-oscillators become the Fourier modes of these fields .
|
warning/0005/hep-ph0005279.html
|
ar5iv
|
text
|
# 1 Virtual gamma–gamma total cross-section by the BFKL Pomeron versus L3 Collaboration data at energies: (a) 183 GeV and (b) 189 GeV of 𝑒⁺𝑒⁻ collisions. Solid curves correspond to NLO BFKL in BLM; dashed: LO BFKL; and dotted: LO Born contribution. Two different curves are for two different choices of the Regge scale: 𝑠₀=𝑄²/2 and 𝑠₀=2𝑄².
CERN-TH/2000-148
IITAP-2000-010
hep-ph/0005279
Parton Scattering at Small-$`x`$ and Scaling Violation
Victor T. Kim<sup>&†</sup>, Grigorii B. Pivovarov<sup>§</sup> and James P. Vary<sup>$‡</sup>
<sup>&</sup> : CERN, CH-1211, Geneva 23, Switzerland
<sup>§</sup> : Institute for Nuclear Research, 117312 Moscow, Russia
<sup>$</sup> : Int. Inst. of Theoretical and Applied Physics, Iowa State University, Ames, IA 50011, USA
: Department of Physics and Astronomy, Iowa State University, Ames, IA 50011, USA
ABSTRACT
Scaling violation of inclusive jet production at small-$`x`$ in hadron scattering with increasing total collision energy is discussed. Perturbative QCD based on the factorisation theorem for hard processes and GLAPD evolution equations predicts a minimum for scaled cross-section ratio that depends on jet rapidity. Studies of such a scaling violation can reveal a vivid indication of new dynamical effects in the high-energy limit of QCD. The BFKL effects, which seem to be seen in recent L3 data at CERN LEP2, should give different results from GLAPD predictions.
Based on the talk by VTK at the Xth Quantum Field Theory and High Energy Physics Workshop (QFTHEP’99), Moscow, Russia, May 27 - June 2, 1999, to appear in the Proceedings.
Permanent address: St.Petersburg Nuclear Physics Institute, 188300 Gatchina, Russia;
e-mail: kim@pnpi.spb.ru.
CERN-TH/2000-148
April 2000
QCD is an essential ingredient of the Standard Model, and it is well tested in hard processes when transferred momentum is of the order of the total collision energy (Bjorken limit: $`Q^2s\mathrm{}`$). The cornerstones of perturbative QCD at this kinematic regime (QCD-improved parton model): factorization of inclusive hard processes and the Gribov–Lipatov–Altarelli–Parisi–Dokshitzer (GLAPD) evolution equation provides a basis for the successful QCD-improved parton model. The factorisation theorem for inclusive hard processes ensures that the inclusive cross section factorises into partonic subprocess(es) and parton distribution function(s). The GLAPD evolution equation governs the $`\mathrm{log}Q^2`$-dependence (at $`Q^2\mathrm{}`$) of the inclusive hard process cross-sections at fixed scaling variable $`x=Q^2/s`$.
Another kinematic domain that is very important at high-energy is given by the (Balitsky–Fadin–Kuraev–Lipatov) BFKL limit \[3–6\], or QCD Regge limit, whereby at fixed $`Q^2\mathrm{\Lambda }_{QCD}^2`$, $`s\mathrm{}`$. In the BFKL limit, the BFKL evolution in the leading order (LO) governs $`\mathrm{log}(1/x)`$ evolution (at $`x0`$) of inclusive processes. Note that the BFKL evolution in the next-to-leading order (NLO) \[7–10\], unlike the LO BFKL \[3–5\], partly includes GLAPD evolution with the running coupling constant of the LO GLAPD, $`\alpha _S(Q^2)=4\pi /\beta _0\mathrm{log}(Q^2/\mathrm{\Lambda }_{QCD}^2)`$.
Therefore, the BFKL and especially the NLO BFKL \[7–10\] are anticipated to be important tools for exploring the high-energy limit of QCD. In particular, this importance arises since the highest eigenvalue, $`\omega ^{max}`$, of the BFKL equation \[3–6, 9, 10\] is related to the intercept of the Pomeron, which in turn governs the high-energy asymptotics of the total cross-sections: $`\sigma (s/s_0)^{\alpha _{IP}1}=(s/s_0)^{\omega ^{max}}`$, where the Regge parameter $`s_0`$ defines the approach to the asymptotic regime. The BFKL Pomeron intercept in the LO turns out to be rather large: $`\alpha _{IP}1=\omega _{LO}^{max}=12\mathrm{log}2(\alpha _S/\pi )0.54`$ for $`\alpha _S=0.2`$; hence, it is very important to analyse recently calculated NLO corrections to the BFKL.
One of the striking features of the NLO BFKL analysis is that the NLO value for the intercept of the BFKL Pomeron, improved by the BLM procedure , has a very weak dependence on the gluon virtuality $`Q^2`$: $`\alpha _{IP}1=\omega _{NLO}^{max}`$ 0.13 – 0.18 at $`Q^2=1`$ – 100 GeV<sup>2</sup>. This agrees with the conventional Regge theory where one expects universal intercept of the Pomeron without any $`Q^2`$-dependence. The minor $`Q^2`$-dependence obtained leads to approximate conformal invariance.
There have recently been a number of papers which analyse the NLO BFKL predictions \[12–16\]. Also, a lot of work should be done to clarify the very important issue of the factorisation properties of the BFKL regime \[17–23\].
As a phenomenological application of the NLO BFKL improved by the BLM procedure, with its effective resummation of the conformal-violating $`\beta _0`$-terms into the running coupling in all orders of the perturbation theory, one can consider the gamma–gamma scattering . This process is attractive because it is theoretically more under control than the hadron–hadron and lepton–hadron collisions, where nonperturbative hadronic structure functions are involved. In addition, for the gamma–gamma scattering the unitarisation (screening) corrections due to multiple Pomeron exchange would be less important than in hadron collisions.
The gamma–gamma cross sections with the BFKL resummation in the LO were considered in . In the NLO BFKL case one should obtain a formula analogous to LO BFKL .
In Fig. 1 we present the comparison of BFKL predictions for LO and NLO BFKL improved by the BLM procedure with data from L3 at CERN LEP. The different curves reflect the uncertainty of the theoretical predictions with the choice of the Regge scale parameter $`s_0`$, which defines the transition to the asymptotic regime. For the present calculations two variants have been chosen $`s_0=Q^2/2`$ and $`s_0=2Q^2`$, where $`Q^2`$ is the virtuality of the photons. One can see from Fig. 1 that the agreement of the NLO BFKL improved by the BLM procedure is reasonably good at energies of LEP2 $`\sqrt{s_{e^+e^{}}}=`$ 183 – 189 GeV. One can notice also that the sensitivity of the NLO BFKL results to $`s_0`$ is much smaller than in the case of the LO BFKL.
It was shown in Refs. that the unitarisation corrections in hadron collisions can lead to a value of the (bare) Pomeron intercept higher than the effective intercept value. Since the hadronic data fit yields about 1.1 for the effective intercept value , the bare Pomeron intercept value should be above it. Therefore, assuming small unitarisation corrections in the gamma–gamma scattering at large $`Q^2`$, one can accommodate the NLO BFKL Pomeron intercept value 1.13 – 1.18 in the BLM optimal scale setting, along with larger unitarisation corrections in hadronic scattering , where they can lead to a smaller effective Pomeron intercept value of about 1.1 for hadronic collisions. The above intercept value of the NLO BFKL Pomeron is in good agreement with the analysis of the diffractive dijet production at the Tevatron.
Another possible application of the BFKL approach can be the collision energy dependence of the inclusive jet production . Unlike the case with the selection of most forward/backward (Mueller-Navelet) jets \[19, 33–38\], the usual inclusive jets \[20–22, 31, 39, 40\] can be more reliable for detectors with the limited acceptance in rapidity.
The advent of the Fermilab Tevatron and the CERN LHC provides a new testing ground for the parton model—the kinematic conditions when the energies of the produced hadrons are large enough to be described by perturbation theory and, at the same time, are much smaller than the total energy of the collision (BFKL semi-hard kinematics). Because the parton model was originally invented and subsequently tested for the hard kinematics, the second condition makes it plausible that a substantial modification of the parton model will be needed to describe this BFKL semi-hard kinematic region.
The range of applicability of the QCD-improved parton model is a subject of controversy at the moment. There are statements (see, e.g., ) that the fitting capacity of the conventional QCD-improved parton model is sufficient to accommodate all the data on deep inelastic scattering (DIS) parton structure functions available at small-$`x`$ kinematics. On the other hand, the same data from HERA on DIS structure functions can be described by NLO BFKL . In addition, the data from HERA and Tevatron on most forward/backward jet production may be interpreted as a manifestation of the BFKL Pomeron , which is beyond the conventional QCD-improved parton model. The situation is further complicated by the observation that the range of applicability of the QCD-improved parton model may be different for different observables. In particular, the cross-sections of processes with specific kinematics exhibit breakdown of the applicability of finite-order perturbative QCD via the development of sensitivity to the choice of the normalisation scale. On the other hand, some dedicated combinations (ratios) of cross sections may be less sensitive to the inclusion of the higher-order corrections. An example is the scaled cross-section ratio \[45–48\] since, as follows from Ref. , it is relatively insensitive to the inclusion of the NLO correction.
Under these circumstances, it is crucial to have qualitative predictions from the conventional QCD-improved parton model (without BFKL resummation of the energy logarithms) for the new kinematic domain. If the predictions would turn out qualitatively incorrect, a generalisation, and a substitute in this kinematical domain of the QCD-improved parton model would become indispensable.
Here we discuss such a prediction, made in Ref. . It is a prediction for the ratio of inclusive single jet production at a smaller energy $`\sqrt{s_N}`$ of the hadron collision to the one at a higher energy $`\sqrt{s_D}`$:
$$R(x,y,s_N,s_D)=\frac{s_Nd\sigma }{dxdy}(x,y,s_N)/\frac{s_Dd\sigma }{dxdy}(x,y,s_D).$$
(1)
Here the cross-section is made dimensionless by the rescaling with the corresponding total invariant energy of the collision squared $`s_N(s_D)`$. The ratio depends on the (pseudo)rapidity $`y=1/2\mathrm{ln}(k_+/k_{})`$, where $`k_\pm =E\pm k_3`$ are the light-cone components of the momentum of the produced jet, and on the fraction of the energy $`x=(k_++k_{})/\sqrt{s_i}`$, $`i=N,D`$ deposited in the jet produced ($`s_N`$ is used for the definition of $`x`$ in the numerator, $`s_D`$ in the denominator, so that $`x`$ varies from zero to unity for both energies). Note that this scaling variable coincides in the centre-of-mass system with $`x_R=E/E_{max}=2E/\sqrt{s}`$ , the radial Feynman variable, or for $`y=0`$, i.e. for $`(\theta _{CMS}=\pi /2)`$, the scaling variable becomes the transverse Feynman variable: $`x=x_{}=2E_{}/\sqrt{s}`$.
The ratio $`R`$ (taken at $`y0`$, i.e. for jets perpendicular to the collision axes) was used in Refs. as a means to test QCD predictions for scaling violations:
$`R(x,y=0,s_N,s_D)`$ $`=`$ $`E_{}^4{\displaystyle \frac{Ed\sigma }{d^3p}}(x_{},y=0,s_N)/E_{}^4{\displaystyle \frac{Ed\sigma }{d^3p}}(x_{},y=0,s_D)`$ (2)
$`=`$ $`E_{}^4{\displaystyle \frac{d\sigma }{dydE_{}^2}}(x_{},y=0,s_N)/E_{}^4{\displaystyle \frac{d\sigma }{dydE_{}^2}}(x_{},y=0,s_D).`$ (3)
Note that without scaling violations the ratio $`R`$ is exactly unity. At fixed $`y`$ and $`x_{}`$ the dependence of $`R`$ from $`s_N`$, $`s_D`$ comes from the presence of the fundamental QCD scale, $`\mathrm{\Lambda }_{QCD}`$, in the running coupling and in the parton distribution functions.
To be more explicit, we remark that the QCD-improved parton model, based on the factorisation theorem for hard processes and the GLAPD $`\mathrm{log}Q^2`$-evolution, presents the inclusive jet scaled cross-section in hadron collsions as
$$E_{}^4\frac{Ed\sigma }{d^3p}(x_{},y=0,s)=_{x_{A,min}}^1_{x_{B,min}}^1dx_Adx_BF_A(x_A,Q^2)F_B(x_B,Q^2)E_{}^4\frac{\widehat{s}}{\pi }\frac{d\widehat{\sigma }}{d\widehat{t}}\delta (\widehat{s}+\widehat{t}+\widehat{u}),$$
(4)
where $`\widehat{s}`$, $`\widehat{t}`$ and $`\widehat{u}`$ are the Mandelstam variables for the partonic subprocess, the scale of the hard partonic subprocess $`\widehat{t}=Q^2E_{}^2x_{}^2s`$, $`F_A`$ and $`F_B`$ are parton distribution functions with the GLAPD evolution following from perturbative $`\alpha _S(Q^2)`$ expansion, and the scaled partonic subprocess is
$$E_{}^4\frac{d\widehat{\sigma }}{d\widehat{t}}\alpha _S^2(Q^2)[1+C_{NLO}\alpha _S+\mathrm{}]=\alpha _S^2(x_{}s^2)[1+C_{NLO}\alpha _S+\mathrm{}].$$
Hence, within the QCD-improved parton model, the scaled cross-section ratio for inclusive jet production at fixed $`x`$ and $`y`$ is the dimensionless function of $`\alpha _S`$. The GLAPD scaling violation due to the interacting QCD partons appears as the logarithmic<sup>1</sup><sup>1</sup>1Indeed, the small-$`x`$ asymptotics of the GLAPD evolution gives a growth of parton structure functions that is faster than any power of a logarithm, but slower than any power — so-called double-logarithmic asymptotics . dependence on the total energy of collision through the coupling constant $`\alpha _S`$. We note here that the BFKL leads to a power-like scaling violation, the strength of which depends on the Regge scale $`s_0`$.
Perturbative QCD calculations of Ref. with hard kinematics ($`Q^2s`$) predict for $`R`$ at $`y=0`$ a steep increase around the value of 1.8 – 1.9 for $`x`$ growing in the range above 0.1 (for the case $`\sqrt{s_N}`$/$`\sqrt{s_D}`$ = 0.63 TeV/1.8 TeV ). For moderate $`x`$, the prediction is in reasonable agreement with CDF data . For $`x<0.1`$, calculations are above the preliminary data of CDF . This was one of the reasons for the conclusion of Ref. that NLO GLAPD with hard kinematics is insufficient to describe the absolute cross section of jets with transverse energy less than 50 GeV within an accuracy of $`10\%`$. It was shown in Ref. that resummation of the energy logarithms, i.e. BFKL, restores the agreement between theory and experiment.
In Ref. we have found the following result: the QCD-improved parton model predicts that $`R`$ is not a monotonic function of its arguments, i.e. the single-jet production cross-section, if measured in the natural units of the same cross section taken at another (higher) energy of the collision, has extrema. Namely, it has minima (“dips”): there is a value of $`x`$ for each $`y`$ with the smallest ratio of jets produced. The reason this fact was overlooked is that for $`y=0`$ (the only value for which the calculations were reported earlier) the minimum is at a value of $`x`$ too small to be inside the acceptance of the existing detectors ($`x_{dip}(y=0)<0.01`$ at the Tevatron).
Fig. 3a presents the ratio for energies 0.63 TeV/1.8 TeV at the Fermilab Tevatron, Fig. 3b for energies 6 TeV/14 TeV at the CERN LHC. Each curve on the plots presents the dependence of the ratio on $`x`$ at different values of the rapidity $`y`$. Each curve ends at a lowest value of $`x`$ where $`\alpha _S(Q^2)`$ has the value of about 0.5 (it corresponds to $`Q=0.7`$ GeV, and $`Q`$ was taken to be half of the transverse energy of the jet produced). For lower values of $`x`$, perturbative theory becomes unreliable because the coupling approaches unity.
There is another lower bound on the values of $`x`$ at which our plots make sense, because there is a lowest energy for which the jet may be resolved. This energy is accepted now to be around 5 GeV <sup>2</sup><sup>2</sup>2At the HERA lepton–hadron collider jets are resolved from $`E_{}=3`$ GeV, and at the Tevatron hadron–hadron collider the CDF Collaboration tags jets from $`E_{}=8`$ GeV . , which corresponds to $`x>0.016`$ for the Tevatron (Fig. 3a), and to $`x>0.0017`$ for the LHC (Fig. 3b).
There is an important issue concerning the accuracy of the present LO calculation. The most important advantage of the scaled cross-section ratio is that this is the ratio of two perturbative series with the same coefficients and with different scales in the running coupling. These scales are defined by the two initial collision energies at fixed scaling variable. One can show that the theoretical accuracy of the ratio in LO of perturbative QCD is not less than the accuracy of NLO calculations for absolute cross-sections.
The minima in Fig. 3 originate from a competition between the running of the parton distribution functions and the running of the coupling constant. Namely, the ratio with frosen parton distribution functions is decreasing monotonously (this tendency is realised at small $`x`$), while the one with frosen coupling constant is growing monotonously (which is realised for $`x`$ larger than the position of the minimum).
We suggest the following potential implications of the minima we have predicted with the parton model: (i) If one observes the minima experimentally, one employs the orthodox QCD-improved parton model and tries to account for observed positions and depths of the minima by taking into account higher-order corrections, in particular, resummation of the energy logarithms. (ii) If one does not observe the minima experimentally, more radical changes are motivated, such as an alternative model of the effective constituents inside the hadrons for the BFKL semi-hard asymptotics. One example might be the colour dipole model .
Finally, we comment on the possibility of searching for the minima at the Fermilab Tevatron and at the CERN LHC: positions of the minima for the Tevatron energies (Fig. 3a) seem to be reached by both D$`\mathrm{}`$ and CDF detectors. The minima of the LHC plot (Fig. 3b) seem to be well inside the acceptance of, for example, the FELIX , the ALICE and the CMS detectors.
We take the ratio of 6 TeV/14 TeV for the LHC, because, in addition to 14 TeV $`pp`$ collisions, lead–lead collisions at the LHC are planned with the collision energy of 6 TeV per nucleon–nucleon collision. Since nuclear collisions bring in nuclear effects, which can distort our predicted curves, we also considered the ratio 6 TeV/100 TeV (see Ref. ).
Further consideration should be given to deciding which pair of energies and value of rapidity are most convenient for an experimental search of the minima. Also, more work is needed to make quantitative predictions on the locations and the shapes of the dips with the NLO corrections taken into account. It is interesting to observe a resemblance of the dips presented here with the nonmonotonic behaviour of parton distribution functions in DIS .
It is also worth noting that, in the case of nuclear collisions, the effects of initial nuclear parton distributions (small-$`x`$ EMC–effect ) and dynamical effects such as quark–gluon plasma, jet quenching, etc. , phase structure of the QCD vacuum will demand special consideration.
Before reaching conclusions, we would like to note that many of the above ideas can be studied also for the case of heavy-quarkonium production, where similar phenomena should be present .
To sum up, we find a new qualitative prediction of the QCD-improved parton model for hadron collisions and suggest its use to test the applicability of the parton model for certain regions of high-energy hadron collisions. Study of the scaling violation of the scaled cross section ratio can reveal such new dynamic effects as the BFKL asymptotics.
We thank G. Altarelli, V. P. Andreev, J. Ellis, A. De Roeck, V. S. Fadin, J. H. Field, M. Kienzle-Focacci, C.-H. Lin, L. N. Lipatov, J.-W. Qiu, V. A. Schegelsky, A. A. Vorobyov and M. Wadhwa for useful discussions, and, also, S. Vascotto for reading of the manuscript. VTK thanks the Organizing Committee of the Xth Quantum Field Theory and High Energy Physics Workshop and the CERN Theory Division for their warm hospitality. This work was supported in part by the Russian Foundation for Basic Research, Grant No. 00-02-17432, the NATO Science Programme, Collaborative Linkage Grant No. PST.CLG.976521, and the U.S. Department of Energy, Contract No. DE-FG02-87ER40371, Division of High Energy and Nuclear Physics.
|
warning/0005/math0005210.html
|
ar5iv
|
text
|
# Quasi-actions on trees Research announcement
### Remark
Theorems 3 and 4 are just samples. The main result on Inhomogeneous Rigidity, Theorem 11, has a multitude of applications. In this research announcement we give only sketches of proofs, and we focus particularly on coarse $`\mathrm{PD}`$ vertex and edge groups, ignoring wider contexts for our results. Full statements in wider contexts, and full proofs, will be found in \[MSW\].
### Remark
As in the homogeneous case, the techniques for the inhomogeneous case can be applied to quasi-isometric classification as well as to quasi-isometric rigidity; although we have not mentioned here any of these classification theorems, the preprint \[MSW\] will include some results. We should also mention \[PW00\], where techniques similar to those of Theorem 11 are used in proving that the quasi-isometric classification of accessible groups completely reduces to the classification of one-ended groups.
## Graphs of groups and Bass-Serre trees: a review
A graph of groups $`\mathrm{\Gamma }`$ of finite type consists of the following data. Start with a finite graph $`\mathrm{\Gamma }`$ with vertex set $`𝒱(\mathrm{\Gamma })`$ and edge set $`(\mathrm{\Gamma })`$. For each edge $`e(\mathrm{\Gamma })`$, the two ends of $`e`$ form a set $`\mathrm{Ends}(e)`$, and each end $`\eta \mathrm{Ends}(e)`$ is attached to some vertex $`v(\eta )𝒱(\mathrm{\Gamma })`$. Associated to each vertex $`v𝒱(\mathrm{\Gamma })`$ there is a finitely generated vertex group $`\mathrm{\Gamma }_v`$, associated to each edge $`e(\mathrm{\Gamma })`$ there is an edge group $`\mathrm{\Gamma }_e`$, and associated to each end $`\eta \mathrm{Ends}(e)`$ there is an edge-to-vertex injection $`\gamma _\eta :\mathrm{\Gamma }_e\mathrm{\Gamma }_{v(\eta )}`$.
Associated to a graph of groups $`\mathrm{\Gamma }`$ is a *graph of spaces* $`B_\mathrm{\Gamma }`$, as follows. Choose path connected, pointed spaces $`B(v)`$, $`v𝒱(\mathrm{\Gamma })`$, and $`B(e)`$, $`e(\mathrm{\Gamma })`$, whose fundamental groups are identified with the associated vertex or edge groups of $`\mathrm{\Gamma }`$. Choose pointed *attaching maps* $`\xi _\eta :B(e)B(v)`$ inducing the injections $`\gamma _\eta `$. For each edge $`e`$ let $`\widehat{e}=\mathrm{int}(e)\mathrm{Ends}(e)`$ denote the end compactification of $`e`$, ($`\widehat{e}`$ is a compact arc, regardless of whether $`e`$ is a compact arc or a loop). Construct a quotient space $`B_\mathrm{\Gamma }`$ from the disjoint union of the set
$$\{B(v),B(e)\times \widehat{e}|v𝒱(\mathrm{\Gamma }),e(\mathrm{\Gamma })\}$$
by gluing $`B(e)\times \eta `$ to $`B(v(\eta ))`$ via the map $`(x,\eta )\xi _\eta (x)`$, for each $`e(\mathrm{\Gamma })`$, $`\eta \mathrm{Ends}(e)`$. The fundamental group $`\pi _1(B_\mathrm{\Gamma })`$ is well-defined up to isomorphism, and it is called the fundamental group of $`\mathrm{\Gamma }`$, denoted $`\pi _1(\mathrm{\Gamma })`$.
The map $`B\mathrm{\Gamma }`$, taking $`B(v)`$ to $`v`$ and projecting $`B(e)\times \widehat{e}`$ to $`e`$, induces a decomposition of $`B`$ into point preimages. In the universal cover $`X=\stackrel{~}{B}`$, taking connected lifts via the universal covering map $`XB`$ of point preimages of $`B\mathrm{\Gamma }`$ gives a decomposition of $`X`$ into path connected sets. This decomposition of $`X`$ is $`\pi _1(\mathrm{\Gamma })`$-equivariant. The quotient space of this decomposition of $`X`$ is a tree $`T`$ on which $`\pi _1(\mathrm{\Gamma })`$ acts, the *Bass-Serre tree* of $`\mathrm{\Gamma }`$. This action is well-defined up to equivariant tree isomorphisms, independent of the choices, and the quotient graph of $`T`$ is canonically identified with $`\mathrm{\Gamma }`$:
The map $`XT`$ is called the *Bass-Serre tree of spaces* associated to the graph of spaces $`B\mathrm{\Gamma }`$. The inverse image of a vertex $`vT`$ is a *vertex space* $`X(v)`$ of $`X`$. The inverse image of the midpoint of an edge $`e`$ of $`T`$ is called an *edge space* $`X(e)`$ of $`X`$. The topological space $`X`$ is constructed from the disjoint union of the set $`\{X(v),X(e)\times e|v𝒱(T),e(T)\}`$, where for each vertex $`v`$ and edge $`e`$ incident to $`v`$ we glue $`X(e)\times v`$ to a subset of $`X(v)`$ via an attaching map $`X(e)X(v)`$ which is a lift of an attaching map for the graph of spaces $`B`$. The image of the attaching map $`X(e)X(v)`$ is called an *incident edge space inside $`X(v)`$*. The set of incident edge spaces inside $`X(v)`$ is called the *edge space pattern inside $`X(v)`$*.
A simple trichotomy holds for Bass-Serre trees: $`T`$ is *bounded*; or $`T`$ is *line-like* meaning that it contains a line as a cobounded subset; or $`T`$ is *bushy* meaning that it has infinitely many ends. A graph of groups $`\mathrm{\Gamma }`$ (and its associated Bass-Serre tree) is said to be *reduced* if for each vertex $`v`$, the number of surjective edge-to-vertex injections $`\gamma _\eta :\mathrm{\Gamma }_e\mathrm{\Gamma }_v`$ is not equal to $`1`$; it can be $`0`$ or $`2`$. When $`T`$ is reduced, the trichotomy simplifies as follows. First, $`T`$ is bounded if and only if $`T`$ and $`\mathrm{\Gamma }`$ are each a point. Second, $`T`$ is line-like if and only if $`T`$ is a line and $`\mathrm{\Gamma }`$ is a *mapping torus*, meaning either a circle with isomorphic edge-to-vertex inclusions all around, or an arc with isomorphic inclusions in the interior and index 2 inclusions at the endpoints. Finally, $`T`$ is bushy if and only if $`T`$ has at least one vertex of valence $`3`$; valence of a vertex in $`T`$ is easily computed in terms of the image vertex in $`\mathrm{\Gamma }`$, as the sum of the indices of the edge groups inside the vertex group.
## Coarse language
Let $`X`$ be a metric space. Given $`AX`$ and $`R0`$, denote $`N_R(A)=\{xX|aA\text{such that}d(a,x)R\}`$. Given subsets $`A,BX`$, let $`A_cB[R]`$ denote $`AN_R(B)`$. Let $`A_cB`$ denote the existence of $`R0`$ such that $`A_cB[R]`$; this is called *coarse containment* of $`A`$ in $`B`$. Let $`A=_cB[R]`$ denote $`A_cB[R]`$ and $`B_cA[R]`$. Let $`A=_cB`$ denote the existence of $`R`$ such that $`A=_cB[R]`$; this is called *coarse equivalence* of $`A`$ and $`B`$.
Given a metric space $`X`$ and subsets $`A,B`$, we say that a subset $`C`$ is a *coarse intersection* of $`A`$ and $`B`$ if for all sufficiently large $`R`$ we have $`N_R(A)N_R(B)=_cC`$. A coarse intersection of $`A`$ and $`B`$ may not exist, but if one does exist then it is well-defined up to coarse equivalence.
Given metric spaces $`X,Y`$, a map $`f:XY`$ is a *uniform embedding* if there exists proper, increasing functions $`g,h:[0,\mathrm{})[0,\mathrm{})`$ such that
$$g(d_X(x,y))d_Y(fx,fy)h(d_X(x,y))$$
If $`X,Y`$ are geodesic metric spaces then the upper bound $`h`$ can always be taken to be an affine function. When $`h(d)=Kd+C`$ and $`g(d)=\frac{1}{K}dC`$ then we say that $`f`$ is a *$`K,C`$ quasi-isometric embedding*. If this is so then we say in addition that $`f`$ is a *$`K,C`$ quasi-isometry from $`X`$ to $`Y`$* if $`f(X)=_cY[C]`$. A *$`C^{}`$-coarse inverse* of $`f`$ is a $`K,C^{}`$ quasi-isometry $`g:YX`$ such that $`x=_cg(f(x))[C^{}]`$ and $`y=_cf(g(y))[C^{}]`$, for all $`xX`$, $`yY`$. A simple fact says that for all $`K,C`$ there exists $`C^{}`$ such that each $`K,C`$ quasi-isometry has a $`C^{}`$-coarse inverse.
Let $`G`$ be a group and $`X`$ a metric space. A *$`K,C`$ quasi-action* of $`G`$ on $`X`$ is a map $`(g,x)gx`$ from $`G\times X`$ to $`X`$, such that: for each $`g`$ the map $`xgx`$ is a $`K,C`$ quasi-isometry; and for each $`xX,g,hG`$ we have
$$g(hx)=_c(gh)x[C]$$
A quasi-action is *cobounded* if there exists a constant $`R`$ such that for each $`xX`$ we have $`Gx=_cX[R]`$. A quasi-action is *proper* if for each $`R`$ there exists $`M`$ such that for all $`x,yX`$, the cardinality of the set $`\{gG|\left(gN(x,R)\right)N(y,R)\mathrm{}\}`$ is at most $`M`$.
A fundamental principle of geometric group theory says that if $`G`$ is a finitely generated group equipped with the word metric, and if $`X`$ is a proper geodesic metric space on which $`G`$ acts properly discontinuously and cocompactly by isometries, then $`G`$ is quasi-isometric to $`X`$.
A partial converse to this result is the *quasi-action principle* which says that if $`G`$ is a finitely generated group with the word metric and $`X`$ is a metric space quasi-isometric to $`G`$ then there is a cobounded, proper quasi-action of $`G`$ on $`X`$; the constants for this action depend only on the quasi-isometry constants between $`G`$ and $`X`$.
The quasi-action principle motivates the following question, which is a common point of departure for many quasi-isometric rigidity problems:
* Given a proper, cobounded quasi-action of a group $`G`$ on a metric space $`X`$, when can we get some action of $`G`$ on $`X`$?
Partial information about this question can sometimes be obtained from the following result (see e.g. \[KL97\]):
###### Proposition 5 (Coboundedness Principle).
Suppose that a finitely generated $`G`$ quasi-acts properly and coboundedly on a metric space $`X`$. Let $``$ be a collection of subsets of $`X`$ which satisfies the following properties: the elements of $``$ are pairwise coarsely inequivalent in $`X`$; there exists $`A0`$ such that for each $`gG`$, $`H`$ there is an $`H^{}`$ such that $`gH=_cH^{}[A]`$; and every metric ball in $`X`$ intersects at most finitely many elements of $``$. Then for each $`H`$, the *stabilizer subgroup* $`\mathrm{Stab}_G(H)=\{gG|gH=_cH\}`$ quasi-acts properly and coboundedly on $`H`$.
We will need some coarse algebraic topology. This subject originated in \[FS96\], by applying Alexander Duality with naturality to fundamental groups of closed aspherical manifolds. These ideas were generalized in \[KK99\] to work for a general class of spaces, the *coarse $`\mathrm{PD}(n)`$ spaces $`X`$*. These are “uniformly acyclic” simplicial complexes which satisfy a coarse, uniform version of the Poincaré duality property of $`𝐑^n`$: there exist chain maps $`C_{}(X)\stackrel{𝐷}{}C_c^n(X)`$ and $`C_c^n(X)\stackrel{\overline{D}}{}C_{}(X)`$ with “uniform distortion” such that $`D\overline{D}`$ and $`\overline{D}D`$ are “uniformly chain homotopic” to the identity; see \[KK99\] for the full development. A coarse $`\mathrm{PD}(n)`$ group $`G`$ is one which acts properly and coboundedly on some coarse $`\mathrm{PD}(n)`$ space.
A subset $`H`$ of a metric space $`X`$ is *deep* if for each $`r`$ there exists $`xX`$ such that the ball in $`X`$ around $`x`$ of radius $`r`$ is a subset of $`H`$.
###### Proposition 6 (Coarse Jordan Separation).
If $`X`$ is a coarse $`\mathrm{PD}(n)`$ space and $`SX`$ is a uniformly embedded coarse $`\mathrm{PD}(n1)`$ space, then for sufficiently large $`A`$ there are exactly two deep components of $`XN_A(S)`$. The coarse intersection of these two components is $`S`$.
For fundamental groups of closed, aspherical $`n`$-manifolds, this result with a weaker conclusion of “at least two deep components” comes from \[FS96\]. The improvement to “exactly two deep components”, and the generalization to coarse $`\mathrm{PD}(n)`$ spaces, comes from \[KK99\].
We will state other coarse algebraic topology results as we need them below.
## Quasi-actions on trees
Suppose that $`\mathrm{\Gamma }`$ is a finite type graph of groups, $`B\mathrm{\Gamma }`$ is an associated graph of spaces, and $`XT`$ is the associated Bass-Serre tree of spaces. We make the additional assumption that each edge and vertex space of $`B`$ is compact; for instance, if all edge and vertex groups are finitely presented then we can take the edge and vertex spaces to be presentation complexes. We impose a geodesic metric on $`B`$, which lifts to a geodesic metric on $`X`$. It follows that $`\pi _1\mathrm{\Gamma }`$ is quasi-isometric to $`X`$. If $`G`$ is any finitely generated group quasi-isometric to $`\pi _1\mathrm{\Gamma }`$, it follows that $`G`$ is also quasi-isometric to $`X`$, and therefore $`G`$ has a cobounded, proper quasi-action on $`X`$. In this situation, the motivating question is:
* Does the quasi-action of $`G`$ on $`X`$ coarsely respect the vertex and edge spaces of $`X`$?
While this question is somewhat vague, there are several distinct ways to make it more precise. Here is the most rigid possible behavior:
* Tree rigidity There is an action of $`G`$ on $`T`$ such that for each $`gG`$, $`v𝒱(T)`$, and $`e(T)`$ we have $`gX(v)=_cX(gv)`$ and $`gX(e)=_cX(ge)`$ (uniformly, i.e. with uniform coarseness constant independent of $`g,v`$).
Tree rigidity immediately implies, for example, that $`G`$ is the fundamental group of a graph of groups whose vertex and edge groups are all quasi-isometric to vertex and edge groups of $`\mathrm{\Gamma }`$.
In general, tree rigidity is a lot to expect. Ignoring edges for the moment, here is a sequence of vertex rigidity properties, from weaker to stronger:
* Weak vertex rigidity For each $`gG`$, $`v𝒱(T)`$ there exists $`w𝒱(T)`$ such that $`gX(v)_cX(w)`$ \[uniformly\].
* Vertex rigidity For each $`gG`$, $`v𝒱(T)`$ there exists $`w𝒱(T)`$ such that $`gX(v)=_cX(w)`$ \[uniformly\].
* Vertex rigidity with uniqueness For each $`gG`$, $`v𝒱(T)`$ there exists a unique $`w𝒱(T)`$ such that $`gX(v)=_cX(w)`$ \[uniformly\].
In general these three properties are not equivalent. In certain situations they become equivalent, e.g. when no vertex space is coarsely contained in another, which occurs if and only if every edge-to-vertex group injection has infinite index (a property which fails completely in the context of Theorem 1).
The utility of vertex rigidity with uniqueness, for instance, is that it allows us to define an action of $`G`$ on $`𝒱(T)`$, where $`gv=w`$ if and only if $`gX(v)=_cX(w)`$. Even then we don’t get any action of $`G`$ on the edges, without knowing something more about edge spaces, such as:
* Strong edge rigidity For all edges $`ee^{}(T)`$, $`X(e_{\mathrm{}}=_cX(e^{})`$.
Vertex rigidity with uniqueness, coupled with strong edge rigidity, implies tree rigidity, for if $`e`$ is an edge with endpoints $`v_1,v_2`$, and if $`e^{}`$ is the edge with endpoints $`gv_1,gv_2`$, then we have
$$gX(e)=_cg(X(v_1)_cX(v_2))=_cgX(v_1)_cgX(v_2)=_cX(gv_1)_cX(gv_2)=_cX(e^{})$$
and so by setting $`ge=e^{}`$ we obtain a well-defined action of $`G`$ on $`T`$ satisfying tree rigidity.
### Example
By the work of Kapovich–Leeb \[KL97\], the torus decomposition of a non-solv Haken 3-manifold satisfies tree rigidity. The heart of their work is a proof of weak vertex rigidity, using asymptotic cones. Tree rigidity follows easily from that, because the torus decomposition has infinite index edge-to-vertex group injections, and obviously satisfies strong edge rigidity.
## Homogeneous graphs of groups: proof of Theorem 1
A graph of groups is *geometrically homogeneous* if it it satisfies any of the following equivalent statements: each edge-to-vertex injection is a quasi-isometry; each edge-to-vertex injection has finite index image; the Bass-Serre tree has bounded valence. It follows that the edge and vertex groups are all quasi-isometric.
Here’s the first step in the proof of Theorem 1:
###### Theorem 7.
Suppose $`\mathrm{\Gamma },\mathrm{\Gamma }^{}`$ are geometrically homogeneous graphs of coarse $`\mathrm{PD}`$ groups with bushy Bass-Serre trees $`T,T^{}`$, and trees of spaces $`X,X^{}`$. If $`h:XX^{}`$ is a quasi-isometry then $`h`$ respects vertex spaces. More precisely, for each $`K,C`$ there exists $`A`$ such that if $`h:XX^{}`$ is a $`K,C`$ quasi-isometry then for each $`v𝒱(T)`$ there exists $`v^{}𝒱(T^{})`$ such that $`h(X(v))=_cX^{}(v^{})[A]`$.
A proof with some extra topological assumptions is found in \[FM00b\], and that proof generalizes almost word-for-word to coarse $`\mathrm{PD}(n)`$ vertex and edge groups. Applying Theorem 7 to each element of a quasi-action we obtain:
###### Corollary 8.
With the same notation as in Theorem 7, if $`H`$ is a finitely generated group quasi-isometric to $`\pi _1\mathrm{\Gamma }`$ then the quasi-action of $`H`$ on $`X`$ satisfies vertex rigidity. It follows that there is a cobounded quasi-action of $`H`$ on $`𝒱(T)`$, with the property that $`hX(v)=_cX(hv)`$ \[uniformly\].
The following result is the technical heart of Theorem 1:
###### Theorem 9.
Let $`T`$ be a bushy tree of uniformly bounded valence. Let $`H`$ be a group quasi-acting coboundedly on $`T`$. Then the quasi-action of $`H`$ on $`T`$ is quasiconjugate to a cobounded action of $`H`$ on a tree $`T^{}`$. That is, there is a quasi-isometry $`f:TT^{}`$ such that $`h(f(v))=_cf(hv)`$ \[uniformly\].
This combines with Corollary 8 to prove Theorem 1, for one can show first that the tree $`T^{}`$ has bounded valence, and so there is a graph of groups $`\mathrm{\Gamma }^{}`$ with fundamental group $`H`$ and Bass-Serre tree $`T^{}`$. Furthermore, the quasi-isometry $`H\pi _1(\mathrm{\Gamma })`$ takes the vertex and edge stabilizers of the action of $`H`$ on $`T^{}`$ to subsets of $`\pi _1\mathrm{\Gamma }`$ coarsely equivalent to the vertex and edge stabilizers of the action of $`\pi _1\mathrm{\Gamma }`$ on $`T`$.
###### Proof of Theorem 9.
The first step of the proof is to construct a vertex set and an action of $`H`$. We do not even have a true action of $`H`$ on $`𝒱(T)`$, however. What we do have is an action of $`H`$ on the space of ends $`\mathrm{Ends}(T)`$. This can be promoted into an action of $`H`$ on a new vertex space as follows. Bushiness and bounded valence of $`T`$ tells us that the space of ends $`\mathrm{Ends}(T)`$ is a Cantor set. Note that each edge of $`T`$ determines a partition of $`\mathrm{Ends}(T)`$ into a pairwise disjoint set of clopens (closed-open subsets). Define a *quasi-edge* of $`T`$ to be any partition of $`\mathrm{Ends}(T)`$ into a pair of disjoint clopens. Clearly $`H`$ acts on the set of quasi-edges of $`T`$. If $`O\mathrm{Ends}(T)`$ is a clopen then the convex hull $`(O)`$ is a subtree of $`T`$. If $`\{O,O^{}\}`$ is a quasi-edge then the intersection of the 1-neighborhoods of the convex hulls $`N_1((O))N_1((O^{}))`$ is bounded; the diameter of this intersection is defined to be the *quasi-edge constant* of $`\{O,O^{}\}`$. Note that a 1-quasi-edge is the same thing as (the partition determined by) a true edge of $`T`$. The action of a $`(K,C)`$ quasi-isometry of $`T`$ on an $`A`$-quasi-edge produces an $`A^{}`$-quasi-edge, where $`A^{}`$ depends only on $`K,C,A`$. It follows that the $`H`$-orbit of a single quasi-edge has uniform quasi-edge constant.
We therefore *define* a vertex set $`Y^0`$ to be the $`H`$-orbit of some arbitrarily chosen quasi-edge of $`T`$. It is not hard, using uniformity of the quasi-edge constant, to attach edges to $`Y^0`$ in a locally finite, $`H`$-equivariant manner, thereby producing a locally finite graph $`Y^1`$ on which $`H`$ acts coboundedly, and a coarse conjugation $`TY^1`$.
Unfortunately, $`Y^1`$ may not be a tree, i.e. in particular it may not be simply connected. The latter problem can be corrected without too much difficulty by attaching 2-cells to $`Y^1`$ in a locally finite, $`H`$-equivariant manner, producing a locally finite, simply connected 2-complex $`Y^2`$ on which $`H`$ acts coboundedly; the coarse conjugation $`TY^1`$ extends to a coarse conjugation $`TY^2`$.
The final step is to get from the 2-complex $`Y^2`$ back to a tree. This is accomplished by using tracks in the sense of Dunwoody \[Dun85\]: a track $`\tau `$ in $`Y^2`$ is a locally separating, connected 1-complex in general position with respect to the skeleta of $`Y^2`$. Since $`Y^2`$ is simply connected, each track $`\tau `$ separates $`Y^2`$ into two components, and $`\tau `$ is *essential* if each component of $`Y\tau `$ is unbounded. Using minimal surface ideas as in \[Dun85\], we construct an $`H`$-equivariant family of pairwise disjoint, essential tracks $`\{\tau _i\}`$ in $`Y^2`$. By a Haken finiteness argument as in \[Dun85\] one shows that a maximal such family $`\{\tau _i\}`$ exists, and that all components of $`Y_i\tau _i`$ are bounded. The final tree $`T^{}`$ has vertex set in 1–1 correspondence with the components of $`Y_i\tau _i`$, and edge set in 1–1 correspondence with the set $`\{\tau _i\}`$. ∎
Theorem 9 has an interesting consequence concerning lattices in locally compact topological groups. Some quasi-isometry classes $`𝒞`$ have the nice property that there is a locally compact group $`H`$ such that each group $`G𝒞`$ has a discrete, cocompact quotient in $`H`$ with finite kernel. This is true, for example, when $`𝒞`$ is the quasi-isometry class of cocompact lattices in a semisimple Lie group. However, this doesn’t work for $`𝒞=\{\text{virtually free of rank }2\}`$:
###### Corollary 10.
There does not exist a locally compact group $`H`$ such that every group which is virtually free of rank $`2`$ has a discrete, cocompact quotient in $`H`$ with finite kernel.
###### Proof.
Suppose $`H`$ exists as stated. Choose a free, discrete, cocompact subgroup $`GH`$. We may give $`H`$ a left-invariant metric so that the inclusion of $`G`$ into $`H`$ is a quasi-isometry. Since $`H`$ is quasi-isometric to the free group $`G`$, it follows that $`H`$ is quasi-isometric to a finite valence tree $`T`$. The left action of $`H`$ on itself is quasi-conjugate to a cobounded quasi-action of $`H`$ on $`T`$. Applying Theorem 9 it follows that $`H`$ has a cobounded quasi-action on a bushy tree $`T^{}`$ of bounded valence. It follows that *every* virtually free group $`G`$ has a cobounded action on $`T^{}`$. Pick a prime $`p`$ larger than the maximal valence of a vertex of $`T^{}`$. The group $`G=𝐙/p𝐙/p`$ is virtually free, and therefore has discrete cocompact image in $`H`$ with finite kernel, and so $`G`$ acts coboundedly on $`T^{}`$. But each of the free factors $`𝐙/p`$ acts trivially on $`T^{}`$, making $`G`$ act trivially, contradicting coboundedness. ∎
## Inhomogeneous graphs of groups
Theorem 3, and part of the conclusion of Theorem 4, will both follow from more general results about geometrically inhomogeneous graphs of groups. Although as in the homogeneous case we have results in wider contexts, for present purposes we restrict to coarse $`\mathrm{PD}(n)`$ vertex and edge groups.
Suppose that $`\mathrm{\Gamma }`$ is a graph of coarse $`\mathrm{PD}(n)`$ groups of various dimensions. An *$`n`$-raft* is a connected subgraph of $`\mathrm{\Gamma }`$ (or of $`T`$) of constant dimension $`n`$, such that each edge incident to the raft but not contained in the raft has dimension $`<n`$. Rafts in $`T`$ are connected lifts of rafts in $`\mathrm{\Gamma }`$. Each raft in $`T`$ is the Bass-Serre tree for the corresponding quotient raft in $`\mathrm{\Gamma }`$.
Assuming $`\mathrm{\Gamma }`$ is reduced, each line-like raft in $`T`$ is actually a line, and the quotient is a mapping torus raft in $`\mathrm{\Gamma }`$. The problem with mapping tori is that they generally fail to satisfy even weak vertex rigidity. For example, a closed hyperbolic 3-manifold fibering over $`S^1`$ has a fundamental group whose quasi-isometry group is all of $`\mathrm{QI}(𝐇^3)=QC(S^2)`$, and weak vertex rigidity fails miserably. A lattice in 3-dimensional solv geometry is a mapping torus of $`𝐙^2`$, and while vertex rigidity is conjectured to hold \[FM00c\], the conjecture remains open. For these and other reasons, in all of our further theorems we must avoid line-like rafts in Bass-Serre trees.
Consider an $`n`$-dimensional point raft $`v`$ of the Bass-Serre tree $`T`$ with associated vertex space $`X(v)X`$. Let $``$ be the edge space pattern inside $`X(v)`$. For each $`m<n`$ define a subset $`_m`$ of $`m`$-dimensional edge spaces inside $`X(v)`$, and let $`_{[m,m^{}]}=_m_{m+1}\mathrm{}_m^{}`$. The Coarse Jordan Separation Theorem implies that each $`E_{n1}`$ coarsely separates $`X(v)`$ into two deep pieces, whose coarse intersection is $`E`$; in this situation we say that a subset of $`X(v)`$ *crosses* $`E`$ if it intersects each of the two deep pieces arbitrarily far from $`E`$. Define the *crossing graph* of $`X(v)`$ to be the graph whose vertex set is $`_{n1}`$, with an edge between $`E,E^{}_{n1}`$ if $`E,E^{}`$ cross each other or if there exists some element of $`_{[1,n2]}`$ which crosses both $`E`$ and $`E^{}`$.
###### Theorem 11.
Let $`\mathrm{\Gamma }`$ be a finite, reduced graph of coarse $`\mathrm{PD}`$ groups satisfying the following *“raft hypotheses”*: the Bass-Serre tree $`T`$ has no line rafts; and for each point raft $`v`$ of $`T`$, the crossing graph of $`X(v)`$ is connected or empty. If $`H`$ is a finitely generated group quasi-isometric to $`\pi _1\mathrm{\Gamma }`$, then there is a finite type, reduced graph of groups $`\mathrm{\Gamma }^{}`$ with $`H\pi _1\mathrm{\Gamma }^{}`$ and with Bass-Serre tree of spaces $`X^{}T^{}`$, and there is a quasi-isometry $`f:X^{}X`$ coarsely conjugating the $`H`$ action on $`X^{}`$ to the $`H`$-quasi-action on $`X`$, such that $`f`$ coarsely respects rafts and their vertex spaces, and $`f`$ coarsely respects all edge spaces. That is: for each raft $`t^{}`$ of $`X^{}`$ there exists a raft $`t`$ of $`X`$ such that $`f(X^{}(t^{}))=_cX(t)`$ \[uniformly\], for each vertex $`v^{}t^{}`$ there exists a vertex $`vt`$ such that $`f(X^{}(v^{}))=_cX(v)`$ \[uniformly\]; and for each edge $`e^{}`$ of $`T^{}`$ there exists an edge $`e`$ of $`T`$ such that $`f(X^{}(e^{}))=_cX(e)`$ \[uniformly\].
Although the conclusion makes no mention of vertices $`v^{}T^{}`$ not on any raft of $`T^{}`$, such information may be derived as follows: since $`v^{}`$ is not on any raft, there exists an edge $`e^{}T^{}`$ incident to $`v^{}`$ such that $`X^{}(v^{})=_cX^{}(e^{})`$, and so $`f(X^{}(v^{}))=_cf(X^{}(e^{}))`$ is coarsely equivalent to some edge space in $`X`$. However, counterexamples show that $`f(X^{}(v^{}))`$ might not be coarsely equivalent to any vertex space in $`X`$.
Before proving the theorem, we apply it to Theorem 3 and part of Theorem 4.
The hypothesis (\*) in Theorem 3 immediately implies the raft hypotheses in Theorem 11, and the conclusion of Theorem 11 immediately implies the conclusion of Theorem 3. Note however that Theorem 11 gives a much stronger conclusion, namely that the new Bass-Serre tree of spaces $`X^{}T^{}`$ maps to original Bass-Serre tree of spaces $`XT`$, preserving the coarse inclusions among vertex and edge spaces. We’ll use this additional information later, combining it with our Abelian Pattern Rigidity Theorem 14 to obtain a stronger rigidity result.
For Theorem 4, the fact that vertex groups have dimension $`3`$ and edge groups have dimension $`1`$ implies that all rafts are point rafts with empty crossing graph, and so Theorem 11 applies. We conclude that all edge groups of $`\mathrm{\Gamma }^{}`$ are quasi-isometric to $`𝐙`$, and therefore commensurable to $`𝐙`$. We also conclude that each vertex group $`\mathrm{\Gamma }_v^{}^{}`$ satisfies one of two possibilities: $`\mathrm{\Gamma }_v^{}^{}`$ is quasi-isometric and so commensurable to $`𝐙`$; or $`\mathrm{\Gamma }_v^{}^{}`$ quasi-isometric to some vertex group $`\mathrm{\Gamma }_v`$, and so to $`𝐇^n`$, and so $`\mathrm{\Gamma }_v^{}^{}`$ is weakly commensurable to *some* closed $`𝐇^n`$ orbifold group. But in the latter case we get more information: the ambient quasi-isometry $`X^{}X`$ takes $`X^{}(v^{})`$ to $`X(v)`$, coarsely mapping the edge space pattern inside $`X^{}(v^{})`$ to the edge space pattern inside $`X(v)`$. We’ll make this more precise later, and combine it with Schwartz’ Geodesic Pattern Rigidity Theorem 13, to get the stronger conclusion of Theorem 4, namely that $`\mathrm{\Gamma }_v^{}^{}`$ is weakly commensurable to $`\mathrm{\Gamma }_v`$ itself.
## Sketch of proof of Theorem 11
Let $`N`$ be the maximal dimension of a vertex in $`T`$, and define a filtration $`T_N\mathrm{}T_{i+1}T_i\mathrm{}T_0=T`$ where $`T_i`$ is the union of all vertices and edges of dimension $`i`$. Note that $`T_N`$ is a disjoint union of $`N`$-rafts. There may be lower dimensional rafts as well: any component of $`T_i`$ which does not contain a component of $`T_{i+1}`$ is an $`i`$-raft. Let $`\mathrm{\Gamma }_i`$ be the image of $`T_i`$ in $`\mathrm{\Gamma }`$, and let $`X_i`$ be the inverse image of $`T_i`$ in $`X`$. Thus, each component of $`X_i`$ maps to a component of $`T_i`$, giving a Bass-Serre tree of spaces for the corresponding component of $`\mathrm{\Gamma }_i`$.
Let $`H`$ be a group quasi-isometric to $`\pi _1\mathrm{\Gamma }`$, and so $`H`$ quasi-acts properly and coboundedly on $`X`$. The method of the proof is to work inductively down from the top dimension, replacing the quasi-action of $`H`$ on the tree of spaces $`X_iT_i`$ by a true action, starting with $`i=N`$. The basis step of the induction depends on the following:
###### Proposition 12 (Vertex rigidity at the top dimension).
The quasi-action of $`H`$ on $`X`$ coarsely respects $`N`$-rafts and their vertex spaces, that is: for each $`N`$-raft $`tT`$ and each $`hH`$ there exists an $`N`$-raft $`t^{}T`$ such that $`hX(t)=_cX(t^{})`$ \[uniformly\], and for each vertex $`vt`$ there exists a vertex $`v^{}t^{}`$ such that $`hX(v)=_cX(v^{})`$ \[uniformly\].
Once this is shown, applying Proposition 5 to the collection of $`N`$-rafts we conclude that the stabilizer of each $`N`$-raft quasi-acts coboundedly and properly, and then applying Theorem 1 we may quasi-conjugate the action of $`\mathrm{Stab}(T)`$ on each $`N`$-raft $`T`$ to a true action of $`\mathrm{Stab}(T)`$ on a new tree $`T^{}`$. Thus we establish the basis step for the inductive proof of Theorem 11.
###### Proof of Proposition 12.
We shall prove:
* There exists $`A`$ such that for each $`vT_N`$ and each $`hH`$ there is an $`N`$-raft $`t`$ with $`hX(v)_cX^{}(t)[A]`$.
To see why this suffices, consider a raft $`t`$ of $`T_N`$ and two vertices $`v_1,v_2𝒱(t)`$, and so $`X(v_1)=_cX(v_2)`$. Applying $`()`$ we get $`hX(v_1)_cX^{}(t_1)`$, $`hX(v_2)_cX^{}(t_2)`$ for $`N`$-rafts $`t_1,t_2`$ of $`T`$, and we want to verify that $`t_1=t_2`$. Since $`X(v_1)=_cX(v_2)`$ it follows that the coarse intersection of $`X^{}(t_1)`$ and $`X^{}(t_2)`$ coarsely contains a coarse $`\mathrm{PD}(N)`$ space, namely $`hX(v_1)=_chX(v_2)`$. If $`t_1t_2`$ then the coarse intersection of $`X^{}(t_1)`$ and $`X^{}(t_2)`$ is coarsely equivalent to some edge space of lower dimension, but a coarse $`\mathrm{PD}(n)`$ space with $`n<N`$ cannot contain a uniformly embedded copy of a coarse $`\mathrm{PD}(N)`$ space. It follows that $`t_1=t_2`$ and $`hX(t)_cX(t_1)`$. Similarly $`h^1X(t_1)_cX(t^{})`$ for some $`N`$-raft $`t^{}`$, and so $`X(t^{})_cX(t)`$; but this is only possible if $`t=t^{}`$. This shows that $`H`$ coarsely respects $`N`$-rafts.
If (\*) is not true then, taking counterexamples for larger and larger values of $`A`$, using coboundedness of the isometry group of $`XT`$, and passing to a limit, we obtain a quasi-isometry $`h:XX`$, a vertex $`v𝒱(T_N)`$ and an edge $`e(T)(T_N)`$, such that $`h(X(v))`$ intersects both components of $`XX(e)`$ arbitrarily deeply. It follows that for sufficiently large $`A`$, the subset
$$S=h^1(N_A(X(e)))X(v)$$
coarsely separates $`X(v)`$ into at least two deep components. But the set $`S`$, with metric restricted from $`X(v)`$, is uniformly equivalent to a subset of $`X(e)`$ with restricted metric.
When $`X(e)`$ has dimension $`N2`$ we obtain a contradiction using arguments of coarse algebraic topology: a coarse $`\mathrm{PD}(N)`$ space cannot be coarsely separated by a subset which is uniformly equivalent to a subset of a coarse $`\mathrm{PD}`$ space of dimension $`N2`$.
If $`X(e)`$ is of dimension $`=N1`$ then in fact $`\pi (h(S))=_cX(e)`$: otherwise, a subset of the coarse $`\mathrm{PD}(N1)`$ space $`X(e)`$ which is not coarsely equivalent to all of $`X(e)`$ would embed uniformly in coarse $`\mathrm{PD}(N)`$ space $`X(v)`$, coarsely separating $`X(v)`$, and that is impossible. We therefore have $`h(S)=_cX(e)`$ in $`X`$. This shows that $`S`$ is a coarse $`\mathrm{PD}(N1)`$ space uniformly embedded in the coarse $`\mathrm{PD}(N)`$ space $`X(v)`$, and so $`X(v)S`$ has exactly two deep components each coarsely containing $`S`$, each contained in a separate deep component of $`XS`$.
Now the argument breaks into cases.
Case 1: $`v`$ is contained in a bushy raft $`t`$. In this case the coarse $`\mathrm{PD}(N1)`$ space $`S`$ coarsely separates $`X(t)`$, a tree of $`\mathrm{PD}(N)`$ spaces, and that is clearly impossible.
Case 2: $`v`$ is a point raft. By hypothesis, the crossing graph of $`X(v)`$ is either connected or empty.
Case 2a: The crossing graph of $`X(v)`$ is connected. Using connectedness of the crossing graph together with some coarse separation arguments, one shows that one of the deep components of $`X(v)S`$ coarsely contains the union of all codimension-1 edge spaces inside $`X(v)`$. But this is absurd, because coboundedness of the $`\pi _1\mathrm{\Gamma }`$ stabilizer subgroup of $`X(v)`$ shows that the union of incident codimension-1 edge spaces intersects every deep subset of $`X(v)`$.
Case 2b: The crossing graph of $`X(v)`$ is empty. Each edge incident to $`v`$ therefore has dimension $`N2`$, and it follows that the inclusion of $`S`$ in $`X`$ has the following “coarse Jordan separation property”: for all sufficiently large $`A0`$ there are *exactly* two deep components of $`XN_A(S)`$ which coarsely contain $`S`$. The inclusion of $`X(e)=_ch(S)`$ into $`X`$ therefore has the same property: for all sufficiently large $`A`$ there are exactly two deep components of $`XN_A(X(e))`$ which coarsely contain $`X(e)`$.
Let $`T_e^{}`$ be the subtree of $`T`$ spanned by all edges in $`T`$ whose edge space is coarsely equivalent to $`X(e)`$; we think of $`T_e^{}`$ as the “edge raft” containing $`e`$, although a priori $`T_e^{}`$ can have $`N`$-dimensional vertices and edges. But by using the coarse Jordan separation property for $`X(e)`$ one shows that $`T_e^{}`$ contains at most one $`N`$-dimensional vertex of $`T`$. Moreover, $`T_e^{}`$ cannot contain exactly one $`N`$-dimensional vertex, for then one would be able to find an $`N1`$ dimensional valence 1 vertex of $`T_e^{}`$ which would violate irreducibility of the Bass-Serre tree $`T`$. It follows that $`T_e^{}`$ in fact consists entirely of $`N1`$ dimensional vertices and edges, and so is an $`N1`$ raft. $`T_e^{}`$ cannot be a bounded raft, for it has at least one edge, namely $`e`$, and again that would violate irreducibility. $`T_e^{}`$ cannot be a line raft, by hypothesis. Finally, it cannot be a bushy raft, because that would violate the coarse Jordan separation property for $`X(e)`$: for larger and larger $`A`$, the number of deep components of $`XN_A(X(e))`$ coarsely containing $`X(e)`$ would approach infinity. ∎
We now continue the proof of Theorem 11 by induction down the dimension. Suppose by induction that we have altered $`X`$ and $`T`$ down to dimension $`n+1`$, producing a filtered tree of spaces
so that $`H`$ quasi-acts properly and coboundedly on $`X_0`$, restricting to a true action on $`X_k^{}T_k^{}`$, $`Nkn+1`$, and we have an $`H`$-quasiconjugation back to the original tree of spaces, restricting to the identity on $`X_0X_{n+1}^{}`$, and satisfying the conclusions of Theorem 11 on $`X_{n+1}^{}`$.
The arguments of the basis step can be applied to any component of $`T_n`$ which is an $`n`$-raft.
Consider now an edge $`e`$ of $`T_nT_{n+1}^{}`$ which does not lie on an $`n`$-raft, and an element $`hH`$; we study the image $`hX(e)`$. Using irreducibility of the original tree of spaces it follows that there are two vertices $`v,wT_{n+1}^{}`$ such that $`X(e)`$ is a coarse intersection of $`X^{}(v)`$ and $`X^{}(w)`$, with constants independent of $`e`$; moreover, $`v,w`$ can be chosen to have a distance bounded above in $`T_{n+1}^{}`$ independent of $`e`$. Now we apply the inductive hypothesis to $`hX(v),hX(w)`$, splitting into subcases depending on whether $`v,w`$ lie on rafts.
Case 1: Suppose $`v,w`$ do lie on rafts. By induction, there exist vertices $`v^{},w^{}`$ such $`hX(v)=_cX(v^{})`$, $`hX(w)=_cX(w^{})`$. It follows that $`hX(e)`$ is a coarse intersection of $`X(v^{})`$ and $`X(w^{})`$. Let $`e_1,\mathrm{},e_K`$ be the simple edge path in $`T_0`$ connecting $`v^{}`$ to $`w^{}`$; $`K`$ is bounded independent of $`v,w`$, depending only on the quasi-isometry constants of $`h`$. The coarse intersection of $`X(v^{})`$ and $`X(w^{})`$ equals the coarse intersection of $`X(e_1),\mathrm{},X(e_K)`$. This implies that $`hX(e)`$ is a coarse intersection of $`X(e_1),\mathrm{},X(e_K)`$, and in particular $`hX(e)`$ is coarsely contained in each of $`X(e_1),\mathrm{},X(e_K)`$.
The first consequence of this is that the dimensions of $`e_1,\mathrm{},e_K`$ are greater than or equal to the dimension of $`e`$, because a coarse $`\mathrm{PD}(n)`$ space cannot uniformly embed in a coarse $`\mathrm{PD}`$ space of lower dimension.
The second consequence is that if $`e_k`$ is $`n`$-dimensional then $`hX(e)=_cX(e_k)`$, because of “packing”: a uniform embedding of a coarse $`\mathrm{PD}(n)`$ space in a coarse $`\mathrm{PD}(n)`$ space must have image coarsely equivalent to the whole space.
To complete the proof in Case 1 it therefore remains to check that the edges $`e_1,\mathrm{},e_K`$ cannot all have dimension $`N+1`$. For this we need the fact that $`X(e)`$ coarsely separates $`X(v)`$ and $`X(w)`$ in $`X`$, and it follows that $`hX(e)`$ coarsely separatex $`X(v^{})`$ and $`X(w^{})`$ in $`X`$. But if $`e_1,\mathrm{},e_K`$ all have dimension $`N+1`$ then the coarse $`\mathrm{PD}(n)`$ space $`hX(e)`$ cannot coarsely separate $`X(v^{})`$ and $`X(w^{})`$.
Remaining cases: If, say, $`v`$ lies on a raft and $`w`$ does not, then $`w`$ is incident to an edge $`e_1`$ of the same dimension. Applying induction, $`hX(v)=_cX(v^{})`$ for some raft vertex $`v^{}`$, and $`hX(w)=_chX(e_1)=_cX(e^{})`$ for some edge $`e^{}`$. Now connect $`v^{}`$ and $`e^{}`$ by a simple edge path and repeat the arguments of Case 1. The other cases are similar.
We have shown that the quasi-action of $`H`$ coarsely respects $`n`$-dimensional edge spaces. It remains to attach additional edges to the $`H`$-forest $`T_{n+1}^{}`$ to make an $`H`$-forest $`T_n^{}`$, so that a newly attached edge $`e^{}`$ between vertices $`v^{},w^{}T_{n+1}^{}`$ has an coarse $`\mathrm{PD}(n)`$ edge space $`X^{}(e^{})X_n^{}`$ taken coarsely to some coarse $`\mathrm{PD}(n)`$ edge space $`X(e)X_n`$. The construction of $`T_n^{}`$ is a relative version of the construction in the basis step, which itself is adapted from the proof of the Homogeneous Theorem 1. First one chooses any $`n`$-dimensional edge $`e`$ of $`T_n`$, between vertices $`v^{},w^{}T_{n+1}^{}`$. Then for each $`hH`$ one attaches an edge $`he`$ between $`hv^{}`$ and $`hw^{}`$. The result is not a forest, but it is quasi-isometrically identified with the forest $`T_n`$. Now attach 2-cells in an $`H`$-equivariant manner, with finitely many $`H`$-orbits, so that each of the edges of a given 2-cell have edge spaces in the same coarse equivalent class. Now we have an $`H`$-complex each of whose components is simply connected. Apply tracks to get an $`H`$-forest containing $`T_{n+1}^{}`$ as a subforest. Each new edge (resp. vertex) of this forest has an $`n`$-dimensional edge space (resp. vertex space) mapping back to a $`n`$-dimensional edge space of $`X_n`$.
## Pattern rigidity
Suppose we are in the setting of Theorem 11. Let $`v^{}`$ be a point raft of $`T^{}`$ and $`v`$ the corresponding point raft of $`T`$. Note that as a consequence of Theorem 11, the quasi-isometry $`X^{}(v^{})X(v)`$ takes the edge space pattern inside $`X^{}(v^{})`$ coarsely to the edge space pattern inside $`X(v)`$. This information can be used to strengthen applications of Theorem 11, by applying “pattern rigidity” results.
For example, in the setting of Theorem 4 we have the following theorem of R. Schwartz \[Sch97\]:
###### Theorem 13 (Geodesic Pattern Rigidity in $`𝐇^n`$).
Let $`G`$ be a discrete, cocompact group of isometries of $`𝐇^n`$, $`n3`$. Let $`A`$ be a nonempty, $`G`$-equivariant set of geodesics in $`𝐇^n`$, with finitely many $`G`$-orbits. Let $`H`$ be a group and let $`H\times 𝐇^n\stackrel{\mathit{\varphi }}{}𝐇^n`$ be a cobounded, proper quasi-action which coarsely respects $`A`$: for each $`hH`$ and $`aA`$ there exists $`a^{}A`$ such that $`\varphi (h,a)=_ca^{}`$. Then there is an isometric action $`\psi :H\mathrm{Isom}(𝐇^n)`$ with finite kernel which strictly respects $`A`$, such that $`\psi `$ is a bounded distance from $`\varphi `$, i.e. $`sup_{h,x}d(\varphi (h,x),\psi (h,x))<\mathrm{}`$.
In this setting of this theorem, let $`\mathrm{Isom}(𝐇^n,A)`$ be the set of isometries of $`𝐇^n`$ respecting $`A`$. This is a discrete, cocompact group of isometries, containing $`G`$ and $`H/\mathrm{Ker}(\mathrm{\Psi })`$ as finite index subgroups. It follows that $`G`$ and $`H/\mathrm{Ker}(\psi )`$ are commensurable, by a commensuration taking the $`A`$-stabilizers in $`G`$ to the $`A`$-stabilizers of $`H/\mathrm{Ker}(\psi )`$. Combined with the discussion above, this completes the proof of Theorem 4.
We can get even stronger conclusions under stronger hypotheses. For instance, consider a graph of groups $`\mathrm{\Gamma }`$ as in Theorem 4. Recall that in a discrete, cocompact group of isometries of $`𝐇^n`$, the set of loxodromic axis stabilizers is identical to the set of maximal, virtually cyclic subgroups. Suppose that we make the following additional assumption:
* Each edge-to-vertex injection $`\xi _\eta :\mathrm{\Gamma }_e\mathrm{\Gamma }_{v(\eta )}`$ has image equal to a maximal virtually cyclic subgoup of $`\mathrm{\Gamma }_{v(\eta )}`$, and two distinct edge-to-vertex injections into $`\mathrm{\Gamma }_{v(\eta )}`$ have distinct images.
This implies that for any vertex $`v`$ of the Bass-Serre tree $`T`$, distinct incident edge spaces inside $`X(v)`$ are all coarsely inequivalent in $`X(v)`$. The Bestvina-Feighn Combination Theorem \[BF92\] implies that $`\pi _1\mathrm{\Gamma }`$ is word hyperbolic.
With this additional assumption, Theorem 11 implies that for any group $`H`$ quasi-isometric to $`\pi _1\mathrm{\Gamma }`$, the proper, cobounded quasi-action of $`H`$ on $`X`$ satisfies tree rigidity. Combining this with Geodesic Pattern Rigidity, it follows that $`H`$ is weakly commensurable to $`\pi _1\mathrm{\Gamma }`$.
In fact, what this argument shows, under the additional assumption, is that the group $`\pi _1\mathrm{\Gamma }`$ has finite index in its quasi-isometry group $`\mathrm{QI}(\pi _1\mathrm{\Gamma })`$; this is the strongest form of quasi-isometric rigidity. We also obtain a computation the abstract commensurator of $`\pi _1\mathrm{\Gamma }`$: it is isomorphic to $`\mathrm{QI}(\pi _1\mathrm{\Gamma })`$, which is isomorphic to the full isometry group of the tree of spaces $`X`$.
Finally we turn to abelian pattern rigidity and a strengthening of Theorem 3.
###### Theorem 14 (Abelian Pattern Rigidity).
Suppose that $`V_1,\mathrm{},V_K𝐄^n`$ and $`W_1,\mathrm{},W_K𝐄^n`$ are affine foliations. Suppose that there exists a quasi-isometry $`f:𝐄^n𝐄^n`$ which maps $`W_k`$ coarsely to $`V_k`$, for each $`k=1,\mathrm{},K`$. Then there exists a linear isomorphism $`F:𝐄^n𝐄^n`$ such that $`F(W_k)=V_k`$ for each $`k`$.
###### Proof.
By passing to asymptotic cones we replace $`f`$ with a bilipschitz homeomorphism taking $`V_k`$ to $`W_k`$ for each $`k`$. Applying the Rademacher Theorem, at almost any point $`x`$ the derivative $`F=D_xf`$ gives the desired conclusion. ∎
If $`S(V)`$ is the linear subspace of $`𝐄^n`$ parallel to the leaves of an affine foliation $`V`$, then a pattern of affine foliations $`V_1,\mathrm{},V_K`$ induces a pattern of linear subspaces $`S(V_1),\mathrm{},S(V_K)𝐄^n`$, which in turn induces a pattern of projective subspaces $`𝐏S(V_1),\mathrm{},𝐏S(V_K)𝐏^{n1}`$, called the *projective pattern* associated to $`V_1,\mathrm{},V_K`$. The Abelian Pattern Rigidity Theorem 14 shows that the associated projective pattern is a quasi-isometry invariant of patterns of affine foliations in $`𝐄^n`$. For example, four distinct 1-dimensional affine foliations of $`𝐄^2`$ have an associated projective pattern of four distinct points in $`𝐏^1`$. The moduli space of such projective patterns is 1-dimensional, parameterized by the cross-ratio, and so the cross-ratio is a pattern preserving quasi-isometry invariant.
We can strengthen Theorem 3 as follows. Suppose that $`\mathrm{\Gamma }`$ is a reduced graph of abelian groups as in Theorem 3, with Bass-Serre tree of spaces $`XT`$. Recall that the hypotheses of Theorem 3 imply that each vertex of $`T`$ is a raft. Our strengthened conclusions say that projective patterns associated to edge space patterns inside vertex spaces of $`X`$ are quasi-isometricially rigid. Here are the details.
Let $`H`$ be a finitely generated group quasi-isometric to $`\pi _1\mathrm{\Gamma }`$. The proof of Theorem 3 gives a graph of groups $`\mathrm{\Gamma }^{}`$ with Bass-Serre tree of spaces $`X^{}T^{}`$, such that $`H\pi _1\mathrm{\Gamma }^{}`$, together with a quasi-isometry $`\varphi :X^{}X`$, such that for each edge $`e^{}`$ of $`T^{}`$ there is an edge $`e`$ of $`T`$ with $`\varphi (X(e))=_cX^{}(e^{})`$, and for each vertex $`v^{}`$ of $`T^{}`$, either $`v^{}`$ is a raft of $`T^{}`$ and $`\varphi (X^{}(v^{}))=_cX(v)`$ for some vertex $`v`$ of $`T`$, or $`v^{}`$ is not a raft of $`T^{}`$ and so $`v^{}`$ has an incident edge $`e^{}`$ with $`X^{}(v^{})=_cX^{}(e^{})`$.
Consider a raft vertex $`v^{}`$ of $`T^{}`$ and the corresponding vertex $`v`$ of $`T`$. Composing the map $`X(v^{})\stackrel{\mathit{\varphi }}{}X^{}`$ with the closest point projection to $`X(v)`$ gives a quasi-isometry still denoted $`\varphi :X^{}(v^{})X(v)`$. Moreover, this quasi-isometry takes the edge space pattern inside $`X^{}(v^{})𝐄^n`$ coarsely to the edge space pattern inside $`X(v)𝐄^n`$. Applying the Abelian Pattern Rigidity Theorem 14, the projective patterns in $`𝐏^{n1}`$ associated to $`X^{}(v^{})`$ and $`X(v)`$ are projectively equivalent.
Lee Mosher:
Dept. of Mathematics and Computer Science
Rutgers University, Newark
Newark, NJ 07102
E-mail: mosher@andromeda.rutgers.edu
Michah Sageev:
Dept. of Mathematics
Technion—Israel Institute of Technology
Haifa 32000, Israel
E-mail: sageevm@techunix.technion.ac.il
Kevin Whyte:
Dept. of Mathematics
Univ. of Utah
Salt Lake City, Utah 84112-0090
E-mail: kwhyte@math.utah.edu
|
warning/0005/math0005129.html
|
ar5iv
|
text
|
# The fourth tautological group of ℳ̄_{𝑔,𝑛} and relations with the cohomology
## 1 Introduction
Let $`\overline{}_{g,n}`$ be the moduli space of $`n`$-pointed complex stable algebraic curves of genus $`g`$.
The existence of some degree $`4`$ relations among tautological classes has been proved with various methods by E. Getzler, C. Faber, R. Pandharipande and P.Belorousski, while other relations are obtained as a consequence of the well known ones in degree $`2`$.
We actually prove that no other relations can arise, and that for genus $`g8`$, the cohomology group $`H^4(\overline{}_{g,n},)`$ coincides with its tautological subgroup. The main results of this paper are formally stated in Theorems 10 and 19.
It turns out that new relations appear only in genus up to $`5`$, whereas for higher genus all possible relations arise only as a consequence of degree $`2`$ ones. The proof of this fact allows us to suggest in Conjecture 18 an upper bound depending on the genus for higher degree new tautological relations.
As for the methods, E. Arbarello and M. Cornalba proposed in \[AC1\] new methods for computing the cohomology groups with rational coefficients of $`\overline{}_{g,n}`$; their strategy is to establish a strict relation between the cohomology of the moduli space and the one of the irreducible components of the boundary, which in turn can be expressed in terms of moduli spaces of curves with lower genus or with lower number of marked points. With similar arguments, we establish inductive procedures on genus and/or number of markings to derive constraints among coefficients in possible relations.
We will therefore be able to give the explicit expression of a new relation in $`H^4\left(\overline{}_{3,2}\right)`$, whose existence was proved by Faber as a consequence of the existence of a tautological relation on the open part $`_{3,2}`$. Furthermore, we will exclude the existence of any relation other than the known ones.
A description of $`H^4(\overline{}_g,)`$, for $`g12`$, has been given by D. Edidin in \[Ed\], and once the tautological group is known, we can adapt his argument to prove that for $`g8`$, it coincides with the cohomology. For this, we make use of the results by Harer (\[Ha\]), Ivanov (\[Iv\]) and Loojenga (\[Lo\]) on the homology of the mapping class group.
This paper is extracted from my Tesi di Perfezionamento at the Scuola Normale Superiore, Pisa. In the present exposition, many of the calculations will be omitted. The interested reader can find them all in the thesis (\[Po\]), available upon request from the author.
I wish to thank my advisor, Enrico Arbarello, as well as Gilberto Bini, Maurizio Cornalba, Carel Faber and Rahul Pandharipande for many extremely useful conversations.
## 2 Stable graphs and tautological classes
To every stable curve $`C`$ of genus $`g`$, with $`P`$ as a set of markings, one can associate a labelled graph $`\mathrm{\Gamma }`$ in the following way:
1. draw a vertex $`v`$ for every irreducible component $`C\left(v\right)`$ of the normalization $`\stackrel{~}{C}`$ of $`C`$, and label it with the genus $`g\left(v\right)`$ of that component,
2. draw an edge between two vertices $`v_1`$, $`v_2`$ (possibly a loop if $`v_1=v_2`$) whenever the normalization map $`\nu :\stackrel{~}{C}`$ $`C`$ identifies two points lying respectively in $`C\left(v_1\right)`$ and $`C\left(v_2\right)`$,
3. draw a half-edge with vertex $`v`$ whenever there is a marking in $`\nu \left(C\left(v\right)\right)`$, and label it with the marking’s name. We denote by $`P(v)`$ the set of these markings.
We call marked half-edges the half-edges constructed in $`3`$. The total set of half-edges is the union of the set of marked half-edges with the set consisting of the halves of the edges constructed in $`2`$.
Let $`r\left(v\right)`$ be the valence of a vertex, namely the number of half-edges with vertex $`v`$. The stability condition translates to: $`2g\left(v\right)+r\left(v\right)3`$, for every vertex $`v`$. The genus of a curve corresponding to the graph $`\mathrm{\Gamma }`$ is $`g\left(\mathrm{\Gamma }\right)=\chi \left(\mathrm{\Gamma }\right)+_vg\left(v\right)`$. Observe that the construction of the graph is only based on the topological type of the curve.
###### Definition 1
A $`P`$-marked stable graph of genus $`g`$ (briefly a $`(g,P)`$ graph), is a connected graph with $`n=|P|`$ marked half-edges, with the following additional data:
1) each vertex $`v`$ is labelled with an integer $`g(v)`$,
2) the valence $`r(v)`$ of any vertex satisfies the stability condition $`2g(v)+r(v)3`$,
3) there is a bijection between marked half-edges and elements in $`P`$,
4) $`g=\chi \left(\mathrm{\Gamma }\right)+_vg\left(v\right)`$.
The codimension of a graph is defined as the number of its edges.
Given a $`P`$-marked stable graph of genus $`g`$ and codimension $`d`$, with set of vertices $`V`$, one can associate to it a closed stratum of codimension $`d`$ in $`\overline{}_{g,P}`$. For every vertex $`vV`$, we let $`S\left(v\right)`$, denote the set of unmarked half-edges with vertex $`v`$.
Let $`\overline{}_\mathrm{\Gamma }:=_{vV}_{g\left(v\right),P\left(v\right)S\left(v\right)}`$
The map
$$\xi _\mathrm{\Gamma }:\overline{}_\mathrm{\Gamma }\overline{}_{g,P}$$
is called a boundary map, and has the closed stratum $`\mathrm{\Delta }_\mathrm{\Gamma }=\xi _\mathrm{\Gamma }\left(\overline{}_\mathrm{\Gamma }\right)`$ as image.
The notation $`\overline{}_\mathrm{\Gamma }`$ will be used also when $`\mathrm{\Gamma }`$ is disconnected: if $`\mathrm{\Gamma }=\mathrm{\Gamma }_1\mathrm{\Gamma }_2`$, then $`\overline{}_\mathrm{\Gamma }=\overline{}_{\mathrm{\Gamma }_1}\times \overline{}_{\mathrm{\Gamma }_2}`$.
Let $`\mathrm{\Gamma }`$ be a $`(g,P)`$-graph.
###### Definition 2
The graph $`G`$ is a $`\mathrm{\Gamma }`$-graph if it is the disjoint union of a collection of $`(g(v),P(v)S(v))`$-graphs.
Look at a $`\mathrm{\Gamma }`$-graph $`G`$. Set $`G=_{vV}G_v`$. We can define the map
$$\overline{}_G=\overline{}_{G_v}\stackrel{\zeta _G}{}\overline{}_\mathrm{\Gamma }$$
as $`\zeta _G=\{\xi _{G_v}\}_{vV}`$.
Let $`\pi `$ be the forgetful map:
$`\pi `$ $`:`$ $`\overline{}_{g,n+1}\overline{}_{g,n}`$
$`[C,p_1,\mathrm{},p_n,p_{n+1}]`$ $``$ $`[C,p_1,\mathrm{},p_n]`$
We will also refer to the map $`\pi `$ as the universal curve, or the projection map.
Let $`\sigma _1,\mathrm{},\sigma _n`$ be the $`n`$ canonical sections of the forgetful map, and let $`D_i`$ be the image of $`\sigma _i`$. Finally, let $`\omega _\pi `$ be the relative dualizing sheaf of $`\pi `$.
We recall the definition of the basic cohomology classes in $`\overline{}_{g,P}`$ (see \[AC2\]):
###### Definition 3
$`\psi _i`$ $`=`$ $`\sigma _i^{}\left(c_1\left(\omega _\pi \right)\right),i=1,\mathrm{},n`$
$`\kappa _a`$ $`=`$ $`\pi _{}\left(\left(c_1\left(\omega _\pi \left({\displaystyle D_j}\right)\right)\right)^{a+1}\right),a=0,\mathrm{}3g3+n`$
The class $`\psi _i`$ can be interpreted as the first Chern class of the orbifold bundle whose fiber over the point $`[C,p_1,\mathrm{},p_n]`$ is the cotangent bundle to the curve $`C`$ evaluated at the point $`p_i`$.
###### Definition 4
A Mumford class in $`H^{}(\overline{}_{g,P},)`$ is a polynomial in the classes $`\psi _i,\kappa _a`$. The Mumford ring is
$$[\psi _1,\mathrm{},\psi _n,\kappa _1,\mathrm{},\kappa _{3g3+n}].$$
The Mumford ring on a product or a disjoint union of moduli spaces is the tensor product or the direct sum of the Mumford rings.
It is worth noticing that the following formula (see Formula 1.7 in \[AC2\]) holds:
$$\kappa _a=\pi _{}(\psi _{n+1}^{a+1}).$$
###### Definition 5
A Mumford class in $`H^{}(_{g,P},)`$ is the pull-back under the inclusion
$$_{g,P}\overline{}_{g,P}$$
of a polynomial in the classes $`\psi _i,\kappa _a`$.
###### Definition 6
A tautological class is the push-forward of a Mumford class via a boundary map. The $`k`$-th tautological group $`T_{g,P}^k`$ is the subspace of $`H^k(\overline{}_{g,P},)`$ generated by these classes.
In Figures 1 and 2 we draw all the graphs of codimension $`1`$ and codimension $`2`$ which we need in our study of $`T_{g,P}^4`$. In each figure we will also write the name of the corresponding graph. Every time half-edges are drawn, one should imagine them labelled with the correspondent markings.
If $`p`$ is a Mumford class, we use the following notation:
$$p|\delta _\mathrm{\Gamma }:=\frac{\xi _\mathrm{\Gamma }\left(p\right)}{|Aut\mathrm{\Gamma }|}.$$
We will often write $`\delta _{irr},\xi _{irr}`$ instead of $`\delta _{\mathrm{\Gamma }_{irr}},\xi _{\mathrm{\Gamma }_{irr}}`$, and $`\delta _{a,A},\xi _{a,A}`$ instead of $`\delta _{\mathrm{\Gamma }_{a,A}},\xi _{\mathrm{\Gamma }_{a,A}}`$.
Degree $`4`$ autological classes are:
1. Pure boundary classes: let $`\mathrm{\Gamma }`$ be a graph of codimension $`2`$, then we define:
$$\delta _\mathrm{\Gamma }:=\frac{\xi _\mathrm{\Gamma }\left(1\right)}{|Aut\mathrm{\Gamma }|}$$
2. Mixed boundary classes: if codim $`\mathrm{\Gamma }=1`$, and $`p`$ is a Mumford class of degree $`2`$ in $`\overline{}_\mathrm{\Gamma }`$, then
$$p|\delta _\mathrm{\Gamma }:=\frac{\xi _\mathrm{\Gamma }(p)}{|Aut\mathrm{\Gamma }|}.$$
We will often use the following simplified notation:
* $`\psi _i\delta _{a,A}=\left(\psi _i1\right)|\delta _{a,A}=\frac{1}{Aut\mathrm{\Gamma }_{a,A}}\xi _{a,A}(\psi _i1)`$,
* $`\psi |\delta _{a,A}=\left(\psi _s1\right)|\delta _{a,A}=\frac{1}{Aut\mathrm{\Gamma }_{a,A}}\xi _{a,A}(\psi _s1)`$,
* $`\delta _{a,A}|\psi =\left(1\psi _t\right)|\delta _{a,A}=\frac{1}{Aut\mathrm{\Gamma }_{a,A}}\xi _{a,A}(1\psi _t)=\psi |\delta _{ga,A^c}`$,
* $`\kappa |\delta _{a,A}=\left(\kappa _11\right)|\delta _{a,A}=\frac{1}{Aut\mathrm{\Gamma }_{a,A}}\xi _{a,A}(\kappa _11)`$,
* $`\delta _{a,A}|\kappa =\left(1\kappa _1\right)|\delta _{a,A}=\frac{1}{Aut\mathrm{\Gamma }_{a,A}}\xi _{a,A}(1\kappa _1)=\kappa |\delta _{g,A^c}`$,
* $`\psi _i\delta _{irr}=\left(\psi _i\right)|\delta _{irr}=\frac{1}{Aut\mathrm{\Gamma }_{irr}}\xi _{irr}(\psi _i)`$,
* $`\psi |\delta _{irr}=\left(\psi _q+\psi _r\right)|\delta _{irr}=\frac{1}{Aut\mathrm{\Gamma }_{irr}}\xi _{irr}(\psi _q+\psi _r)`$,
* $`\kappa _1\delta _{irr}=\kappa _1|\delta _{irr}=\frac{1}{Aut\mathrm{\Gamma }_{irr}}\xi _{irr}(\kappa _1)`$.
3. Mumford classes : these are simply monomials in Mumford classes (considered as push-forward via the map corresponding to the trivial graph).
In the mixed boundary classes we intentionally used ambiguous notation. Some of the classes ($`\psi _i\delta _{a,A},\psi _i\delta _{irr},\kappa _1\delta _{irr}`$) turn out to be written as a product of a codimension $`1`$ boundary class with a Mumford class. In the proof of the next Proposition we will show that the above notation is unambiguous.
###### Proposition 7
The image of the map:
$`H^2\left(\overline{}_{g,P}\right)\times H^2\left(\overline{}_{g,P}\right)H^4\left(\overline{}_{g,P}\right)`$
$`(\alpha ,\beta )\alpha \beta `$
lies in $`T_{g,P}^4`$.
Proof. Recall that $`H^2\left(\overline{}_{g,P}\right)=T_{g,P}^2`$. Two irreducible codimension $`1`$ boundary classes either coincide or intersect transversally. In the latter case, it is trivial to check that their intersection is a linear combination of tautological pure boundary classes. The product of two Mumford classes is clearly a Mumford class.
Finally, using the push-pull formula, one is able to express the product of a Mumford class and a boundary class, and the square of a boundary class, as linear combination of tautological classes:
$`\psi _i\delta _{a,A}=\psi _i|\delta _{a,A}`$ $`\psi _i\delta _{irr}=\psi _i|\delta _{irr}`$
$`\kappa _1\delta _{a,A}=\kappa _1|\delta _{a,A}+\delta _{a,A}|\kappa _1`$ $`\kappa _1\delta _{irr}=\kappa _1|\delta _{irr}`$
$`\delta _{a,A}^2=\psi |\delta _{a,A}\delta _{a,A}|\psi `$ $`+\{\begin{array}{cc}\frac{2}{|Aut\mathrm{\Gamma }_{a,A}|}\delta _{G(ga,\mathrm{},2ag,P)}\hfill & ifA=P\hfill \\ \frac{2}{|Aut\mathrm{\Gamma }_{a,A}|}\delta _{G(a,\mathrm{},g2a,P)}\hfill & ifA=\mathrm{}\hfill \end{array}`$
$`\delta _{irr}^2`$ $`={\displaystyle \frac{1}{2}}\xi _{irr}(\psi _q+\psi _r)+2\delta _F+2{\displaystyle \delta _{E(a,A)}}`$
$`\psi |\delta _{irr}+2\delta _F+2{\displaystyle \delta _{E(a,A)}}`$
We compute explicitely one sample case. Since
$`\xi _{irr}^{}(\delta _{irr})`$ $`=`$ $`\delta _{irr}+{\displaystyle \delta _{a,A\left\{q\right\}}}\psi _q\psi _r,`$
then
$`2\delta _{irr}^2=\xi _{irr}\xi _{irr}^{}(\delta _{irr})`$ $`=`$ $`{\displaystyle \frac{1}{2}}\xi _{irr}\stackrel{~}{\xi }_{irr}(1)+{\displaystyle \xi _{irr}\stackrel{~}{\xi }_{a,A\left\{q\right\}}(1)}\xi _{irr}(\psi _q+\psi _r),`$
where the symbol $`\stackrel{~}{\xi }`$ is used for boundary maps of $`\overline{}_{g1,P\{q,r\}}`$. In fact, from now on, when composing two boundary maps, we will append the second one with the twiddle.
We easily compute: $`\frac{1}{2}\xi _{irr}\stackrel{~}{\xi }_{irr}(1)=\frac{1}{2}\xi _F(1)=4\delta _F`$, and then observe that $`\xi _{irr}\stackrel{~}{\xi }_{a,A\left\{q\right\}}=\xi _{E(a,A)}`$ and that the corresponding graph has automorphism order $`2`$, unless $`P=\mathrm{},a=g/2`$, when the order is $`4`$. Moreover, $`\xi _{irr}\stackrel{~}{\xi }_{a,A\left\{q\right\}}(1)=\xi _{irr}\stackrel{~}{\xi }_{ga,A^C\left\{q\right\}}(1)=|Aut\mathrm{\Gamma }_{E(a,A)}|\delta _{E(a,A)}`$. Whenever $`|Aut\mathrm{\Gamma }_{E(a,A)}|=4`$, then by symmetry only one of the summands above does appear, hence we can write
$`\delta _{irr}^2`$ $`={\displaystyle \frac{1}{2}}\xi _{irr}(\psi _q+\psi _r)+2\delta _F+2{\displaystyle \delta _{E(a,A)}}`$
$`\psi |\delta _{irr}+2\delta _F+2{\displaystyle \delta _{E(a,A)}}.`$
$`\mathrm{}`$
## 3 Essential tautological classes
It is well known that, for genus up to $`2`$, there are some relations between degree $`2`$ tautological classes; thus, certain tautological classes could be expressed as linear combination of other ones; they are: $`\kappa _1`$ and $`\psi _i,iP`$ for genera $`g=0,1`$, $`\kappa _1`$ for genus $`g=2`$.
Moreover, there are Keel’s relations among boundary classes in genus $`0`$.
All these relations reproduce themselves in every genus. The reason is quite clear: every time there is a relation among tautological classes in the second cohomology group of a codimension $`1`$ boundary component, we can push it forward to $`H^4\left(\overline{}_{g,P}\right)`$.
In this section we will choose a set of degree $`4`$ tautological classes which generate $`T_{g,P}^4`$, by eliminating the above relations. We will call these classes the essential tautological classes. The set of essential tautological classes will be denoted by $`_{g,P}^4`$ and it is obtained from the set of all tautological classes by removing the unessential classes which we are presently going to list.
The unessential tautological classes are:
| $`\psi |\delta _{0,A}=\xi _{0,A}(\psi _s1)`$ | $`\psi _i\delta _{0,A}=\xi _{0,A}(\psi _i1)`$ | $`\kappa |\delta _{0,A}=\xi _{0,A}(\kappa _11)`$ | for any $`g`$, |
| --- | --- | --- | --- |
| $`\psi |\delta _{1,A}=\xi _{1,A}(\frac{\psi _s1}{|Aut\mathrm{\Gamma }_{1,A}|})`$ | $`\psi _i\delta _{1,A}=\xi _{1,A}(\psi _i1)`$ | $`\kappa |\delta _{1,A}=\xi _{1,A}(\frac{\kappa _11}{|Aut\mathrm{\Gamma }_{1,A}|})`$ | for any $`g`$, |
| $`\kappa |\delta _{2,A}=\xi _{2,A}(\frac{\kappa _11}{|Aut\mathrm{\Gamma }_{2,A}|})`$ | | | for any $`g`$, |
| $`\psi |\delta _{irr}=\xi _{irr}(\frac{\psi _q+\psi _r}{2})`$ | | | for $`g=1,2`$ |
| $`\psi _i\delta _{irr}=\xi _{irr}(\frac{\psi _i}{2})`$ | | | for $`g=1,2`$, |
| $`\kappa |\delta _{irr}=\xi _{irr}(\frac{1\kappa _1}{2})`$ | | | for $`g=1,2,3`$, |
| $`\psi _i^2,\psi _i\psi _j`$ | $`\kappa _1^2,\kappa _1\psi _i`$ | | for $`g=0,1`$, |
| | $`\kappa _1^2,\kappa _1\psi _i`$ | | for $`g=2`$. |
Moreover, some classes $`\delta _{G(0,A,0,B)}`$ are unessential (see below); in fact, in genus $`0`$ there are Keel’s relations (\[Ke\]) among boundary classes: we can push them forward by means of the maps
$$H^2\left(\overline{}_{0,A\left\{s\right\}}\right)\stackrel{\varphi _{0,A}}{}H^4\left(\overline{}_{g,P}\right)$$
to obtain the following relations :
$$\underset{\begin{array}{c}x,yB,\\ z,wC,\\ BC=A\end{array}}{}\delta _{G(0,B,0,C)}+\delta _{G(0,C,0,B)}=\underset{\begin{array}{c}x,zB,\\ y,wC,\\ BC=A\end{array}}{}\delta _{G(0,B,0,C)}+\delta _{G(0,C,0,B)},$$
$$\underset{\begin{array}{c}x,yB,\\ zC,\\ BC=A\end{array}}{}\delta _{G(0,B,0,C)}=\underset{\begin{array}{c}x,zB,\\ yC,\\ BC=A\end{array}}{}\delta _{G(0,B,0,C)}.$$
We now describe a subset of essential classes of this type; if we fix an ordering in $`P,`$ this induces an ordering of every subset $`A`$; a basis for $`H^2\left(\overline{}_{0,A\left\{s\right\}}\right)`$ consists of classes $`\delta _{0,\left\{s\right\}C}`$, with $`B=A\backslash C`$, $`|B|3`$, or $`\left|B\right|=2`$ and $`b<c`$ $`bB,cC`$. This implies that we are going to consider only classes $`\delta _{G(0,B,0,C)}`$, with $`|B|3`$, or $`\left|B\right|=2`$ and $`b<c`$ $`bB`$, $`cC`$ .
## 4 Pull-back formulas
In this section we show how to pull back tautological classes to the codimension $`1`$ boundary components and to the universal curve. Let $`A`$ be a stable $`(g,P)`$-graph of codimension $`1`$, as defined in the introduction, and let $`\mathrm{\Gamma }`$ be a stable connected $`(g,P)`$-graph of codimension $`2`$.
We fix our attention on a class of the form $`p|\delta _\mathrm{\Gamma }=\frac{1}{|Aut\mathrm{\Gamma }|}\xi _\mathrm{\Gamma }(p).`$ We want to describe the boundary components of $`\overline{}_A`$ on which the pull-back $`\xi _A^{}(p|\delta _\mathrm{\Gamma })`$ is supported.
Given any stable $`A`$-graph $`G`$, let $`j_{s,t}\left(G\right)`$ be the graph obtained by gluing the half edges $`s`$ and $`t`$, and let $`f_{s,t}\left(G\right)`$ be the graph obtained from $`j_{s,t}\left(G\right)`$ by collapsing the new edge. Via the operation $`j_{s,t}`$ we are either creating a node on an irreducible component, or joining two irreducible components at a point. In either case we are creating a node. Via the operation $`f_{s,t}`$ we are smoothing the new node.
We claim that the boundary components we are looking for correspond to $`A`$-graphs $`G`$ such that $`j_{s,t}\left(G\right)=\mathrm{\Gamma }`$ or $`f_{s,t}\left(G\right)=\mathrm{\Gamma }`$. It is very simple to produce graphs $`G`$ of this sort.
Either $`\mathrm{\Delta }_\mathrm{\Gamma }\mathrm{\Delta }_A`$, or $`\mathrm{\Delta }_\mathrm{\Gamma }`$ and $`\mathrm{\Delta }_A`$ intersect transversally. If $`\mathrm{\Delta }_\mathrm{\Gamma }`$ and $`\mathrm{\Delta }_A`$ intersect transversally there must be at least a vertex $`v`$ of $`\mathrm{\Gamma }`$ and a simple Feynman move based at $`v`$ making $`\mathrm{\Gamma }`$ a degeneration of $`A`$. Cutting into a half the edge produced by the Feynman move, and calling the two new half edges $`s`$ and $`t`$, creates a stable $`A`$-graph $`G`$ having the property that $`f_{s,t}\left(G\right)=\mathrm{\Gamma }`$.
Suppose, on the other hand, that $`\mathrm{\Delta }_\mathrm{\Gamma }`$ is contained in $`\mathrm{\Delta }_A`$. This simply means that there is at least one edge of $`\mathrm{\Gamma }`$ cutting which produces two half edges $`s`$ and $`t`$ and a stable $`A`$-graph $`G`$ with the property that $`j_{s,t}\left(G\right)=\mathrm{\Gamma }`$.
Furthermore we can say that $`\mathrm{\Delta }_\mathrm{\Gamma }\mathrm{\Delta }_A`$ if and only if there exist a graph $`G`$ such that $`j_{s,t}\left(G\right)=\mathrm{\Gamma }`$.
In conclusion, whatever the position of $`\mathrm{\Delta }_\mathrm{\Gamma }`$ is with respect to $`\mathrm{\Delta }_A`$, we can build a diagram:
for any graph $`G`$ such that $`j_{s,t}\left(G\right)=\mathrm{\Gamma }`$ or $`f_{s,t}\left(G\right)=\mathrm{\Gamma }`$. The maps $`\xi _A`$ and $`\xi _\mathrm{\Gamma }`$ are boundary maps, the map $`\zeta _G`$ has been defined in section 2, and the map $`\eta _G`$ consists in joining the two half-edges $`s`$ and $`t`$ of the graph $`G`$.
Observe that some of these maps could be the identity: e.g if $`\mathrm{\Gamma }=A=\mathrm{\Gamma }_{irr}`$, then the trivial $`A`$ \- graph $`G`$ satisfies $`j_{s,t}(G)=\mathrm{\Gamma }`$, and the map $`\zeta _G`$ is the identity.
###### Proposition 8
Let $`\mathrm{\Gamma }`$ be any stable graph, of codimension $`2`$. Let $`A`$ be any graph of codimension $`1`$. Then the following formula holds:
$$\frac{\xi _A^{}(\xi _\mathrm{\Gamma }(p))}{Aut\mathrm{\Gamma }}=\underset{f_{s,t}(G)=\mathrm{\Gamma }}{}\frac{\zeta _G(\eta _G^{}(p))}{AutG}+\underset{j_{s,t}(G)=\mathrm{\Gamma }}{}\frac{\zeta _G(\eta _G^{}(p))}{AutG}c_1(N_{\xi _A}),$$
where we denote by $`N_{\xi _A}`$ the normal bundle to the map $`\xi _A`$.
As usual, we will adopt the simplified notation:
$$\xi _A^{}(p|\delta _\mathrm{\Gamma })=\underset{f_{s,t}(G)=\mathrm{\Gamma }}{}(\eta _G^{}(p))|\delta _G+\underset{j_{s,t}(G)=\mathrm{\Gamma }}{}(\eta _G^{}(p))|\delta _Gc_1(N_{\xi _A}).$$
Proof. As we already explained, the two cycles $`\mathrm{\Delta }_\mathrm{\Gamma }`$ and $`\mathrm{\Delta }_A`$ do not intersect transversally in $`\overline{}_{g,P}`$ if and only if there exist a graph $`G`$ such that $`j_{s,t}\left(G\right)=\mathrm{\Gamma }`$. In this case, we consider a tubular neighborhood $`T`$ of the divisor with normal crossing $`\mathrm{\Delta }_A\overline{}_{g,P}`$.
Consider the diagram:
and the normal bundle $`N_{f_A}`$ to the map $`f_A`$. Also observe that $`g_A^{}N_{f_A}=N_{\xi _A}`$.
Introduce a metric in $`N_{f_A}`$, construct a tubular neighborhood $`\stackrel{~}{T}`$ of its zero section, and extend $`f_A`$ in the obvious way to a $`C^{\mathrm{}}`$ map
$$\stackrel{~}{f}_A:\stackrel{~}{T}T.$$
Take then a sufficiently generic $`C^{\mathrm{}}`$ section $`s`$ of $`N_{f_A}`$ lying in $`\stackrel{~}{T}`$. The composition $`\stackrel{~}{f}_Asg_A`$ yields a $`C^{\mathrm{}}`$ map
$$s_A:\overline{}_A\overline{}_{g,P}$$
homotopic to $`\xi _A`$.
As Poincarè duality holds for smooth compact orbifolds, we may pull back cycles from $`\overline{}_{g,P}`$ to $`\overline{}_A`$. If $`\mathrm{\Delta }`$ is any irreducible boundary component, then because of our generic choice of the sections, we have, by transverse intersection,
$$s_A^{}([\mathrm{\Delta }])=\underset{i}{}\left[\mathrm{\Delta }_i\right]$$
(2)
where the sum ranges over the irreducible components $`\mathrm{\Delta }_i`$ of the preimage of $`\mathrm{\Delta }`$ in $`\overline{}_A`$.
The first step is to describe the irreducible components $`\mathrm{\Delta }_i`$. We claim that they are of two types, which can combinatorially described as follows. The first one is simply a cycle $`\mathrm{\Delta }_G\overline{}_A`$ for each graph $`G`$ such that $`f_{s,t}\left(G\right)=\mathrm{\Gamma }`$. If $`\mathrm{\Delta }_A`$ and $`\mathrm{\Delta }_\mathrm{\Gamma }`$ intersect transversally, these are the only components $`\mathrm{\Delta }_i`$ appearing in the above expression. If not, the remaining $`\mathrm{\Delta }_i`$’s are all of the form
$$\frac{\xi _G\xi _G^{}(c_1(N_{\xi _A}))}{AutG},$$
where $`G`$ is a graph such that $`j_{s,t}\left(G\right)=\mathrm{\Gamma }`$.
Once this is established, we get the Proposition for the case $`p=1`$, that is:
$$\xi _A^{}(\delta _\mathrm{\Gamma })=\underset{f_{s,t}(G)=\mathrm{\Gamma }}{}\delta _G+\underset{j_{s,t}(G)=\mathrm{\Gamma }}{}\delta _Gc_1(N_{\xi _A}).$$
Instead of proving our assertion about the $`\mathrm{\Delta }_i`$’s in general, we shall restrict ourselves to some typical examples. The first example is $`\mathrm{\Gamma }=A=\mathrm{\Gamma }_{b,B}`$, with $`B\mathrm{}`$, $`B^c\mathrm{}`$. There is only one $`\mathrm{\Delta }_i`$, which is the zero locus of a section of the normal bundle to the map $`\xi _A`$. One may notice that $`\mathrm{\Delta }_i`$ corresponds to the trivial $`A`$-graph $`G`$, drawn on the right, and that one has that
$$\xi _{b,B}^{}(\delta _{b,B})=(\eta _G^{}(1))|\delta _Gc_1(N_{\xi _A})=c_1(N_{\xi _A}).$$
This is the standard situation of excess intersection, and there is no surprise in finding this term in the general formula of Proposition 8 we are discussing.
The opposite situation occurs for example in the formula for
$$\xi _{irr}^{}(\delta _{b,B})=\xi _{irr}^{}(1|\delta _{b,B})$$
where we further assume that $`b1,gb1`$. There are two components $`\mathrm{\Delta }_i`$, corresponding to the $`A`$-graphs $`G_1`$ and $`G_2`$ having the property that $`f_{s,t}\left(G_i\right)=\mathrm{\Gamma }_{b,B}`$. In this case
$$\xi _{irr}^{}(\delta _{b,B})=\delta _{G_1}+\delta _{G_2}.$$
This is the standard situation of transverse intersection.
What is somewhat unexpected in the formula we are discussing, is the mixture between terms related to excess intersection and terms related to transverse intersection. To illustrate this phenomenon, let us consider the case
$$\xi _{irr}^{}(\delta _F).$$
The formula in the statement tells us that
$$\xi _{irr}^{}\left(\delta _F\right)=\left(\psi _q+\psi _r\right)\delta _{irr}+\delta _F+\delta _{E(a,A\left\{q\right\})}+\left(\delta _{H(a,A\left\{q\right\})}+\delta _{H(a,A\left\{r\right\})}\right),$$
where the two sums range over all the possible graphs of the corresponding type.
The first term is clear: it comes from excess intersection, and corresponds to the only graph $`G`$ such that $`j_{s,t}(G)=F`$, i.e. the graph with one vertex of genus $`g2`$, one loop, and half-edges with labels in $`P\{s,t\}`$.
As a sample case, let us explain the presence of the term $`\delta _F`$. The presence of the other terms can be justified by similar arguments. Draw a picture of $`\mathrm{\Delta }_{irr}`$ in a neighborhood of a generic point of the cycle $`\mathrm{\Delta }^{}`$ corresponding to the locus of irreducible curves with at least three nodes (Figure 6). We cut it with a codimension three generic subspace, in order to draw the picture. The cycle $`\mathrm{\Delta }^{}`$ is drawn as a triple point of $`\mathrm{\Delta }_{irr}`$, which is locally the union of three planes, intersecting each other in the three lines belonging to $`\mathrm{\Delta }_F`$.
Now we “move ” a little bit $`\mathrm{\Delta }_{irr}`$ (Figure 7), we call it $`\stackrel{~}{\mathrm{\Delta }}_{irr}`$, and draw it with a dotted line. There are three points of transverse intersection between $`\stackrel{~}{\mathrm{\Delta }}_{irr}`$ and $`\mathrm{\Delta }_F`$. This shows that $`s_A^{}(\delta _F)`$ contains, with multiplicity $`1`$, the codimension $`2`$ cocycle in $`\overline{}_{g1,P\{s,t\}}`$ corresponding to the locus of irreducible two-noded curves, which by abuse of notation is again denoted by $`\mathrm{\Delta }_F`$.
The formula in the statement, in the case $`p=1`$,
$$\xi _A^{}(\delta _\mathrm{\Gamma })=\underset{f_{s,t}(G)=\mathrm{\Gamma }}{}\delta _G+\underset{j_{s,t}(G)=\mathrm{\Gamma }}{}\delta _Gc_1(N_{\xi _A})$$
is now completely justified.
To prove the general formula we make the following preliminary remark; we seek a formula for the pull-back under a $`\xi _A`$ map of one of the following classes:
* pure boundary classes, hence orbifold Poincaré duals of cycles;
* $`\psi `$-mixed classes, hence orbifold Chern classes of bundles supported on cycles;
* $`\kappa `$-mixed classes. These are linear combinations of the above two types. In fact, we recall Mumford theorem
$$\kappa _1=12\lambda _1+\psi _i\delta _G,$$
where the second sum ranges over the set of stable graphs of codimension $`1`$ , and $`\lambda _1`$ is the first Chern class of the Hodge bundle; this implies that $`\kappa _1`$ is a linear combination of Poincaré duals of cycles and of Chern classes of bundles;
* pure Mumford classes, hence polynomials in classes of the above types.
In order to pull-back a tautological class, we first decompose it into a linear combination of Mumford classes supported on cycles, and then pull back each summand separately.
We therefore seek a formula for
$$\xi _A^{}(\frac{\xi _\mathrm{\Gamma }(c_1(F))}{Aut\mathrm{\Gamma }})$$
where $`F`$ is a line bundle on $`\overline{}_\mathrm{\Gamma }`$.
Suppose first that $`\mathrm{\Delta }_\mathrm{\Gamma }`$ and $`\mathrm{\Delta }_A`$ intersect transversally. Take a sufficiently generic $`C^{\mathrm{}}`$ section $`\sigma _F`$ of the line bundle $`F`$. For every graph $`G`$ such that $`f_{s,t}(G)=\mathrm{\Gamma }`$, we denote by $`F_G`$ the bundle $`\eta _G^{}(F)`$, and by $`\sigma _{F_G}`$ its section $`\eta _G^{}(\sigma _F)`$.
By Poincaré duality, we can pull back cycles. We claim that
$$\xi _A^{}(\frac{\xi _\mathrm{\Gamma }([\{\sigma _F=0\}])}{Aut\mathrm{\Gamma }})=\underset{f_{s,t}(G)=\mathrm{\Gamma }}{}\frac{\xi _G^{}([\{\sigma _{F_G}=0\}])}{AutG}.$$
Let $`\mathrm{\Delta }`$ be a cycle in in $`\overline{}_{g,P}`$ such that
$$[\mathrm{\Delta }]=\frac{\xi _\mathrm{\Gamma }([\{\sigma _F=0\}])}{Aut\mathrm{\Gamma }};$$
we can pick
$$\mathrm{\Delta }=\{x\overline{}_{g,P}x=\xi _\mathrm{\Gamma }(y),\sigma _F(y)=0\}$$
with orbifold multiplicity $`1`$. Because of transverse intersection of $`\mathrm{\Delta }_\mathrm{\Gamma }`$ and $`\mathrm{\Delta }_A`$, Formula 2 applies in this case too. $`\mathrm{\Delta }`$ is a cycle contained in $`\mathrm{\Delta }_\mathrm{\Gamma }`$. We therefore seek the irreducible components $`\mathrm{\Delta }_i`$ inside the irreducible components of the preimage of $`\mathrm{\Delta }_\mathrm{\Gamma }`$ in $`\overline{}_A`$, that is, inside the $`\mathrm{\Delta }_G`$’s, where $`f_{s,t}(G)=\mathrm{\Gamma }`$. One can easily check that
$`\mathrm{\Delta }_G\xi _A^1(\mathrm{\Delta })`$ $`=`$ $`\{z\overline{}_A\mathrm{\Delta }_G\xi _A(z)=\xi _\mathrm{\Gamma }(y)\text{ for some }y\text{ such that }\sigma _F(y)=0\}`$
$`=`$ $`\{z\overline{}_Az=\zeta _G(w)\text{ for some }w\text{}\xi _A(z)=\xi _\mathrm{\Gamma }(y)\text{ for some }y\text{ such that }\sigma _F(y)=0\}`$
$`=`$ $`\{z\overline{}_Az=\zeta _G(w)\text{ for some }w\text{ such that }\sigma _{F_G}(w)=0\},`$
again with orbifold multiplicity $`1`$.
Suppose, on the other hand, that $`\mathrm{\Delta }_\mathrm{\Gamma }\mathrm{\Delta }_A`$. We need formulas for degree $`4`$ classes, hence the only new and significant situation occurs when $`\mathrm{\Delta }_\mathrm{\Gamma }=\mathrm{\Delta }_A`$, and $`\mathrm{\Gamma }=A`$ is a graph of codimension $`2`$.
From the construction of the map $`s_A`$, we see that the diagram
commutes only up to homotopy. To explain the presence of the transverse intersection terms in the pull-back formula,
$$\underset{f_{s,t}(G)=\mathrm{\Gamma }}{}\frac{\zeta _G(\eta _G^{}(c_1(F)))}{AutG},$$
we observe that the induced diagram in cohomology commutes, hence, if one chooses suitable sections $`\sigma _{F_G}`$’s of the bundles $`\eta _G^{}(F)`$, one can proceed as in the transverse intersection case. We now pass to justify the self-intersection term. In our specific situation this term is
$$\eta _G^{}(c_1(F))c_1(N_{\xi _A}),$$
in fact, since $`\mathrm{\Gamma }=A`$, the only $`A`$-graph $`G`$ such that $`j_{s,t}(G)=\mathrm{\Gamma }`$ is the trivial $`A`$-graph and the map $`\zeta _G`$ is the identity. The corresponding component in the preimage of $`\mathrm{\Delta }_\mathrm{\Gamma }`$ under the map $`s_A`$ is the Poincaré dual to $`c_1(N_{\xi _A})`$. Take a section of such bundle, call it $`\tau `$. The component we are looking for is the Poincaré dual of
$$\{x\overline{}_A\sigma _{F_G}(x)=0,\tau (x)=0\},$$
that is, the first Chern class of the bundle
$$\eta _G^{}(F)N_{\xi _A},$$
as we claimed.
$`\mathrm{}`$
### 4.1 Formulas for $`\pi ^{}`$
Let
$$\pi _A:\overline{}_{g,PA}\overline{}_{g,P}$$
be the map forgetting the $`A`$ markings. We first recall pull-back formulas for degree $`2`$ classes (see \[AC1\] and \[AC2\]).
| $`\pi _A^{}\left(\delta _{c,C}\right)=_{BA}\delta _{c,CB}`$ | | $`\pi _A^{}\left(\psi _i\right)=\psi _i_{BA}\delta _{0,B\left\{i\right\}}`$ |
| --- | --- | --- |
| $`\pi _A^{}\left(\delta _{irr}\right)=\delta _{irr}`$ | | $`\pi _A^{}\left(\kappa _1\right)=\kappa _1_{iA}\psi _i+_{BA}\delta _{0,B}`$ |
The pull-back formulas for Mumford classes are recursively deduced from Formula (1.10) in \[AC2\] and Lemma (1.2) in \[AC1\]; if $`\pi :\overline{}_{0,n}\overline{}_{0,n1}`$ is the forgetful map, then
$$\psi _i=\pi ^{}\left(\psi _i\right)+\delta _{0,\{i,n\}},$$
(3)
and
$$\kappa _i=\pi ^{}\left(\kappa _i\right)+\psi _n^i.$$
(4)
Let us now come to degree $`4`$ classes.
Mumford classes are pulled back via formulas 3 and 4:
$`\pi _A^{}\left(\psi _i^2\right)`$ $`=`$ $`\psi _i^2{\displaystyle \underset{BA}{}}\delta _{0,B\left\{i\right\}}|\psi +\text{type }G\text{ classes,}`$
$`\pi _A^{}\left(\psi _i\psi _j\right)`$ $`=`$ $`\psi _i\psi _j\psi _j{\displaystyle \underset{BA}{}}\delta _{0,B\left\{i\right\}}\psi _i{\displaystyle \underset{BA}{}}\delta _{0,B\left\{j\right\}}+\text{type }G\text{ classes,}`$
$`\pi _A^{}\left(\kappa _1\psi _i\right)`$ $`=`$ $`\kappa _1\psi _i\psi _i{\displaystyle \underset{jA}{}}\psi _j{\displaystyle \underset{BA}{}}\delta _{0,B\left\{i\right\}}|\kappa +{\displaystyle \underset{BA,jA\backslash B}{}}\psi _j\delta _{0,B\left\{i\right\}}+{\displaystyle \underset{BA}{}}\psi _i\delta _{0,B}`$
$`+\text{type }G\text{ classes,}`$
$`\pi _A^{}\left(\kappa _1^2\right)`$ $`=`$ $`\kappa _1^22{\displaystyle \underset{iA}{}}\kappa _1\psi _i+{\displaystyle \underset{iA}{}}\psi _i^2+2{\displaystyle \underset{i,jA,ij}{}}\psi _i\psi _j`$
$`+2{\displaystyle \underset{BA}{}}\delta _{0,B}|\kappa 2{\displaystyle \underset{BA,iA\backslash B}{}}\psi _i\delta _{0,B}{\displaystyle \underset{BA}{}}\delta _{0,B}|\psi +\text{ type }G\text{ classes,}`$
$$\pi _A^{}\left(\kappa _2\right)=\kappa _2\underset{iA}{}\psi _i^2+\underset{BA}{}\delta _{0,B}|\psi +\text{type }G\text{ classes;}$$
this last formula is computed by induction on $`|A|`$.
With arguments similar to the ones used in Proposition 8, one can easily prove the following:
###### Proposition 9
The following formulas hold:
| $`\pi _A^{}\left(p|\delta _{irr}\right)=\left(\stackrel{~}{\pi }_A^{}\left(p\right)\right)|\delta _{irr},`$ | | $`\pi _A^{}\left(\delta _{E(c,C)}\right)=_{BA}\delta _{E(c,CB)},`$ |
| --- | --- | --- |
| $`\pi _A^{}\left(p|\delta _{c,C}\right)=_{BA}\left(\stackrel{~}{\pi }_B^{}\left(p\right)\right)|\delta _{c,CB}`$ | | $`\pi _A^{}\left(\delta _{H(c,C)}\right)=_{BA}\delta _{H(c,CB)}`$ |
| $`\pi _A^{}\left(\delta _F\right)=\delta _F`$ | | $`\pi _A^{}\left(\delta _{G(c,C,d,D)}\right)=_{\left(BB^{}\right)A}\left(\delta _{G(c,CB,d,DB^{})}\right)`$ |
where
$`\stackrel{~}{\pi }_A`$ $`:`$ $`\overline{}_{g1,PA\{q,r\}}\overline{}_{g1,P\{q,r\}}`$
$`\stackrel{~}{\pi }_B`$ $`:`$ $`\overline{}_{c,CB\left\{s\right\}}\times \overline{}_{gc,\left(P\backslash C\right)(A\backslash B)\left\{t\right\}}\overline{}_{c,C\left\{s\right\}}\times \overline{}_{gc,\left(P\backslash C\right)\left\{t\right\}}\text{.}`$
$`\mathrm{}`$
## 5 Relations in degree $`4`$
New relations arising in degree $`4`$ appear in $`\overline{}_{g,n}`$ for $`g5`$ and for suitable $`n`$, and can be pulled back with formulas in 4.1. They have been computed with different techniques by E. Getzler, R. Pandharipande, P. Belorousski, and C. Faber. Most of them can be found in the literature, and we will give below the precise reference. The existence of some of them follows from \[Fa5\], as a consequence of the existence of tautological relations on $`_{g,n}`$, while their explicit expression on $`\overline{}_{g,n}`$ has been recently computed by C. Faber and privately communicated to the author (\[Fa4\]). The only exception is the new relation in $`\overline{}_{3,2}`$, whose coefficients will be determined in section 6 by the “pull-back to the boundary” techniques.
### 5.1 Genus $`0`$
The only new result is that
$$\kappa _2=0\text{ in }H^4\left(\overline{}_{0,4}\right)$$
for dimension reasons.
### 5.2 Genus $`1`$
As above,
$$\kappa _2=0\text{ in }H^4\left(\overline{}_{1,1}\right).$$
Moreover, as observed by Faber in \[Fa3\],
$$\delta _{irr}^2=0\text{.}$$
There are other relations: the first one originates in $`H^4\left(\overline{}_{1,2}\right):`$
$$\delta _{E(0,\left\{i\right\})}\delta _{H(0,\mathrm{})}=0,$$
as the push-forward of Keel relation with the map $`\xi _{irr}:\overline{}_{0,4}\overline{}_{1,2}`$. The second one originates in $`H^4\left(\overline{}_{1,4}\right)`$:
$`0`$ $`=`$ $`12{\displaystyle \underset{i}{}}\delta _{G(0,\{1,i\},1,\mathrm{})}12{\displaystyle \underset{i}{}}\delta _{G(1,\left\{i\right\},0,\{\})}2{\displaystyle \underset{i,j}{}}\delta _{G(1,\mathrm{},0,\{i,j\})}`$
$`+6{\displaystyle \underset{i}{}}\delta _{G(1,\mathrm{},0,\left\{i\right\})}2{\displaystyle \underset{i}{}}\delta _{E(\{1,i\})}+{\displaystyle \underset{i}{}}\delta _{H(\left\{i\right\})}+\delta _{H\left(\mathrm{}\right)}.`$
This was discovered by Getzler (\[G1\]), while Pandharipande (\[Pa\]) then proved it is algebraic.
### 5.3 Genus $`2`$
Following Mumford (\[Mu\]),
$$60\kappa _2=\delta _F+6\delta _{H(0,\mathrm{})}$$
in $`H^4\left(\overline{}_{2,0}\right)`$. Faber proves that in $`H^4\left(\overline{}_{2,1}\right)`$
$$\psi _i^2=\frac{1}{120}\delta _F+\frac{1}{5}\delta _{E(1,\mathrm{})}+\frac{13}{120}\delta _{H(0,i)}\frac{1}{120}\delta _{H(0,\mathrm{})}+\frac{7}{5}\delta _{G(1,\mathrm{},0,i)}.$$
Getzler proves in (\[G2\]) that, in $`H^4\left(\overline{}_{2,2}\right)`$,
$`\psi _i\psi _j`$ $`=`$ $`3\psi |\delta _{2,\mathrm{}}+{\displaystyle \frac{1}{72}}\delta _F+{\displaystyle \frac{7}{15}}\delta _{E(1,\mathrm{})}+{\displaystyle \frac{1}{15}}\left(\delta _{E(1,i)}+\delta _{E(1,j)}\right)`$
$`+{\displaystyle \frac{23}{120}}\delta _{H(0,ij)}+{\displaystyle \frac{1}{24}}\left(\delta _{H(0,i)}+\delta _{H(0,j)}\right){\displaystyle \frac{1}{40}}\delta _{H(0,\mathrm{})}{\displaystyle \frac{1}{15}}\delta _{H(1,\mathrm{})}`$
$`+{\displaystyle \frac{13}{5}}\delta _{G(1,\mathrm{},0,ij)}+{\displaystyle \frac{4}{5}}\left(\delta _{G(1,i,0,j)}+\delta _{G(1,j,0,i)}\right){\displaystyle \frac{4}{5}}\delta _{G(0,ij,1,\mathrm{})}.`$
A new algebraic relation was discovered by Belorousski and Pandharipande (\[BP\]) in $`H^4\left(\overline{}_{2,3}\right)`$:
$`0`$ $`=`$ $`12\psi |\delta _{2,\mathrm{}}6{\displaystyle \underset{i=1}{\overset{3}{}}}\psi |\delta _{2,i}+6{\displaystyle \underset{i=1}{\overset{3}{}}}\psi _i\delta _{2,i}+{\displaystyle \frac{6}{5}}\delta _{E(1,\mathrm{})}{\displaystyle \frac{6}{5}}{\displaystyle \underset{i=1}{\overset{3}{}}}\delta _{E(1,i)}+{\displaystyle \frac{2}{5}}{\displaystyle \underset{i=1}{\overset{3}{}}}\delta _{E(0,i)}`$
$`+{\displaystyle \frac{1}{10}}\delta _{H(0,123)}{\displaystyle \frac{3}{10}}{\displaystyle \underset{i=1}{\overset{3}{}}}\delta _{H(0,jk)}+{\displaystyle \frac{3}{10}}{\displaystyle \underset{i=1}{\overset{3}{}}}\delta _{H(0,i)}{\displaystyle \frac{1}{10}}\delta _{H(0,\mathrm{})}{\displaystyle \frac{3}{5}}\delta _{H(1,\mathrm{})}{\displaystyle \frac{1}{5}}{\displaystyle \underset{i=1}{\overset{3}{}}}\delta _{H(1,i)}`$
$`12\delta _{G(2,\mathrm{},0,)}+{\displaystyle \frac{12}{5}}\delta _{G(1,\mathrm{},0,123)}{\displaystyle \frac{12}{5}}{\displaystyle \underset{i=1}{\overset{3}{}}}\delta _{G(1,i,0,jk)}+{\displaystyle \frac{24}{5}}{\displaystyle \underset{i=1}{\overset{3}{}}}\delta _{G(1,\mathrm{},0,i)}`$
$`{\displaystyle \frac{36}{5}}{\displaystyle \underset{i=1}{\overset{3}{}}}\delta _{G(1,,0,i)}{\displaystyle \frac{36}{5}}{\displaystyle \underset{i=1}{\overset{3}{}}}\delta _{G(1,\mathrm{},1,\mathrm{})}+{\displaystyle \frac{18}{5}}{\displaystyle \underset{i=1}{\overset{3}{}}}\delta _{G(1,i,1,\mathrm{})}{\displaystyle \frac{12}{5}}{\displaystyle \underset{i=1}{\overset{3}{}}}\delta _{G(1,\mathrm{},1,i)}.`$
Here, and from now on, every time we write the symbol $``$ instead of a marking’s name, we mean that any marking which does not appear elsewhere in the notation could replace the $``$.
### 5.4 Genus $`3`$
In $`H^4\left(\overline{}_{3,0}\right)`$(\[Fa4\] and \[Fa1\]):
$`\kappa _1^2`$ $`=`$ $`{\displaystyle \frac{5}{7}}\psi |\delta _{irr}{\displaystyle \frac{89}{7}}\psi |\delta _{2,\mathrm{}}{\displaystyle \frac{2}{35}}\delta _F{\displaystyle \frac{94}{35}}\delta _{E(1,\mathrm{})}+{\displaystyle \frac{103}{84}}\delta _{H(0,\mathrm{})}{\displaystyle \frac{2}{7}}\delta _{H(1,\mathrm{})}{\displaystyle \frac{22}{35}}\delta _{G(1,\mathrm{},1,\mathrm{})},`$
$`\kappa _2`$ $`=`$ $`{\displaystyle \frac{5}{42}}\psi |\delta _{irr}{\displaystyle \frac{41}{21}}\psi |\delta _{2,\mathrm{}}+{\displaystyle \frac{1}{630}}\delta _F{\displaystyle \frac{11}{35}}\delta _{E(1,\mathrm{})}+{\displaystyle \frac{41}{252}}\delta _{H(0,\mathrm{})}+{\displaystyle \frac{2}{105}}\delta _{H(1,\mathrm{})}+{\displaystyle \frac{8}{35}}\delta _{G(1,\mathrm{},1,\mathrm{})},`$
whereas in $`H^4\left(\overline{}_{3,1}\right)`$ a new relation involving $`\kappa _1\psi _i`$ appears, and the three of them could be written as follows (\[Fa4\]):
$`\kappa _1\psi _i`$ $`=`$ $`5\psi _i^2{\displaystyle \frac{1}{7}}\psi _i\delta _{irr}{\displaystyle \frac{1}{42}}\psi |\delta _{irr}{\displaystyle \frac{5}{7}}\psi _i\delta _{2,i}{\displaystyle \frac{16}{21}}\psi |\delta _{2,i}{\displaystyle \frac{40}{21}}\psi \delta _{2,\mathrm{}}{\displaystyle \frac{1}{630}}\delta _F`$
$`+{\displaystyle \frac{13}{21}}\delta _{E(0,i)}{\displaystyle \frac{9}{35}}\delta _{E(1,i)}+{\displaystyle \frac{61}{252}}\delta _{H(0,i)}{\displaystyle \frac{2}{105}}\delta _{H(1,i)}+{\displaystyle \frac{4}{105}}\delta _{H(1,\mathrm{})}+{\displaystyle \frac{4}{63}}\delta _{H(0,\mathrm{})}`$
$`+{\displaystyle \frac{16}{35}}\delta _{G(1,i,1,\mathrm{})}+{\displaystyle \frac{61}{21}}\delta _{G(1,\mathrm{},0,i)}{\displaystyle \frac{8}{35}}\delta _{G(1,\mathrm{},1,i)},`$
$`\kappa _1^2`$ $`=`$ $`9\psi _i^2{\displaystyle \frac{2}{7}}\psi _i\delta _{irr}{\displaystyle \frac{16}{21}}\psi |\delta _{irr}{\displaystyle \frac{10}{7}}\psi _i\delta _{2,i}{\displaystyle \frac{299}{21}}\psi |\delta _{2,i}{\displaystyle \frac{347}{21}}\psi \delta _{2,\mathrm{}}{\displaystyle \frac{19}{315}}\delta _F`$
$`+{\displaystyle \frac{83}{3}}\delta _{E(0,i)}{\displaystyle \frac{16}{5}}\delta _{E(1,i)}+{\displaystyle \frac{431}{252}}\delta _{H(0,i)}{\displaystyle \frac{34}{105}}\delta _{H(1,i)}{\displaystyle \frac{22}{105}}\delta _{H(1,\mathrm{})}+{\displaystyle \frac{341}{252}}\delta _{H(0,\mathrm{})}`$
$`+{\displaystyle \frac{2}{7}}\delta _{G(1,i,1,\mathrm{})}+{\displaystyle \frac{389}{21}}\delta _{G(1,\mathrm{},0,i)}{\displaystyle \frac{38}{35}}\delta _{G(1,\mathrm{},1,i)},`$
$`\kappa _2`$ $`=`$ $`\psi _i^2{\displaystyle \frac{5}{42}}\psi |\delta _{irr}{\displaystyle \frac{41}{21}}\psi |\delta _{2,i}{\displaystyle \frac{347}{21}}\psi \delta _{2,\mathrm{}}+{\displaystyle \frac{1}{630}}\delta _F`$
$`+{\displaystyle \frac{5}{21}}\delta _{E(0,i)}{\displaystyle \frac{11}{35}}\delta _{E(1,i)}+{\displaystyle \frac{41}{252}}\delta _{H(0,i)}+{\displaystyle \frac{2}{105}}\delta _{H(1,i)}+{\displaystyle \frac{2}{105}}\delta _{H(1,\mathrm{})}+{\displaystyle \frac{41}{252}}\delta _{H(0,\mathrm{})}`$
$`+{\displaystyle \frac{8}{35}}\delta _{G(1,i,1,\mathrm{})}+{\displaystyle \frac{41}{21}}\delta _{G(1,\mathrm{},0,i)}+{\displaystyle \frac{8}{35}}\delta _{G(1,\mathrm{},1,i)}.`$
Finally, in $`H^4\left(\overline{}_{3,2}\right)`$, we have:
$`0`$ $`=`$ $`\psi _a^2+\psi _b^2{\displaystyle \frac{6}{5}}\psi _a\psi _b\kappa |\delta _{3,\mathrm{}}+5\psi |\delta _{3,\mathrm{}}{\displaystyle \frac{40}{21}}\psi |\delta _{2,\mathrm{}}+{\displaystyle \frac{5}{3}}(\psi |\delta _{2,a}+\psi |\delta _{2,b})`$
$`{\displaystyle \frac{6}{7}}\left(\psi _a\delta _{2,a}+\psi _b\delta _{2,b}\right){\displaystyle \frac{16}{21}}\psi |\delta _{2,ab}+{\displaystyle \frac{12}{35}}\left(\psi _a\delta _{2,ab}+\psi _b\delta _{2,ab}\right){\displaystyle \frac{1}{42}}\psi |\delta _{irr}`$
$`+{\displaystyle \frac{1}{35}}\left(\psi _a\delta _{irr}+\psi _b\delta _{irr}\right){\displaystyle \frac{1}{630}}\delta _F+{\displaystyle \frac{13}{21}}\delta _{E(2,\mathrm{})}{\displaystyle \frac{4}{15}}\left(\delta _{E(2,a)}+\delta _{E(2,b)}\right)`$
$`{\displaystyle \frac{9}{35}}\delta _{E(1,\mathrm{})}{\displaystyle \frac{34}{105}}\delta _{E(1,a)}+{\displaystyle \frac{1}{7}}\delta _{H(2,\mathrm{})}{\displaystyle \frac{2}{105}}\delta _{H(1,ab)}+{\displaystyle \frac{4}{105}}\delta _{H(1,\mathrm{})}`$
$`+{\displaystyle \frac{1}{105}}\left(\delta _{H(1,a)}+\delta _{H(1,b)}\right)+{\displaystyle \frac{4}{63}}\delta _{H(0,\mathrm{})}+{\displaystyle \frac{10}{63}}\delta _{H(0,ab)}{\displaystyle \frac{5}{36}}\left(\delta _{H(0,a)}+\delta _{H(0,b)}\right)`$
$`+{\displaystyle \frac{40}{21}}\delta _{G(2,\mathrm{},0,ab)}\delta _{G(2,\mathrm{},1,\mathrm{})}+{\displaystyle \frac{16}{35}}\delta _{G(1,ab,1,\mathrm{})}{\displaystyle \frac{8}{35}}\delta _{G(1,\mathrm{},1,ab)}`$
$`{\displaystyle \frac{5}{3}}\left(\delta _{G(2,b,0,a)}+\delta _{G(2,a,0,b)}\right){\displaystyle \frac{40}{21}}\left(\delta _{G(2,\mathrm{},0,a)}+\delta _{G(2,\mathrm{},0,b)}\right).`$
### 5.5 Genus $`4`$
In $`H^4\left(\overline{}_{4,\mathrm{}}\right)`$(\[Fa4\] and \[Fa2\]):
$`0`$ $`=`$ $`{\displaystyle \frac{45}{2}}\kappa _1^2240\kappa _27\kappa _1\delta _{irr}+{\displaystyle \frac{35}{2}}\psi |\delta _{irr}39\kappa |\delta _{3,\mathrm{}}+{\displaystyle \frac{315}{2}}\psi |\delta _{3,\mathrm{}}+{\displaystyle \frac{45}{2}}\psi |\delta _{2,\mathrm{}}`$
$`+\delta _F+13\delta _{E(2,\mathrm{})}{\displaystyle \frac{105}{8}}\delta _{H(0,\mathrm{})}+2\delta _{H(1,\mathrm{})}+5\delta _{H(2,\mathrm{})}+24\delta _{G(1,\mathrm{},1,\mathrm{})}+21\delta _{G(1,\mathrm{},2,\mathrm{})},`$
and since another relation appears in $`H^4\left(\overline{}_{4,1}\right)`$ (\[Fa4\]), we get there the following two relations:
$`0`$ $`=`$ $`5\kappa _1^230\kappa _240\kappa _1\psi _i+245\psi _i^2\kappa _1\delta _{irr}+7\psi _i\delta _{irr}2\kappa |\delta _{3,i}+44\psi _i\delta _{3,i}`$
$`35\psi |\delta _{3,i}32\kappa |\delta _{3,\mathrm{}}+175\psi |\delta _{3,\mathrm{}}30\psi _i\delta _{2,i}+95\psi |\delta _{2,i}85\psi |\delta _{2,\mathrm{}}`$
$`36\delta _{E(0,i)}+24\delta _{E(1,i)}12\delta _{E(2,i)}+{\displaystyle \frac{35}{12}}\delta _{H(0,\mathrm{})}+\delta _{H(1,\mathrm{})}+5\delta _{H(2,\mathrm{})}`$
$`{\displaystyle \frac{175}{12}}\delta _{H(0,i)}+\delta _{H(1,i)}\delta _{H(2,i)}18\delta _{G(1,i,1,\mathrm{})}+28\delta _{G(1,i,2,\mathrm{})}`$
$`+12\delta _{G(2,i,1,\mathrm{})}175\delta _{G(3,\mathrm{},0,i)}10\delta _{G(2,\mathrm{},0,i)}+12\delta _{G(1,\mathrm{},1,i)}4\delta _{G(1,\mathrm{},2,i)}`$
$`0`$ $`=`$ $`{\displaystyle \frac{25}{2}}\kappa _1^2180\kappa _2+35\kappa _1\psi _i{\displaystyle \frac{455}{2}}\psi _i^25\kappa _1\delta _{irr}+{\displaystyle \frac{35}{2}}\psi |\delta _{irr}7\psi _i\delta _{irr}`$
$`35\kappa |\delta _{3,i}49\psi _i\delta _{3,i}+{\displaystyle \frac{455}{2}}\psi |\delta _{3,i}+25\kappa |\delta _{3,\mathrm{}}{\displaystyle \frac{385}{2}}\psi |\delta _{3,\mathrm{}}+60\psi _i\delta _{2,i}{\displaystyle \frac{335}{2}}\psi |\delta _{2,i}+{\displaystyle \frac{385}{2}}\psi |\delta _{2,\mathrm{}}`$
$`+\delta _F+37\delta _{E(0,i)}35\delta _{E(1,i)}+37\delta _{E(2,i)}{\displaystyle \frac{455}{24}}\delta _{H(0,\mathrm{})}5\delta _{H(2,\mathrm{})}+{\displaystyle \frac{385}{24}}\delta _{H(0,i)}+7\delta _{H(2,i)}`$
$`+60\delta _{G(1,i,1,\mathrm{})}35\delta _{G(1,i,2,\mathrm{})}+{\displaystyle \frac{385}{2}}\delta _{G(3,\mathrm{},0,i)}25\delta _{G(2,\mathrm{},0,i)}+49\delta _{G(1,\mathrm{},2,i)}.`$
### 5.6 Genus $`5`$
Finally, in $`H^4\left(\overline{}_{5,0}\right)`$(\[Fa4\]):
$`0`$ $`=`$ $`{\displaystyle \frac{25}{2}}\kappa _1^2180\kappa _25\kappa _1\delta _{irr}+{\displaystyle \frac{35}{2}}\psi |\delta _{irr}35\kappa |\delta _{4,\mathrm{}}+{\displaystyle \frac{455}{2}}\psi |\delta _{4,\mathrm{}}+25\kappa |\delta _{3,\mathrm{}}`$
$`{\displaystyle \frac{385}{2}}\psi |\delta _{3,\mathrm{}}{\displaystyle \frac{385}{2}}\delta _{3,\mathrm{}}|\psi +\delta _F+37\delta _{E(1,\mathrm{})}35\delta _{E(2,\mathrm{})}{\displaystyle \frac{455}{24}}\delta _{E(0,\mathrm{})}`$
$`5\delta _{H(2,\mathrm{})}+7\delta _{H(3,\mathrm{})}35\delta _{G(1,\mathrm{},2,\mathrm{})}+49\delta _{G(1,\mathrm{},3,\mathrm{})}+25\delta _{G(2,\mathrm{},1,\mathrm{})}.`$
## 6 Degree $`4`$ relations in the tautological group
###### Theorem 10
For $`g6`$, $`_{g,P}^4`$ is a basis for $`T_{g,P}^4`$. For $`2g5`$, the relations among elements of $`_{g,P}^4`$ are the ones listed in section 5.
We will prove this Theorem by induction on $`g`$. We start with a sketchy exposition of an argument which covers the cases $`g6`$, once the previous ones are established. Unfortunately, this argument fails to extend to the low genus cases. We will therefore give a second, less direct argument. The initial cases require more involved computations, because of the presence of many relations among tautological classes. We will work out two sample cases in Lemmas 14 and 16, and recover the coefficients of the new relation in $`\overline{}_{3,2}`$ in Proposition 15.
###### Proposition 11
Suppose that Theorem 10 holds for $`g=5`$. Then it holds for every genus $`g6`$.
Proof. For the first proof we make an induction on $`g`$. Consider the boundary maps:
$$\xi _{a,A}:\overline{}_{a,A\left\{s\right\}}\times \overline{}_{ga,A^C\left\{t\right\}}\overline{}_{g,P},$$
on varying $`(a,A)`$ in such a way that $`a3,ga3`$. Consider the composition of the induced pull-back map with the projection on $`H^2H^2`$:
$$g_{a,A}:H^4\left(\overline{}_{g,P}\right)H^2\left(\overline{}_{a,A\left\{s\right\}}\right)H^2\left(\overline{}_{ga,A^C\left\{t\right\}}\right).$$
We need a few remarks:
* Under the above hypotheses on genera, there are no relation among tautological classes in $`H^2\left(\overline{}_{a,A\left\{s\right\}}\right)H^2\left(\overline{}_{ga,A^C\left\{t\right\}}\right)`$.
* Every class of the standard basis in $`H^2\left(\overline{}_{a,A\left\{s\right\}}\right)H^2\left(\overline{}_{ga,A^C\left\{t\right\}}\right)`$ (by the standard basis we mean the one described in \[AC1\]), appears, with the suitable sign, as a summand in the pull-back of at most one tautological class of $`H^4\left(\overline{}_{g,P}\right)`$, with the exception of $`\psi _s\psi _t`$, which is a summand both of $`\xi _{a,A}^{}\left(\psi |\delta _{a,A}\right)`$ and $`\xi _{a,A}^{}\left(\delta _{a,A}|\psi \right)`$. This is a combinatorial remark which follows from the description of pull-backs of section 4. In particular, one should look at the description of the operations on graphs denoted by $`f_{s,t}`$ and $`j_{s,t}`$.
* Almost every essential tautological class $`\alpha `$ in $`_{g,P}^4`$ satisfies $`g_{a,A}\left(\alpha \right)0`$ for at least one $`(a,A)`$ satisfying the hypotheses. This is also a combinatorial remark, and it is based on the relative position of boundary cycles in $`\overline{}_{g,P}`$. The exceptions are:
$`\kappa _2,\psi _x\text{ for every }xP,\delta _{E(b,B)},`$
$`\delta _{G(c,C,d,D)}\text{, if }c+d2.`$
Suppose there is a relation among essential tautological classes in $`H^4\left(\overline{}_{g,P}\right)`$. Applying all the maps $`g_{a,A}`$, one obtains that many coefficients have to vanish. The relation should then be:
$$c\kappa _2+\underset{xP}{}c_x\psi _x^2+c_{b,B}\delta _{E(b,B)}+\underset{c+d2}{}c_{c,C,d,D}\delta _{G(c,C,d,D)}=0$$
We pull it back with the map
$$\xi ^{}:H^4\left(\overline{}_{g,P}\right)H^4\left(\overline{}_{g1,P\{q,r\}}\right),$$
and get
$`c\kappa _2+{\displaystyle \underset{xP}{}}c_x\psi _x^2+{\displaystyle c_{b,B}\left(\delta _{E(b,B)}+\delta _{E(b1,B\{q,r\})}+\mathrm{}\right)}`$
$`+{\displaystyle \underset{c+d2}{}}c_{c,C,d,D}\left(\delta _{G(c1,C\{q,r\},d,D)}+\delta _{G(c,C,d1,D\{q,r\})}+\delta _{G(c,C,d,D)}\right)`$ $`=`$ $`0`$
By induction hypothesis, the coefficients $`c,c_x,c_{b,B}`$ all have to vanish. Every type $`G`$ class appears at most once as a summand in the image of at type $`G`$ class. If we call “critical” the classes corresponding to graphs $`G(0,A,0,B)`$, i.e. the possibly unessential ones, we observe that every non-critical class has at least one non-critical summand in its pull-back. On the other hand, if we extend the ordering of $`P`$ to an ordering for $`P\{q,r\}`$ imposing $`\{q,r\}`$ to be the last two elements, then a basis of critical classes maps to a set of linearly independent critical classes. Thus, the coefficients $`c_{c,C,d,D}`$ vanish.
$`\mathrm{}`$
The main tool used in the second proof is the map:
$$\xi ^{}:H^4\left(\overline{}_{g,P}\right)H^4\left(\overline{}_{g1,P\{q,r\}}\right).$$
The combinatorics of tautological classes and pull-back formulas becomes rather intricate, but nevertheless it suggests a partition of $`_{g,P}^4`$, corresponding to any given partition of $`P`$, which, inductively, turns out to give a direct sum decomposition of the tautological group.
###### Definition 12
1. Pure boundary classes of type $`E`$ and $`F`$
are essential pure boundary classes corresponding to graphs $`F`$ and $`E(a,A)`$.
They generate the subspace $`𝐖_{EF}`$ of $`T_{g,P}^4`$.
2. Pure boundary classes of type $`H`$ and $`G`$
are essential pure boundary classes corresponding to graphs $`H(a,A)`$ and $`G(a,A,b,B)`$.
They generate the subspace $`𝐖_{GH}`$ of $`T_{g,P}^4`$.
3. $`\mathrm{\Psi }`$\- mixed classes
are essential mixed boundary classes $`\psi |\delta _{irr}`$ and $`\psi |\delta _{a,A}`$, generating $`𝐖_\mathrm{\Psi }.`$
4. $`\mathrm{\Psi }_I`$-mixed classes
are essential mixed boundary classes $`\psi _i\delta _{irr}`$ and $`\psi _i\delta _{a,A}`$, with $`iIA`$, generating $`𝐖_{\mathrm{\Psi }I}.`$
5. $`K`$\- mixed classes
are essential mixed boundary classes $`\kappa _1\delta _{irr}`$ and $`\kappa |\delta _{a,A}`$, generating $`𝐖_K.`$
6. Mumford $`K`$ classes
are essential classes $`\{\begin{array}{c}\kappa _1^2,\kappa _2\text{, for }g6\\ \kappa _1^2\text{, for }g=5\text{ and }g=4,P=\mathrm{}\\ \mathrm{}\text{, for }g=4,P\mathrm{}\text{, and }g3\end{array}`$, and generate $`𝐊`$.
7. Mumford $`\mathrm{\Psi }_I`$ classes
are essential classes $`\{\begin{array}{c}\kappa _1\psi _i,\psi _i^2,\psi _i\psi _j\text{, for }g4\\ \psi _i^2,\psi _i\psi _j\text{, for }g=3\\ \mathrm{}\text{, for }g2\end{array}`$, with $`i,jI`$, and generate $`𝚿_I`$.
8. Mumford $`\mathrm{\Psi }_{IJ}`$ classes
are essential classes $`\{\begin{array}{c}\psi _i\psi _j\text{, for }g3\\ \mathrm{}\text{, for }g2\end{array}`$, with $`iI`$, $`jJ`$, and generate $`𝚿_{IJ}`$.
###### Proposition 13
Suppose that Theorem 10 holds for $`g=6`$. Then it holds for every genus $`g6`$.
Proof. Let $`O=\{q,r\}`$, so that $`P\{q,r\}=PO`$. Following formulas of section 4, we describe how the above subspaces of $`T_{g,P}^4`$ behave with respect to the map
$$\xi ^{}:H^4\left(\overline{}_{g,P}\right)H^4\left(\overline{}_{g1,P\{q,r\}}\right).$$
We write down the behavior for genus $`g4`$. When no confusion will arise, we will denote by the same letter the subspaces of the same type in $`H^4\left(\overline{}_{g,P}\right)`$ and $`H^4\left(\overline{}_{g1,P\{q,r\}}\right)`$.
$`K`$ $``$ $`K\text{, for }g7`$
$`\mathrm{\Psi }_P`$ $``$ $`\mathrm{\Psi }_P\text{, for }g5`$
$`W_K`$ $``$ $`W_K+W_\mathrm{\Psi }+W_{\mathrm{\Psi }P}+W_{EF}+W_{GH}+W_{\mathrm{\Psi }O}+\mathrm{\Psi }_O`$
$`W_\mathrm{\Psi }`$ $``$ $`W_\mathrm{\Psi }+W_{GH}+\mathrm{\Psi }_O`$
$`W_{\mathrm{\Psi }P}`$ $``$ $`W_{\mathrm{\Psi }P}+W_{GH}+\mathrm{\Psi }_{OP}`$
$`W_{EF}`$ $``$ $`W_{EF}+W_{GH}+W_{\mathrm{\Psi }O}`$
$`W_{GH}`$ $``$ $`W_{GH}+W_{\mathrm{\Psi }O}.`$
We prove the Proposition by induction on $`g`$. Suppose that
$$T_{g1,P\{q,r\}}^4=W_{EF}W_{GH}W_\mathrm{\Psi }W_{\mathrm{\Psi }P}W_{\mathrm{\Psi }O}W_K\mathrm{\Psi }_P\mathrm{\Psi }_O\mathrm{\Psi }_{OP}$$
and that every summand is freely generated by essential tautological classes. We write down in block form the matrix of the map
$$\xi ^{}:H^4\left(\overline{}_{g,P}\right)H^4\left(\overline{}_{g1,P\{q,r\}}\right).$$
$$\begin{array}{ccccccccccc}& K\hfill & \mathrm{\Psi }_P\hfill & W_K\hfill & W_\mathrm{\Psi }\hfill & W_{\mathrm{\Psi }P}\hfill & W_{EF}\hfill & W_{GH}\hfill & W_{\mathrm{\Psi }O}\hfill & \mathrm{\Psi }_O\hfill & \mathrm{\Psi }_{OP}\hfill \\ K\hfill & A\hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill \\ \mathrm{\Psi }_P\hfill & 0\hfill & B\hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill \\ W_K\hfill & 0\hfill & 0\hfill & C\hfill & \mathrm{}\hfill & \mathrm{}\hfill & 0\hfill & \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill & 0\hfill \\ W_\mathrm{\Psi }\hfill & 0\hfill & 0\hfill & 0\hfill & D\hfill & 0\hfill & 0\hfill & \mathrm{}\hfill & 0\hfill & \mathrm{}\hfill & 0\hfill \\ W_{\mathrm{\Psi }P}\hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill & E\hfill & 0\hfill & \mathrm{}\hfill & 0\hfill & 0\hfill & \mathrm{}\hfill \\ W_{EF}\hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill & F\hfill & \mathrm{}\hfill & \mathrm{}\hfill & 0\hfill & 0\hfill \\ W_{GH}\hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill & 0\hfill & G\hfill & \mathrm{}\hfill & 0\hfill & 0\hfill \end{array}.$$
We claim that the elements of $`_{g,P}^4`$ form a basis for $`T_{g,P}^4`$. Because of the form of the above matrix, it is sufficient to check that every subset generating each subspace consists of independent classes. For this, we look at blocks $`A,\mathrm{},G`$, and check that each of them has maximal rank, equal to the number of rows. It is easy to see that $`A`$ and $`B`$ are both the identity matrix, whereas from
$`\kappa _1\delta _{irr}\kappa _1\delta _{irr}+\mathrm{}\text{, if}g5`$
$`\kappa |\delta _{a,A}\{\begin{array}{c}\kappa |\delta _{a,A}+\kappa |\delta _{a1,A\{q,r\}}\text{, if }ga1\text{}a4\\ \kappa |\delta _{a1,A\{q,r\}}\text{, if }g=a4\\ \kappa |\delta _{a,A}\text{, if }ga1,a3\\ 0\text{, if }g=a3\end{array}`$
we observe that $`C`$ has maximal rank for $`g5`$.
Similarly, $`D`$ and $`E`$ have maximal rank for $`g3`$, wheres $`F`$ has maximal rank for $`g2`$.
As for the block $`G`$, from
$`\delta _{H(a,A)}\{\begin{array}{c}\delta _{H(a,A)}+\delta _{H(a1,A\{q,r\})}+\mathrm{}\text{, if }g1a1\text{}a2\\ \delta _{H(a,A)}+\frac{5}{6}\delta _{H(a1,A\{q,r\})}+\mathrm{}\text{, if }g1a1\text{}a=1\\ \delta _{H(a1,A\{q,r\})}+\mathrm{}\text{, if }g=a+13\\ \frac{5}{6}\delta _{H(a1,A\{q,r\})}+\mathrm{}\text{, if }g=a+1=2\\ \delta _{H(a,A)}+\mathrm{}\text{, if }g1a1\text{}a=0\\ 0\text{, if }g=a+1=1\end{array}`$
$`\delta _{G(a,A,b,B)}\delta _{G(a,A,b,B)}+\delta _{G(a1,A\{q,r\},b,B)}+\delta _{G(a,A,b1,B\{q,r\})};`$
we observe that type $`H`$ classes are independent, and independent from type $`G`$ ones. For the type $`G`$ class, the argument used in the proof of Proposition 11 works in this case as well. One can write the block $`G`$ in a triangular form, and see that it has maximal rank for $`g3`$.
$`\mathrm{}`$
###### Lemma 14
Suppose that Theorem 10 holds for $`g=5`$. Then it holds for genus $`g=6.`$
Proof. The same proof of Proposition 13 can be repeated to prove that $`_{6,P}^4\backslash \left\{\kappa _2\right\}`$ is a set of linearly independent classes. Thus, if a relation does exist, it should be of the form
$$\kappa _2+\mathrm{}=0;$$
since $`\xi ^{}\left(\kappa _2+\mathrm{}\right)=\kappa _2+\mathrm{}=0`$, then the relation should be a pull-back of the relation in $`H^4\left(\overline{}_{5,0}\right)`$ (see section 5):
$$\kappa _2\frac{1}{180}\delta _F+\frac{37}{180}\delta _{E\left(1\right)}+\mathrm{}=0,$$
and hence it should be of the form
$$\kappa _2\frac{1}{180}\delta _F+\frac{37}{180}\left(\delta _{E(1,q)}+\delta _{E(1,r)}\right)+\mathrm{}=0;$$
but one can easily observe that classes $`\delta _F`$ and $`\delta _{E(1,q)}+\delta _{E(1,r)}`$ do only appear in the pull-back $`\xi ^{}\left(\delta _F\right)=\delta _F+\left(\delta _{E(1,q)}+\delta _{E(1,r)}\right)+\mathrm{}`$, hence cannot have different coefficients. This leads to a contradiction.
$`\mathrm{}`$
###### Proposition 15
There is a unique new relation in $`\overline{}_{3,2}`$, and it is the one described in section 5.
Proof.
We know from \[Fa1\] and \[Fa4\] the relations arising in $`H^4\left(\overline{}_{3,0}\right)`$ and $`H^4\left(\overline{}_{3,1}\right)`$, and further we know that a new relation does exist in $`H^4\left(\overline{}_{3,2}\right)`$, involving pure Mumford classes $`\psi _a`$, $`\psi _b`$, $`\psi _a\psi _b`$. We need to prove that the relation has exactly the form described in section 5, and that no other relation appears. We also recall that the group $`H^4\left(\overline{}_{2,2}\right)`$ has been computed in \[G2\].
The relations in $`H^4\left(\overline{}_{3,0}\right)`$ and $`H^4\left(\overline{}_{3,1}\right)`$ can be all used to write classes $`\kappa _1^2`$, $`\kappa _2`$, $`\kappa _1\psi _i`$ in terms of other boundary classes, when $`|P|2`$.
Therefore, a possible new relation in $`H^4\left(\overline{}_{3,2}\right)`$, can be written as follows:
$$\underset{\mathrm{\Gamma }}{}c_\mathrm{\Gamma }\delta _\mathrm{\Gamma }+\underset{p}{}c_{p\left(irr\right)}p|\delta _{irr}+\underset{p,(a,A)}{}c_{p(a,A)}p|\delta _{a,A}+c_i\psi _i^2+c_{ab}\psi _a\psi _b=0.$$
(7)
The first constraints on coefficients in (7) are derived by writing down explicitly the non-vanishing pull-backs of tautological classes under the map
$$\overline{}_{3,s}\overline{}_{3,ab},$$
which glues a fixed rational tail marked by $`Pt`$ by identifying $`t`$ and $`s`$, and observing that the pull-back in $`H^4\left(\overline{}_{3,s}\right)`$ of (7) must be a multiple of Faber’s relation involving $`\kappa _1\psi _s`$ (see section 5). They are:
| $`c_F=\frac{1}{630}k`$ | $`c_{H(2,\mathrm{})}=\frac{1}{7}k`$ | $`c_{G(1,P,1,\mathrm{})}=\frac{16}{35}k`$ | $`c_{\psi \left(irr\right)}=\frac{1}{42}k`$ |
| --- | --- | --- | --- |
| $`c_{E(1,P)}=\frac{9}{35}k`$ | $`c_{H(1,\mathrm{})}=\frac{4}{105}k`$ | $`c_{G(1,\mathrm{},1,P)}=\frac{8}{35}k`$ | $`c_{\psi (3,\mathrm{})}=5k`$ |
| $`c_{E(0,P)}=\frac{13}{21}k`$ | $`c_{H(1,P)}=\frac{2}{105}k`$ | $`c_{G(1,\mathrm{},2,\mathrm{})}=\frac{5}{7}k`$ | $`c_{\psi (2,\mathrm{})}=\frac{40}{21}k`$ |
| | $`c_{H(0,\mathrm{})}=\frac{4}{63}k`$ | $`c_{G(2,\mathrm{},0,P)}=\frac{40}{21}k`$ | $`c_{\psi (2,P)}=\frac{16}{21}k`$ |
| | $`c_{H(0,P)}=\frac{10}{63}k`$ | $`c_{G(2,\mathrm{},1,\mathrm{})}=k`$ | $`c_{\kappa (3,\mathrm{})}=k`$. |
To determine the coefficient of some classes of type $`H`$ and $`G`$ we also need to use the map
$$H^4\left(\overline{}_{3,P}\right)H^2\left(\overline{}_{2,s}\right)H^2\left(\overline{}_{1,Pt}\right).$$
We then know by \[Fa4\] and \[Fa5\] that a new relation does actually exist, and therefore we fix the value of the constant $`k`$ to be $`1`$.
We consider the following maps:
$`H^4\left(\overline{}_{3,ab}\right)`$ $``$ $`H^2\left(\overline{}_{2,s}\right)H^2\left(\overline{}_{1,abt}\right)`$
$`H^4\left(\overline{}_{3,ab}\right)`$ $``$ $`H^2\left(\overline{}_{2,as}\right)H^2\left(\overline{}_{1,bt}\right)`$
$`H^4\left(\overline{}_{3,ab}\right)`$ $``$ $`H^4\left(\overline{}_{2,as}\right);`$
the constraints on the coefficient derived by pulling back (7) force all of them to be the ones indicated in section 5.
$`\mathrm{}`$
###### Lemma 16
Theorem 10 holds for genus $`g=2.`$
Proof. The cases $`n=0,1`$ are well known (see \[Mu\]); the cases $`n=2,3`$ are entirely described in \[G2\] and \[BP\]. Recall that a new relation appears in $`H^4\left(\overline{}_{2,3}\right)`$ (see section 5).
For every set $`\{i,j,k\}P`$, only the relation pulled back from $`\overline{}_{2,\{i,j,k\}}`$ contains the summand:
$$\psi _i\delta _{2,P\backslash \{j,k\}}+\psi _j\delta _{2,P\backslash \{i,k\}}+\psi _k\delta _{2,P\backslash \{i,j\}};$$
we fix an ordering on $`P`$, and use the relation in $`H^4\left(\overline{}_{2,\{i,j,k\}}\right)`$ to express $`\psi _i\delta _{2,P\backslash \{j,k\}}`$, for $`i<j,i<k,`$ as linear combination of other classes.
Let $`𝒞_{2,P}^4`$ be the set obtained from the set of essential classes $`_{2,P}^4`$ after having eliminated the relations arising in degree $`4`$, that is, after having removed all pure Mumford classes, and the classes $`\psi _i\delta _{2,P\backslash \{j,k\}}`$, for $`i<j,i<k.`$ Observe that the definition of $`𝒞_{2,P}^4`$ depends on the choice of an ordering on $`P`$.
If $`n=4`$, there is no new relation among essential tautological classes; we postpone the proof of this fact. If $`n5`$, let $`F_{2,P}^4`$ be the free vector space generated by classes in $`𝒞_{2,P}^4`$. One can define every pull-back map on $`F_{2,P}^4`$, following formulas in section 4. Our claim is that the map
$$f=\left\{f_{ij}^{}\right\}:F_{2,P}^4_{\{i,j\}P}F_{2,P\backslash \{i,j\}\left\{s\right\}}^4$$
is injective for $`|P|5`$. This implies, by induction, that no new relation among tautological classes can appear for $`n5`$: any new one should map to zero with $`f`$.
We use a decomposition of $`F_{2,P}^4`$ similar to the one described at the beginning of this section.
* $`W_F`$ is generated by $`\delta _F,`$
* $`W_E`$ is generated by classes $`\delta _{E(1,A)},`$
* $`W_{H\left(0\right)}`$ is generated by classes $`\delta _{H(0,A)},`$
* $`W_{H\left(1\right)}`$ is generated by classes $`\delta _{H(1,A)},`$
* $`W_{G(2,0)}`$ is generated by classes $`\delta _{G(2,A,0,B)},`$
* $`W_{G(0,2)}`$ is generated by classes $`\delta _{G(0,A,2,B)},`$
* $`W_{G(1,1)}`$ is generated by classes $`\delta _{G(1,A,1,B)},`$
* $`W_{G(1,0)}`$ is generated by classes $`\delta _{G(1,A,0,B)},`$
* $`W_\psi `$ is generated by classes $`\psi |\delta _{2,A},`$
* $`W_{\psi _I}`$ is generated by classes $`\psi _i\delta _{2,A},`$ with $`iI.`$
In the space $`_{\{i,j\}P}F_{2,P\backslash \{i,j\}\left\{s\right\}}^4`$, we denote by $`W_X=_{ij}`$ $`W_X^{ij}`$ the direct sum of subspaces $`W_X^{ij}F_{2,P\backslash \{i,j\}\left\{s\right\}}^4`$. The matrix of the map $`f`$ can be written in triangular block form (we omit all zeroes):
| | $`W_F`$ | $`W_E`$ | $`W_{H\left(0\right)}`$ | $`W_{G(1,0)}`$ | $`W_{G(0,2)}W_{\psi _S}`$ | $`W_{G(2,0)}`$ | $`W_{H\left(1\right)}`$ | $`W_{G(1,1)}`$ | $`W_\psi `$ | $`W_{\psi _P}`$ |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| $`W_F`$ | $`A`$ | | | | | | | | | |
| $`W_E`$ | | $`B`$ | | | | | | | | |
| $`W_{H\left(0\right)}`$ | | | $`C`$ | | | | | | | |
| $`W_{G(1,0)}`$ | | | | $`D`$ | | | | | | |
| $`W_{G(0,2)}`$ | | | | | $`E`$ | | | | | |
| $`W_{G(2,0)}`$ | | | | | | $`F`$ | | | | |
| $`W_{H\left(1\right)}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | | | | $`G`$ | | | |
| $`W_{G(1,1)}`$ | | | $`\mathrm{}`$ | $`\mathrm{}`$ | | | | $`H`$ | | |
| $`W_\psi `$ | $`\mathrm{}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | | $`\mathrm{}`$ | | | $`I`$ | |
| $`W_{\psi _P}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | $`\mathrm{}`$ | $`L`$ |
We just need to check that the blocks on the diagonal have maximal rank. This is completely trivial for the blocks $`A,B,C,D,E`$. We check block $`G`$, and observe that blocks $`H`$ and $`I`$ present a very similar combinatorics. $`G`$ is of the form
$$\left(\begin{array}{ccccc}G^{12}& G^{13}& \mathrm{}& G^{ij}& \mathrm{}\end{array}\right),$$
where $`G^{ij}`$ is a block of the matrix of the map $`f_{ij}^{}.`$ We can write $`G^{ij}`$ as
| | $`\delta _{H(1,B\left\{s\right\})}`$ | $`\delta _{H(1,B)}`$ |
| --- | --- | --- |
| $`\delta _{H(1,A)},\{i,j\}A`$ | $`Id`$ | $`0`$ |
| $`\delta _{H(1,A)},\{i,j\}A^C`$ | $`0`$ | $`Id`$ |
| $`\delta _{H(1,P\backslash \{i,j\})}`$ | $`0`$ | $`\mathrm{}`$ |
| $`\delta _{H(1,A)},|\{i,j\}A|=1`$ | $`0`$ | $`0`$ |
We consider the matrix $`G^{}`$ obtained removing the second column of blocks from each $`G^{ij}`$, except for the columns corresponding to $`\delta _{H(1,\mathrm{})},\delta _{H(1,x)}.`$ Finally, we can extract such a triangular matrix
δH(1,B),|B|1
δH(1,B{s})δH(1,A),|A|1Id0δH(1,A),|A|2Id.fragments
δH(1,B),|B|1
δH(1,B{s})fragmentsδ𝐻1𝐴,|A|1fragmentsId0fragmentsδ𝐻1𝐴,|A|2fragmentsId\begin{tabular}[]{|l|l|l|}\hline\cr&$\delta_{H\left(1,B\right)},|B|\leq 1$&$\delta_{H\left(1,B\cup\left\{s\right\}\right)}$\\
\hline\cr$\delta_{H\left(1,A\right)},|A|\leq 1$&$Id$&$0$\\
\hline\cr$\delta_{H\left(1,A\right)},|A|\geq 2$&$...$&$Id$\\
\hline\cr\end{tabular}.
Observe that we just need the weaker assumption $`|P|4`$.
As for the block $`L`$, observe that any essential class maps to essential classes, except for
$$\psi _i\delta _{2,P\backslash \{j,k\}}\stackrel{f_{jk}^{}}{}\psi _i\psi _s=\underset{|C^C|3}{}\psi _i\delta _{2,C}\underset{x<i}{}\psi _i\delta _{2,P\backslash \{x,s\}}+\underset{x>i}{}\psi _x\delta _{2,P\backslash \{i,s\}};$$
but this doesn’t prevent us from extracting a non-degenerate matrix
ψiδ2,B,|B|1
ψiδ2,B{s}ψiδ2,A,|A|1Id0ψiδ2,A,|A|2Id.fragments
ψiδ2,B,|B|1
ψiδ2,B{s}fragmentsψ𝑖δ2𝐴,|A|1fragmentsId0fragmentsψ𝑖δ2𝐴,|A|2fragmentsId\begin{tabular}[]{|l|l|l|}\hline\cr&$\psi_{i}\delta_{2,B},|B|\leq 1$&$\psi_{i}\delta_{2,B\cup\left\{s\right\}}$\\
\hline\cr$\psi_{i}\delta_{2,A},|A|\leq 1$&$Id$&$0$\\
\hline\cr$\psi_{i}\delta_{2,A},|A|\geq 2$&$...$&$Id$\\
\hline\cr\end{tabular}.
With the same argument, one can write a sub-block of $`F`$ of the form
| | $`\delta _{G(2,A,0,D\left\{s\right\})},\delta _{G(2,A,0,D)}`$ | $`\delta _{G(2,C\left\{s\right\},0,D)}`$ |
| --- | --- | --- |
| $`\delta _{G(2,A,0,B)},|A|1`$ | $`K`$ | $`0`$ |
| $`\delta _{G(2,A,0,B)},|A|2`$ | $`\mathrm{}`$ | $`Id`$ |
The set $`P\backslash \{i,j\}\left\{s\right\}`$ inherits an ordering from $`P`$, assuming $`s`$ to be the last point; therefore the second column of blocks gives no problem. As for the matrix $`K`$, write it in sub-blocks $`K_A`$, where $`K_A`$ involves classes $`\delta _{G(2,A,0,B)}`$. These classes are all obtained pushing forward from $`H^2\left(\overline{}_{0,A^C\left\{z\right\}}\right)`$, and so are the relations among them in $`F_{2,P}^4`$. The combinatorics of the map corresponding to the block $`K_A`$ is then exactly the same of the map
$$H^2\left(\overline{}_{0,A^C\left\{z\right\}}\right)_{\{i,j\}A^c}H^2\left(\overline{}_{0,A^C\backslash \{i,j\}\{z,s\}}\right)$$
which will be proved in lemma 17 to be injective for $`|A^C|4`$. Therefore each $`K_A`$, and consequently $`K`$, has maximal rank.
As for the case $`n=4`$, we first prove by using the pull-back map
$$H^4\left(\overline{}_{2,\{i,j,k,l\}}\right)H^2\left(\overline{}_{2,\{i,s\}}\right)H^2\left(\overline{}_{0,\{j,k,l,t\}}\right);$$
that in a possible new relation, the coefficients of $`\psi `$-mixed classes and of classes of type $`G`$ vanish.
We now restrict the map $`f`$ to the free vector space generated by the classes with non vanishing coefficient in a possible new relation in $`T_{2,4}^4`$. By the same arguments used for the general case, the new map $`f`$ is injective, and the proof of our Lemma is complete.
$`\mathrm{}`$
###### Lemma 17
For $`|P|5`$, the map
$$H^2\left(\overline{}_{0,P}\right)_{\{x,y\}P\backslash \left\{h\right\}}H^2\left(\overline{}_{0,P\backslash \{x,y\}\left\{s\right\}}\right)$$
is injective.
Proof. The case $`|P|=5`$ is trivial.
We can consider
$$\varphi _A^{}:H^2\left(\overline{}_{0,P}\right)H^2(\overline{}_{0,A\left\{s\right\}}\times \overline{}_{0,A^C\left\{t\right\}})$$
as the sum of the two maps
$$f_A^{}:H^2\left(\overline{}_{0,P}\right)H^2\left(\overline{}_{0,A\left\{s\right\}}\right),$$
$$f_{A^C}^{}:H^2\left(\overline{}_{0,P}\right)H^2\left(\overline{}_{0,A^c\left\{t\right\}}\right),$$
where the two maps are the pull-back of the map that glues any fixed rational tail to the extra marked point. For any such $`A`$, there exist $`\{x,y\}P`$ such that $`AP\backslash \{x,y\}`$. For a suitable choice of the rational tail to glue, we can write a commutative diagram
so that from the induced diagram on $`H^2`$ we read: $`\mathrm{ker}f_{P\backslash \{x,y\}}^{}\mathrm{ker}f_A^{}`$. Therefore, by proposition 2.8 in \[AC1\],
$$\left(_{\{x,y\}P}\mathrm{ker}f_{P\backslash \{x,y\}}^{}\right)\left(_{AP,}\mathrm{ker}f_A^{}\right)=0.$$
The statement is proved by induction on $`|P|`$ . Suppose that $`x\left(_{\{x,y\}P\backslash \left\{h\right\}}\mathrm{ker}f_{P\backslash \{x,y\}}^{}\right)`$, but there exist $`kP\backslash \left\{h\right\}`$, such that $`yf_{P\backslash \{h,k\}}^{}\left(x\right)0`$. By the commutativity of
we see that $`y\left(_{\{x,y\}P\backslash \{h,k\}}\mathrm{ker}f_{P\backslash \{x,y\}}^{}\right)`$, hence by induction hypothesis, $`y=0`$, and we are done.
$`\mathrm{}`$
Proof of Theorem 10. The induction on the genus starts with Lemma 16; then one can perform the next few steps by arguments similar to the one used in Lemma 14, and get the result for genus up to $`6`$. The procedure is then completed with Proposition 13.
$`\mathrm{}`$
## 7 A conjecture on higher degree tautological relations
At this point it is natural to formulate a conjecture which is suggested by the proof of Proposition 11.
This conjecture agrees with Harer’s and Ivanov’s stability theorems (see \[Ha\] and \[Iv\]), and with Faber’s results and conjectures concerning the tautological ring of the open part $`_{g,n}`$ (see \[Fa5\]).
###### Conjecture 18
There are no relations between essential tautological classes in $`H^{2k}\left(\overline{}_{g,P}\right)`$ whenever $`g3k`$.
We justify our conjecture. We first need to extend some definitions. A tautological class of degree $`2k`$ is a push-forward of a degree $`2l`$ Mumford class from a codimension $`kl`$ boundary component; a degree $`2k`$ class is unessential if it can be eliminated by means of a relation among tautological classes arising in degree $`<2k`$.
Then we need to build new pull-back formulas, but we can give conjectural ones starting from the ones we proved for degree $`4`$. In particular, we claim that they preserve the tautological group.
Under the above hypotheses, there are plenty of boundary components $`\overline{}_{\mathrm{\Gamma }_i}`$ in $`\overline{}_{g,P}`$ such that
$$H^{2k}(\overline{}_{\mathrm{\Gamma }_i})=_{{\scriptscriptstyle i}=2k}(H^i\left(\overline{}_{g_i,P_i}\right))$$
contains at least one summand $`H^{2j}\left(\overline{}_{g_j,P_j}\right)`$ with $`g_j3j`$, and $`g_j<g`$.
Write a generic linear combination of tautological classes in $`H^{2k}\left(\overline{}_{g,P}\right)`$, and suppose it is equal to $`0`$; by pulling back these relation to the above components we can prove that many coefficients do vanish: in fact, inductively, there are no relations among essential classes in these summands of the cohomology. We also conjecture that the pull-back maps in higher degree still satisfy the property that each class is generically a summand in the pull-back of at most one class.
It is then hard to believe that a new relation holds among the few classes whose coefficient has not yet been showed to be zero.
## 8 Generators of the cohomology group
###### Theorem 19
$`H^4(\overline{}_{g,P},)`$ is generated by tautological classes for all $`g8`$.
Proof. We are following Edidin’s scheme of Proof (\[Ed\]).
In the proof of this Proposition we plan to give an upper bound for the dimension of the cohomology group, and then to use the knowledge of the tautological group and of the homology of the mapping class group to prove that, this bound is achieved.
Let $`n=3g3+|P|`$ be the complex dimension of $`\overline{}_{g,P}`$. We write a part of the exact homology sequence of the pair $`(\overline{}_{g,P},\overline{}_{g,P}\backslash _{g,P})`$ :
$$\mathrm{}H_{2n4}\left(\overline{}_{g,P}\backslash _{g,P}\right)\stackrel{j_{}}{}H_{2n4}\left(\overline{}_{g,P}\right)H_{2n4}(\overline{}_{g,P},\overline{}_{g,P}\backslash _{g,P})\mathrm{}$$
hence, using Poincaré duality for smooth orbifolds:
$`dimH^4\left(\overline{}_{g,P}\right)`$ $`=`$ $`dimH_{2n4}\left(\overline{}_{g,P}\right)`$
$``$ $`dimj_{}H_{2n4}\left(\overline{}_{g,P}\backslash _{g,P}\right)+dimH^4\left(_{g,P}\right)`$
We refer to the description of the stratified structure of $`\overline{}_{g,P}`$ which has been explained in section 2. For any stable graph $`\mathrm{\Gamma }`$, we further denote by $`\mathrm{\Delta }_\mathrm{\Gamma }^0`$ the open stratum $`\xi _\mathrm{\Gamma }(_\mathrm{\Gamma })`$.
Let
$$_{g,P}=\overline{}_{g,P}\backslash _{g,P}.$$
We recall that $`_{g,P}=_i\mathrm{\Delta }_{\mathrm{\Gamma }_i}`$, where the $`\mathrm{\Delta }_{\mathrm{\Gamma }_i}`$’s are the codimension $`1`$ boundary components.
We denote by $`_{g,P}`$ the union of the codimension two boundary components, and write the homology exact sequence for the pair $`(_{g,P},_{g,P}):`$
$$\mathrm{}H_{2n4}\left(_{g,P}\right)\stackrel{i_{}}{}H_{2n4}\left(_{g,P}\right)H_{2n4}(_{g,P},_{g,P})\mathrm{}$$
Let us look at the relative term. By Lefschetz Theorem (\[Sp\]) we have:
$$H_{2n4}(_{g,P},_{g,P})H^2\left(_{g,P}\backslash _{g,P}\right).$$
The space
$$_{g,P}\backslash _{g,P}$$
consists of the disjoint union of the interior parts of the codimension $`1`$ boundary components, the $`\mathrm{\Delta }_i^0`$’s.
We have a precise description of these $`\mathrm{\Delta }_i^0`$’s as quotients of moduli spaces of smooth curves:
$$_{g,P}\backslash _{g,P}_{a,A}(_{a,A\{s\}}\times _{ga,A^c\{t\}})/Aut\mathrm{\Gamma }_{a,A}_{g1,P\{qr\}}/Aut\mathrm{\Gamma }_{irr}$$
The rational cohomology of such quotients satisfies:
$$H^k(_{\mathrm{\Gamma }_i}/Aut\mathrm{\Gamma }_i,)H^k\left(_{\mathrm{\Gamma }_i}\right)^{Aut\mathrm{\Gamma }_i}$$
where we denote by $`H^k\left(_{\mathrm{\Gamma }_i}\right)^{Aut\mathrm{\Gamma }_i}`$ the invariants with respect to the induced $`Aut\mathrm{\Gamma }_i`$ action on the cohomology. In the case $`k=2`$, these invariants can be precisely described. The cohomology group $`H^2\left(_{\mathrm{\Gamma }_i}\right)`$ is generated by Mumford classes of degree $`2`$. The class $`\kappa _1`$ is fixed by the automorphism group of any graph, whereas the $`\psi _i`$ classes, for $`i`$ a special point , are permuted by the group action in the obvious way.
We then get
$`H^2(_{g,P}\backslash _{g,P})`$ $``$ $`_{a,A}H^2(_{a,A\{s\}}\times _{ga,A^c\{t\}})^{Aut\mathrm{\Gamma }_{a,A}}H^2(_{g1,P\{q,r\}})^{Aut\mathrm{\Gamma }_{irr}}.`$
At this point, the bound for the dimension of the cohomology group is:
$`dimH^4\left(\overline{}_{g,P}\right)`$ $``$ $`+{\displaystyle \underset{a,A}{}}dimH^2(_{a,A\{s\}}\times _{ga,A^c\{t\}})^{Aut\mathrm{\Gamma }_{a,A}}`$
$`+dimH^2(_{g1,P\{q,r\}})^{Aut\mathrm{\Gamma }_{irr}}+dimH^4(_{g,P})+dimi_{}j_{}H_{2n4}\left(_{g,P}\right)`$
The space $`_{g,P}`$ is the union of the codimension two boundary components, which we will call $`\mathrm{\Theta }_i`$’s. Their complex dimension is $`n2`$. An easy application of the Maier-Vietoris exact sequence, shows that the obvious map
$$k:_i\mathrm{\Theta }_i_i\mathrm{\Theta }_i=_{g,P}$$
from the disjoint union into the union of these components induces the following isomorphism in homology:
$$_iH_{2n4}(\mathrm{\Theta }_i)H_{2n4}(_{g,P}).$$
Observe that for dimension reasons, $`dimH_{2n4}(\mathrm{\Theta }_i)=1`$.
We claim that
$$dimi_{}j_{}H_{2n4}\left(_{g,P}\right)r$$
where $`r`$ equals the number of essential pure boundary classes. This number differs from the number of codimension two boundary components because of the presence of Keel’s relations in genus $`0`$. These relations live in the second homology group of $`\overline{}_{0,n}`$.
The push-forward induced by the map
$$\overline{}_{0,A\{s\}}\overline{}_{g,P}$$
determines homological equivalences among codimension $`2`$ boundary components of $`\overline{}_{g,P}`$.
Let
$$\varphi :_i\mathrm{\Theta }_i\overline{}_{g,P}$$
be the collection of the inclusion maps of the codimension $`2`$ boundary components. By what we said above, the image of the map
$$\varphi _{}:H_{2n4}(_i\mathrm{\Theta }_i)H_4(\overline{}_{g,P})$$
has dimension less or equal than $`r`$. Since $`\varphi =kij`$, and $`k_{}`$ is an isomorphism, this implies that
$$dimi_{}j_{}H_{2n4}\left(_{g,P}\right)r.$$
Our final bound is:
$`dimH^4\left(\overline{}_{g,P}\right)`$ $``$ $`{\displaystyle \underset{a,A}{}}dimH^2(_{a,A\{s\}}\times _{ga,A^c\{t\}})^{Aut\mathrm{\Gamma }_{a,A}}`$
$`+dimH^2(_{g1,P\{qr\}})^{Aut\mathrm{\Gamma }_{irr}}`$
$`+dimH^4(_{g,P})+r`$
By Ivanov (\[Iv\]), Harer (\[Ha\]), and Loojenga’s (\[Lo\]) stability theorems for the homology of the mapping class group, $`H^4\left(_{g,P}\right)`$ is freely generated by Mumford classes, for $`g8`$.
Instead of computing the dimension of all the cohomology groups involved in (8), we proceed more indirectly. We show that there is a bijection between the following two sets. On one hand, the set $`_{g,P}^4`$, on the other, the set whose elements are the $`r`$ pure boundary classes in $`_{g,P}^4`$ and the vectors belonging to the natural bases of the cohomology vector spaces appearing on the right hand side of the above inequality (8). The upper bound for the dimension of the cohomology group is therefore achieved, and consequentely the tautological classes generate the cohomology group.
The bijection directly follows from the definition of essential tautological classes:
* pure Mumford classes in $`_{g,P}^4`$ correspond to a basis for
$$H^4(_{g,P}),$$
* mixed boundary classes in $`_{g,P}^4`$ correspond to a basis for
$$_{a,A}H^2(_{a,A\{s\}}\times _{ga,A^c\{t\}})^{Aut\mathrm{\Gamma }_{a,A}}H^2(_{g1,P\{qr\}})^{Aut\mathrm{\Gamma }_{irr}},$$
* pure boundary classes in $`_{g,P}^4`$ are exactly $`r`$.
This completes the proof.
$`\mathrm{}`$
Department of Mathematics
California Institute of Technology
91125 Pasadena, CA
polito@caltech.edu
|
warning/0005/math0005146.html
|
ar5iv
|
text
|
# Unirationality of cubic hypersurfaces
A remarkable result of \[Segre43\] says that a smooth cubic surface over $``$ is unirational iff it has a rational point. \[Manin72, II.2\] observed that similar arguments work for higher dimensional cubic hypersurfaces satisfying a certain genericity assumption over any infinite field. \[CT-S-SD87, 2.3.1\] extended the result of Segre to any normal cubic hypersurface (other than cones) over a field of characteristic zero. It is also clear that the result should hold for all sufficiently large finite fields, though the details were not worked out in general. \[Manin72, IV.8\] settles the cubic surface case for finite fields with at least 34 elements. The aim of this note is to observe that a variant of the Segre–Manin method works for all fields and for all cubics:
###### Theorem 1.
Let $`k`$ be a field and $`X^{n+1}`$ a smooth cubic hypersurface of dimension $`n2`$ over $`k`$. Then the following are equivalent:
1. $`X`$ is unirational (over $`k`$).
2. $`X`$ has a $`k`$-point.
Similar results hold for singular cubic hypersurfaces, with a few exceptions.
###### Theorem 2.
Let $`k`$ be a perfect field and $`X^{n+1}`$ an irreducible cubic hypersurface of dimension $`n2`$ over $`k`$ which is not a cone over an $`(n1)`$-dimensional cubic. Then the following are equivalent:
1. $`X`$ is unirational (over $`k`$).
2. $`X`$ has a $`k`$-point.
3. $`X`$ has a smooth $`k`$-point.
###### 3Nonperfect fields.
Over nonperfect fields of characteristic 3, there are nonsingular cubic hypersurfaces of arbitrary dimension which are not unirational but do have a $`k`$ point (17). I have not been able to find examples in characteristic 2.
###### Question 4.
Unirationality of varieties is very poorly understood in general and there are very basic open questions. We do not even have a list of unirational surfaces and very few examples are known in higher dimensions. For instance, let $`X`$ be a smooth projective variety over $`k`$ such that $`X`$ is unirational over $`\overline{k}`$. Assume for simplicity that $`k`$ is infinite and consider the following properties:
1. $`X`$ is unirational (over $`k`$).
2. $`X(k)`$ is dense in $`X`$.
3. $`X`$ has a $`k`$-point.
It is clear that each property implies the next. They are equivalent for cubic hypersurfaces by (1). It is extremely unlikely that they are always equivalent, but no counter examples are known.
###### 5Proof of (2.1) $``$ (2.2) $``$ (2.3).
(2.1) $``$ (2.2) is clear for infinite fields. For finite fields it follows from \[Nishimura55\].
Assume (2.2) and let $`xX`$ be a $`k`$-point. We are done if $`X`$ is smooth at $`x`$. Otherwise $`x`$ is a double point and we can choose affine coordinates such that $`x=(0,\mathrm{},0)`$ and $`X`$ is given by an equation $`q(x_1,\mathrm{},x_n)+c(x_1,\mathrm{},x_n)=0`$ where $`q`$ is quadratic and $`c`$ is cubic. Assume that there is a point $`(p_1,\mathrm{},p_n)k^n`$ such that $`q(p_1,\mathrm{},p_n)0`$. Then the line connecting the origin and $`(p_1,\mathrm{},p_n)`$ intersects $`X`$ in a single point outside the origin and this is a smooth $`k`$-point of $`X`$. Thus we are done unless $`q`$ vanishes everywhere on $`k^n`$.
However, if a homogeneous polynomial $`f`$ of degree $`d`$ vanishes on $`k^n`$ and $`|k|d`$ then $`f`$ is identically zero. ∎
The intersting part is to show unirationality starting with a smooth $`k`$-point. The construction is presented in 3 stages, successive version working in greater and greater generality. At least in retrospect, all of this is only a slight modification of the works of Segre.
###### 6First unirationality construction.
Let $`X^{n+1}`$ be a cubic and $`pX`$ a point. Let $`C_p`$ denote the intersection of $`X`$ with the tangent plane at $`p`$. We expect that usually $`C_p`$ is an irreducibe cubic with a double point at $`p`$. If this is indeed the case then the inverse of the projection from $`p`$ gives a birational map $`\pi _p:^{n1}C_p`$. If $`pX(k)`$ then $`C_p`$ is birational to $`^{n1}`$ over $`k`$.
Assume next that we have two points $`p,qX`$ and $`C_p,C_q`$ are both irreducible with a double point at $`p`$ (resp. $`q`$). Define the “3rd intersection point” map
$$\varphi :C_p\times C_qX$$
as follows. Take $`uC_p,vC_q`$. If the line connecting $`u,v`$ is not contained in $`X`$, it has a unique 3rd intersection point with $`X`$; call it $`\varphi (u,v)`$. Under very mild genericity assumptions (15) this is a well defined dominant map. Thus we get that $`X`$ is unirational via
$$\mathrm{\Phi }:^{n1}\times ^{n1}\stackrel{\pi _p\times \pi _q}{}C_p\times C_q\stackrel{\varphi }{}X.$$
###### Definition 7 (Restriction of scalars).
Let $`L/K`$ be a finite degree field extension. Restriction of scalars (or Weil restriction) is way to associate to an $`L`$-variety $`U`$ a $`K`$-variety $`_{L/K}U`$ such that there is a natural identification of the $`L`$-points of $`U`$ with the $`K`$-points of $`_{L/K}U`$. This dictates that $`dim_{L/K}U=\mathrm{deg}(L/K)dimU`$; see \[BLR90, 7.6\] for details.
This can be done very explicitly in the affine case as follows. Let $`U𝔸^n`$ be an affine variety defined over $`L`$. Choose equations $`f_i(x_1,\mathrm{},x_n)`$ for $`U`$ and let $`e_1,\mathrm{},e_dL`$ be a $`K`$-basis. Choose new coordinates $`y_{ij}:i=1,\mathrm{},n,j=1,\mathrm{},d`$ and set $`x_i=_je_jy_{ij}`$. We can then write
$$f_k(x_1,\mathrm{},x_n)=\underset{\mathrm{}}{}e_{\mathrm{}}f_k\mathrm{}(y_{ij})\text{where }f_k\mathrm{}K[y_{ij}]\text{.}$$
Let $`_{L/K}U`$ be the subvariety of $`𝔸^{nd}`$ defined by the equations $`f_k\mathrm{}=0`$.
In particular we see that $`_{L/K}(𝔸^n)𝔸^{dn}`$. In the projective case, $`_{L/K}(^n)`$ is inconvenient to describe by explicit equations but we at least get that $`_{L/K}(^n)`$ is birational to $`^{dn}`$ over $`K`$. (They are not isomorphic for $`d>1`$.)
###### 8Second unirationality construction.
Let $`X^{n+1}`$ be a cubic defined over $`k`$ and $`k^{}k`$ a quadratic extension. Let $`pX(k^{})`$ be a point and $`\overline{p}X(k^{})`$ its conjugate. (Let us ignore that $`k^{}k`$ may be inseparable in characteristic 2.) We have conjugate birational maps $`\pi _p:^{n1}C_p`$ and $`\pi _{\overline{p}}:^{n1}C_{\overline{p}}`$. If $`u^{n1}(k^{})`$ then $`\pi _p(u)`$ and $`\pi _{\overline{p}}(\overline{u})`$ are conjugate points of $`X`$, thus the line connecting them is defined over $`k`$. Hence $`\varphi (u,\overline{u})X(k)`$. Putting this invariantly, we obtain a rational map (defined over $`k`$)
$$\mathrm{\Phi }:_{k^{}/k}^{n1}X$$
which is dominant under mild genericity assumptions.
###### 9Final unirationality construction.
Assume now that $`X`$ is a cubic defined over $`k`$ and $`xX`$ is a smooth $`k`$-point. Let $`L`$ be a line through $`x`$. If $`L`$ is not contained in $`X`$ then it intersects $`X`$ in a point pair $`\{p,q\}`$. These points are usually not in $`k`$, but they are conjugate over $`k`$ and lie in a quadratic extension $`k^{}=k^{}(L)`$ of $`k`$. Hence, under some genericity assumptions, we obtain a dominant map
$$\mathrm{\Phi }:_{k^{}/k}^{n1}X$$
which shows that $`X`$ is unirational. There are very few problems if $`k`$ is infinite, since then a general choice of $`L`$ should work. (Extra work is needed in characteristic 2.) The situation is less clear over finite fields since there may not be enough room to choose $`L`$ general; see, for example, (16).
To avoid this difficulty, we do not choose any line, rather we work with all lines simultaneously. We should obtain a map
$$\mathrm{\Psi }:_{xL^{n+1}}_{k^{}(L)/k}^{n1}X.$$
We are in good shape if we can indentify the left hand side with a product $`^n\times ^{n1}\times ^{n1}`$, at least birationally. Once this problem is settled, it is enough to check dominance over the algebraic closure where the previous arguments work. It seems best to give an explicit algebraic description.
###### 10Algebraic description of $`\mathrm{\Psi }`$.
We work in affine coordinates, assuming that the origin is a smooth point of $`X`$. Thus the equation of $`X`$ can be written as
$$F=L(x_1,\mathrm{},x_{n+1})+Q(x_1,\mathrm{},x_{n+1})+C(x_1,\mathrm{},x_{n+1})$$
where $`L`$ is linear, $`Q`$ is quadratic and $`C`$ is cubic. We may assume that $`F/x_{n+1}`$ is not indetically zero (for instance we can even assume that $`L=x_{n+1}`$).
We write down a rational map
$$\mathrm{\Psi }:𝔸^{3n2}(u_1,\mathrm{},u_n,v_1,\mathrm{},v_{n1},w_1,\mathrm{},w_{n1})X.$$
Later we check that it is dominant with a few exceptions.
Consider the universal line through the origin $`(\tau u_1,\mathrm{},\tau u_n,\tau )`$. It intersects $`X`$ in two further points which correspond to the roots of the quadratic equation
$$L(u_1,\mathrm{},u_n,1)+\tau Q(u_1,\mathrm{},u_n,1)+\tau ^2C(u_1,\mathrm{},u_n,1)=0.$$
The equation is irreducible if $`X`$ is irreducible. Let its roots be $`t_1,t_2\overline{k(u_1,\mathrm{},u_n)}`$.
The equation of the tangent space of $`X`$ at $`𝐩=(p_1,\mathrm{},p_{n+1})X`$ is
$$\frac{F}{x_1}(𝐩)(x_1p_1)+\mathrm{}+\frac{F}{x_{n+1}}(𝐩)(x_{n+1}p_{n+1})=0.$$
Thus the universal tangent line at $`(t_1u_1,\mathrm{},t_1u_n,t_1)`$ can be described parametrically as
$$\begin{array}{c}x_1=t_1u_1+\sigma (v_1+t_1w_1),\mathrm{},x_{n1}=t_1u_{n1}+\sigma (v_{n1}+t_1w_{n1})\hfill \\ x_n=t_1u_n+\sigma ,\hfill \\ x_{n+1}=t_1\sigma \left(\frac{F}{x_{n+1}}(t_1𝐮,t_1)\right)^1_{i=1}^n\frac{F}{x_i}(t_1𝐮,t_1)(v_i+t_1w_i),\hfill \end{array}$$
where we set $`v_n=1,w_n=0`$. Substituting the above parametric representation into $`F`$, we obtain a cubic equation in $`\sigma `$
$$\underset{j=0}{\overset{3}{}}\sigma ^jH_j\text{where}H_jk(𝐮,𝐯,𝐰,t_1).$$
$`H_0=H_1=0`$ since we have a tangent line, thus the 3rd intersection point corresponds to the value $`\sigma =H_2/H_3`$. Thus we obtain a point
$$Q_1k(𝐮,𝐯,𝐰,t_1)^{n+1}.$$
Replacing $`t_1`$ by its conjugate $`t_2`$ we obtain another point $`Q_2`$. The line connecting $`Q_1`$ and $`Q_2`$ can be given parametrically as
$$L(\lambda )=\frac{\lambda t_2}{t_1t_2}Q_1+\frac{\lambda t_1}{t_2t_1}Q_2,$$
and this is a parametrization over $`k(𝐮,𝐯,𝐰)`$. Evaluating $`F`$ on the line we have that $`t_1,t_2`$ are roots, so
$$F(L(\lambda ))=(A\lambda +B)(C\lambda ^2+Q\lambda +L).$$
Thus if we expand
$$F(L(\lambda ))=\underset{j=0}{\overset{3}{}}\lambda ^jG_j,\text{then}G_jk(𝐮,𝐯,𝐰),$$
and the 3rd root is
$$\frac{B}{A}=\frac{G_2}{G_3}+\frac{Q}{C}.$$
Substituting this into the parametrization of the line gives
$$\mathrm{\Psi }(𝐮,𝐯,𝐰)X(k(𝐮,𝐯,𝐰)).$$
Depending on our definition of unirationality, we also need to check the following:
###### Lemma 11.
For a $`k`$-variety $`X`$ the following are equivalent:
1. There is a dominant map $`\varphi _m:𝔸^mX`$ for some $`m`$.
2. There is a dominant map $`\varphi _m:𝔸^mX`$ for $`m=dimX`$.
Proof. Assume that $`m>dimX`$. There is a dense open set $`U𝔸^m`$ such that $`\varphi _m|_U`$ is open with $`mdimX`$ dimensional fibers. Let $`𝐮U`$ be a point. If $`𝐮Z𝔸^m`$ is a hypersurface which does not contain the irreducible component of the fiber of $`\varphi _m|_U`$ through $`𝐮`$, then $`\varphi _m|_Z:ZX`$ is dominant.
If $`k`$ is infinite and $`m>dimX`$ then we can choose $`Z`$ to be a general hyperplane.
Assume next that $`k`$ is finite. Fix a prime $`\mathrm{}\mathrm{char}k`$ and let $`k^{}`$ be the composite of all algebraic extensions of degree $`\mathrm{}^s`$ of $`k`$. $`k^{}`$ is infinite, hence we can choose a point $`𝐮=(u_1,\mathrm{},u_m)U(k^{})`$. By permuting the coordinates we may assume that $`\mathrm{deg}k(u_m)/k\mathrm{deg}k(u_1)/k`$, or, equivalently, $`k(u_m)k(u_1)`$. This implies that $`u_m`$ can be written as a polynomial of $`u_1`$, hence the ideal $`I(𝐮)k[x_1,\mathrm{},x_m]`$ contains an polynomial of the form $`x_mp(x_1)`$. This implies that $`I(𝐮)`$ is generated by polynomials of the form $`x_mP(x_1,\mathrm{},x_{m1})`$. Thus we can choose $`Z=(x_m=P(x_1,\mathrm{},x_{m1}))`$ for suitable $`P`$.∎
###### 12Proof of (2.3) $``$ (2.1).
Let $`X`$ be an irreducible cubic hypersurface. The set of all triple points of $`X(\overline{k})`$ is a linear space and it is defined over $`k`$ if $`k`$ is perfect. Thus $`X`$ is a cone over a cubic hypersurface without triple points. Therefore, $`X`$ has no triple points over $`\overline{k}`$.
Assume next that $`X`$ is not normal. The nonnormal locus has dimension $`(n1)`$ and the linear space spanned by it is in $`X`$. Thus the nonnormal locus is a linear space $`L^{n1}^{n+1}`$ which is defined over $`k`$ if $`k`$ is perfect. Projecting form $`L`$ realizes $`X`$ as a $`^{n1}`$-bundle over $`^1`$, hence rational.
For the rest of the proof assume that $`X`$ is normal. We need to check three conditions.
First we prove that $`C_x`$ is irreducible with a double point at $`x`$ for general $`xX(\overline{k})`$. This is done in (14).
Second, we need to check that the 3rd intersection point map $`\varphi :C_p\times C_qX`$ is dominant. It is, however, not enough to check this for a general pair $`p,q`$. In our construction $`p,q`$ are the two intersection points of a line through $`x`$, hence dependent. Assume that $`\pi _x:X^n`$, the projection from $`x`$, is separable. Then for a generic line $`xL`$ we get 2 distinct intersection points and both intersections are transverse. In particular, the tangent space of $`X`$ at one point does not contain the other point. In (15) we see that this is sufficient to guarantee that $`\varphi :C_p\times C_qX`$ is dominant.
Third, we need to consider the case when the projection $`\pi _x:X^n`$ is inseparable. This can happen only in characteristic 2. Over a perfect field a purely inseparable map induces a purely inseparable map in the reverse direction, hence in this case $`X`$ is (purely inseparably) unirational. Nonetheless, we check in (20) that we can always choose a smooth $`k`$-point such that projection from it is separable.∎
By looking at the proof we obtain the following for nonperfect fields. We check in (18) that (13.3) is satisfied for $`X`$ smooth. This shows that our proof covers all smooth hypersurfaces.
###### Proposition 13.
Let $`k`$ be a nonperfect field and $`X^{n+1}`$ a cubic hypersurface which is not a cone. Then the 3 parts of (2) are equivalent if one of the following conditions holds:
1. $`\mathrm{char}k5`$.
2. $`\mathrm{char}k=3`$ and $`X`$ has no triple points over $`\overline{k}`$.
3. $`\mathrm{char}k=2`$ and there is a smooth $`k`$ point $`pX`$ such that projection from $`p`$ is separable.∎
###### Proposition 14.
Let $`k`$ be an algebraically closed field and $`X^{n+1}`$ a normal cubic hypersurface over $`k`$ without triple points. Then $`C_x`$ is irreducible with a double point at $`x`$ for general $`xX`$.
Proof. Let $`xX`$ be arbitrary. If $`C_x`$ is irreducible with a triple point at $`x`$ then $`C_x`$ is a cone, hence there is an $`(n2)`$-dimensional family of lines through $`x`$. If $`C_x`$ is reducible then either $`C_x`$ contains an $`(n1)`$-dimensional linear space through $`x`$ or a quadric cone with vertex at $`x`$. In either case, there is an $`(n2)`$-dimensional family of lines through $`x`$. Thus it is enough to prove that for a general $`xX`$ the family of lines in $`X`$ through $`x`$ has dimension at most $`n3`$. This is equivalent to proving that a general surface section through $`x`$ has no lines through $`x`$.
$`X`$ has no triple points, hence by Bertini, a general surface section of $`X`$ is also normal with no triple points.
If $`S`$ is a normal cubic surface without triple points then there are only finitely many lines through each double point. (Choose affine coordinates such that the equation becomes $`q(x_1,x_2,x_3)+c(x_1,x_2,x_3)=0`$. The lines through $`(0,0,0)`$ correspond to the solutions of $`(q=c=0)^2`$. If there are infinitely many solutions, then $`q`$ and $`c`$ have a common factor, thus the surface is reducible.) Each line in the smooth locus has selfintersection $`1`$, hence rigid. Thus $`S`$ has only finitely many lines.∎
###### Lemma 15.
Let $`X^{n+1}`$ be an irreducible cubic hypersurface. Let $`x,yX`$ be smooth points and $`C_x,C_y`$ the corresponding intersections with the tangent hyperplanes. Assume that
1. $`C_x`$ and $`C_y`$ are irreducible.
2. $`xC_y`$ and $`yC_x`$.
Then the 3rd intersection point map $`\varphi :C_x\times C_yX`$ is dominant.
Proof. Let us see first that $`\varphi `$ is indeed defined. Pick a point $`uC_x`$ which is a smooth point of $`X`$. Pick $`vC_y`$ which is a smooth point of $`X`$ such that $`v`$ does not lie on $`T_uX`$. If we now choose a general $`wC_x`$ then $`v`$ does not lie on $`T_wX`$ and $`w`$ does not lie on $`T_vX`$. Thus the line connecting $`u,w`$ has a unique third intersection point with $`X`$. This shows that $`\varphi `$ is defined at the pair $`(v,w)`$.
In order to prove dominance, we need to show that $`\varphi `$ has at least one fiber of dimension $`n2`$. Pick a point $`zX`$ which is not on $`C_xC_y`$ and let $`\pi :^{n+1}T_yX`$ denote the projection from $`z`$. Then $`\varphi ^1(z)`$ is the set of pairs $`(v,w)`$ such that $`\pi (v)=w`$. Thus we are done if
$$dim((C_y\pi (C_x))(C_yC_x))=n2.$$
For this it is sufficient to find one projection $`\pi ^{}:^{n+1}T_yX`$ where this holds. Then the same holds for a general projection and a general projection always corresponds to a point of $`X`$. Pick any smooth point $`vC_y`$ and let $`\pi ^{}`$ be a projection such that $`\pi ^{}(x)=v`$. Then $`\pi ^{}(C_x)`$ and $`C_y`$ intersect at $`v`$ but they have different multiplicty there. Hence their intersection has dimension $`n2`$.∎
###### Example 16.
Let $`S`$ be the cubic surface $`(x_0^3+x_1^3+x_2^3+x_3^3=0)`$. By \[Hirschfeld81\] over the fields $`𝔽_2,𝔽_4,𝔽_{16}`$ all the points are on the 27 lines. Hence the second unirationality construction does not work over $`𝔽_2`$ and $`𝔽_4`$. (It does work over $`𝔽_{16}`$.)
###### Example 17.
Let $`k`$ be a field of characteristic 3 and $`t_i`$ algebraically independent over $`k`$. Set $`K=k(t_1,\mathrm{},t_n)`$ and
$$Y:=(y^3yz^2=\underset{i=1}{\overset{n}{}}t_ix_i^3)^{n+1}.$$
1. $`Y`$ is non–singular.
2. Over $`\overline{K}`$, $`Y`$ is a cone over a cuspidal cubic curve.
3. $`Y(K)=\{(0,1,0,\mathrm{},0),(1,1,0,\mathrm{},0),(1,1,0,\mathrm{},0)\}`$.
4. $`Y`$ is not unirational (over $`K`$).
Proof. $`Y`$ is the generic fiber of the smooth variety
$$(y^3yz^2=\underset{i}{}t_ix_i^3)𝔸_{(t_1,\mathrm{},t_n)}^n\times _{(y,z,x_1,\mathrm{},x_n)}^{n+1}$$
over $`𝔸^n`$, thus $`Y`$ is non–singular. (2) holds since over $`\overline{K}`$ we can write our equation as
$$(y\underset{i}{}\sqrt[3]{t_i}x_i)^3yz^2=0.$$
In order to see (3) we may as well assume that $`k`$ is algebraically closed. Assume that we have relatively prime polynomials $`f,g,h_ik[t_1,\mathrm{},t_n]`$ such that
$$f^3fg^2=\underset{i}{}t_ih_i^3.$$
We are done if $`h_1=\mathrm{}=h_n=0`$. Otherwise, we can make a substitution $`t_i=c_it`$ for $`i=1,\mathrm{},n`$ and general $`c_i`$ to get a solution of
$$f(fg)(f+g)=th^3\text{with }f,g,hk[t]\text{ and }h0\text{.}$$
We may assume that $`f`$ and $`g`$ are relatively prime. Thus 2 of the factors $`f,fg,f+g`$ are cubes and the third is $`t`$ times a cube. However, $`f+(fg)+(f+g)=0`$, hence if 2 are cubes then so is their sum which is minus the 3rd factor. This is a contradiction.
Since $`Y`$ has only 3 points in $`K`$, it does not contain any rational curves and so it is definitely not unirational. ∎
###### Lemma 18.
Assume that $`\mathrm{char}k=2`$. Let $`V^{n+1}`$ be the linear span of all points $`pX(k)`$ such that projection from $`p`$ is a purely inseparable map $`X^n`$. Let $`(y_i=0)`$ be equations of $`V`$ and $`x_j`$ coordinates on $`V`$. Then the equation of $`X`$ can be written as
$$f:=\underset{j}{}\mathrm{}_j(𝐲)x_j^2+g(𝐲)$$
where the $`\mathrm{}_j`$ are linear and $`g`$ is cubic. If $`V\mathrm{}`$ then $`X`$ is not smooth.
Proof. We can choose coordinates such that the points
$$p_1=(1:0:\mathrm{}:0),\mathrm{},p_m=(0:\mathrm{}:\stackrel{mth}{1}:0:\mathrm{}:0)$$
are in $`X(k)`$ and projection from $`p_i`$ is a purely inseparable map for $`i=1,\mathrm{},m`$. $`p_i`$ is inseparable iff $`x_i`$ occurs in the equation of $`X`$ always with even exponent. This gives the above equation.
$`f/x_j`$ is zero, and the equations $`f/y_i=0`$ have a common solution. Since $`f=3f=_i(f/y_i)`$, these give singular points of $`X(\overline{k})`$. ∎
###### Lemma 19.
Let $`k`$ be a perfect field of characteristic 2. Let $`X`$ be a cubic of dimension at least 2 given by an equation
$$f(𝐱,𝐲):=\underset{j}{}\mathrm{}_j(𝐲)x_j^2+g(𝐲).$$
Then $`X`$ has a smooth $`k`$-point with nonzero $`y`$-coordinate.
Proof. Assume first that we have at least two $`x`$-variables. If $`\mathrm{}_1=c\mathrm{}_2`$ then
$$\mathrm{}_1x_1^2+\mathrm{}_2x_2^2=\mathrm{}_1(x_1+\sqrt{c}x_2)^2$$
thus we can change coordinates to eliminate one $`x`$-variable. Otherwise we can pick $`𝐲_0`$ such that $`\mathrm{}_1(𝐲_0)0`$ and $`\mathrm{}_2(𝐲_0)=0`$. Then
$$(\sqrt{g(𝐲_0)/\mathrm{}_1(𝐲_0)},x_2,𝐲_0)X(k)$$
for every $`x_2`$. Since
$$\frac{f}{y_i}=x_2^2+\frac{g}{y_i}$$
the above point is smooth for suitable choice of $`x_2`$.
Thus assume that there is only one $`x`$-variable and write the equation as $`y_1x_1^2+g(𝐲)`$.
Take any $`(p_1,\mathrm{},p_n)k^n`$. If $`p_1=0`$ and $`g(p_1,\mathrm{},p_n)=0`$ then $`(x_1:p_1:\mathrm{}:p_n)X(k)`$ for any $`x_1`$ and one of them is a smooth by looking at $`f/y_1`$.
If $`p_10`$ then $`p_0:=\sqrt{g(p_1,\mathrm{},p_n)/p_1}k`$ and $`(p_0:p_1:\mathrm{}:p_n)`$ is a smooth point unless
$$g(p_1,\mathrm{},p_n)p_1\frac{g}{y_1}(p_1,\mathrm{},p_n)=0.$$
Thus we are done unless the following holds:
1. $`gy_1(g/y_1)`$ is nonzero for $`y_1=0`$, and
2. $`gy_1(g/y_1)`$ is zero for $`y_10`$.
Write $`g=y_1^ig_{3i}(y_2,\mathrm{},y_n)`$. Then
$$gy_1(g/y_1)=y_1^2g_1(y_2,\mathrm{},y_n)+g_3(y_2,\mathrm{},y_n).$$
$`g_1`$ is a linear form thus it has a nontrivial zero $`(p_2,\mathrm{},p_n)`$. If $`g_3(p_2,\mathrm{},p_n)=0`$ then set $`p_1=0`$ and if $`g_3(p_2,\mathrm{},p_n)0`$ then set $`p_1=1`$.∎
Combining the above lemmas we obtain:
###### Corollary 20.
Let $`k`$ be a perfect field of characteristic 2 and $`X^{n+1}`$ a cubic with a smooth $`k`$-point. Assume that $`n2`$. Then there is a smooth point $`xX(k)`$ such that the projection from $`x`$ is separable.∎
###### Acknowledgments .
I thank J.-L. Colliot-Thélène and J. Ellenberg for helpful comments and references. Partial financial support was provided by the NSF under grant number DMS-9970855.
Princeton University, Princeton NJ 08544-1000
```
kollar@math.princeton.edu
```
|
warning/0005/physics0005060.html
|
ar5iv
|
text
|
# Sub-Doppler resolution with double coherently driving fields
## I Introduction
Multi-photon transitions are embedded in a variety of optical effects based on correlated absorption and emission of two or more photons. Step-wise and multi-photon processes can be distinguished by their frequency-correlation properties . Sub-Doppler resolution for an inhomogeneously broadened medium based on multi-photon processes has attracted much attention in recent decades. The physical essence for this method lies in the fact that a detuning for multi-photon transition stipulated by Doppler frequency shift can be reduced or eliminated by adopting appropriate light propagation direction, since the detuning is the sum or difference of multiple single-photon transition detunings. Sub-Doppler nonlinear optical resonance as appearance of quantum coherence and interference in the context of frequency-correlation properties of the coherent components have been proposed in . A widely used method is Doppler-free two- or multi-photon absorption where the velocity dependent detuning could be removed in the case that the wave vectors of the interacting beams sum down to zero (for a review, see ). This type of sub-Doppler spectrum is characterized by a large detuning from the intermediate resonance, measuring the fluorescence from the upper level as well as by small cross sections for these multi-photon processes. Enhanced by intermediate resonance sub-Doppler processes based on the use of strong driving fields and on a change of frequency-correlation properties of resonant multi-photon process in strong resonant fields, have been proposed in and further developed in (also in ). Substantial enhancement in absorption, gain and fluorescence controlled with the auxiliary appropriately propagating electro-magnetic wave has been predicted. An interference nature of the spectrum modification has been stressed, which implies that along the growth of the absorption (gain) in certain spectral intervals, integral over the frequency may even decrease. A role of increase of intensity of a probe field as well as features of Doppler-free lasers were explored too.
Control of atomic response with intense coupling lasers has been a subject of many intensive studies in the context of electromagnetically induced transparency (EIT) , amplification without inversion (AWI) , enhancement of the refractive index without absorption and so on (for review see ). Sub-Doppler resolution by using intense coherently driving field(s) at other transition(s) has been recently further explored in . In Ref. probe weak field absorption spectrum in three-level schemes was considered. It was pointed out that the linewidth of one of the two Autler-Townes absorption peaks can be reduced by match of the coupling field intensity and frequency. In a later work , an experimental observation of sub-Doppler linewidth in a Doppler-broadened $`\mathrm{\Lambda }`$-type Rb atomic system was reported. Related works to reduce Doppler broadening with atomic coherence effects could be found in Ref. and multiple sub-Doppler lines have been shown to achieve with a strong coupling field and a saturation effects in a three-level system.
In the recent papers , cancellation of Doppler broadening by applying different frequency fields was proposed with one or two coherently driving lasers. As it has been outlined, sub-Doppler structures can be induced without population redistribution. At that not only absorption and gain, but four-wave mixing output can be enhanced through Doppler-free coupling controlled with auxiliary co- or counter-propagating driving electromagnetic radiation.
In this paper, we propose a four-level scheme, where both upper and lower levels of a probe transition are coupled to other levels by strong coherent fields. In this scheme, four absorption peaks could be found on account of the fact that both the upper and lower states are split into Autler-Townes doublets. Sub-Doppler resolution is achievable because the modified two-photon transition occurs contributing in the probe process and two-photon detunings, resulting from the Doppler shift, could be reduced by choosing appropriate optical geometry. A crucial role of the ratio of the frequencies of the coupled transitions is outlined. The results obtained on the basis of the developed theory are accompanied with numerical simulation addressed to real experimental schemes. Despite the decreased integral intensity, we predict enhanced laser-induced sub-Doppler absorption peaks along with transparency windows. The features are interpreted in the terms of quantum coherence and interference processes.
The paper is organized into four sections. In the next section, we present a four-level model and the density-matrix equations describing the system. The solution in the limit of a weak probe and steady-state condition is obtained and the absorption spectrum is then analyzed, first without Doppler broadening. In Sec. III, taking into consideration Doppler broadening, we demonstrate that the multiple sub-Doppler lines as well as enhanced absorption at the resonance could be obtained in this scheme. From numerical illustrations of the effect in a practical medium, the conditions for and features of Doppler-free resonance, stipulated by compensation of Doppler shifts with light shifts are discussed. In Sec. IV, we summarize the results.
## II Model and absorption spectrum at homogeneously broadened probe transition
We consider a closed four-level scheme shown in Fig.1(a). In this scheme two coherent driving fields $`E_c`$ and $`E_d`$ with coupling Rabi frequencies $`\mathrm{\Omega }_c=𝐄_𝐜𝐝_{\mathrm{𝟏𝟑}}/2\mathrm{}`$ and $`\mathrm{\Omega }_d=𝐄_𝐝𝐝_{\mathrm{𝟐𝟒}}/2\mathrm{}`$ interact with the transitions labeled $`|4|2`$ and $`|3|1`$, respectively. The transition $`|4|1`$ is probed by the weak nonperturbating field $`E_p`$. Absorption index $`\alpha _p`$ for this field, reduced by it’s value $`\alpha _{p0}^0`$ at $`\omega _p=\omega _{41}`$, $`\mathrm{\Omega }_c=\mathrm{\Omega }_d=0`$ is found as $`\alpha _p/\alpha _{p0}^0=Re\{\chi _p/\chi _{p0}^0\}`$, where $`\chi _p`$ is corresponding dressed susceptibility. The later is convenient to calculate with aid of dressed nonlinear polarization $`P^{NL}`$ and density matrix $`\rho _{ij}`$ as $`P^{NL}(\omega _p)=N\chi _pE_p/2=N\rho _{ij}d_{ji}`$ ($`d_{ji}`$ is transition electrodipole moment, N - number density of atoms). In the framework of semi- classical theory and using the standard density matrix formalism with the rotating wave approximation, the description equations of this scheme can be written in general form as:
$`L_{nn}\rho _{nn}=q_ni[V,\rho ]_{nn}+{\displaystyle \underset{m>n}{}}\gamma _{mn}\rho _{mm};L_{14}\rho _{14}=L_p\rho _p=i[V,\rho ]_{14}(etc.),`$ (1)
where $`L_{ij}=d/dt+\mathrm{\Gamma }_{ij};V_{14}=\mathrm{\Omega }_p\mathrm{exp}\{i\mathrm{\Delta }_pt\};\mathrm{\Omega }_p=𝐄_𝐩𝐝_{\mathrm{𝟏𝟒}}/2\mathrm{}`$; $`\mathrm{\Delta }_p=\omega _p\omega _{41}`$ (etc.) are frequency detunings from the corresponding resonance; $`\mathrm{\Gamma }_{ij}`$ \- homogeneous half-widths of transitions (in absence of collisions $`\mathrm{\Gamma }_{mn}=(\mathrm{\Gamma }_m+\mathrm{\Gamma }_n)/2`$); $`\mathrm{\Gamma }_n=_j\gamma _{nj}`$ \- inverse lifetimes of levels; $`\gamma _{mn}`$ \- rate of relaxation from the level $`m`$ to $`n`$, $`q_n=_jw_{nj}r_j`$ \- rate of incoherent exitation to a state $`n`$ from the underlying levels.
We represent density matrix elements as $`\rho _{14}=r_p\mathrm{exp}\{i\mathrm{\Delta }_pt\},\rho _{33}=r_3`$ etc. Then in a steady-state regime a set of density-matrix equation may be reduced to the set of algebraic equations for the amplitudes $`r_{ij}`$. Analytical solution of this system of 10 coupled equation both for closed and open four-level scheme is given in . In the case under consideration incoherent exitation of the upper levels as well as relaxation from the level 3 to 2 are supposed to be negligible small. Then $`r_2=r_4=r_c=0`$ and equation system (1) reduces to
$`P_pr_p=i\mathrm{\Omega }_pr_1+ir_{12}\mathrm{\Omega }_ci\mathrm{\Omega }_dr_{34},P_{12}r_{12}=ir_p\mathrm{\Omega }_c^{}i\mathrm{\Omega }_dr_{32},`$ (2)
$`P_{34}r_{34}=i\mathrm{\Omega }_d^{}r_p+ir_d^{}\mathrm{\Omega }_p+ir_{32}\mathrm{\Omega }_c,P_{23}r_{23}=i\mathrm{\Omega }_cr_{43}+ir_{21}\mathrm{\Omega }_d.`$ (3)
$`P_dr_d=i\mathrm{\Omega }_d(r_1r_3),\mathrm{\Gamma }_3r_3=2Re\{i\mathrm{\Omega }_d^{}r_d\},r_1=1r_3.`$ (4)
Here $`P_{14,13}P_{p,d}=\mathrm{\Gamma }_{p,d}+i\mathrm{\Delta }_{p,d}`$, $`P_{12}=\mathrm{\Gamma }_{12}+i(\mathrm{\Delta }_p\mathrm{\Delta }_c)`$, $`P_{34}=\mathrm{\Gamma }_{34}+i(\mathrm{\Delta }_p\mathrm{\Delta }_d)`$, $`P_{23}=\mathrm{\Gamma }_{23}+i(\mathrm{\Delta }_c\mathrm{\Delta }_p+\mathrm{\Delta }_d)`$. For atom moving with speed $`\mathrm{v}`$, Doppler shift of resonances must be taken into account by substituting $`\mathrm{\Delta }_j`$ for $`\mathrm{\Delta }_j^{^{}}=\mathrm{\Delta }_j𝐤_j𝐯`$.
Under considered approximation solution for two level system $`|1|3`$ is found apart of the other elements.
$`r_d=i{\displaystyle \frac{\mathrm{\Omega }_dP_d^{}}{\mathrm{\Gamma }_d^2(1+\text{æ}_d)+\mathrm{\Delta }_d^2}},r_3={\displaystyle \frac{\mathrm{\Gamma }_d^2\text{æ}/2}{\mathrm{\Gamma }_d^2(1+\text{æ})+\mathrm{\Delta }_d^2}},\text{æ}={\displaystyle \frac{4|\mathrm{\Omega }_d|^2}{\mathrm{\Gamma }_3\mathrm{\Gamma }_d}}.`$ (5)
Dressed susceptibility for the probe field is found in the notations of as:
$`{\displaystyle \frac{\chi _p}{\chi _{p0}^0}}={\displaystyle \frac{\mathrm{\Gamma }_p}{P_p}}R_p,R_p={\displaystyle \frac{r_1(1+g_5+v_5)(r_1r_3)(1+g_5v_6)g_1}{1+g_5+v_5+g_4(1+g_5v_6)+v_4(1+v_5g_6)}},`$ (6)
$$g_1=\frac{|\mathrm{\Omega }_d|^2}{P_dP_{34}},g_4=\frac{|\mathrm{\Omega }_d|^2}{P_pP_{34}},g_5=\frac{|\mathrm{\Omega }_d|^2}{P_{12}P_{23}^{}},g_6=\frac{|\mathrm{\Omega }_d|^2}{P_{34}P_{23}^{}},v_4=\frac{|\mathrm{\Omega }_c|^2}{P_pP_{12}},v_5=\frac{|\mathrm{\Omega }_c|^2}{P_{34}P_{23}^{}},v_6=\frac{|\mathrm{\Omega }_c|^2}{P_{12}P_{23}^{}}.$$
At $`\mathrm{\Omega }_d=0`$ all $`g_i=0`$ and the equation (6) reduces to that describing $`\mathrm{\Lambda }`$ scheme:
$`{\displaystyle \frac{\alpha _p}{\alpha _{0p}^0}}=Re{\displaystyle \frac{\mathrm{\Gamma }_p[\mathrm{\Gamma }_{12}+i(\mathrm{\Delta }_p\mathrm{\Delta }_c)]}{(\mathrm{\Gamma }_p+i\mathrm{\Delta }_p)[\mathrm{\Gamma }_{12}+i(\mathrm{\Delta }_p\mathrm{\Delta }_c)]+|\mathrm{\Omega }_c|^2}}=Re{\displaystyle \frac{\mathrm{\Gamma }_p[\mathrm{\Gamma }_{12}+i(\mathrm{\Delta }_p\mathrm{\Delta }_c)]}{(\mathrm{\Delta }_p\delta _1)(\mathrm{\Delta }_p\delta _2)}},`$ (7)
$`\delta _{1,2}={\displaystyle \frac{\mathrm{\Delta }_c+i(\mathrm{\Gamma }_{12}+\mathrm{\Gamma }_p)}{2}}\sqrt{{\displaystyle \frac{\mathrm{\Delta }_c+i(\mathrm{\Gamma }_{12}\mathrm{\Gamma }_p)}{2}}+|\mathrm{\Omega }_c|^2}.`$ (8)
Following to , we introduce frequency-correlation factor
$`M_{1,2}={\displaystyle \frac{d\delta _{1,2}}{d\mathrm{\Delta }_c}}={\displaystyle \frac{1}{2}}\left[1{\displaystyle \frac{\mathrm{\Delta }_c}{\sqrt{4|\mathrm{\Omega }_c|^2+\mathrm{\Delta }_c^2}}}\right].`$ (9)
The denominator in (7) displays two resonances. At $`\mathrm{\Omega }_c0`$ we obtain $`\delta _10+i\mathrm{\Gamma }_p,\delta _2\mathrm{\Delta }_c+i\mathrm{\Gamma }_{12}`$. This indicates one resonance at $`\mathrm{\Delta }_p=0`$, the HWHM is $`\mathrm{\Gamma }_p`$, which corresponds to one-photon resonance with no correlation with $`\omega _c`$ ($`M_1=d\delta _1/d\mathrm{\Delta }_c=0`$). The second resonance at $`\mathrm{\Delta }_p=\mathrm{\Delta }_c`$ is of HWHM $`\mathrm{\Gamma }_{12}`$ that corresponds to two-photon resonance, fully correlated with $`\omega _c`$ ($`M_2=d\delta _2/d\mathrm{\Delta }_c=1`$). With growth of the coupling Rabi frequency $`\mathrm{\Omega }_c`$ the resonance becomes split in two component, their frequency-correlation properties modify, and $`M_1M_21/2`$ at $`|\mathrm{\Omega }_c|^2\{|\mathrm{\Delta }_c|^2,\mathrm{\Gamma }_{p,12}^2\}`$, which neither correspond to one- nor to two-photon processes. HWHM of the resonances becomes also nearly equal to each other and to $`(\mathrm{\Gamma }_{12}+\mathrm{\Gamma }_p)/2`$. Note that all the effects are determined by the coherence $`\rho _{12}`$ induced in the transition $`|1|2`$ by two coupled fields, and that always $`M_1+M_2=1`$. More detailed discussion can be found in .
In the alternative case $`\mathrm{\Omega }_c=0`$ all $`v_i=0`$ and the equation (6) reduces to that describing $`V`$ scheme:
$`{\displaystyle \frac{\alpha _p}{\alpha _{p0}^0}}=Re\left\{\mathrm{\Gamma }_p{\displaystyle \frac{[\mathrm{\Gamma }_{34}+i(\mathrm{\Delta }_p\mathrm{\Delta }_d)](r_1r_3)|\mathrm{\Omega }_d|^2/(\mathrm{\Gamma }_p+i\mathrm{\Delta }_p)}{(\mathrm{\Gamma }_p+i\mathrm{\Delta }_p)[\mathrm{\Gamma }_{34}+i(\mathrm{\Delta }_p\mathrm{\Delta }_d)]+|\mathrm{\Omega }_d|^2}}\right\}.`$ (10)
The structure of the denominator is similar to (7) and determined by the coherence $`\rho _{34}`$. Additional term in the nominator is stipulated by the coherence induced at the transition $`|1|3`$ with not zero populations of the levels unlike the transition $`|2|4`$. Alongside with the coherence $`\rho _{34}`$ this term is a source of the nonlinear interference effects (NIEF) in absorption (gain) and refraction. Specific features of NIEF in coupled Doppler broadened $`\mathrm{\Lambda }`$, $`V`$ and ladder schemes were explored in . It has been outlined that frequency integrated absorption index is proportional to $`r_1r_4`$ only, i.e. it’s change is determined only by the population change. As it was first outlined in , indeed NIEF give rise to difference in the line shapes of pure absorption and emission spectra. As the consequence of this effect, the appearance of amplification without inversion on the base of NIEF has been predicted. Amplification without inversion was introduced and its features were explicitly analyzed and illustrated for the model of neon transitions in the early publication .
Specific feature of the case under consideration in this paper is that both of the levels of the probe transition can be driven independently. This gives rise to multiple resonance structure, corresponding to the roots of denominator in (6). If resonance splitting is much greater than their widths, we can set $`\mathrm{\Gamma }_{ij}=0`$, and the resonance values of $`\mathrm{\Delta }_p`$ are described by the equation:
$`[\mathrm{\Delta }_p\mathrm{\Delta }_d\mathrm{\Delta }_c][\mathrm{\Delta }_p(\mathrm{\Delta }_p\mathrm{\Delta }_d)(\mathrm{\Delta }_p\mathrm{\Delta }_c)(\mathrm{\Delta }_p\mathrm{\Delta }_c)\mathrm{\Omega }_d^2(\mathrm{\Delta }_p\mathrm{\Delta }_d)\mathrm{\Omega }_c^2]`$
$`\mathrm{\Delta }_p[(\mathrm{\Delta }_p\mathrm{\Delta }_d)\mathrm{\Omega }_d^2+(\mathrm{\Delta }_p\mathrm{\Delta }_c)\mathrm{\Omega }_c^2]+(\mathrm{\Omega }_d^2\mathrm{\Omega }_c^2)^2=0.`$ (11)
Basically, this equation possesses four roots, that indicates appearance of four nonlinear resonances, determined by splitting of each of the levels $`|1`$ and $`|4`$ into two quasi levels.
In the case $`\mathrm{\Delta }_d=\mathrm{\Delta }_c=\mathrm{\Delta }`$ resonance positions are given by the equation:
$`\mathrm{\Delta }_p^2=(\mathrm{\Delta }^2+2\mathrm{\Omega }_d^2+2\mathrm{\Omega }_c^2)/2\pm \sqrt{[(\mathrm{\Delta }^2+2\mathrm{\Omega }_d^2+2\mathrm{\Omega }_c^2)/2]^2(\mathrm{\Omega }_d^2\mathrm{\Omega }_c^2)^2}.`$ (12)
At $`|\mathrm{\Omega }_d|=|\mathrm{\Omega }_c|=|\mathrm{\Omega }|`$ two resonances merge in one not shifted resonance at $`\mathrm{\Delta }_p^{(1)}=0`$, the other two resonances are given by
$`\mathrm{\Delta }_p^{(2,3)}=\pm \sqrt{\mathrm{\Delta }^2+4\mathrm{\Omega }^2}.`$ (13)
In the case $`\mathrm{\Delta }_d=\mathrm{\Delta }_c=0`$ resonance positions are found as:
$`\mathrm{\Delta }_p^{(1,2)}=\pm (|\mathrm{\Omega }_d||\mathrm{\Omega }_c|);\mathrm{\Delta }_p^{(3,4)}=\pm (|\mathrm{\Omega }_d|+|\mathrm{\Omega }_c|).`$ (14)
We shall illustrate major outcomes of the paper with numerical simulations addressed to the conditions of the experiment . Transitions $`|1|3|2|4|1`$ in Fig. 1a are attributed to those of the sodium dimers $`Na_2`$: $`X^{}\mathrm{\Sigma }_g^+(v\mathrm{"}=0,J\mathrm{"}=45)A^{}\mathrm{\Sigma }_u^+(6,45)(\lambda _d=655`$ $`nm)X^1\mathrm{\Sigma }_g^+(14,45)(\lambda _{32}=756`$ $`nm)B^1\mathrm{\Pi }_u(5,45)(\lambda _c=532`$ $`nm)X^{}\mathrm{\Sigma }_g^+(0,45)(\lambda _p=480`$ $`nm)`$. The Doppler widths of the transition at wavelength $`\lambda _p=480`$ $`nm`$ at the temperature about 450 C is $`2D_p1.7`$ GHz. Then the Boltzmann’s population of the level $`|2`$ makes about 1.5% from that of the level $`|1`$. The following relaxation parameters are used: $`\mathrm{\Gamma }_4=\mathrm{\Gamma }_3=120`$, $`\mathrm{\Gamma }_2=\mathrm{\Gamma }_1=20`$, $`\gamma _{42}=5`$, $`\gamma _{41}=10`$, $`\gamma _{32}=4`$, $`\gamma _{31}=7`$, $`\mathrm{\Gamma }_{12}=20`$, $`\mathrm{\Gamma }_{34}=120`$, $`\mathrm{\Gamma }_{23}=\mathrm{\Gamma }_c=\mathrm{\Gamma }_d=\mathrm{\Gamma }_p=70`$, (all in $`10^6`$ $`s^1`$). For numerical simulation we have used the full set of the equations from , accounting for various relaxation transitions and not zero population of the level $`|2`$.
We begin our discussion with the case when the Doppler-broadening is not included. Performing numerical calculation for Eq. (6), we depict the absorption profiles in Fig. 2. From Fig. 2 a and b, it is easy to find the typical EIT and absorption spectrum with $`\left|\mathrm{\Omega }_d\right|=0`$ (Fig. 2a) and $`\left|\mathrm{\Omega }_c\right|=0`$ (Fig. 2b). Under increasing Rabi frequency $`\left|\mathrm{\Omega }_c\right|`$ ($`\left|\mathrm{\Omega }_d\right|`$), the four-peak spectrum forms, when the driving fields become comparable to or larger than the spectral width $`\mathrm{\Gamma }_p`$ of the atomic transition. The emergence of the four-peak spectrum is attributed to the fact that levels $`|1`$ and $`|4`$ have been driven into two dressed states (Fig. 1). As $`\left|\mathrm{\Omega }_d\right|`$ approaches $`\left|\mathrm{\Omega }_c\right|`$, the four-peak spectrum degenerates into a three-peak spectrum since two dressed-state transitions have the same resonant frequency and contribute to the same central component. As already known, a three-peak spectrum occurs in resonance fluorescence (RF) in the strong field limit of a two-level system (for a review, see ). Even though the dressed-state diagrams are similar in the two cases, the four-peak spectrum can not appear in RF because in the case of RF the dynamic Stark splitting for the two doublets must be the same, while in our system the two doublet splittings are controlled separately by different driving field intensity and frequency. We would like to stress that the multiple EIT windows can appear simultaneously with the multiple absorption peaks. In Fig. 2c, the absorption amplitude $`A`$ at $`\mathrm{\Delta }_p=0`$ is plotted against the Rabi frequencies for two resonant driving fields. An important feature is that the maximum absorption occurs under the condition that the two Rabi frequencies are equal, which is due to the fact that the two dressed-state transitions simultaneously contribute to the central component. This result may provide a way to measure the field intensity as well as atomic parameters.
Fig.3 depicts absorption spectrum for the molecules with the velocity projection on the propagation direction of the coupled fields $`\mathrm{v}=0`$. The plots illustrate dependencies, described by the equations (12) - (14). The appearance of the four-peak spectrum is attributed to the fact that levels $`|1`$ and $`|4`$ have been driven into two dressed states. Direct computing shows that frequency integral absorption for the plot 3 is the same as for 5, whereas for the other plots it is about 0.5 of that for the graph 5. This is due to strong saturation of population difference at the transition $`|1`$$`|3`$.
## III Coherence induced resonances in Doppler-broadened medium of sodium dimers
In the previous sections, Doppler broadening has not been accounted for. In order to consider this effect, Doppler shifts must be introduced in the three interacting field resonance detunings in Eq. (6).
Below we will proceed with analysis of the absorption spectrum under effects of Doppler-broadening. Figure 5 displays same dependencies, but for velocity averaged absorption index and co-propagating probe and driving fields. While number of peaks and their positions are identical to the previous case, the specific feature of the inhomogeneous broadening is that plot 1 displays enhanced absorption in the central component despite the fact that integral intensity of this graph is about 2 times less than that for the plot 4 (due to the strong saturation of the populations at the transition $`|1`$$`|3`$). In absence of driving fields only small fraction of the molecules (about $`\mathrm{\Gamma }_p/D_p10^3`$) can be coupled by probe field. Overlap of the two dressed transitions gives rise to increased amount of the molecules coupled with the probe field and consequently to enhanced absorption as compared to not perturbed absorption in the center of Doppler broadened transition.
As it was shown in , resonance requirements for molecules at velocity $`\mathrm{v}`$ and at $`max\{|\mathrm{\Omega }_d|^2,\mathrm{\Delta }_d|^2\}k_du`$ ($`u`$ is thermal velocity, $`k_d`$ \- wave number) are described by the equation
$$\mathrm{\Delta }_p𝐤_p𝐯=\delta _{1,2}M_{1,2}𝐤_d𝐯.$$
(15)
The equation indicates that cancellation of Doppler shifts is possible even at $`k_pk_d`$, if $`k_p<k_d`$. This is due to the interplay of the Doppler and ac Stark shifts giving rise to variation of $`M_1`$ in the interval 0…1/2, whereas $`M_2`$ – in the interval 1…1/2. Possible enhancements of cross-section of optical processes through concurrent coupling of molecules from wide velocity interval by means of the above outlined approach are shown in Fig. 5. The frequencies of the probe and $`E_d`$ fields are interchanged (see Fig. 1b), all other relaxation parameters remain the same. Plot 1 displays enhanced sub-Doppler resonances with the FWHM comparable with the natural linewidth (see insets (a),(b))for co-propagating wave. Dramatic change as compared to similar graphs in Fig. 5 is explicitly seen. The other plots show evolution of the structures with variation of the intensity of one of the driving field. As plot 4 indicates, the sub-Doppler structures exist only under certain ratio of the driving field intensities. Even more dramatic change occurs while driving fields (especially shorter wavelength one) are counter propagating (see inset (c)).
## IV Conclusion
In conclusion, we have explored spectral properties of the molecular transitions in the case that both upper and lower levels are coupled to other atomic levels and Doppler effects play an important role. The features of the four-peak or three-peak spectra induced by the driving fields are investigated. Despite the dominated Doppler-broadening, two or more peaks may possess the sub-Doppler resolusion. A crucial role of the ratio of the frequencies of the coupled transition is shown. With manipulation the detunings and/or intensities of two coupling fields, one can improve the spectral resolution. The enhanced absorption in the sub-Doppler peak is shown to be realized, whereas the peak can be induced in the center of the inhomogeneously broadened transition. The predicted effects are attributed to the quantum coherence and interference processes, while frequency-correlation properties of multi-photon processes experience substantial modification with the growth of the driving fields intensities, which leads to corresponding dramatic changes of the role of Doppler effects.
The similar sub-Doppler technique can also be implemented for other atomic coherence effects, such as EIT, LWI and for enhancing various resonant nonlinear optical phenomena in inhomogeneously broadened media. In one of our previous papers , the inversionless amplification has been considered with the important feature that incoherent excitation to the upper level is not necessary in the same scheme as considered in this letter. Similar scheme is often used in quantum optics. Most recently, electromagnetically induced absorption was studied in a similar four-level system . We believe that, by utilizing this sub-Doppler technique, the above effects in inhomengeneously broadened media can be further enhanced and manipulated.
We want to stress that in optically thick media the processes similar to optical parametric amplification and involved in the schemes Fig. 1a,b may play a crucial role, that imposes dramatic consequences on the features of the output probe signal .
## V Acknowledgements
One of the authors, Jin-Yue Gao, would like to acknowledge the support from NSF in China, the research Fund for the Doctoral Program of Higher Education in China, Education Department of China, and DFG in Germany. A. K. Popov acknowledges support from the Russian (grant 99-02-39003) and Krasnoyarsk Regional Foundations for Basic Research, from the Center on Fundamental Natural Sciences at St. Petersburg University (grant 97-5.2-61) and INTAS (INTAS-99-19). He gratefully acknowledge support of his visit in Germany from Institute of Quantum Optics of Hannover University and along with Jin-Yue Gao would like to thank B. Wellegehausen for stimulating discussions. The authors thank S. A. Myslivets for valuable assistance in numerical simulation.
|
warning/0005/quant-ph0005013.html
|
ar5iv
|
text
|
# How entangled can two couples get?11footnote 1See ref. [8].
## 1 Introduction
In a system of two qubits the state
$$|C_2=\frac{1}{\sqrt{2}}\left(|00+|11\right)$$
(1.1)
is, on all counts, the most entangled of all pure states. It gives the greatest violation of Bell inequalities, it has the largest entropy of entanglement, and its one-party reduced states are both maximally mixed. All of these properties determine it uniquely up to local unitary transformations.
A pure state of three qubits with similar properties is the GHZ state
$$|C_3=\frac{1}{\sqrt{2}}(|000+|111).$$
(1.2)
This state has the maximum value of pure 3-party entanglement, as measured by Wootters’s 3-tangle , and its one-particle reduced density matrices are all maximally mixed. Like the two-qubit state $`|C_2`$, it is characterised uniquely, up to local unitary transformations, by the latter property .
The obvious $`n`$-party generalisation is the “Schrödinger cat” state
$$|C_n=\frac{1}{\sqrt{2}}(|00\mathrm{}0+|11\mathrm{}1).$$
(1.3)
Like $`|C_2`$ and $`|C_3`$, this state has the property that its one-party reduced states are all maximally mixed. On the strength of this, $`|C_n`$ is sometimes called the “maximally entangled” pure state of $`n`$ qubits. For $`n>3`$, however, not all states with this property are locally equivalent, and it is not clear that $`|C_n`$ is really the most entangled of them. Here we examine the case $`n=4`$, show that there are 4-qubit states which are more entangled than $`|C_4`$, and attempt to find the most entangled among them.
A four-qubit system can be regarded, in three different ways, as a system of two pairs of qubits, and one can ask how entangled are these pairs. In the state $`|C_4`$ this entanglement is not maximal: each pair $`X,Y`$ of the four qubits $`A,B,C,D`$ is not in the maximally mixed state but exhibits correlations, all two-qubit density matrices being
$$\rho _{XY}=\frac{1}{2}\left(|0000|+|1111|\right).$$
(1.4)
But there do exist pure states of four qubits in which four of the six two-qubit reduced density matrices are maximally mixed. (We note that these density matrices come in pairs: if the reduced state of one pair is maximally mixed, so is that of the complementary pair.) For example, the state
$$|\mathrm{\Psi }=\frac{1}{2}\left(|0000+|0111+|1001+|1110\right)$$
(1.5)
has two-qubit density matrices
$`\rho _{AB}`$ $`=\rho _{AC}=\rho _{BD}=\rho _{CD}=\frac{1}{4}\mathrm{𝟏},`$
$`\rho _{AD}`$ $`=\frac{1}{2}(|++++|+||),`$
$`\rho _{BC}`$ $`=\frac{1}{2}\left(|0000|+|1111|\right),`$
where
$$|\pm =\frac{1}{\sqrt{2}}\left(|0\pm |1\right).$$
(1.6)
The one-qubit reduced density matrices of $`|\mathrm{\Psi }`$ are all maximally mixed. But $`|\mathrm{\Psi }`$ is not locally equivalent to the cat state $`|C_4`$, for it has different values from that state of the entanglement entropies $`E_{AB},E_{AC},E_{BD},E_{CD}`$, which are invariants under local unitary transformations. (We write
$$E_{XY}=\mathrm{tr}(\rho _{XY}\mathrm{log}_2\rho _{XY}),\rho _{XY}=\mathrm{tr}_{ZW}|\mathrm{\Psi }\mathrm{\Psi }|$$
(1.7)
where $`\{W,X,Y,Z\}`$ is a permutation of $`\{A,B,C,D\}`$). In fact , the state $`|\mathrm{\Psi }`$ cannot be asymptotically reversibly converted into any collection of the states $`|C_n`$ for $`n=2,3,4`$.
It is natural to wonder whether there is a still more entangled state in which the reduced density matrices of all six pairs of qubits are maximally entangled.<sup>2</sup><sup>2</sup>2This question has been studied by Gisin and Bechmann-Pasquinucci for the general case of $`n`$ qubits, but under the restriction that the states are symmetric under permutations of the qubits. This is an unnatural requirement in this context, since it is not invariant under local unitary transformations. In Section 2 we will show that there is no such state of four qubits. The situation changes, however, if particles with larger state spaces are considered, and we will show that a system of four four-state particles does have pure states with this property.
In Section 3 we ask what is the maximum two-pair entanglement possible for a system of four qubits. Taking as a measure the average entanglement entropy $`E_2=\frac{1}{3}(E_{AB}+E_{AC}+E_{AD})`$, we exhibit a pure state $`|M_4`$ with a greater value of this quantity than the state (1.5); we show that $`E_2`$ has a stationary value at $`|M_4`$, and present evidence that this is a global maximum.
## 2 Non-existence of maximal entanglement <br>between all pairs of pairs
###### Theorem 1.
There is no pure state of four qubits whose two-qubit density matrices are all multiples of the identity.
###### Proof.
Write the four-qubit state $`|\mathrm{\Psi }`$ as
$$|\mathrm{\Psi }=\underset{i,j,k,l=0,1}{}t^{ijkl}|i|j|k|l$$
The tensor $`t^{ijkl}`$ can be regarded as a $`4\times 4`$ matrix in three different ways:
$$t^{ijkl}=\frac{1}{2}(U_1)_{kl}^{ij}=\frac{1}{2}(U_2)_{jl}^{ik}=\frac{1}{2}(U_3)_{jk}^{il}$$
and the requirement of maximal entanglement is that $`U_1,U_2`$ and $`U_3`$ should all be unitary matrices.
We show first that by local unitary transformations we can arrange that the coordinates of $`|\mathrm{\Psi }`$ satisfy
$$t^{1000}=t^{0100}=t^{0010}=t^{0001}=0.$$
To do this, we find normalised states $`|\alpha ,|\beta ,|\gamma ,|\delta `$ which maximise $`N=|\mathrm{\Psi }|\alpha |\beta |\gamma |\delta |^2`$. Such states certainly exist, since $`N(\alpha ,\beta ,\gamma ,\delta )`$ is a continuous function on the compact space $`S^3\times S^3\times S^3\times S^3`$. Change basis in each of the single-qubit spaces so that $`|\alpha |\beta |\gamma |\delta `$ becomes the basis state $`|0000`$; then $`N=|t^{0000}|^2`$. But if $`t^{1000}`$ were non-zero we could change basis for states of the first qubit so as to increase $`|t^{0000}|`$; since we have maximised it, we must have $`t^{1000}=0`$. Similarly $`t^{0100}=t^{0010}=t^{0001}=0`$.
The matrices $`U_1,U_2,U_3`$ which have to be unitary are now
$$\left(\begin{array}{cccc}t^{0000}& 0& 0& t^{0011}\\ 0& t^{0101}& t^{0110}& t^{0111}\\ 0& t^{1001}& t^{1010}& t^{1011}\\ t^{1100}& t^{1101}& t^{1110}& t^{1111}\end{array}\right),\left(\begin{array}{cccc}t^{0000}& 0& 0& t^{0101}\\ 0& t^{0011}& t^{0110}& t^{0111}\\ 0& t^{1001}& t^{1100}& t^{1101}\\ t^{1010}& t^{1011}& t^{1110}& t^{1111}\end{array}\right),\left(\begin{array}{cccc}t^{0000}& 0& 0& t^{0110}\\ 0& t^{0011}& t^{0101}& t^{0111}\\ 0& t^{1010}& t^{1100}& t^{1110}\\ t^{1001}& t^{1011}& t^{1101}& t^{1111}\end{array}\right).$$
Hence the following three conditions must hold:
1. $`t^{0011}=0`$ or $`t^{0111}=t^{1011}=0`$;
2. $`t^{0101}=0`$ or $`t^{0111}=t^{1101}=0`$;
3. $`t^{0110}=0`$ or $`t^{0111}=t^{1110}=0`$.
Following the branching consequences of these, and requiring that the matrices are unitary, leads to a conclusion of the form that three $`2\times 2`$ matrices
$$\left(\begin{array}{cc}a& b\\ c& d\end{array}\right),\left(\begin{array}{cc}b& e\\ f& c\end{array}\right),\left(\begin{array}{cc}a& e\\ f& d\end{array}\right)$$
must all be unitary. In this situation all the matrix elements must be non-zero: if, for example, $`a=0`$, then $`d=0`$ and $`b,c,e,f`$ would all have unit modulus, which is impossible if the second matrix is to be unitary. Now unitarity gives
$$a=\frac{b\overline{d}}{\overline{c}}=\frac{e\overline{d}}{\overline{f}},$$
so
$$b\overline{f}=e\overline{c}.$$
But
$$b\overline{f}+e\overline{c}=0,$$
so $`b\overline{f}=e\overline{c}=0`$, which is impossible. ∎
Four-party states with this property do exist, however, if the individual parties have more independent states. For example, suppose each party has a 4-dimensional state space; then the (non-normalised) state with coordinates
$$t^{ijkl}=\{\begin{array}{cc}1\hfill & \text{ if }i=j=k=l\text{ or }(i,j,k,l)\text{ is an even permutation of }(1,2,3,4)\hfill \\ 0\hfill & \text{ otherwise }\hfill \end{array}$$
has every two-qubit density matrix equal to a multiple of the identity.
We note that the method introduced in this proof can be used to define a canonical form for any multipartite pure state, in which a maximal number of coordinates are zero. This can be regarded as a generalisation of the Schmidt decomposition of a bipartite state. Details can be found in .
## 3 A Maximally Entangled Four-qubit State?
Given that a four-qubit state cannot have maximal entropy of entanglement for every two-qubit subset, we now ask what is the greatest possible average for such entropies, i.e. we seek to maximise
$`E_2`$ $`=\frac{1}{6}(E_{AB}+E_{AC}+E_{AD}+E_{BC}+E_{BD}+E_{CD})`$
$`=\frac{1}{3}(E_{AB}+E_{AC}+E_{AD}).`$
The second equality holds because complementary pairs have equal entropy.
We have not been able to solve this problem analytically. We will adopt a heuristic approach, using an (indefensible) analogy with the two-qubit system to obtain a candidate maximally entangled state and then showing that $`E_2`$ is indeed stationary at this state and appears to be maximal.
Maximally entangled states of two qubits like (1.1) are often referred to as “singlet” states, since an example of such a state is the state of two spin-$`\frac{1}{2}`$ particles with zero total angular momentum. This description is not invariant under local unitary transformations, which need not preserve the total angular momentum. Nevertheless, let us take this as a hint in investigating four-qubit states. A system of four spin-$`\frac{1}{2}`$ particles has two independent singlet states. The most symmetric combination of these, in the sense that all of its two-qubit density matrices are unitarily equivalent, is locally equivalent to
$$|M_4=\frac{1}{\sqrt{6}}\left[|0011+|1100+\omega (|1010+|0101)+\omega ^2(|1001+|0110)\right],$$
where $`\omega =e^{2\pi i/3}`$, or to $`|\overline{M}_4`$, the complex conjugate of $`|M_4`$ in the computation basis. Note that this state is not symmetric between the qubits but belongs to a two-dimensional representation of the permutation group, the other state in the representation being $`|\overline{M}_4`$.
The two-qubit density matrices of $`|M_4`$ are
$$\rho _{AB}=\rho _{AC}=\rho _{AD}=\frac{1}{6}\left[|0000|+|1111|+|\mathrm{\Phi }_+\mathrm{\Phi }_+|\right]+\frac{1}{2}|\mathrm{\Phi }_{}\mathrm{\Phi }_{}|$$
(3.1)
where
$$|\mathrm{\Phi }_\pm =\frac{1}{\sqrt{2}}(|10\pm |01)$$
(in angular momentum terms, the two-qubit reduced states are equal mixtures of a singlet and a maximally mixed triplet). Hence the entanglement entropies are
$$E_{AB}=E_{AC}=E_{AD}=1+\frac{1}{2}\mathrm{log}_23.$$
Comparing with the cat state $`|C_4`$, for which
$$E_{AB}=E_{AC}=E_{AD}=1,$$
and the state of (1.5), for which
$$E_{AB}=E_{AC}=2,E_{AD}=1,$$
we see that $`|M_4`$ has a greater value of $`E_2`$ than either.
We will now show that the function $`E_2`$ is stationary at $`|M_4`$. To simplify the calculation, we consider the functions $`E_{XY}`$ defined by (1.7) for all state vectors $`|\mathrm{\Psi }`$, though these functions coincide with the entropies only when $`|\mathrm{\Psi }`$ is normalised. Suppose $`|\mathrm{\Psi }`$ changes by $`|\delta \mathrm{\Psi }`$. Because the trace makes the operators behave as if they commuted, the consequent change in $`E_{XY}`$ is given by
$$\delta E_{XY}=\mathrm{tr}\left[\rho _{XY}\mathrm{log}_2(\text{e}\rho _{XY})\right].$$
At $`|\mathrm{\Psi }=\sqrt{6}|M_4`$ we find
$`\delta E_{AB}`$ $`=\mathrm{Re}\left[\mathrm{\Psi }|\delta \mathrm{\Psi }\mathrm{log}_2(3\text{e}^2)+\overline{\mathrm{\Psi }}|\delta \mathrm{\Psi }\mathrm{log}_23\right],`$
$`\delta E_{AC}`$ $`=\mathrm{Re}\left[\mathrm{\Psi }|\delta \mathrm{\Psi }\mathrm{log}_2(3\text{e}^2)+\omega \overline{\mathrm{\Psi }}|\delta \mathrm{\Psi }\mathrm{log}_23\right],`$
$`\delta E_{AD}`$ $`=\mathrm{Re}\left[\mathrm{\Psi }|\delta \mathrm{\Psi }\mathrm{log}_2(3\text{e}^2)+\omega ^2\overline{\mathrm{\Psi }}|\delta \mathrm{\Psi }\mathrm{log}_23\right],`$
where $`|\overline{\mathrm{\Psi }}`$ is the complex conjugate of $`|\mathrm{\Psi }`$ in the computation basis. Hence the gradient of the average $`E_2`$ in the real 32-dimensional space of all state vectors is in the direction of $`|\mathrm{\Psi }`$:
$$\delta E_2=\mathrm{log}_2(3\text{e}^2)\mathrm{Re}\mathrm{\Psi }|\delta \mathrm{\Psi }.$$
Thus $`E_2`$ is stationary at $`|\mathrm{\Psi }`$ for variations on the sphere of vectors with the same norm as $`|\mathrm{\Psi }`$. But each $`E_{XY}`$, as given by (1.7), is linearly related to the true entropy of entanglement, so the average two-party entropy is stationary at $`|M_4`$.
We have searched numerically for states which maximise $`E_2`$ starting from several arbitrarily chosen states. All the states we have obtained in this manner are locally equivalent to $`|M_4`$ or $`|\overline{M}_4`$.
## 4 Robustness of the entanglement
One of the criteria for maximal entanglement suggested in is that the entanglement should be maximally fragile: a measurement on any one of the qubits destroys the entanglement between the remaining qubits. This is true of the cat state (1.3) if the measurement projects onto the computation basis $`\{|0,|1\}`$. Projection onto the basis $`\{|+,|\}`$ defined in (1.6), however, leaves the remaining qubits in a cat state.
The state constructed in the previous section behaves in the opposite way to the cat state: measurement of one qubit leaves the other three qubits in an entangled state, and the amount of entanglement afterwards is independent of what measurement is performed. To be precise, projection of one qubit onto a computation basis state leaves the remaining qubits in one of the entangled states
$`{\displaystyle \frac{1}{\sqrt{3}}}\left(|011+\omega |101+\omega ^2|110\right)`$ (4.1)
or $`{\displaystyle \frac{1}{\sqrt{3}}}\left(|100+\omega |010+\omega ^2|001\right).`$ (4.2)
These states, which are clearly locally equivalent, have recently been identified as having maximal average two-qubit entanglement in a certain sense . Projection onto any other state of the first qubit leaves the other three qubits in another state which is locally equivalent to the above. This follows from the fact that the state $`|M_4`$ is an SU$`(2)`$ singlet, so that the unitary transformation from the computation basis to any other basis of one qubit can be undone by performing the same transformation on the other three qubits.
Thus the entanglement of the state $`|M_4`$ is *robust*. Any carelessness by the holder of any one of the qubits, resulting in an uncontrolled decoherence of that qubit, does not completely destroy the entanglement of the remaining qubits, and always leaves them with the same amount of entanglement.
## Acknowledgements
AS is grateful for the hospitality of the Physics Department at Imperial College, London, where this work was started, and to Martin Plenio for a helpful conversation.
|
warning/0005/astro-ph0005233.html
|
ar5iv
|
text
|
# What Does It Take to Stabilize Gravitational Clustering?
## 1 Introduction
Gravitational clustering is the fundamental physical process responsible for the formation and evolution of structure in the universe. Galaxy-hosting dark matter halos are the products of the strongly nonlinear stage of this process, but a detailed understanding of this important regime has mostly eluded us because few numerical simulations have the dynamic range to explore comfortably such a small length scale (e.g., Moore et al. 1999; Bullock et al. 1999). The one analytic handle has been the stable clustering assumption (Peebles 1974; Davis & Peebles 1977; Peebles 1980), which has allowed us to extrapolate into nonlinear regimes beyond the reach of numerical simulations. Predictions of the stable clustering hypothesis are: (1) the two-point correlation function of the density field $`\xi (r)`$ and the power spectrum $`P(k)`$ are power laws, with $`\xi (r)r^\gamma `$ and $`\mathrm{\Delta }(k)4\pi k^3P(k)k^\gamma `$, where $`\gamma =(9+3n)/(5+n)`$ for a primordial spectral index $`n`$; (2) the higher order $`N`$-point correlation functions $`\xi _N(r)`$ scale as $`\xi _N(r)r^{\gamma _N}`$, with $`\gamma _N=(N1)\gamma `$, or $`\xi _N\xi ^{N1}`$; and (3) the pairwise peculiar velocity exactly cancels the Hubble expansion on small scales, $`v/Hr=1`$, so that bound, high-density halos maintain a fixed physical size. Some of these predictions have agreed well enough with numerical simulation results (Jain 1997 and references therein) to have obtained general acceptance.
In two recent papers we presented evidence from high-resolution $`N`$-body simulations that gravitational clustering does not necessarily follow the scaling $`\xi _3\xi ^2`$ required by the stability condition for the two- and three-point functions $`\xi `$ and $`\xi _3`$ (Ma & Fry 2000a). For a deeper understanding of the strongly nonlinear regime beyond simulations, we proceeded to construct an analytic halo model (Ma & Fry 2000b) in which mass is distributed in spherical dark matter halos with phenomenological mass distribution functions, density profiles, and halo-halo correlations. We showed that this halo model can reproduce analytically the two- and three-point correlation functions measured in numerical simulations, and that it also makes useful predictions beyond the limited range of validity of simulations.
In this Letter we investigate the extent to which the analytic halo model is consistent with stable gravitational clustering in the strongly nonlinear regime. In particular, we derive the halo model predictions for the asymptotic nonlinear behavior of the $`N`$-point correlation functions (§2) and the pairwise peculiar velocities (§3) in terms of dark matter halo properties. We then obtain the conditions that must be satisfied by the halo mass function and density profile in order to reproduce results of the stable clustering hypothesis.
## 2 $`N`$-Point Correlation Functions
The analytic halo model of Ma & Fry (2000b) for the $`N`$-point correlation functions of the mass density takes three inputs: the halo mass distribution function $`dn/dM`$; the halo density profile $`\rho (r)/\overline{\rho }=\overline{\delta }u(r/R_s)`$ where $`R_s`$ is a scale radius and $`\overline{\delta }`$ is the normalization; and halo-halo correlations. In this model, the two-point function contains contributions from particle pairs residing in a single halo and in two separate halos, $`\xi (r)=\xi _{1h}(r)+\xi _{2h}(r)`$, where the subscripts “$`1h`$” and “$`2h`$” denote the respective “1-halo” and “2-halo” contributions (see also Seljak 2000). Similarly, we write the bispectrum $`B(k_1,k_2,k_3)`$ as three separate terms, representing contributions from particle triplets that reside in a single halo and in two and three distinct halos. As shown in Ma & Fry (2000b), the dominant contribution in the nonlinear regime on small length scales is from the “1-halo” terms for particles that reside in the same halos. This makes intuitive sense because closely spaced particles are most likely to be found in the same halo. These “1-halo” terms for the two- and $`N`$-point statistics are
$`\xi _{1h}(r)={\displaystyle 𝑑M\frac{dn}{dM}R_s^3\overline{\delta }^2\lambda _2(r/R_s)},P_{1h}(k)={\displaystyle 𝑑M\frac{dn}{dM}[R_s^3\overline{\delta }\stackrel{~}{u}(kR_s)]^2},`$
$`\xi _{N,1h}(r_1,\mathrm{},r_N)={\displaystyle 𝑑M\frac{dn}{dM}R_s^3\overline{\delta }^N\lambda _N(r_1/R_s,\mathrm{},r_N/R_s)},`$ (1)
$`P_{N,1h}(\text{k}_1,\mathrm{},\text{k}_N)={\displaystyle 𝑑M\frac{dn}{dM}[R_s^3\overline{\delta }\stackrel{~}{u}(k_1R_s)]\mathrm{}[R_s^3\overline{\delta }\stackrel{~}{u}(k_NR_s)]}.`$
where $`\lambda _N(x_1,\mathrm{},x_N)d^3yu(y)u(|\text{y}+\text{x}_1\text{x}_2|)\mathrm{}u(|\text{y}+\text{x}_1\text{x}_N|)`$; $`\stackrel{~}{u}`$ is the Fourier transform of $`u(x)`$; the $`N`$-point $`P_N`$ is defined for $`\text{k}_i=0`$, and $`P_3=B`$ is the bispectrum.
For our goal of studying the $`N`$-point functions in the strongly nonlinear regime, we need to consider only the halo mass function and halo density profile and not halo-halo correlations because the latter make negligible contributions on small length scales and are therefore irrelevant for stable clustering (see Ma & Fry 2000b for the full expressions for $`P`$ and $`B`$). For the halo mass function, we write it as
$$\frac{dn}{dM}M^2\nu ^\alpha e^{\nu ^2/2},\nu =\frac{\delta _c}{\sigma (M)},$$
(2)
where $`\alpha `$ parameterizes the uncertainty in the logarithmic slope at the low-mass end, $`\delta _c1.68`$ characterizes the linear overdensity at the onset of gravitational collapse, and $`\sigma (M)`$ is the linear rms mass fluctuations in spheres of radius $`R`$, where $`M=4\pi \overline{\rho }R^3/3`$. Values in use for $`\alpha `$ vary from 1 (Press & Schechter 1974) to about 0.4 (Sheth & Tormen 1999; Jenkins et al. 2000). For the halo density profile, we assume a “universal” shape, $`\rho /\overline{\rho }=\overline{\delta }u(r/R_s)`$, where $`u=1/[x(1+x)^2]`$ is proposed by Navarro et al. (1997), while $`u(x)=1/(x^{3/2}+x^3)`$ is suggested by Moore et al. (1999). Both the scale radius $`R_s`$ and the normalization $`\overline{\delta }`$ can be expressed in terms of a concentration parameter $`c(M)R_{200}/R_s`$, where $`R_{200}`$ is the virial radius within which the average halo density is 200 times the mean density of the universe. Specifically, $`R_sM^{1/3}/c`$, and $`\overline{\delta }c^3g(c)`$, where $`g(c)`$ is a logarithmic factor with $`g=1/[\mathrm{ln}(1+c)c/(1+c)]`$ for Navarro et al. (1997) and $`g=1/\mathrm{ln}(1+c^{3/2})`$ for Moore et al. (1999). High resolution simulations of individual halos show a roughly power-law $`c(M)`$ with some, perhaps substantial, scatter (e.g., Cole & Lacey 1996; Tormen et al 1997; Navarro et al. 1997; Jing & Suto 2000), so we will write
$$c\frac{R_{200}}{R_s}\left(\frac{M}{M^{}}\right)^{\beta _0},$$
(3)
where $`M^{}`$ is defined by $`\sigma (M^{})=1`$, and $`\beta _0`$ is used to parameterize the uncertainty in the mass dependence, which has been found to be in the range $`0\beta _01/2`$.
We now derive analytically the small-$`r`$ behavior of equation (1), in which the factors $`dn/dM`$, $`R_s`$, and $`\overline{\delta }`$ are all functions of the halo mass $`M`$. For small $`r`$ or high $`k`$, the integral for $`\xi _{1h}(r)`$ or $`P_{1h}(k)`$ converges before the exponential cutoff in the mass function, and the dominant contribution to the integral comes from masses near the scale for which $`r/R_s=1`$ or $`kR_s=1`$. We then have $`P(k)P_{1h}(k)𝑑M\nu ^\alpha \stackrel{~}{u}^2(kR_s)g^2(c)`$ at high $`k`$. Changing variables to $`y=kR_skM^{\beta _0+1/3}`$, we obtain the asymptotic behavior
$$\xi (r)r^\gamma ,\mathrm{\Delta }(k)4\pi k^3P(k)k^\gamma ,\gamma =\frac{18\beta \alpha (n+3)}{2(3\beta +1)},$$
(4)
where we have $`\beta =\beta _0`$ if the logarithmic factor $`g(c)`$ in $`\overline{\delta }`$ is ignored, or more accurately, we have $`\beta 0.8\beta _0`$ if the effect of $`g(c)`$ is approximated by a power law. From equation (4), we see that in order to reproduce the two-point stable clustering result for a given spectral index $`n`$, the indices $`\alpha `$ and $`\beta `$ must satisfy
$$\frac{18\beta \alpha (n+3)}{2(3\beta +1)}=\frac{9+3n}{5+n}.$$
(5)
For $`n=2`$, for example, this is satisfied if $`\beta =0.25`$ or $`\beta _00.31`$ for $`\alpha =1`$, and $`\beta =0.2`$ or $`\beta _00.25`$ for $`\alpha =0.4`$.
For the three-point function $`\xi _3`$ and the bispectrum $`B`$, the integrals in equation (1) converge if $`ϵ<3/2`$ for an inner halo density of $`r^ϵ`$. For $`\xi _3`$, $`B`$, and the commonly defined three-point amplitude $`QB/P^2`$, we find
$$\xi _3(r)r^{\gamma _3},k^6B(k)k^{\gamma _3},Q(k)k^{\gamma _32\gamma },\gamma _3=\frac{36\beta \alpha (n+3)}{2(3\beta +1)}.$$
(6)
Stable clustering requires a scale-independent $`Q(k)`$, and
$$\gamma _32\gamma =\frac{\alpha (n+3)}{2(3\beta +1)}=0,$$
(7)
which holds only if $`\alpha =0`$ or $`n=3`$. Combining equations (5) and (7), we find that stable clustering requires $`\alpha =0`$ and $`\beta =(3+n)/6`$. All work on halo mass functions performed thus far has reported a much larger value for $`\alpha `$; stable clustering therefore does not appear to be followed, which is consistent with the scale-dependent $`Q(k)`$ reported in Ma & Fry (2000a).
We illustrate these results by showing in Figure 1 the power spectrum and bispectrum computed from the analytic halo model for three sets of ($`\alpha `$, $`\beta `$) for the $`n=2`$ scale-free model. The values of $`\alpha `$ and $`\beta `$ are all chosen to give the same slope $`\gamma =1`$ for the nonlinear two-point function; the corresponding power spectra are therefore nearly identical at high $`k`$. The three-point $`Q`$, however, have very different high-$`k`$ behavior for the three cases, and only the case $`\alpha =0`$ produces an approximately scale-independent $`Q`$. ($`Q`$ is not perfectly flat for $`(\alpha ,\beta )=(0,1/6)`$ because as mentioned above, the simple expressions in eqs. (5) and (7) are derived by approximating the logarithmic factor $`g(c)`$ in the halo profile as a power law.)
For the four-point and higher correlation functions, provided the integral is convergent without the exponential cutoff in the mass function, the same change of variable gives
$$\xi _Nr^{\gamma _N},k^{3(N1)}P_Nk^{\gamma _N},\gamma _N=\frac{18(N1)\beta \alpha (n+3)}{2(3\beta +1)}.$$
(8)
Stable clustering requires
$$\gamma _N(N1)\gamma =(N2)\frac{\alpha (n+3)}{2(3\beta +1)}=0,$$
(9)
which again is satisfied by $`\alpha =0`$ or $`n=3`$. So if the three-point correlation follows equation (7), stable clustering continues to hold for all orders for which the integral converges. At large $`k`$, we are probing the inner cores of halos, where the inner profile goes as $`r^ϵ`$. The Fourier amplitude $`\stackrel{~}{u}`$ thus falls off as $`\stackrel{~}{u}(k)(kR_s)^{(3ϵ)}`$, and convergence requires $`N[1(3ϵ)(3\beta +1)/3]1+\alpha (n+3)/6<0`$. For $`ϵ=1`$ (Navarro et al.) and $`n>2`$, this holds for all $`N`$, but for $`ϵ=\frac{3}{2}`$ (Moore et al.) and $`n=2`$, it holds only as far as $`N=3`$.
## 3 Pairwise Peculiar Velocities
Besides the scaling properties of the $`N`$-point correlation functions discussed above, stable clustering also predicts that the mean peculiar velocity $`v`$ between pairs of objects of separation $`r`$ must cancel the Hubble expansion $`Hr`$ on small scales so that bound halos maintain a fixed physical size. The velocity ratio $`v/Hr`$ is related to the two-point correlation function $`\xi `$ by the pair conservation equation (Peebles 1980)
$$\frac{v}{Hr}=\frac{\overline{\xi }}{3(1+\xi )}\frac{\mathrm{ln}\overline{\xi }}{\mathrm{ln}a},$$
(10)
where $`a`$ is the expansion factor and $`\overline{\xi }(x)3x^3_0^x𝑑x^{}x^2\xi (x^{})`$ is the mean two-point correlation function interior to $`x`$. Stable clustering requires $`v/Hr=1`$.
The analytic halo model of Ma & Fry (2000b) makes specific predictions for the behavior of $`v/Hr`$. As eq. (10) indicates, the main task is to compute the time dependence of $`\xi `$. Equation (1) shows that the time dependence of $`\xi `$ enters through the rms density fluctuation $`\sigma `$ in $`dn/dM`$ of eq. (2), and the halo concentration parameter $`c(M)`$ of eq. (3). We obtain
$$\frac{\xi _{1h}}{\mathrm{ln}a}=𝑑M\left[f(\nu ^2\alpha )+(2\overline{\delta }^{}3+\lambda ^{})c^{}\right]\frac{dn}{dM}R_s^3\overline{\delta }^2\lambda (r/R_s),$$
(11)
where $`fd\mathrm{ln}\sigma /d\mathrm{ln}a\mathrm{\Omega }^{0.6}`$ is the standard growth rate of the density field, $`\overline{\delta }^{}d\mathrm{ln}\overline{\delta }/d\mathrm{ln}c`$, $`\lambda ^{}d\mathrm{ln}\lambda /d\mathrm{ln}x`$, and $`c^{}d\mathrm{ln}c/d\mathrm{ln}a`$. Specifically, for the Navarro et al. profile, $`\overline{\delta }(c)=(200c^3/3)/[\mathrm{ln}(1+c]c/(1+c)]`$ and $`\lambda (x)=8\pi [(x^2+2x+2)\mathrm{ln}(1+x)/(x^2+2x)1]/(x^3+2x^2)`$; and similar expressions for the Moore et al. profile can be found in Ma & Fry (2000b). The term $`c^{}`$ for the time dependence of the concentration parameter is somewhat uncertain. Numerical simulation results of Bullock et al. (1999) give a stronger redshift evolution, $`ca`$, than Navarro et al. (1997), while more complicated forms have also been proposed (e.g., Cooray et al. 2000). We note that in equation (3), $`c(M)`$ has an implicit dependence on $`a`$ through $`M^{}`$ which roughly agrees with Bullock et al. when $`\beta =(3+n)/6`$.
The prediction of the halo model for $`v/Hr`$ can be calculated from equations (1) and (11). Before doing so, it is useful to derive first the asymptotic behavior of $`v/Hr`$ in the deeply nonlinear regime at small $`r`$. On such small scales, $`\xi 1`$, and equation (4) gives a power law $`\xi r^\gamma `$ and $`\overline{\xi }=3\xi /(3\gamma )`$. For a scale-free model, the self similarity solution gives $`\overline{\xi }/\mathrm{ln}a=2/(3+n)\overline{\xi }/\mathrm{ln}r`$. We find that equation (10) then reaches
$$\frac{v}{Hr}\stackrel{r0}{}\frac{2}{n+3}\frac{18\beta \alpha (n+3)}{6+\alpha (n+3)}.$$
(12)
It can be easily checked that $`v/Hr`$ approaches unity, as it should, for any pair $`(\alpha ,\beta )`$ that satisfies the stable clustering condition of equation (5) for the two-point function.
Figure 2 shows our results for the general behavior of $`v/Hr`$ calculated from the analytic halo model. An $`n=2`$ scale-free power spectrum is used for illustration. Although $`\xi `$ is dominated by the 1-halo term $`\xi _{1h}`$ of equation (1) at $`\xi 1`$, which is the range of interest for this paper, we improve the accuracy at $`\xi <1`$ by adding the linear theory $`\xi _lr^{(3+n)}r^1`$ to $`\xi _{1h}`$. As Figure 2 shows, the resulting $`v/Hr`$ indeed follows the linear perturbative prediction at $`\xi 1`$, and equation (12) at $`\xi 1`$. We find $`v/Hr1`$ for the same pairs of $`(\alpha ,\beta )`$ shown in Figure 1, which were chosen to illustrate the stable clustering scaling of the two-point function. However, $`v/Hr`$ reaches a smaller value at $`\xi 1`$ for the $`(\alpha ,\beta )`$ suggested in the literature, e.g., $`\alpha =1`$ or 0.4 for the mass function (Press & Schechter 1974; Sheth & Tormen 1997; Jenkins et al. 2000) and $`\beta =1/6`$ (Cole & Lacey 1996; Navarro et al. 1997), implying that the pairwise peculiar velocities are too small to cancel exactly the Hubble expansion. As a consistency check, we have computed the derivative $`\overline{\xi }/a`$ using both eq. (11) and $`[\overline{\xi }(a+\mathrm{\Delta }a)\overline{\xi }(a)]/\mathrm{\Delta }a`$. The two methods give identical results, but eq. (11) offers a more physical insight into the various contributing terms. For example, the broad peak in $`v/Hr`$ for $`1\xi 100`$ is the result of the term in $`\nu ^2`$ in eq. (11), which reflects that the integral on these scales is dominated by mass scales near $`M^{}`$ where $`\sigma (M^{})=1`$.
## 4 Summary and Discussion
In this Letter, we have derived the asymptotic nonlinear behavior of the $`N`$-point correlation functions $`\xi _N`$ (eqs. , , ) and the pairwise peculiar velocities $`v/Hr`$ (eq. ) in the framework of the recently proposed analytic halo model of Ma & Fry (2000b). We have shown that their small scale behavior is consistent with the stable clustering hypothesis only if dark matter halos satisfy certain criteria. The two halo parameters whose variation we have explored are the logarithmic slope $`\alpha `$ of the halo mass distribution $`dn/dM`$ at the low mass end in eq. (2) and the slope $`\beta _0`$ for the mass dependence of the halo concentration parameter $`c(M)`$ in eq. (3) (recall $`\beta 0.8\beta _0`$). These two parameters are highlighted because results to date from numerical simulations have indicated a significant uncertainty, with $`0.4\alpha 1`$ and $`0\beta 1/2`$ being plausible values.
From the derived asymptotic nonlinear behavior, we have obtained analytically the relations that must be satisfied by $`\alpha `$ and $`\beta `$ in order for stable clustering to occur. Equations (5), (7), and (9) summarize the results for the $`N`$-point correlation function $`\xi _N`$. These equations and Figure 1 indicate that although the two-point stable clustering condition alone is satisfied by an infinite set of $`(\alpha ,\beta )`$ for a given $`n`$, the three- and higher-point stable clustering conditions, $`\xi _N\xi ^{N1},N3`$, imply $`\alpha `$ must be zero. Achieving stable clustering to all orders therefore requires stringent conditions: $`\alpha =0`$ and $`\beta =(3+n)/6`$. For non-scale-free models (e.g., cold dark matter), $`n`$ is the index at the high-$`k`$ end of the power spectrum.
Compared with the hierarchical model of Davis & Peebles (1977), results from this paper and Ma & Fry (2000b) show that the analytic halo model makes more general and physical predictions for the behavior of the $`N`$-point correlation functions and the pairwise velocities. Imposing the stability condition $`\alpha =0`$ and $`\beta =(3+n)/6`$ in equations (4), (6), and (8), we indeed recover their results of $`\xi r^\gamma `$ with $`\gamma =(9+3n)/(5+n)`$ at small $`r`$ and $`\xi _N\xi ^{N1}`$. The condition $`\beta =(3+n)/6`$ can be understood further by examining individual halos. Using $`cM_{}^{\beta _0}`$, $`M_{}a^{6/(n+3)}`$, and $`\beta _0\beta `$ up to a logarithmic factor, we see that $`\beta =(3+n)/6`$ gives approximately $`ca(t)`$, which is consistent with the simulation results in Bullock et al. (2000). The halo scale radius $`R_s`$ is then approximately constant in time in physical coordinates, as expected for a stable system, and the physical density $`\rho `$ of individual halos is indeed very close to constant in time, more so for very small halos that form early ($`\nu 1`$, $`c1`$) than for those that form recently ($`\nu 1`$, $`c1`$).
Yano & Gouda (1999) have found that for halo density $`\rho r^ϵ`$ with $`ϵ3/2`$ (such as Navarro et al. and Moore et al.), the velocity parameter $`v/Hr`$ approaches 0 and the two-point correlation function $`\xi `$ approaches a constant (i.e. $`\gamma =0`$) in the nonlinear limit. Neither behavior is seen in our Figures because their derivation ignores the mass distribution function $`dn/dM`$ and is therefore valid only for equal mass halos.
Our earlier paper (Ma & Fry 2000a) has already shown signs of departure from the stable clustering hypothesis in high resolution $`N`$-body simulations. In this paper, the specific values $`\alpha =0`$ and $`\beta =(3+n)/6`$ required for stable clustering provide additional evidence that this long-cherished hypothesis may not be applicable in all situations. In such cases, one consequence is that the frequently used fitting formulas for the nonlinear power spectrum (Hamilton et al 1991; Jain et al. 1995; Peacock & Dodds 1996; Ma 1998) will need modifications at high $`k`$, e.g., $`k/k_{nl}50`$ for the $`n=2`$ model and $`k20h`$ Mpc<sup>-1</sup> for cold dark matter models with a cosmological constant (see Figs. 3 and 4 of Ma & Fry 2000b). We believe the analytic halo model is a powerful tool that is providing new insight into the nonlinear regime of gravitational clustering, the most fundamental process in cosmology.
C.-P. M thanks Jim Peebles and Uros Seljak for useful discussions. She acknowledges support of an Alfred P. Sloan Foundation Fellowship, a Cottrell Scholars Award from the Research Corporation, a Penn Research Foundation Award, and NSF grant AST 9973461.
|
warning/0005/quant-ph0005071.html
|
ar5iv
|
text
|
# Robustness and diffusion of pointer states
## Abstract
Classical properties of an open quantum system emerge through its interaction with other degrees of freedom (decoherence). We treat the case where this interaction produces a Markovian master equation for the system. We derive the corresponding distinguished local basis (pointer basis) by three methods. The first demands that the pointer states mimic as close as possible the local non-unitary evolution. The second demands that the local entropy production be minimal (predictability sieve). The third imposes robustness on the inherent quantum and emerging classical uncertainties. All three methods lead to localized Gaussian pointer states, their formation and diffusion being governed by well-defined quantum Langevin equations.
The programme of decoherence has been very successful in explaining the classical appearance of macroscopic or mesoscopic quantum systems, both theoretically and experimentally deco . The interaction of a quantum system with its environment leads in these cases to a delocalization of phase relations in the full configuration space of system plus environment, preventing them to be observed locally, i.e. at the system itself. Thereby the environment distinguishes a certain preferred basis for the system, which can be used to describe the apparent classical behavior (pointer basis Zur81 ). An important question is how the pointer basis can be determined and whether it is uniquely fixed. It would then be possible to give a unique decomposition of the reduced density matrix into an apparent ensemble of wave functions.
No unique rules have so far been adopted to calculate the pointer states in the general case. In Zeh73 the suggestion was made that the pointer basis (there called collection of “memory states”) is characterized by its robustness (there called “dynamical stability”). A first quantitative measure to investigate the dynamical stability is the rate of de-separation introduced in Kue73 – it measures how fast a quantum system becomes entangled with environmental degrees of freedom. In a model consisting of harmonic oscillators, it was shown that coherent states are the most stable states and therefore can be considered as pointer states Kue73 ; deco . A different measure for robustness is the “predictability sieve” put forward in Zur93 . The pointer basis is there distinguished by the property of having the least production rate for local entropy during the coupling to the environment. In the case of harmonic oscillators, this again leads to the coherent states as the pointer basis ZurHabP93 . At least for such simple systems, rate of de-separation and predictability sieve are roughly equivalent measures deco .
On the other hand, the theory and formalism of quantum state diffusion (QSD) were put forward to attribute random wave functions for local systems, which satisfy an appropriate Langevin equation qsd . These wave functions are known to be related to possible continuous measurements Dio88a as well as to decoherent histories DioGisHP95 of the given local system. But even if we take them as wave functions of mere formal meaning (since subsystems do not, in general, possess their own pure states), the question arises whether there is any connection between these states and the pointer states, in cases where the local system exhibits classical properties. We shall show that there is in fact such a connection – pointer basis and QSD basis are substantially the same. For this aim, we shall also present below a new, alternative, derivation of QSD.
In the following we shall consider the dynamics of the reduced density matrix, $`\widehat{\rho }(t)`$, of a system interacting with a certain decohering environment. Ideal pointer states (described by a fixed set of projectors $`\widehat{P}_n`$) are characterized by the fact that $`\widehat{\rho }(t)`$ can be decomposed as
$$\widehat{\rho }(t)\underset{n}{}f_n\widehat{P}_n,tt_D,$$
(1)
for a generic initial state $`\widehat{\rho }(0)`$, where $`t_D`$ is the decoherence time. The weights $`\{f_n\}`$ correspond to a normalized probability distribution. The pointer states $`\{\widehat{P}_n;n=1,2,\mathrm{}\}`$ form in this case an orthogonal system. For macroscopic systems, $`t_D`$ is extremely short deco ; JooZeh85 . More generally, one would expect
$$\widehat{\rho }(t)f(\mathrm{\Gamma })\widehat{P}(\mathrm{\Gamma })\text{d}\mathrm{\Gamma },tt_D,$$
(2)
where $`f(\mathrm{\Gamma })`$ is a probability distribution over the pointer states $`\widehat{P}(\mathrm{\Gamma })`$ which project now on an overcomplete set of pure states (the above-mentioned coherent states provide an example for this). The pointer states in (1) or (2) result after an explicit interaction with the environment is taken into account deco .
In many cases, the effects of decoherence can be described by the following Markovian master equation deco ,
$`{\displaystyle \frac{\text{d}\widehat{\rho }(t)}{\text{d}t}}\widehat{\rho }(t)`$ $`=`$ $`{\displaystyle \frac{\text{i}}{2m}}[\widehat{p}^2,\widehat{\rho }(t)]{\displaystyle \frac{D}{2}}[\widehat{x},[\widehat{x},\widehat{\rho }(t)]],`$ (3)
where $`D`$ describes the strength of the interaction with the environment. Such an equation arises, for example, in situations where environmental degrees of freedom scatter off a macroscopic object and localize it by carrying away quantum correlations with the object JooZeh85 ; deco . Applying the concept of predictability sieve would mean to minimize the local production of “linear entropy” $`S(t)=1\text{tr}\widehat{\rho }^2(t)`$. This does not give a unique answer, since the result depends explicitly on $`t`$. In the oscillator case, one has therefore calculated the time-integrated entropy production ZurHabP93 ; deco . If one considers the initial entropy production rate $`\dot{S}(0)`$, assuming the initial state $`\widehat{\rho }(0)`$ of the subsystem to be a pure state $`\widehat{P}(\mathrm{\Gamma })`$, one finds from (3) that $`\dot{S}(0)=D\left(\widehat{x}^2\widehat{x}^2\right)D\sigma ^2,`$ where the expectation value refers to the initial state. This rate would be minimized if the pointer wave functions were delta functions. However, their spread $`\sigma `$ increases dynamically due to the kinetic term in (3). Therefore, very narrow wave functions do not produce minimum entropy on a finite time scale and thus cannot correspond to pointer states. The (coherent) unitary spreading and the (incoherent) non-unitary localizing terms of the master equation (3) are competing with each other. For a wave function of characteristic width $`\sigma `$, the above two effects are approximately balanced for the “equilibrium width” JooZeh85 ; Dio87a ; deco :
$$\sigma _0\left(Dm\right)^{1/4}.$$
(4)
It is then reasonable to conjecture that, in the spirit of predictability sieve and of dynamical robustness, $`\sigma _0`$ will be the characteristic width of the pointer states. This is in fact what we shall show by using three different methods, all invoking a principle of robustness.
The first method goes as follows. Let us allow for the pointer state $`\widehat{P}(\mathrm{\Gamma })`$ a certain natural time dependence such that it may initially evolve as close as possible along the true state $`\widehat{\rho }`$ satisfying the master equation (3), and then reach a stationary state. We introduce the “speed” $`v`$ describing the departure of the states $`\widehat{P}(\mathrm{\Gamma })`$ from $`\widehat{\rho }`$ in the Hilbert-Schmidt norm:
$$v^2=\text{tr}\left[\frac{\text{d}}{\text{d}t}\widehat{P}(\mathrm{\Gamma })\widehat{P}(\mathrm{\Gamma })\right]^2.$$
(5)
The smaller $`v`$, the greater is the robustness of the pointer states $`\widehat{P}(\mathrm{\Gamma })`$. Hence one defines the optimum drift of the pure pointer state $`\widehat{P}(\mathrm{\Gamma })\psi \psi ^{}`$ by minimizing $`v`$ GisRig95 . This is given by the nonlinear equation Dio86 ; GisRig95 :
$$\dot{\psi }=(\psi \psi ^{})\psi \psi \psi ^{}\psi .$$
(6)
This result is valid for all kinds of Markovian subdynamics. In our special case (3), it yields the following nonlinear wave equation:
$$\dot{\psi }=\frac{\text{i}}{2m}\widehat{p}^2\psi \frac{D}{2}\left[(\widehat{x}<\widehat{x}>)^2\sigma ^2\right]\psi .$$
(7)
As shown in Dio87b , this equation has a stationary solution which is unique up to Galilean transformations. The wave function of the fiducial stationary state is the complex Gaussian wave packet
$$\psi _0(x)=\left(\alpha _R/2\pi \right)^{1/4}\mathrm{exp}(\alpha x^2/4),$$
(8)
with parameter
$$\alpha \alpha _R+\text{i}\alpha _I=(1\text{i})\sqrt{2Dm}.$$
(9)
The principle of “Hilbert-Schmidt robustness” has thus singled out unique Gaussian pointer states as the robust pure states closest to the true non-unitary local dynamics. The exact width $`\sigma 1/\sqrt{\alpha _R}`$ confirms the heuristic estimate (4). Accordingly, we restrict our further discussion to pointer states $`\widehat{P}(\mathrm{\Gamma })`$ with Gaussian wave functions and make the ansatz
$$\psi _\mathrm{\Gamma }(x)=\left(\alpha _R/2\pi \right)^{1/4}\mathrm{exp}\left(\alpha (x\overline{x})^2/4+\text{i}\overline{p}(x\overline{x})\right),$$
(10)
where $`\mathrm{\Gamma }(\overline{x},\overline{p})^T`$ has been understood. For later purposes, we calculate the correlation matrix $`𝐂`$
$$𝐂\psi _0|\left(\begin{array}{cc}\widehat{x}^2& (\widehat{x}\widehat{p}+h.c.)/2\\ (\widehat{x}\widehat{p}+h.c.)/2& \widehat{p}^2\end{array}\right)|\psi _0=\frac{1}{\alpha _R}\left(\begin{array}{cc}1& \alpha _I/2\\ \alpha _I/2& |\alpha |^2/4\end{array}\right)$$
(11)
of the quantum uncertainties of the canonical observables in the pointer state itself, where $`\psi _0`$ denotes the fiducial state (8). It was shown in JooZeh85 that the states diagonalizing $`\widehat{\rho }`$ exactly are the harmonic oscillator eigenfunctions which are very broad, while narrow eigenfunctions are apparently obtained only for discrete systems. In contrast to these, the above pointer states are well localized. For example, in the situation of a small dust particle ($`m=10^{14}\text{g}`$) scattered by air molecules one has $`D10^{32}\text{cm}^2\text{s}^1`$ JooZeh85 and therefore $`\sigma _0(Dm)^{1/4}10^{11}\text{cm}`$ and $`t_D\sqrt{m/D}10^{10}\text{s}`$.
We now come to the second method. As a preparation, we shall discuss the reduced dynamics of the local system in the basis given by (10). We allow temporarily the parameter $`\alpha `$ to take an arbitrary complex value, and then derive again a distinguished value. If one allows a “natural” time dependence for the probability distribution $`f(\mathrm{\Gamma };t)`$ of the pointer, the asymptotic condition (2) can be turned into an exact identity:
$$\widehat{\rho }(t)=f(\mathrm{\Gamma };t)\widehat{P}(\mathrm{\Gamma })\text{d}\mathrm{\Gamma },t>t_D,$$
(12)
where $`\text{d}\mathrm{\Gamma }\text{d}\overline{x}\text{d}\overline{p}/2\pi `$. This important fact will be proven elsewhere DioKie00 . It generalizes the corresponding statement made in Dio87b for the specific value (9) of $`\alpha `$ as well as the asymptotic statement proved in HalZou95 ; HalZou97 .
From (3) and (12) one can derive an evolution equation for $`f(\mathrm{\Gamma };t)`$:
$$\frac{\text{d}f(\mathrm{\Gamma };t)}{\text{d}t}=\frac{\overline{p}}{m}_{\overline{x}}f(\mathrm{\Gamma };t)+\frac{1}{2}\left[D_{pp}_{\overline{p}\overline{p}}^2+D_{xx}_{\overline{x}\overline{x}}^2+2D_{px}_{\overline{p}\overline{x}}^2\right]f(\mathrm{\Gamma };t),$$
(13)
where the elements of the diffusion matrix are given by
$$𝐃\left(\begin{array}{cc}D_{xx}& D_{xp}\\ D_{px}& D_{pp}\end{array}\right)=\left(\begin{array}{cc}\alpha _I/m\alpha _R& |\alpha |^2/4m\alpha _R\\ |\alpha |^2/4m\alpha _R& D\end{array}\right).$$
(14)
To find a formal solution of (13), we use the Fourier representation $`\stackrel{~}{f}(\stackrel{~}{\mathrm{\Gamma }};t)=f(\mathrm{\Gamma };t)\mathrm{exp}\left[\text{i}(\stackrel{~}{x}\overline{p}\stackrel{~}{p}\overline{x})\right]\text{d}\mathrm{\Gamma }`$ with $`\stackrel{~}{\mathrm{\Gamma }}=(\stackrel{~}{x},\stackrel{~}{p})^T`$. Eq. (13) then leads to
$$\frac{\text{d}\stackrel{~}{f}(\stackrel{~}{\mathrm{\Gamma }};t)}{\text{d}t}=\frac{\stackrel{~}{p}}{m}_{\stackrel{~}{x}}\stackrel{~}{f}(\stackrel{~}{\mathrm{\Gamma }};t)\frac{1}{2}|𝐃|\left[\stackrel{~}{\mathrm{\Gamma }}^T𝐃^1\stackrel{~}{\mathrm{\Gamma }}\right]\stackrel{~}{f}(\stackrel{~}{\mathrm{\Gamma }};t),$$
(15)
where $`|𝐃|`$ denotes the determinant of $`𝐃`$. The solution takes the form
$$\stackrel{~}{f}(\stackrel{~}{\mathrm{\Gamma }};t)=\mathrm{exp}\left[\frac{t}{2}\stackrel{~}{\mathrm{\Gamma }}^T𝐆(t)\stackrel{~}{\mathrm{\Gamma }}\right]\stackrel{~}{f}(\stackrel{~}{x}\stackrel{~}{p}t/m,\stackrel{~}{p};0).$$
(16)
By substitution into (15) one obtains explicitly the matrix of time-dependent coefficients:
$$𝐆(t)=\left(\begin{array}{cc}D_{pp}& D_{xp}+D_{pp}t/2m\\ D_{xp}+D_{pp}t/2m& D_{xx}D_{xp}t/m+D_{pp}t^2/3m^2\end{array}\right).$$
(17)
Eq. (13) can be interpreted as a Fokker-Planck equation provided the diffusion matrix $`𝐃`$ is non-negative. Then the weight function $`f(\overline{x},\overline{p};t)`$ of the pointer states $`\widehat{P}(\overline{x},\overline{p})`$ will drift according to the free-particle dynamics. At the same time the state of the system will diffuse over the pointer states $`\widehat{P}(\overline{x},\overline{p})`$. We can now implement the predictability sieve and minimize the production rate for linear entropy by minimizing the width of the Gaussian pointer states. In other words, we maximize $`\alpha _R`$ under the condition that the diffusion matrix be non-negative. The condition that $`𝐃`$ has a non-negative determinant leads to the condition $`\alpha _R^4+2\alpha _R^2\alpha _I^2+16Dm\alpha _R\alpha _I+\alpha _I^40`$ which can only be fulfilled if $`\alpha _I<0`$, since $`\alpha _R>0`$ for (10) to be normalizable. Introducing dimensionless polar coordinates $`R,\varphi `$ by $`\alpha \sqrt{Dm}R\mathrm{exp}(\text{i}\varphi )`$, this condition reads $`R^2+8\mathrm{sin}2\varphi 0.`$ The maximum for $`\alpha _R=\sqrt{Dm}R\mathrm{cos}\varphi `$ is reached if the equality sign holds, since one could otherwise increase $`R`$ by holding $`\varphi `$ fixed and thus increase $`\alpha _R`$. Maximizing $`\alpha _R`$ under the condition $`R^2=8\mathrm{sin}2\varphi `$ then yields for $`\alpha `$ the following value $`\alpha _s`$ distinguished by the predictability sieve:
$$\alpha _s=3^{1/4}(\sqrt{3}\text{i})\sqrt{Dm},$$
(18)
which coincides, up to a small deviation in the numerical coefficients, with the value (9) following from the criterion of Hilbert-Schmidt robustness. This above slight departure of $`\alpha _s`$ might be related to the fact that the given form of predictability sieve predicts a degenerate diffusion matrix $`𝐃`$. We think, however, that the emerging incoherent uncertainties due to the pointer state diffusion must be made proportional to the quantum uncertainties already present in the pointer states itself. The “robustness of uncertainties” demands that the matrix $`𝐂`$ (11) of quantum correlations be proportional to the diffusion matrix $`𝐃`$ (14) of the corresponding classical coordinates for the pointer. From the condition that $`𝐂=const\times 𝐃`$ we then obtain again the standard value (9) for $`\alpha `$, while $`𝐂=m/2D\times 𝐃`$.
We shall now discuss our last method to determine the pointer basis, which will involve quantum state diffusion. As we see from (12) and (13), the quantum state of the system, when expanded as a mixture of pointer states, performes diffusion after the decoherence time has elapsed. This diffusion will, by construction, preserve the shape (8) of the Gaussian wave packet, and only its center will walk randomly. It is then natural to ask, whether there is a generic QSD process which, first, applies to generic pure initial states and, second, tends to the above specific diffusion process for $`tt_D`$.
As is well known, the Fokker-Planck equation (13) is equivalent to the Itô-Langevin equation Ar
$$\text{d}\mathrm{\Gamma }=V\text{d}t+\text{d}X,$$
(19)
where $`V=(\overline{p}/m,0)`$, and $`\text{d}X=(\text{d}\xi ,\text{d}\pi )`$ is the increment of a zero-mean Gaussian white noise with correlation matrix $`𝐃\text{d}t`$. In case of phase-space diffusion the use of the Itô-Langevin formalism instead of the Fokker-Planck formalism is a matter of taste. But the diffusion of the corresponding pointer states $`\psi _\mathrm{\Gamma }`$ would be quite awkward in the Fokker-Planck formalism. We thus choose the Itô-formalism and apply (19) to the Gaussian pointer states (10). This leads to, substituting $`\widehat{x}=\overline{x}`$ and $`\widehat{p}=\overline{p}`$,
$$\text{d}\psi =\frac{\text{i}}{2m}\widehat{p}^2\psi \text{d}t\frac{D}{2}\left(\widehat{x}\widehat{x}\right)^2\psi \text{d}t+(\widehat{x}\widehat{x})\psi \text{d}z,$$
(20)
where the index $`\mathrm{\Gamma }`$ has been skipped. The deterministic part of the evolution is governed, up to normalization, by the same nonlinear wave equation (7) which we had obtained from the Hilbert-Schmidt robustness, while the random part is driven by the complex Gaussian white noise
$$\text{d}z=\frac{\alpha }{2}\text{d}\xi +\text{i}\text{d}\pi .$$
(21)
Since the correlation matrix of $`\text{d}X=(\text{d}\xi ,\text{d}\pi )`$ is $`𝐃\text{d}t`$, (14,21) yields
$$M[\text{d}z\text{d}z^{}]=D\text{d}t$$
(22)
for the mean of the Hermitian correlation, independent of $`\alpha `$. But the correlation $`M[\text{d}z\text{d}z]`$ still depends on $`\alpha `$. A most remarkable feature of (20) is that any reference to the phase-space variables $`\mathrm{\Gamma }=(\overline{x},\overline{p})`$ has been cancelled. For this reason we have omitted the subscript $`\mathrm{\Gamma }`$ from $`\psi `$ and extend the validity of the equation to arbitrary initial state vectors. It is possible to prove Dio88b ; GatGis91 ; HalZou95 ; Kol95 that, starting from whatever initial state $`\psi (0)`$, the random solution $`\psi (t)`$ will tend to be the Gaussian pointer state $`\psi _{\mathrm{\Gamma }(t)}`$, where $`\mathrm{\Gamma }(t)`$ is governed by the diffusion process (19). Eq. (20) is called the Itô-Schrödinger equation of QSD.
Since the free parameter $`\alpha `$ still appears in $`M[\text{d}z\text{d}z]`$, we are left with the non-uniqueness problem of the QSD equations. If one, however, chooses the distinguished value (9) of $`\alpha `$, one finds, using (21), the simple result
$$M[\text{d}z\text{d}z]=M[\text{d}z^{}\text{d}z^{}]=0,$$
(23)
distinguishing a unique QSD. Historically, this unique QSD was in the Fokker-Planck formalism singled out by certain invariance considerations Dio88c ; Per90 . The Itô-Schrödinger equation (20) with the complex Gaussian white noise (22,23) has become the dominating formalism of standard QSD theory qsd extended for arbitrary Markovian reduced dynamics. Applying exact forms of robustness criteria we have thus obtained a unique QSD which leads to stationary Gaussian pointer states for $`tt_D`$, whose centers undergo a diffusion process. With heuristic forms of robustness one could have chosen other QSD equations like in Dio88b ($`\alpha _R=2\sqrt{Dm}`$) or HalZou95 ($`\alpha _R=\sqrt{Dm}/2\sqrt{2}`$). The recent proposal of “maximal survival probability” from WV differs from our first method and does not lead to (23) DioKie00 .
In conclusion, we have demonstrated that three different methods of dynamical robustness lead to an essentially unique local pointer basis in case of Markovian local dynamics. The corresponding pointer states follow the classical trajectories up to a tiny random diffusion. Well-defined stochastic differential equations, known from the theory of quantum state diffusion, govern both the formation and the diffusion of pointer states. These states can thus be used to characterize local quasiclassical properties. The pointer states are not an absolute property of the system in itself, but only characterize certain stability properties with respect to interactions with the environment: They are least sensitive to quantum entanglement, which is why interference terms between them cannot be noticed by local observers. They possess thus meaning with respect to an observer-related branch of the total wave function or a component corresponding to a potential fundamental collapse deco ; Zeh99 , while the interaction with the environment is encoded in the choice of our master equation (3).
We thank the Institute for Advanced Study Berlin and the Isaac Newton Institute, Cambridge, for their kind hospitality while this work was begun. L.D. thanks the ESF QIT program and the Hungarian OTKA Grant 032640 for financial support. Critical comments by E. Joos, H. Wiseman, and H.D. Zeh are gratefully acknowledged.
|
warning/0005/cond-mat0005493.html
|
ar5iv
|
text
|
# Adsorbed states of a long - flexible polymer chain
## Abstract
A phase diagram for a surface-interacting long flexible polymer chain in a two-dimensional poor solvent where the possibility of collapse exists is determined using exact enumeration method. A model of a self-attracting self avoiding walk (SASAW) on a square lattice was considered and up to 28 steps in series were evaluated. A new adsorbed state having the conformation of a surface attached globule is found. Four phases (i) desorbed expanded, (ii) desorbed collapsed, (iii) adsorbed expanded and (iv) surface attached globule are found to meet at a point on the adsorption line.
64.60.-i,68.35.Rh,5.50.+q
The behaviour of a long flexible polymer chain near an impenetrable surface is a subject of considerable experimental and theoretical importance . This is because such a chain may exhibit a phase diagram characterized by many different universality domains of critical behaviour. The subtle competition between the gain of internal energy and corresponding loss of entropy at the surface may lead to the possibility of the coexistence of different regimes and multicritical behaviour. An attracting surface may lead adsorption - desorption transition from the state when the chain is mostly attached to the surface, to the state of detachment when the temperature is increased. This behaviour finds applications in lubrication, adhesion, surface protection etc.
The essential physics associated with the behaviour of a surface interacting polymer chain in a good solvent where monomer-solvent attraction is greater than the monomer-monomer attraction, is derived from a model of self-avoiding walk (SAW) on a semi-infinite lattice. If the surface is attractive, it contributes an energy $`ϵ_a`$ ($`<0`$) for each step of the walk along the lattice boundary. This leads to an increased probability characterized by the Boltzmann factor $`\omega =\mathrm{exp}(ϵ_a/k_\beta T)`$ of making a step along the wall, since for $`ϵ_a<0`$, $`\omega >1`$ for any finite temperature $`T`$ ($`k_\beta `$ is the Boltzmann constant). Because of this the polymer chain becomes adsorbed at low temperatures on the surface while at high temperatures all polymer conformations have almost same weight and non adsorbed (or desorbed) behaviour prevails. The transition between these two regimes is marked by a critical adsorption temperature $`T_a`$, with a desorbed phase for $`T>T_a`$ and adsorbed phase for $`T<T_a`$. At $`T=T_a`$ one may define the crossover exponent $`\varphi `$, as $`MN^\varphi `$, where $`N`$ is the total number of steps and $`M`$ the number of steps on the surface. The transition point $`T_a`$ is a tricritical point . Both the surface and the bulk critical exponents have been calculated using renormalization group methods , exact enumeration methods and Monte-Carlo simulations . For a two dimensional system exact values of the exponents have been found by using conformal invariance .
The situation, is, however, different when the surface interacting polymer chain is in a poor solvent where monomer-monomer attraction dominates over the monomer-solvent attraction. As is well known, a long-flexible polymer chain in a poor solvent exhibits a transition from a compact globule (collapsed state) to a expanded state when the temperature is increased. Above the critical $`\theta `$ temperature (often referred to as the $`\theta `$-point) the chain behaves as it would in a good solvent and below this temperature it behaves like a compact globule. At the $`\theta `$-point the chain behaviour is described by a tricritical point of the $`O(n)`$ ($`n0`$) spin system . However, when the chain is in the vicinity of an impenetrable surface the competition between the monomer-monomer attraction and the surface-monomer interaction gives rise to many new features. Attempts have been made to study these features using several approaches .
For two-dimensions the transfer matrix method has been used for a directed polymer chain whereas for the nondirected (isotropic) version the exact enumeration method has been used . In both cases, three phases, desorbed expanded, desorbed collapsed and a single adsorbed phase have been predicted. However, the true nature of the phase diagram remained unknown. For three dimensions, the Monte-Carlo simulations method has been used for a finite length ($``$ 100) chain which led to a phase diagram containing four phases; desorbed expanded (DE), desorbed collapsed (DC), adsorbed expanded (AE) and adsorbed collapsed (AC). The phase diagram shows a phase boundary between the AE and DC phases leading to two points on the phase diagram where three phases coexists(triple point). However, the phase diagram found by the exact enumeration technique has many features which are different from that found in ref. . This indicates the possibility of a richer phase diagram than has been realized so far. In view of this we reexamine the problem of simultaneous adsorption and collapse of a linear polymer chain on a square lattice and investigate the phase diagram and critical parameters using the exact enumeration technique. We prefer this technique because in this case the scaling corrections are correctly taken into account by a suitable extrapolation scheme. As shown by Grassberger and Hegger , to achieve the same accuracy with the Monte Carlo method one has to consider a polymer chain of about two orders of magnitude longer than in the exact enumeration method.
We consider SASAW on a square lattice restricted to half space $`Z0`$ (impenetrable hard wall). Walk starts from the middle of the surface. Let $`C_{N,N_s,N_m}`$ be the number of SAWs with $`N`$ steps, having $`N_s`$ $`(N)`$ step on the surface and $`N_m`$ nearest neighbor. We have obtained $`C_{N,N_s,N_m}`$ for $`N28`$ for square lattice by exact enumeration method.
Now we consider the interaction energy $`ϵ_a`$ associated with each walk on the surface and $`ϵ_m`$ for monomer-monomer interaction. Partition function of the attached chain is
$$Z_N(\omega ,u)=\underset{N_s,N_m}{}C_{N,N_s,N_m}\omega ^{N_s}u^{N_m}$$
(1)
where $`\omega =e^{ϵ_a/kT}`$ and $`u=e^{ϵ_m/kT}`$. $`\omega >1`$ and $`u>1`$ for attractive force. Reduced free energy for the chain can be written as
$$G(\omega ,u)=\underset{N\mathrm{}}{lim}\frac{1}{N}\mathrm{log}Z_N(\omega ,u)$$
(2)
In general it is appropriate to assume that as $`N\mathrm{}`$
$$Z_N(\omega ,u)N^{\gamma 1}\mu (\omega ,u)^N$$
(3)
where $`\mu (\omega ,u)`$ is the effective coordination number and $`\gamma `$ is the universal configurational exponents for walks with one end attached to the surface. The value of $`\mu (\omega ,u)`$ can be estimated using ratio method with associated Neville table. From equations (2) and (3) we can write
$$\mathrm{log}\mu (\omega ,u)=\underset{N\mathrm{}}{lim}\frac{1}{N}\mathrm{log}Z_N(\omega ,u)=G(\omega ,u)$$
(4)
$`Z_N(\omega ,u)`$ is calculated from the data of $`C_{N,N_s,N_m}`$ using equation (1) for a given $`\omega `$ and $`u`$. From this we construct linear and quadratic extrapolants of the ratio of $`Z_N(\omega ,u)`$ for the adjacent values of $`N`$ as well as the alternate one. Results for alternate $`N`$ give better convergence. When $`u=1`$ and $`\omega =1`$ the value of $`\mu `$ is found to be 2.638 which is in very good agreement with the value given in ref. .
The value of $`\omega _c(u)`$ at which polymer gets adsorbed for a given value of $`u`$ is found from the $`(i)`$ plot of $`G(\omega ,u)`$ which remains fairly constants until $`\omega =\omega _c`$ and increases consistently as a function of $`\omega `$, for $`\omega \omega _c`$ $`(ii)`$ from the plot of $`^2G(\omega ,u)/ϵ_{a}^{}{}_{}{}^{2}`$ at constant $`u`$ and $`(iii)`$ from the plot of $`\gamma ^0\gamma _1`$ ( see Eq.(6) and discussions which follow it) as a function of $`\omega `$ for different $`N`$. The value of $`\omega _c`$ found from the plot of $`G(\omega ,u)`$ is slightly lower than the peak value of $`^2G/ϵ_{a}^{}{}_{}{}^{2}`$. It is, however, observed that as $`N`$ is increased from $`22`$ to $`28`$ the peak value shifts to smaller $`\omega `$ and appears to converge on the value of $`\omega `$ found from $`G(\omega ,u)`$ plot. We therefore choose the value of $`\omega _c`$ found from the plot of $`G(\omega ,u)`$ and determine lines $`\omega _c(u)`$ and $`\omega _{c1}(u)`$ (see Fig.1) by this method. For $`u=1`$, the value of $`\omega _c`$ is 2.050 which is in very good agreement with the value (= $`2.044\pm 0.002`$) reported in ref. . Similarly the phase boundary separating the extended and collapsed phases is calculated from the plot of $`G(\omega ,u)`$ as a function of $`u`$ for a given $`\omega `$. However, transition point $`u_c`$ is located more accurately from the peak of $`^2G(\omega ,u)/ϵ_{m}^{}{}_{}{}^{2}`$ at constant $`\omega `$. For $`\omega =1`$, the value of $`u_c`$ is 1.93 which is in good agreement with the value found by Foster et al and the Monte Carlo results (= 1.94 $`\pm `$ 0.004) . The method is found to work for all values of $`\omega `$ i.e in both the bulk and the adsorbed regimes. However, as $`\omega `$ is increased the values of $`G(\omega ,u)`$ do not remain as smooth as at lower values of $`\omega `$, therefore introducing some inaccuracy in the value of $`u_c`$. The estimate of this inaccuracy is of the order of $`5\%`$ for $`\omega >4`$. We therefore conclude that the $`u_c`$ and $`\omega _{c2}`$ (Fig.1) lines are determined with reasonable accuracy.
The surface critical exponent $`\gamma _1^{N,k}`$ can be calculated using the relation:
$$\gamma _1^{N,k}=\frac{\mathrm{log}(Z_N/Z_{N2})k\mathrm{log}(\mu )}{\mathrm{log}(N/(Nk)}+1$$
(5)
where subscript “1” indicates the corresponding quantity of the surface (with one end of the polymer is attached to the surface). If we assume that $`\mu `$ does not depend on $`\omega `$ and polymer is in desorbed phase, then we can calculate the quantity $`\gamma ^0\gamma _1`$ from the above equation
$$\gamma ^0\gamma _1=\frac{\mathrm{log}(Z_N^0Z_{N2}/Z_{N2}^0Z_N)}{\mathrm{log}(N/N2)}$$
(6)
where superscript “0” indicates the corresponding quantity of the bulk (i.e without surface). In this case one calculates $`\gamma ^0\gamma _1`$ for different $`N`$ using above equation and plot it as a function of $`\omega `$. The location of adsorption point $`\omega _c`$ can be determined from the intersection of successive approximation to $`\gamma ^0\gamma _1`$ in the limit $`N\mathrm{}`$. This method, however, fails beyond the $`\theta `$-point and reproduces closely the adsorption phase boundary below the $`\theta `$-point found by the method discussed above.
The phase diagram shown in Fig. 1 has four phase boundaries instead of three as reported in earlier work . The $`u_c`$ line separates the expanded and collapsed phases. This line remains straight and parallel to $`\omega `$-axis in the bulk. This result is in agreement with that of ref. . The special adsorption line $`\omega _c`$ separates adsorbed expanded (AE) phase from that of desorbed expanded (DE). Beyond $`\theta `$-point we have two boundaries $`\omega _{c1}`$ and $`\omega _{c2}`$. The line $`\omega _{c1}`$ separates the desorbed collapsed (DC) bulk phase from that of an adsorbed globule (AG) state, whereas the boundary $`\omega _{c2}`$ separates the AG phase from the AE phase. The point where $`u_c`$ line meets the special adsorption line $`\omega _c`$, all the four phases AE, DE, DC and AG coexist. The AG phase which exists between the boundaries $`\omega _{c1}`$ and $`\omega _{c2}`$ for $`u>u_c`$ is essentially a two-dimensional globule sticking to the surface in the same way as a liquid drop may lie on a surface. The existance of such a phase, to the best of our knowledge, is shown for the first time.
FIG.1: The phase diagram of a surface interacting linear polymer in 2-D space. $`\omega `$ and $`u`$ axes represent, respectively the Boltzmann factor of surface interaction and monomer-monomer attraction. Regions marked by AE, AG, DE and DC represent, respectively, the adsorbed polymer in expanded ( swollen) state and globular state, desorbed polymer in expanded and collapsed state.
To find the monomer density as a function of distance normal to the surface for different regimes of the phase diagram we first found the terms which make most significant contribution in Eq.(1) for the walks of 26 steps. For these walks we evaluated the monomers (number of visited sites) on different lattice layers. The results are shown in Fig. 2 where we plot the fraction of monomers lying on different layers for many values of $`u`$ and $`\omega `$ corresponding to different parts of the phase diagram. Since the walks always start from the surface, there is at least one monomer on the surface for all the cases. In Fig. 2(a) $`u`$ is taken equal to 1.5 which is less than $`u_c(1.93)`$ and therefore it corresponds to expanded state. The values of $`\omega =1.0`$ and 3.5 correspond, respectively to DE, and AE phases, while $`\omega =2.20`$ lies on the special adsorption line $`w_c`$ for $`u=1.5`$. Fig. 2(b) shows the change in monomer density as $`\omega `$ is increased for $`u=2.5`$. In this case $`\omega =1.0`$, 3.4 and 3.59 correspond to DC, AG and AE phases, respectively. For $`u=2.5`$ $`\omega _{c1}=2.75`$ and $`\omega _{c2}=3.50`$, thus the value $`\omega =3.59`$ is just above the $`\omega _{c2}`$ line. The large change in monomer density distribution when $`\omega `$ value is changed from 3.4 (slightly below $`\omega _{c2}`$ line) to 3.59 (just above $`\omega _{c2}`$ line) is evident and confirms the existence of new (AG) phase.
In Fig. 3 we plot the quantities $`<n_s>`$ and $`<n_p>`$ giving the average fraction of monomers on the surface and number of pairs and defined as
$$<n_s>=\underset{N\mathrm{}}{lim}\frac{\mathrm{ln}G}{\omega }_u,<n_p>=\underset{N\mathrm{}}{lim}\frac{\mathrm{ln}G}{u}_\omega $$
FIG.2: Fraction of monomers (n) on different layer.
The transition points on each curve is marked by dot. While the dot on a curve of $`<n_s>`$ (in Fig. 3(a)) indicate the adsorption transition (i.e point on $`\omega _c`$ and $`\omega _{c1}`$ lines depending on the values of $`u`$), the dot on a curve $`<n_p>`$ (in Fig. 3(b)) indicates the transition to collapsed state (i.e point on line $`u_c`$ and $`\omega _{c2}`$ depending on the values of $`\omega `$). We may note that the value of $`<n_p>`$ for AG phase is comparable to that in the DC phase. We have found that along the line $`\omega _c`$ and $`\omega _{c1}`$, $`<n_s>=0.07\pm 0.004`$ and along the line $`u_c`$ and $`\omega _{c2}`$, $`<n_p>=0.5\pm 0.003`$. This shows that the error involved in determining the phase boundaries is very small.
It is obvious from these results that when the chain gets adsorbed from the expanded bulk state (i.e for $`u<u_c`$), it acquires a conformation at $`\omega \omega _c(u)`$ such that a small fraction ($`10\%`$) of monomers get attached to the surface (see Fig. 3(a)) and others are still in the bulk. Though the chain has formed a layer parallel to the surface there are considerable fluctuations in direction normal to the surface layer. As $`\omega `$ is increased for the same value of $`u`$, the fluctuations along the normal to the surface get suppressed and at large $`\omega (>>\omega _c)`$ the chain lies on the surface with very little fluctuations (see Fig. 2(a)). On the other hand, when the adsorbing chain was in collapsed bulk state, then at $`\omega =\omega _{c1}`$ the collapsed chain gets attached to the surface. Here again the number of monomers getting attached to the surface are about $`10\%`$ (see Fig. 3(a)). For $`\omega _{c1}\omega \omega _{c2}(u)`$ the chain remains in the form of globule attached to the surface. In this range the monomer-monomer attraction remains effective in holding the monomers in the neighbourhood of each other than the surface-monomer attraction whose tendency is to spread the chain on the surface (see Fig. 2(b)). For $`\omega >\omega _{c2}(u)`$ the globule conformation becomes unstable as surface-monomer attraction becomes more effective than the monomer-monomer attraction and therefore the chain spreads over the surface (just like a liquid spreads over a wetting surface).
FIG.3: (a) Average fraction of monomers on the surface ($`<n_s>`$) as a function of $`\omega `$ and (b) number of pairs ($`<n_p>`$) as a function of $`u`$. Dot on a curve indicates the transition point; in (a) adsorption - desorption and in (b) expanded - collapsed transitions.
Summarizing we studied a SASAW in the presence of an attracting impenetrable wall and obtained the phase boundaries separating different phases of the polymer chain from data obtained by exact enumerations. We report a new adsorbed state which has conformation of a compact globule sticking to a surface in same way as a liquid drop may lie on a non wetting surface. The monomer density distribution, the number of monomers on the surface and the number of nearest neighbours in different regimes of the phase diagram are obtained.
The work was supported by the Department of Science and Technology (India) through a project grant.
|
warning/0005/math-ph0005021.html
|
ar5iv
|
text
|
# 1 Introduction
## 1 Introduction
The purpose of this paper is to provide a complete description of the non-dynamical, constant r-matrices of the standard Calogero-Moser models associated with degenerate potential functions, which can be obtained by gauge transformations of their usual Lax representation. A preliminary account of a part of this work is contained in .
The Calogero-Moser type many particle systems (for a review, see ) have been much studied recently due to their fascinating mathematics and applications ranging from solid state physics to Seiberg-Witten theory. The definition of these models involves a root system and a potential function depending on the inter-particle ‘distance’. The potential is given either by the Weierstrass $`𝒫`$-function or one of its (hyperbolic, trigonometric or rational) degenerations. The classical equations of motion of the models admit Lax representations,
$$\dot{L}=[L,M],$$
(1)
which underlie their integrability. A Lax representation of the Calogero-Moser models based on the root systems of the classical Lie algebras was found by Olshanetsky and Perelomov using symmetric spaces. Recently new Lax representations for these systems as well as their exceptional Lie algebraic analogues and twisted versions have been constructed .
In general , Liouville ingtegrability can be understood as a consequence of the Poisson brackets of the Lax matrix having the r-matrix form,
$$\{L_1,L_2\}=\{L^\mu ,L^\nu \}T_\mu T_\nu =[r_{12},L_1][r_{21},L_2],$$
(2)
where $`r_{12}=r^{\mu \nu }T_\mu T_\nu `$ ($`r_{21}=r^{\mu \nu }T_\nu T_\mu `$) with some constant matrices $`T_\mu `$. The components $`L^\mu `$ of the Lax matrix $`L=L^\mu T_\mu `$ encode the phase space variables, and the components $`r^{\mu \nu }`$ of the classical r-matrix may in general depend on the same variables as $`L_1=L1`$ and $`L_2=1L`$. Of course, $`L`$ and $`r`$ depend also on a spectral parameter in general, but this does not occur for the systems of our interest, and thus is suppressed in (2). When the r-matrix really does depend on the phase space variables, one says that it is ‘dynamical’.
The classical r-matrix has been calculated first for the standard Lax representation of the $`gl_n`$ Calogero-Moser systems associated with degenerate potentials , and then for Krichever’s spectral parameter dependent Lax matrix in the elliptic case . The r-matrices found in these papers are dynamical, but depend only on the coordinates of the particles. These r-matrices have been re-derived by means of Hamiltonian reduction in , and in a very recent paper they have been generalized explicitly for the $`BC_n`$ system as well as for all classical Lie algebras. In the physically most interesting $`gl_n`$ case, dynamical r-matrices have also been found for the relativistic deformations of the Calogero-Moser models introduced by Ruijsenaars and Schneider . Then the quantization of the non-relativistic and the relativistic models has been investigated in a new framework based on quantum dynamical R-matrices.
The above developments have close connections with the new theory of dynamical r-matrices and associated quantized structures reviewed in . However, since the present understanding of most integrable systems involves constant (i.e. ‘non-dynamical’) r-matrices, which form a direct link to Poisson-Lie groups and quantum groups , it is natural to ask if the Lax representation of the Calogero-Moser models can be chosen in such a way to exhibit non-dynamical r-matrices. The obvious way to search for new Lax representations with this property is to perform gauge transformations on the usual Lax representations. In the elliptic case of the standard $`gl_n`$ models a new Lax representation associated with Belavin’s constant elliptic r-matrix has recently been found in this way . To be more precise, the results of are already contained in a somewhat less explicit form in the seminal paper by Hasegawa , where the commuting Ruijsenaars operators have been interpreted as commuting transfer matrices based on a realization of the $`RLL=LLR`$ relation with Belavin’s elliptic R-matrix and certain difference $`L`$-operators. In fact, the dynamical twisting and the classical and non-relativistic limits of the $`L`$-operator leading to Krichever’s Lax matrix for the elliptic Calogero-Moser model are indicated in (see also ). Then in the paper some delicate limit procedures have been considered, whereby non-dynamical R-matrices can be obtained for the trigonometric degenerations of the Ruijsenaars-Schneider and Calogero-Moser models. The resulting R-matrix was found to be non-unique, one possibility being the spectral parameter independent Cremmer-Gervais R-matrix discovered in a different context in .
It is clear from the above that Lax representations for the degenerate Calogero-Moser models with non-dynamical r-matrices can be obtained by taking limits of Hasegawa’s $`RLL=LLR`$ relation. However, the details of the admissible limiting procedures appear rather complicated and the starting point requires familiarity with quite advanced results. In this circumstance, it might be worthwhile to understand the possible non-dynamical r-matrices also from an elementary viewpoint. This is the objective of the present paper, where we aim to perform a self-contained, systematic analysis of the gauge transformations of the usual Lax representation of the degenerate Calogero-Moser models that lead to constant r-matrices.
The organization and the main results of our work are as follows. First, we describe the most general momentum independent dynamical r-matrices for the standard Lax representation in section 2. This amounts to a slight but necessary generalization of the Avan-Talon r-matrix as given by Theorem 1. Second, we select those dynamical r-matrices that become constant by a gauge transformation (defined by eq. (18)) and determine the corresponding ‘gauge potentials’$`A_k(q)`$. This is the content of section 3, in particular Proposition 2 and Theorem 3. Third, in section 4 we compute explicitly the gauge transformations $`g(q)`$ (from eq. (19)) and the resulting most general constant r-matrix, which is given by Theorem 6. It turns out that in the rational case the constant r-matrix is conjugate to the antisymmetric solution of the classical Yang-Baxter equation that belongs to the Frobenius subalgebra of $`gl_n`$ consisting of the matrices with vanishing last row . In the hyperbolic/trigonometric cases the $`sl_n`$-part of the most general $`gl_ngl_n`$-valued constant r-matrix (see Proposition 7) is equivalent to a multiple of the Cremmer-Gervais classical r-matrix , and it can also be made equal to it by a choice of the gauge transformation. This identification of the constant Calogero-Moser r-matrices is presented in section 5. The main results are summarized once more in the conclusion, which occupies section 6. Except for the notations introduced in section 2, this final section is self-contained and it may be useful to consult it before reading the main text. The details of some proofs are contained in three appendices.
The outcome of our direct analysis of the degenerate Calogero-Moser models is consistent with the previous results . In addition to the advantage that our analysis is elementary, we also clarify the extent to which the constant r-matrix is unique in the degenerate cases. In principle, this uniqueness question cannot be answered by studying the limits of the elliptic case, even though in the final analysis it follows that all our constant r-matrices can be regarded as various degenerations (see also ) of Belavin’s elliptic r-matrix.
## 2 Momentum independent dynamical r-matrices
The standard (degenerate) Calogero-Moser-Sutherland models are defined by the Hamiltonian
$$h=\frac{1}{2}\underset{k=1}{\overset{n}{}}p_k^2+\underset{k<l}{}v(q_kq_l),$$
(3)
where $`v`$ is given as
$$v(x)=\{\begin{array}{cc}x^2,& \text{rational case}\\ a^2\mathrm{sinh}^2(ax),& \text{hyperbolic case}\\ a^2\mathrm{sin}^2(ax),& \text{trigonometric case.}\end{array}$$
(4)
One has the canonical Poisson brackets $`\{p_k,q_l\}=\delta _{k,l}`$, the coordinates are restricted to a domain in $`𝐑^n`$ where $`v(q_kq_l)<\mathrm{}`$, and $`a>0`$ is a parameter.
Let us fix the following notation for elements of the Lie algebra $`gl_n`$:
$$H_k:=e_{kk},E_\alpha :=e_{kl},H_\alpha :=(e_{kk}e_{ll}),K_\alpha :=(e_{kk}+e_{ll})\text{for}\alpha =\lambda _k\lambda _l\mathrm{\Phi }.$$
(5)
Here $`\mathrm{\Phi }=\{(\lambda _k\lambda _l)|kl\}`$ is the set of roots of $`gl_n`$, $`\lambda _k`$ operates on a diagonal matrix, $`H=\mathrm{diag}(H_{1,1},\mathrm{},H_{n,n})`$ as $`\lambda _k(H)=H_{k,k}`$, and $`e_{kl}`$ is the $`n\times n`$ elementary matrix whose $`kl`$-entry is $`1`$. Moreover, we denote the standard Cartan subalgebra of $`sl_ngl_n`$ as $`_n`$, and put $`p=_{k=1}^np_kH_k`$, $`q=_{k=1}^nq_kH_k`$, $`\mathrm{𝟏}_n=_{k=1}^nH_k`$.
From the list of known Lax representations we consider the original one for which $`L`$ is the $`gl_n`$ valued function
$$L(q,p)=p+\sqrt{1}\underset{\alpha \mathrm{\Phi }}{}w(\alpha (q))E_\alpha ,$$
(6)
where the real function $`w`$ is chosen according to
$$w(x)=\{\begin{array}{c}x^1\\ a\mathrm{sinh}^1(ax)\\ a\mathrm{sin}^1(ax).\end{array}$$
(7)
Then the function
$$F:=\frac{w^{}}{w}$$
(8)
enjoys the important identities
$$F^{}=w^2,$$
(9)
$$F(x)+F(y)=\frac{w(x)w(y)}{w(x+y)},$$
(10)
$$F(xy)\left(F(x)F(y)\right)+F(x)F(y)=,$$
(11)
where, respectively to the cases above,
$$=\{\begin{array}{c}0\\ a^2\\ a^2.\end{array}$$
(12)
For any real function $`f`$ (like $`v`$, $`w`$ or $`F`$), we introduce the functions $`f_k`$ and $`f_\alpha `$ of $`q`$ as
$$f_k(q):=f(q_k),f_\alpha (q)=f(\alpha (q)),$$
(13)
and sometimes write $`f_{kl}`$ for $`f_\alpha `$ if $`\alpha =(\lambda _k\lambda _l)`$. As an $`n\times n`$ matrix $`L_{k,l}=p_k\delta _{k,l}+\sqrt{1}(1\delta _{k,l})w(q_kq_l)`$, but $`L`$ can also be used in any other representation of $`gl_n`$. The r-matrix corresponding to this $`L`$ was studied by Avan and Talon , who found that it is necessarily dynamical, and may be chosen so as to depend on the coordinates $`q_k`$ only. We next describe a slight generalization of their result.
Theorem 1. The most general $`gl_ngl_n`$-valued r-matrix that satisfies (2) with the Lax matrix in (6) and depends only on $`q`$ is given by
$$r(q)=\underset{\alpha \mathrm{\Phi }}{}F_\alpha (q)E_\alpha E_\alpha +\frac{1}{2}\underset{\alpha \mathrm{\Phi }}{}w_\alpha (q)(C_\alpha (q)K_\alpha )E_\alpha +\mathrm{𝟏}_nQ(q),$$
(14)
where the $`C_\alpha (q)`$ are $`_nsl_n`$ valued functions subject to the conditions
$$C_\alpha (q)=C_\alpha (q),\beta (C_\alpha (q))=\alpha (C_\beta (q))\alpha ,\beta \mathrm{\Phi }$$
(15)
and $`Q(q)`$ is an arbitrary $`gl_n`$-valued function.
Remarks. The functions $`C_\alpha `$ can be given arbitrarily for the simple roots, and are then uniquely determined for the other roots by (15). The r-matrix found by Avan and Talon is recovered from (14) with $`C_\alpha 0`$; and we refer to $`r(q)`$ in (14) as the Avan-Talon r-matrix in its general form. Given that this holds for the Avan-Talon r-matrix, the fact that $`r(q)`$ above satisfies (2) with any $`Q(q)`$ and $`C_\alpha (q)`$ subject to (15) is easy to verify. Theorem 1 can be proved by a careful calculation along the lines of . For the details, see appendix A.
## 3 Is $`r(q)`$ gauge equivalent to a constant?
A gauge transformation of a given Lax representation (1) has the form
$$LL^{}=gLg^1,MM^{}=gMg^1\frac{dg}{dt}g^1,$$
(16)
where $`g`$ is an invertible matrix function on the phase space. If $`L`$ satisfies (2), then $`L^{}`$ will have similar Poisson brackets with a transformed r-matrix $`r^{}`$. The question now is whether one can remove the $`q`$-dependence of any of the r-matrices in (14) by a gauge transformation. It is natural to assume this gauge transformation to be $`p`$-independent, i.e. defined by some function $`g:qg(q)GL_n`$. In this case we find that
$$\{L_1^{},L_2^{}\}=[r_{12}^{},L_1^{}][r_{21}^{},L_2^{}]$$
(17)
holds with
$$r^{}(q)=\left(g(q)g(q)\right)\left(r(q)+\underset{k=1}{\overset{n}{}}A_k(q)H_k\right)\left(g(q)g(q)\right)^1,$$
(18)
$$A_k(q):=g^1(q)_kg(q),_k:=\frac{}{q_k}.$$
(19)
The meaning of this formula is that if $`r(q)`$ is the most general $`p`$-independent r-matrix for which $`L`$ (6) satisfies (2), then $`r^{}(q)`$ has the analogous property in relation to $`L^{}`$.
We wish to find $`r(q)`$ and $`g(q)`$ such that $`_kr^{}=0`$. On account of (18) this is equivalent to
$$_k(r+\underset{l=1}{\overset{n}{}}A_lH_l)+[r+\underset{l=1}{\overset{n}{}}A_lH_l,A_k\mathrm{𝟏}_n+\mathrm{𝟏}_nA_k]=0.$$
(20)
By using (19), whereby
$$_kA_l_lA_k+[A_l,A_k]=0,$$
(21)
it is useful to rewrite (20) as
$$_kr+\underset{l=1}{\overset{n}{}}_lA_kH_l+[r,A_k\mathrm{𝟏}_n+\mathrm{𝟏}_nA_k]+\underset{l=1}{\overset{n}{}}A_l[H_l,A_k]=0.$$
(22)
Our strategy is to first find $`A_k(q)`$ and $`r(q)`$ from eqs. (21), (22), and then determine $`g(q)`$ and the resulting constant r-matrix. For this we now parametrize $`A_k`$ as
$$A_k(q)=\underset{l=1}{\overset{n}{}}A_k^l(q)H_l+\underset{\alpha \mathrm{\Phi }}{}A_k^\alpha (q)E_\alpha ,$$
(23)
and expand the r-matrix from Theorem 1 in the form
$$r(q)=\underset{\alpha }{}F_\alpha (q)E_\alpha E_\alpha +\underset{i,\alpha }{}r_i^\alpha (q)H_iE_\alpha +\underset{i}{}Q^i(q)\mathrm{𝟏}_nH_i.$$
(24)
We here have
$$r_i^\alpha (q)=Q^\alpha (q)+\frac{1}{2}w_\alpha (q)\mathrm{t}r\left(H_i(C_\alpha (q)K_\alpha )\right),$$
(25)
$$Q(q)=\underset{i=1}{\overset{n}{}}Q^i(q)H_i+\underset{\alpha \mathrm{\Phi }}{}Q^\alpha (q)E_\alpha ,$$
(26)
where $`Q(q)`$, $`C_\alpha (q)`$ and $`K_\alpha `$ appear in (14).
With reference to the conventions (5), we define the structure constants $`c_{\alpha ,\beta }^{\alpha +\beta }`$ by writing $`[E_\alpha ,E_\beta ]=c_{\alpha ,\beta }^{\alpha +\beta }E_{\alpha +\beta }`$ if $`\alpha ,\beta ,(\alpha +\beta )`$ all belong to $`\mathrm{\Phi }`$, and $`c_{\alpha ,\beta }^{\alpha +\beta }:=0`$ otherwise. Then (21) yields
$`_lA_k^i_kA_l^i={\displaystyle \underset{\alpha \mathrm{\Phi }}{}}\alpha _iA_l^\alpha A_k^\alpha ,i,k,l,`$ (27)
$`_lA_k^\alpha _kA_l^\alpha ={\displaystyle \underset{i=1}{\overset{n}{}}}\alpha _i(A_l^iA_k^\alpha A_k^iA_l^\alpha )+{\displaystyle \underset{\gamma \mathrm{\Phi }}{}}c_{\gamma ,\alpha \gamma }^\alpha A_l^\gamma A_k^{\alpha \gamma },\alpha ,k,l.`$ (28)
The $`H_iH_j`$ and $`H_iE_\alpha `$ components of (22) require that
$`_kQ^j+_jA_k^i+{\displaystyle \underset{\alpha \mathrm{\Phi }}{}}\alpha _jr_i^\alpha A_k^\alpha =0,i,j,k,`$ (29)
$`_kr_i^\alpha \alpha _iF_\alpha A_k^\alpha +{\displaystyle \underset{j=1}{\overset{n}{}}}\alpha _jQ^jA_k^\alpha {\displaystyle \underset{j=1}{\overset{n}{}}}\alpha _jA_k^jr_i^\alpha +{\displaystyle \underset{\gamma \mathrm{\Phi }}{}}c_{\gamma ,\alpha \gamma }^\alpha r_i^\gamma A_k^{\alpha \gamma }+{\displaystyle \underset{j=1}{\overset{n}{}}}\alpha _jA_j^iA_k^\alpha =0`$ (30)
$`i,k,\alpha `$. From the $`E_\alpha H_i`$ and $`E_\alpha E_\beta `$ components of (22) we find that
$`_iA_k^\alpha +\alpha _iF_\alpha A_k^\alpha =0,i,k,\alpha ,`$ (31)
$`\delta _{\beta ,\alpha }\alpha _kw_\alpha ^2c_{\alpha ,\beta }^{\alpha +\beta }{\displaystyle \frac{w_\alpha w_\beta }{w_{\alpha +\beta }}}A_k^{\alpha +\beta }+{\displaystyle \underset{j=1}{\overset{n}{}}}\alpha _jr_j^\beta A_k^\alpha +{\displaystyle \underset{j=1}{\overset{n}{}}}\beta _jA_j^\alpha A_k^\beta =0`$ (32)
$`k,\alpha ,\beta `$. Note that to derive (32) we have used the identities (9), (10) and the symmetry properties of the structure constants.
It is convenient to focus first on the last two equations, since they do not contain the Cartan components of $`A_k`$. Eq. (31) obviously implies that
$$A_k^\alpha (q)=w_\alpha (q)b_k^\alpha ,b_k^\alpha :\text{some constants}.$$
(33)
The constants are then determined as follows.
Proposition 2. Eq. (32) admits solution for the constants $`b_k^\alpha `$ only for those two families of $`r(q)`$ in (14) for which the $`C_\alpha `$ are chosen according to
$$\text{case I}:C_\alpha =H_\alpha \alpha \mathrm{\Phi },\text{or}\text{case II:}C_\alpha =H_\alpha \alpha \mathrm{\Phi }.$$
(34)
For $`\alpha =\lambda _m\lambda _l`$, the $`b_k^\alpha `$ are respectively given by
$$b_k^{(\lambda _m\lambda _l)}=\delta _{km}+\mathrm{\Omega }\text{in case I,}\text{and}b_k^{(\lambda _m\lambda _l)}=\delta _{kl}+\mathrm{\Omega }\text{in case II,}$$
(35)
where $`\mathrm{\Omega }`$ is an arbitrary constant.
Proof. The statement is obtained by an elementary, but rather lengthy inspection of eq. (32). This is contained in appendix B. Q.E.D.
It is easy to explain why we got two series of solutions in the above. Namely, they arise due to the fact that $`L`$ in (6) is a self-adjoint matrix. Indeed, $`L^{}=L`$ implies that if $`r(q)`$ solves (2) then $`r^{}(q)`$ also solves it, where $`(u_1u_2)^{}=u_1^{}u_2^{}`$. Furthermore, if $`r(q)`$ is gauge transformed to a constant $`r^{}`$ by $`g(q)`$, then $`r^{}(q)`$ is transformed to $`(r^{})^{}`$ by $`(g^{})^1`$. The two series of solutions described in Proposition 2 are exchanged by this symmetry. It is thus enough to consider only one of these series, and from now on we concentrate on case I.
As the main result of this section, we now give the most general ‘gauge potential’ $`A_k`$ and $`r(q)`$ for which $`r^{}`$ (18) will be constant.
Theorem 3. The most general solution of eqs. (21), (22) for $`A_k`$ and $`Q`$ in case I of Proposition 2 can be described as follows. The root part of $`A_k`$ is determined by Proposition 2, while its Cartan part has the form<sup>*</sup><sup>*</sup>*Note that $`F_{\lambda _l\lambda _l}=0`$ by the definition of $`F_{\lambda _l\lambda _k}`$ in (13).
$$A_k^l=F_{\lambda _l\lambda _k}+\mathrm{\Omega }\underset{m(ml)}{}F_{\lambda _l\lambda _m}+_k\theta (k,l=1,\mathrm{},n),$$
(36)
where $`\theta (q)`$ is arbitrary smooth function. The function $`Q(q)gl_n`$ is given by
$$Q=\underset{k=1}{\overset{n}{}}A_k^kH_k\mathrm{\Omega }\underset{\alpha \mathrm{\Phi }}{}w_\alpha E_\alpha +g^1Q^{}g,$$
(37)
where $`g(q)GL_n`$ denotes a solution of $`_kg=gA_k`$ and $`Q^{}gl_n`$ is an arbitrary constant.
Proof. The main steps of the proof can be outlined as follows. After choosing case I of Proposition 2, the right hand side of (27) can be calculated. The general solution of (27) for the unknowns $`A_k^l`$ is then found to be
$$A_k^l=F_{\lambda _l\lambda _k}+\mathrm{\Omega }\underset{m(ml)}{}F_{\lambda _l\lambda _m}+_k\theta ^l(k,l=1,\mathrm{},n),$$
(38)
where the $`\theta ^l`$ are arbitrary smooth functions of $`q`$. Next, it is verified that (38) solves (28) if and only if
$$\theta ^1=\theta ^2=\mathrm{}=\theta ^n:=\theta .$$
(39)
At this point we have the general solution for $`A_k`$ and remaining task is to solve (29), (30) for $`Q`$. By using also (25) with $`C_\alpha =H_\alpha `$, these are inhomogeneous linear differential equations for $`Q`$. It is an easy matter to check that (37) with $`Q^{}=0`$ gives a particular solution, and that the difference $`\delta Q`$ of two solutions must satisfy the equations
$$_k(\delta Q)+[\delta Q,A_k]=0(k=1,\mathrm{},n).$$
(40)
The proof is completed by remarking that the last equation is equivalent to $`_k(g(\delta Q)g^1)=0`$ with $`_kg=gA_k`$. Q.E.D.
We wish to make some observations on the above result. Firstly, note that if $`r^{}`$ is the constant r-matrix obtained from (18) in the case
$$\theta =0,Q^{}=0,$$
(41)
then in the general case of Theorem 3 the same formula yields
$$r^{}+\mathrm{𝟏}_nQ^{}.$$
(42)
This means that the free parameters $`\theta `$ and $`Q^{}`$ in (36), (37) are irrelevant. Henceforth they will be set to zero. An additional convenience of this choice is that it guarantees the antisymmetry of $`r^{}`$ (18). In fact, one can compute the symmetric part of $`(r+_kA_kH_k)`$ and finds it to be zero if $`Q^{}=0`$. Secondly, it is worth pointing out that
$$r^{}sl_nsl_n\mathrm{\Omega }=\frac{1}{n}.$$
(43)
Indeed, the condition $`r^{}sl_nsl_n`$ is clearly equivalent to $`(r+_kA_kH_k)sl_nsl_n`$, and this is easily calculated to hold if and only if $`Q^{}=0`$ and $`\mathrm{\Omega }=\frac{1}{n}`$. Since for a given $`A_k`$ the solution of (19) for $`g(q)GL_n`$ is unique up to a constant,
$$g(q)g_0g(q),g_0GL_n,$$
(44)
we can also conclude that if the condition $`r^{}sl_nsl_n`$ is imposed, then $`r^{}`$ is necessarily antisymmetric and is uniquely determined up to an automorphism of $`sl_n`$.
Finally, let us observe that our $`r(q)`$ and $`A_k(q)`$ for which $`r^{}`$ will be a constant admit the interesting decompositions
$$r=\stackrel{~}{r}\mathrm{\Omega }\mathrm{𝟏}_n𝒜,A_k=\stackrel{~}{A}_k+\mathrm{\Omega }𝒜,$$
(45)
where
$$𝒜=\underset{l,m(lm)}{}\left(F_{\lambda _l\lambda _m}H_l+w_{\lambda _l\lambda _m}E_{\lambda _l\lambda _m}\right).$$
(46)
Here $`r(q)`$, $`A_k`$ are given by Theorem 3 together with (41). In the rest of the paper we shall determine the corresponding constant r-matrices from (18). It will be convenient to consider first the $`\mathrm{\Omega }=0`$ special case, for which $`r`$, $`A_k`$, $`r^{}`$ become $`\stackrel{~}{r}`$, $`\stackrel{~}{A}_k`$, $`\stackrel{~}{r}^{}`$, respectively.
## 4 Constant r-matrices from gauge transformation
If $`A_k`$ is given so that (21) holds then the gauge transformation $`g(q)`$ can be determined from the differential equation in (19). By taking $`A_k`$ and $`r(q)`$ from Theorem 3 with (41), this $`g`$ will transform the dynamical r-matrix $`r(q)`$ into an antisymmetric constant (18). Here we shall determine $`g(q)`$ and $`r^{}`$ explicitly. For an antisymmetric constant $`r^{}`$ the (modified) classical Yang-Baxter equation is a sufficient condition for the Jacobi identity $`\{\{L_1^{},L_2^{}\},L_3^{}\}+\mathrm{cycl}.=0`$, which will be seen to hold for the r-matrices found below.
### 4.1 The case of $`\mathrm{\Omega }=0`$
Now we calculate the gauge transformation and the resulting constant r-matrix in the special case of Theorem 3 for which $`\mathrm{\Omega }=0`$ and (41) hold. In agreement with (45), the various quantities will carry a tilde in this case. We shall use the notation
$$I_k^n:=\{1,\mathrm{},n\}\{k\},k=1,\mathrm{},n,$$
(47)
and write the elements of $`gl_n`$ as matrices. Then $`\stackrel{~}{r}(q)`$ and $`\stackrel{~}{A}_k(q)`$ take the following form:
$$\stackrel{~}{r}=\underset{1kln}{}\left(F_{kl}e_{kl}e_{lk}+w_{kl}e_{kk}e_{kl}\right),\stackrel{~}{A}_k=\underset{lI_k^n}{}\left(w_{kl}e_{kl}+F_{lk}e_{ll}\right).$$
(48)
Let us start by defining the matrix function $`\phi `$ of $`q`$ as follows: $`\phi _{nk}:=1`$ for any $`k=1,\mathrm{},n`$ and
$$\phi _{jk}:=\underset{\begin{array}{c}PI_k^n\\ |P|=nj\end{array}}{}\left(\underset{lP}{}F_l\right)k,1jn1,$$
(49)
where $`|P|`$ denotes the number of the elements of $`P`$. Moreover, let $`\chi `$ be the $`n\times n`$ matrix function of $`q`$ given by
$$\chi _{jk}=\delta _{jk}\underset{lI_k^n}{}\frac{1}{w_l}.$$
(50)
These formulae yield invertible matrices on the admissible domain of $`q`$, where $`v(q)`$ is finite. This is obvious for the diagonal matrix $`\chi `$. By using the identity
$$\underset{l=1}{\overset{n}{}}(F_i)^{l1}\phi _{lj}=\underset{\tau I_j^n}{}(F_\tau F_i),$$
(51)
we can also find the inverse of $`\phi `$ explicitly
$$\left(\phi ^1\right)_{jk}=(F_j)^{k1}\underset{lI_j^n}{}\frac{1}{(F_lF_j)}.$$
(52)
Proposition 4. A gauge transformation $`\stackrel{~}{g}(q)GL_n`$ that satisfies
$$_k\stackrel{~}{g}(q)=\stackrel{~}{g}(q)\stackrel{~}{A}_k(q)$$
(53)
with $`\stackrel{~}{A}_k`$ in (48) is given by $`\stackrel{~}{g}(q)=\phi (q)\chi (q)`$, where $`\phi `$ and $`\chi `$ are defined by (49) and (50).
Proof. The componentwise form of (53) with $`\stackrel{~}{A}_k`$ in (48) reads
$$_k\stackrel{~}{g}_{ik}=0,i,k\{1,\mathrm{},n\},$$
(54)
$$_k\stackrel{~}{g}_{ij}=\stackrel{~}{g}_{ij}F_{jk}\stackrel{~}{g}_{ik}w_{kj},i,j,k\{1,\mathrm{},n\},jk.$$
(55)
We notice that the matrix
$$\stackrel{~}{g}_{ij}(q)=\underset{lI_j^n}{}\frac{1}{w(q_l+c_i)},i,j\{1,\mathrm{},n\},$$
(56)
where the $`\{c_i\}_{i=1}^n`$ are pairwise distinct constants, yields a solution of these equations. Indeed, (54) holds obviously, while (55) is checked with the aid of the identity (10). Using (10) again, we can rewrite the matrix $`\stackrel{~}{g}(q)`$ defined by (56) in the product form
$$\stackrel{~}{g}(q)=𝐂\phi (q)\chi (q),$$
(57)
where $`𝐂`$ is the invertible constant matrix given by
$$𝐂_{ij}=\frac{1}{w(c_i)^{n1}}\left(F(c_i)\right)^{j1}.$$
(58)
Since equation (53) determines $`\stackrel{~}{g}`$ up to multiplication by a constant matrix form the left, the required statement follows. Q.E.D.
We can now calculate the gauge transformed r-matrix from (18). The result turns out to be an antisymmetric, constant solution of the (modified) classical Yang-Baxter equation,
$$[r_{12}^{},r_{13}^{}]+[r_{12}^{},r_{23}^{}]+[r_{13}^{},r_{23}^{}]=\widehat{},$$
(59)
where $``$ appears in (12) and $`\widehat{}(gl_n)^3`$ is given by
$$\widehat{}:=\underset{i,j,k,l,r,s=1}{\overset{n}{}}_{ij,kl}^{rs}e_{ji}e_{lk}e_{rs}\text{with}[e_{ij},e_{kl}]=\underset{r,s=1}{\overset{n}{}}_{ij,kl}^{rs}e_{rs}.$$
(60)
Proposition 5. The gauge transform of $`\stackrel{~}{r}(q)`$ in (48) by $`\stackrel{~}{g}(q)`$ in Proposition 4 is given by
$$\stackrel{~}{r}^{}=\underset{(a,b,c,d)S}{}\left(e_{ab}e_{cd}e_{a+1,b}e_{c+1,d}\right),$$
(61)
$$S=\{(a,b,c,d)𝐍^4|a+c+1=b+d,1ba<n,bc<n,1dn\}.$$
This formula defines an antisymmetric solution of (59).
Proof. The first statement is verified by a direct calculation, which is described in appendix C. The fact that $`\stackrel{~}{r}^{}`$ solves (59) can also be checked directly. Alternatively, it follows from the identification of $`\stackrel{~}{r}^{}`$ in terms of certain well-known solutions of (59), which is presented in section 5. Q.E.D.
It is clear from (59) that the two terms in (61) must separately satisfy the classical Yang-Baxter equation,
$$[b_{12},b_{13}]+[b_{12},b_{23}]+[b_{13},b_{23}]=0.$$
(62)
In fact, this holds since the first term
$$b_{gl_n}:=\underset{(a,b,c,d)S}{}e_{ab}e_{cd}$$
(63)
is nothing but the classical r-matrix associated with the Frobenius subalgebra of $`gl_n`$ spanned by the matrices with vanishing last row, which is described as an example in . More explicitly, it reads as
$$b_{gl_n}=\underset{k=1}{\overset{n1}{}}\underset{j=1}{\overset{nk}{}}e_{jj}e_{nk,n+1k}+\underset{1i<jn}{}\underset{m=1}{\overset{ji1}{}}e_{n+1im,n+1j}e_{n+mj,n+1i}.$$
(64)
The second term is a transform of the first one according to
$$\underset{(a,b,c,d)S}{}e_{a+1,b}e_{c+1,d}=(\sigma \sigma )b_{gl_n},$$
(65)
where $`\sigma :gl_ngl_n`$ is the inner automorphism
$$\sigma :e_{ij}e_{n+1i,n+1j}.$$
(66)
Finally, we note for later purpose that
$$\stackrel{~}{r}^{}=b_{gl_n}+(\sigma \sigma )b_{gl_n}=\stackrel{~}{r}_{sl_n}^{}+X\mathrm{𝟏}_n,$$
(67)
where $`\stackrel{~}{r}_{sl_n}^{}sl_nsl_n`$ and
$$X=\frac{1}{n}\underset{k=1}{\overset{n1}{}}(nk)e_{k+1,k}\frac{}{n}\underset{k=1}{\overset{n1}{}}ke_{k,k+1}.$$
(68)
Of course, $`\stackrel{~}{r}_{sl_n}^{}`$ satisfies the same equation (59) as $`\stackrel{~}{r}^{}`$.
### 4.2 The case of an arbitrary $`\mathrm{\Omega }`$
Now we tackle the general case by making use of the decompositions of $`r(q)`$ and $`A_k`$ in (45).
It is natural to look for $`g(q)`$ as a product
$$g(q)=h(q)\stackrel{~}{g}(q),$$
(69)
where $`\stackrel{~}{g}(q)`$ is given in Proposition 4. Then the equation $`_kg=(\stackrel{~}{A}_k+\mathrm{\Omega }𝒜)g`$ reduces to
$$_kh=h\stackrel{~}{𝒜}\mathrm{\Omega }\text{with}\stackrel{~}{𝒜}:=\stackrel{~}{g}𝒜\stackrel{~}{g}^1,$$
(70)
where $`𝒜`$ is given in (46). By using also the decomposition of $`r(q)`$ in (45) we obtain from (18) that
$$r^{}=(h(q)h(q))\left(\stackrel{~}{r}^{}+\mathrm{\Omega }\stackrel{~}{𝒜}(q)\mathrm{𝟏}_n\right)(h(q)h(q))^1,$$
(71)
where $`\stackrel{~}{r}^{}`$ is given by (61). The fact that $`r^{}`$ and $`\stackrel{~}{r}^{}`$ are both constant permits us to prove the following result without further explicit calculation.
Theorem 6. With the above notations and $`\stackrel{~}{r}^{}`$, $`X`$ defined in (61), (67), we have
$$h(q)=g_0\mathrm{exp}\left(Xn\mathrm{\Omega }\underset{i=1}{\overset{n}{}}q_i\right),$$
(72)
where $`g_0GL_n`$ is an arbitrary constant, and
$$r^{}=(g_0g_0)\left(\stackrel{~}{r}_{sl_n}^{}+(n\mathrm{\Omega }+1)X\mathrm{𝟏}_n\right)(g_0g_0)^1$$
(73)
is the most general constant r-matrix resulting from gauge transformation.
Proof. By substituting (67), we can rewrite (71) as the sum $`r^{}=r_{sl_n}^{}+r_{\mathrm{r}est}^{}`$ with
$$r_{sl_n}^{}=(h(q)h(q))\stackrel{~}{r}_{sl_n}^{}(h(q)h(q))^1$$
(74)
and
$$r_{\mathrm{r}est}^{}=\left(h(q)(\mathrm{\Omega }\stackrel{~}{𝒜}(q)+X)h^1(q)\right)\mathrm{𝟏}_n.$$
(75)
Since $`r^{}`$ is constant, these two terms must be constant separately. Recall now that $`\stackrel{~}{𝒜}(q)`$ is independent of $`\mathrm{\Omega }`$ by its definition (70) and that for $`\mathrm{\Omega }=\frac{1}{n}`$ we must have $`r^{}sl_nsl_n`$ (43). This implies that $`(X\frac{1}{n}\stackrel{~}{𝒜}(q))`$ must vanish, whereby
$$\stackrel{~}{𝒜}=nX.$$
(76)
Hence we obtain (72) from the differential equation in (70). But then the fact that $`r_{sl_n}^{}`$ is constant shows that the relation
$$[X\mathrm{𝟏}_n+\mathrm{𝟏}_nX,\stackrel{~}{r}_{sl_n}^{}]=0,$$
(77)
which is equivalent to
$$r_{sl_n}^{}=(g_0g_0)\stackrel{~}{r}_{sl_n}^{}(g_0g_0)^1,$$
(78)
must be valid. By substituting these results back into (71) we arrive at (73). Q.E.D.
Incidentally, we have also verified by explicit calculation that (76) and (77) are indeed satisfied, which represents a reassuring check on the foregoing considerations in the paper.
## 5 Identification of the constant r-matrices
The constant r-matrix (73) is a solution of (59). For the rational Calogero-Moser model, $`=0`$, this is the classical Yang-Baxter equation. In this case the identification of the r-matrix in terms of a Frobenius subalgebra of $`gl_n`$ has already been mentioned (67). In the hyperbolic/trigonometric cases (59) is the modified classical Yang-Baxter equation, whose antisymmetric solutions have been classified by Belavin and Drinfeld for the complex simple Lie algebras. A well-known solution for the Lie algebra $`sl_n`$, with the normalization
$$[\rho _{12},\rho _{13}]+[\rho _{12},\rho _{23}]+[\rho _{13},\rho _{23}]=\widehat{},$$
(79)
is the so-called Cremmer-Gervais classical r-matrix, which we quote from as
$$r_{CG}=\underset{1i<jn}{}e_{ij}e_{ji}+2\underset{1i<jn}{}\underset{m=1}{\overset{ji1}{}}e_{i,jm}e_{j,i+m}+\frac{1}{n}\underset{1i<jn}{}(n+2(ij))e_{ii}e_{jj}.$$
(80)
Note that $`r_{CG}sl_nsl_n`$ and $`\widehat{}`$ (60) belongs to $`(sl_n)^3`$. Below we show that for $`0`$ the $`sl_n`$-part of the constant Calogero-Moser r-matrix (73) is equivalent to $`r_{CG}`$.
We shall need the following properties of $`r_{CG}`$. As in , first introduce $`J_0,J_\pm sl_n`$ by
$$J_0=\frac{1}{2}\underset{k=1}{\overset{n}{}}(n+12k)e_{kk},J_+=\underset{k=1}{\overset{n1}{}}(nk)e_{k,k+1},J_{}=\sigma (J_+)=\underset{k=1}{\overset{n1}{}}ke_{k+1,k}.$$
(81)
They generate the principal $`sl_2`$ subalgebra of $`sl_n`$,
$$[J_0,J_\pm ]=\pm J_\pm ,[J_+,J_{}]=2J_0.$$
(82)
Then define the elements $`b_{CG}^\pm :=\frac{1}{2}[J_\pm \mathrm{𝟏}_n+\mathrm{𝟏}_nJ_\pm ,r_{CG}]sl_nsl_n`$. Explicitly,
$$b_{CG}^+=\underset{k=1}{\overset{n1}{}}d_ke_{k,k+1}+\underset{1i<jn}{}\underset{m=1}{\overset{ji1}{}}e_{i,jm+1}e_{j,i+m},d_k:=\underset{j=1}{\overset{k}{}}e_{jj}\frac{k}{n}\mathrm{𝟏}_n.$$
(83)
On account of $`(\sigma \sigma )r_{CG}=r_{CG}`$, with $`\sigma `$ defined in (66), $`b_{CG}^{}=(\sigma \sigma )b_{CG}^+`$. It has been pointed out in that the subspace of $`sl_nsl_n`$ spanned by $`r_{CG}`$ and $`b_{CG}^\pm `$ is an irreducible representation of the principal $`sl_2`$ subalgebra. In fact, for the operators
$$𝒥_{0,\pm }(Y):=[J_{0,\pm }\mathrm{𝟏}_n+\mathrm{𝟏}_nJ_{0,\pm },Y]Ygl_ngl_n,$$
(84)
one has the relations:
$$𝒥_0\left(\begin{array}{c}b_{CG}^+\\ r_{CG}\\ b_{CG}^{}\end{array}\right)=\left(\begin{array}{c}b_{CG}^+\\ 0\\ b_{CG}^{}\end{array}\right),𝒥_+\left(\begin{array}{c}b_{CG}^+\\ r_{CG}\\ b_{CG}^{}\end{array}\right)=\left(\begin{array}{c}0\\ 2b_{CG}^+\\ r_{CG}\end{array}\right),𝒥_{}\left(\begin{array}{c}b_{CG}^+\\ r_{CG}\\ b_{CG}^{}\end{array}\right)=\left(\begin{array}{c}r_{CG}\\ 2b_{CG}^{}\\ 0\end{array}\right).$$
(85)
It follows from these relations that $`b_{CG}^\pm `$ satisfy the classical Yang-Baxter equation , and the identification of $`b_{CG}^\pm `$ in terms of Frobenius subalgebras of $`sl_n`$ is also described in this reference.
Now we are prepared to establish the connection between $`r_{CG}`$ and the r-matrix $`r^{}`$ (73). The key observation is the following identity:
$$(TT)\stackrel{~}{r}_{sl_n}^{}=b_{CG}^++b_{CG}^{},$$
(86)
where $`T:gl_ngl_n`$ denotes matrix transposition. This can be checked directly from the formulas (67), (64) (83). It permits us to transform $`\stackrel{~}{r}_{sl_n}^{}`$ into a multiple of $`r_{CG}`$ in a simple manner. To treat the hyperbolic/trigonometric cases together, we introduce the parameter
$$a^{}=\{\begin{array}{cc}a,& \text{hyperbolic case,}\\ \sqrt{1}a,& \text{trigonometric case,}\end{array}$$
(87)
whose square $`=(a^{})^2`$ appears in (59). By using (85) it is not difficult to verify that
$$(u_{}u_+u_{}u_+)\left(TT\stackrel{~}{r}_{sl_n}^{}\right)(u_{}u_+u_{}u_+)^1=a^{}r_{CG}$$
(88)
with
$$u_{}:=\mathrm{exp}\left(\frac{a^{}}{2}J_{}\right),u_+:=\mathrm{exp}\left(\frac{1}{a^{}}J_+\right).$$
(89)
According to (88) the $`sl_n`$-part of $`r^{}`$ is equivalent to $`a^{}r_{CG}`$ under a Lie algebra automorphism.
In the end, notice from (68) and (81) that
$$X=\frac{1}{n}(J_+^T+J_{}^T).$$
(90)
This allows us to present the r-matrix associated with
$$L^{}(q,p)=g_0h(q)\stackrel{~}{g}(q)L(q,p)(g_0h(q)\stackrel{~}{g}(q))^1$$
(91)
in a ‘standard form’. Here $`h(q)`$ and $`\stackrel{~}{g}(q)`$ are the same as in Theorem 6, and our final result is formulated as follows.
Proposition 7. Consider the hyperbolic/trigonometric Calogero-Moser models. If in Theorem 6 the constant $`g_0`$ is chosen to be
$$g_0=\mathrm{exp}\left(\frac{a^{}}{2}J_{}^T\right)\mathrm{exp}\left(\frac{1}{a^{}}J_+^T\right),$$
(92)
then the r-matrix (73) becomes
$$r^{}=a^{}(TT)(r_{CG}+2(\mathrm{\Omega }+\frac{1}{n})J_0\mathrm{𝟏}_n).$$
(93)
Proof. By means of the $`sl_2`$ algebra (82) and (90) it is easy to check that $`g_0Xg_0^1=\frac{2a^{}}{n}J_0^T`$. The statement is obtained by combining this with (88). Q.E.D.
This proposition describes the precise relationship between the most general constant r-matrices of the hyperbolic/trigonometric Calogero-Moser models and the standard Cremmer-Gervais classical r-matrices.
## 6 Conclusion
In this paper we have determined the most general constant r-matrices that may be obtained by coordinate dependent gauge transformations of the standard Lax representation (6) of the degenerate Calogero-Moser models associated with $`gl_n`$. Up to automorhisms of $`gl_n`$ (i.e. up to conjugation by constants $`g_0GL_n`$ and transpose) and addition of an irrelevant term $`\mathrm{𝟏}_nQ^{}`$ with any constant $`Q^{}gl_n`$, the most general such r-matrix turned out to have the form
$$r^{}=\underset{(a,b,c,d)S}{}\left(e_{ab}e_{cd}e_{a+1,b}e_{c+1,d}\right)+n\mathrm{\Omega }X\mathrm{𝟏}_n,$$
(94)
where
$$X=\frac{1}{n}\underset{k=1}{\overset{n1}{}}(nk)e_{k+1,k}\frac{}{n}\underset{k=1}{\overset{n1}{}}ke_{k,k+1},$$
$``$ is given according to (12) in correspondence with the rational, hyperbolic and trigonometric potential functions (4), $`\mathrm{\Omega }`$ is an arbitrary constant scalar, and
$$S=\{(a,b,c,d)𝐍^4|a+c+1=b+d,1ba<n,bc<n,1dn\}.$$
We have seen that $`r^{}`$ solves the classical (modified) Yang-Baxter equation (59), and have identified it in terms of well-known solutions of this equation. In particular, we have shown that in the hyperbolic and trigonometric cases the above $`r^{}`$ with $`\mathrm{\Omega }=\frac{1}{n}`$ is equivalent to a multiple of the Cremmer-Gervais classical r-matrix under an automorphism of $`gl_n`$. We obtained these results by an explicit determination of the gauge transformations $`g(q)GL_n`$ for which the Poisson brackets of $`L^{}(q,p)=g(q)L(q,p)g^1(q)`$, where $`L`$ is the standard Lax matrix (6), can be written in the form (2) with a constant r-matrix. The gauge transformation $`g(q)`$ for which the Poisson brackets of $`L^{}`$ are encoded by $`r^{}`$ in (94) was found as the product
$$g(q)=\mathrm{exp}\left(Xn\mathrm{\Omega }\underset{i=1}{\overset{n}{}}q_i\right)\phi (q)\chi (q),$$
(95)
where the matrices $`\phi (q)`$ and $`\chi (q)`$ are defined by (49) and (50), with the notations fixed by equations (7), (8), (13) in section 2.
The outcome of our direct analysis of the degenerate Calogero-Moser models is consistent with the results obtained in by different means. We hope to present a more detailed comparison with the elliptic case as well as an analogous study for other Lie algebras elsewhere.
Acknowledgments. We wish to thank J. Balog and the referees for useful comments on the manuscript. This work has been supported in part by the Hungarian Ministry of Education under FKFP 0596/1999 and by the National Science Fund (OTKA) under T025120, M028418.
## A Proof of Theorem 1
The proof given below relies on the general analysis of the momentum independent Calogero-Moser r-matrices presented by Braden and Suzuki in . We first specialize the relevant results of to our case and then further elaborate them to obtain the statement of Theorem 1.
Consider the Lax matrix in (6) with a function $`w`$ in (7). Our task is to find the most general momentum independent r-matrix, $`r(q)`$, which satisfies equation (2), i.e.,
$$\{L_1,L_2\}(q,p)=[r_{12}(q),L_1(q,p)][r_{21}(q),L_2(q,p)].$$
(A.1)
Obviously, $`r(q)=r_{12}(q)gl_ngl_n`$ can be expanded in the form
$$r(q)=\underset{i,j=1}{\overset{n}{}}r^{i,j}(q)H_iH_j+\underset{\alpha \mathrm{\Phi }}{}\underset{i=1}{\overset{n}{}}\left(r^{i,\alpha }(q)H_iE_\alpha +r^{\alpha ,i}(q)E_\alpha H_i\right)+\underset{\alpha ,\beta \mathrm{\Phi }}{}r^{\alpha ,\beta }(q)E_\alpha E_\beta .$$
(A.2)
Since the functions $`w`$ in (7) are odd (and thus $`w_\alpha (q)=w_\alpha (q)`$), we can use the results of the third and fourth chapters of , where it has been shown that under our conditions the following equations hold:
$$r^{\alpha ,i}(q)=0(i\{1,\mathrm{},n\},\alpha \mathrm{\Phi }),$$
(A.3)
$$r^{\alpha ,\beta }(q)=\frac{w_\alpha ^{}(q)}{w_\alpha (q)}\delta _{\alpha ,\beta }=F_\alpha (q)\delta _{\alpha ,\beta }(\alpha ,\beta \mathrm{\Phi }).$$
(A.4)
Moreover, according to , the remaining requirements on $`r(q)`$ reduce to the equations
$$\underset{i=1}{\overset{n}{}}\alpha _ir^{i,j}(q)=0(\alpha \mathrm{\Phi },j\{1,\mathrm{},n\}),$$
(A.5)
and
$$\underset{i=1}{\overset{n}{}}(\alpha _ir^{i,\beta }w_\alpha \beta _ir^{i,\alpha }w_\beta )=c_{\alpha ,\alpha +\beta }^\beta (r^{\alpha ,\alpha }w_{\alpha +\beta }+r^{\beta ,\beta }w_{\alpha +\beta })(\alpha ,\beta \mathrm{\Phi }).$$
(A.6)
We here use the basis of $`gl_n`$ introduced in (5), $`\alpha _i:=\alpha (H_i)`$, the structure constants $`c_{\alpha ,\beta }^{\alpha +\beta }`$ satisfy $`[E_\alpha ,E_\beta ]=c_{\alpha ,\beta }^{\alpha +\beta }E_{\alpha +\beta }`$ if $`\alpha ,\beta ,(\alpha +\beta )`$ all belong to $`\mathrm{\Phi }`$, and $`c_{\alpha ,\beta }^{\alpha +\beta }:=0`$ otherwise.
Now consider equation (A.5) for $`\alpha :=(\lambda _k\lambda _l)\mathrm{\Phi }`$. From this we see that $`r^{k,j}(q)r^{l,j}(q)=0`$ $`(kl,j)`$, which means that the general solution of (A.5) is
$$r^{i,j}(q)=M^j(q)(i,j\{1,\mathrm{},n\}),$$
(A.7)
where the $`M^j`$ are arbitrary smooth functions of $`q`$. Let us next solve (A.6) for $`r^{i,\alpha }(q)`$. By substituting (A.4) into (A.6) and using the identity (10) and the symmetry properties of the structure constants we obtain
$$\underset{i=1}{\overset{n}{}}\left(\alpha _i\widehat{r}^{i,\beta }(q)\beta _i\widehat{r}^{i,\alpha }(q)\right)=c_{\alpha ,\beta }^{\alpha +\beta }(\alpha ,\beta \mathrm{\Phi }),$$
(A.8)
where we define $`\widehat{r}^{i,\gamma }:=\frac{r^{i,\gamma }}{w_\gamma }`$ for any $`\gamma \mathrm{\Phi }`$. By introducing the notations
$$\widehat{r}_S^\alpha :=\underset{i=1}{\overset{n}{}}(\widehat{r}^{i,\alpha }+\widehat{r}^{i,\alpha })H_i,\widehat{r}_A^\alpha :=\underset{i=1}{\overset{n}{}}(\widehat{r}^{i,\alpha }\widehat{r}^{i,\alpha })H_i,$$
(A.9)
we have
$$\widehat{r}^\alpha :=\underset{i=1}{\overset{n}{}}\widehat{r}^{i,\alpha }H_i=\frac{1}{2}(\widehat{r}_S^\alpha +\widehat{r}_A^\alpha ),\widehat{r}_S^\alpha (q)=\widehat{r}_S^\alpha (q),\widehat{r}_A^\alpha (q)=\widehat{r}_A^\alpha (q).$$
(A.10)
We now consider equation (A.8) for the pairs of roots $`(\alpha ,\beta )`$ and $`(\alpha ,\beta )`$. By adding these two equations we get
$$\alpha (\widehat{r}_S^\beta (q))=c_{\alpha ,\beta }^{\alpha +\beta }+c_{\alpha ,\beta }^{\alpha \beta }(\alpha ,\beta \mathrm{\Phi }).$$
(A.11)
It follows from the definition of $`K_\alpha `$ in (5) that $`\alpha (K_\beta )=(c_{\alpha ,\beta }^{\alpha +\beta }+c_{\alpha ,\beta }^{\alpha \beta })`$ for any $`\alpha ,\beta \mathrm{\Phi }`$. Therefore the general solution of (A.11) is given by
$$\widehat{r}_S^\alpha (q)=K_\alpha +\tau _S^\alpha (q)\mathrm{𝟏}_n(\alpha \mathrm{\Phi }),$$
(A.12)
where $`\tau _S^\alpha (q)=\tau _S^\alpha (q)`$ are arbitrary smooth functions. On the other hand, by substituting (A.11) and the decomposition in (A.10) into (A.8) we obtain the relation
$$\alpha (\widehat{r}_A^\beta (q))=\beta (\widehat{r}_A^\alpha (q))(\alpha ,\beta \mathrm{\Phi }).$$
(A.13)
Obviously, there exists the decomposition
$$\widehat{r}_A^\alpha (q)=C_\alpha (q)+\tau _A^\alpha (q)\mathrm{𝟏}_n,$$
(A.14)
where $`C_\alpha (q)_nsl_n`$ and $`\tau _A^\alpha (q)`$ are smooth functions. The antisymmetry of $`\widehat{r}_A^\alpha (q)`$ in $`\alpha `$ and (A.13) can be rewritten as
$$C_\alpha (q)=C_\alpha (q),\alpha (C_\beta (q))=\beta (C_\alpha (q)),\tau _A^\alpha (q)=\tau _A^\alpha (q).$$
(A.15)
By the above, we have parametrized the most general $`r(q)`$ in terms of the functions $`M^j`$, $`\tau _A^\alpha `$, $`\tau _S^\alpha `$ and $`C_\alpha `$. If we now introduce the notation
$$Q(q):=\underset{i=1}{\overset{n}{}}M^i(q)H_i+\frac{1}{2}\underset{\alpha \mathrm{\Phi }}{}\left(\tau _S^\alpha (q)+\tau _A^\alpha (q)\right)w_\alpha (q)E_\alpha ,$$
(A.16)
then $`r(q)`$ in (A.2) takes precisely the form stated by Theorem 1, which completes the proof.
## B Proof of Proposition 2
In this appendix we prove Proposition 2 by analyzing equation (32),
$$\alpha _kw_\alpha ^2\delta _{\beta ,\alpha }c_{\alpha ,\beta }^{\alpha +\beta }\frac{w_\alpha w_\beta }{w_{\alpha +\beta }}A_k^{\alpha +\beta }+(\alpha r^\beta )A_k^\alpha +(\beta A^\alpha )A_k^\beta =0(k=1,\mathrm{},n),$$
(B.1)
whereby we determine the constants $`b_k^\alpha `$ that appear in $`A_k^\alpha =w_\alpha b_k^\alpha `$ (33). We here use the notation $`\alpha r^\beta =_{i=1}^n\alpha _ir_i^\beta `$, $`\beta A^\alpha =_{i=1}^n\beta _iA_i^\alpha `$ and similarly for all quantities with Cartan indices. For later reference, note from (25) that
$$\beta r^\alpha =\frac{1}{2}w_\alpha \beta (C_\alpha K_\alpha ),\alpha ,\beta \mathrm{\Phi },$$
(B.2)
where $`K_\alpha `$ is defined in (5) and $`C_\alpha =_{i=1}^nC_\alpha ^iH_i`$ enjoys the properties in (15).
If we fix $`\alpha \mathrm{\Phi }`$, then (B.1) for the pairs of roots $`(\alpha ,\beta )`$ given by $`(\alpha ,\alpha )`$, $`(\alpha ,\alpha )`$, $`(\alpha ,\alpha )`$ and $`(\alpha ,\alpha )`$ leads respectively to the following relations:
$$(\alpha r^\alpha +\alpha A^\alpha )A_k^\alpha =0,$$
(B.3)
$$(\alpha r^\alpha +\alpha A^\alpha )A_k^\alpha =0,$$
(B.4)
$$\alpha _kw_\alpha ^2+(\alpha r^\alpha )A_k^\alpha (\alpha A^\alpha )A_k^\alpha =0,$$
(B.5)
$$\alpha _kw_\alpha ^2+(\alpha r^\alpha )A_k^\alpha (\alpha A^\alpha )A_k^\alpha =0.$$
(B.6)
Since $`\alpha r^\alpha =\alpha r^\alpha `$ by (B.2), these relations imply that
$$\alpha A^\alpha =\alpha A^\alpha =\alpha r^\alpha .$$
(B.7)
On account of (B.7) and (B.2), (B.5) can be written as
$$\alpha _kw_\alpha ^2=(\alpha A^\alpha )(A_k^\alpha +A_k^\alpha ).$$
(B.8)
This expression shows that
$$b_k^\alpha b_k^\alpha =\epsilon ^\alpha \alpha _k$$
(B.9)
with some constants $`\epsilon ^\alpha `$. We then find from the above that
$$\alpha b^\alpha =\epsilon ^\alpha $$
(B.10)
and the $`\epsilon ^\alpha `$ must satisfy
$$\epsilon ^\alpha =\epsilon ^\alpha ,(\epsilon ^\alpha )^2=1.$$
(B.11)
Now it is convenient to introduce $`\mathrm{\Pi }_k^\alpha :=(b_k^\alpha +b_k^\alpha )`$, which results in
$$b_k^\alpha =\frac{1}{2}\epsilon ^\alpha \alpha _k+\frac{1}{2}\mathrm{\Pi }_k^\alpha ,\alpha \mathrm{\Phi }.$$
(B.12)
Let us put $`\mathrm{\Pi }_k^{ij}:=\mathrm{\Pi }_k^{(\lambda _i\lambda _j)}`$. Then the relations $`\mathrm{\Pi }_k^\alpha =\mathrm{\Pi }_k^\alpha `$ and $`\alpha \mathrm{\Pi }^\alpha =0`$ (by (B.7)) give
$$\mathrm{\Pi }_k^{ij}=\mathrm{\Pi }_k^{ji},\mathrm{\Pi }_i^{ij}=\mathrm{\Pi }_j^{ij},k,ij.$$
(B.13)
Consider now such roots $`\alpha =(\lambda _i\lambda _j)`$ and $`\beta =\pm (\lambda _l\lambda _m)\mathrm{\Phi }`$ that $`\{i,j\}\{l,m\}=\mathrm{}`$. In this case (B.1) yields
$$(\alpha \widehat{r}^\beta )b_k^\alpha +(\beta b^\alpha )b_k^\beta =0,$$
(B.14)
$$(\alpha \widehat{r}^\beta )b_k^\alpha (\beta b^\alpha )b_k^\beta =0,$$
(B.15)
where we use the notation $`\widehat{r}^\gamma :=\frac{r^\gamma }{w_\gamma }`$ for any $`\gamma \mathrm{\Phi }`$. Adding these two equations, and using (B.7) and (B.12), we can easily get that now
$$\beta b^\alpha =0,\beta \mathrm{\Pi }^\alpha =0.$$
(B.16)
The general form of $`\mathrm{\Pi }_k^{ij}`$ which obeys (B.13) and (B.16) is in fact the following:
$$\mathrm{\Pi }_k^{ij}=\eta ^\alpha (\delta _{ki}+\delta _{kj})+2\mathrm{\Omega }^\alpha ,$$
(B.17)
where $`\eta ^\alpha `$, $`\mathrm{\Omega }^\alpha `$ are constants. Notice that for $`\alpha =(\lambda _i\lambda _j)`$ that element $`K_\alpha =_{k=1}^nK_\alpha ^kH_k`$ defined in (5) has precisely the components $`K_\alpha ^k=\delta _{ki}+\delta _{kj}`$.
Now, let $`\alpha ,\beta ,\alpha +\beta \mathrm{\Phi }`$ be roots. In this case $`\alpha \beta =\alpha +(\beta )\mathrm{\Phi }`$. Hence (B.1) for the $`(\alpha ,\beta )`$ and the $`(\alpha ,\beta )`$ pairs reads as
$$c_{\alpha ,\beta }^{\alpha +\beta }b_k^{\alpha +\beta }=(\alpha \widehat{r}^\beta )b_k^\alpha +(\beta b^\alpha )b_k^\beta ,$$
(B.18)
$$0=(\alpha \widehat{r}^\beta )b_k^\alpha (\beta b^\alpha )b_k^\beta .$$
(B.19)
By adding these two equations making use of (B.2) and (B.9), we obtain
$$c_{\alpha ,\beta }^{\alpha +\beta }b_k^{\alpha +\beta }=(\alpha K_\beta )b_k^\alpha +\epsilon ^\beta (\beta b^\alpha )\beta _k.$$
(B.20)
If $`\alpha =(\lambda _i\lambda _j)`$, $`\beta =(\lambda _j\lambda _l)`$ are chosen, then $`\alpha K_\beta =1`$ and $`c_{\alpha ,\beta }^{\alpha +\beta }=1`$. Let us then substitute (B.12) with (B.17) into (B.20) and consider the resulting equation for $`k\{i,j,l\}`$ and for $`k\{i,j,l\}`$. In this way we obtain the requirementsWe here implicitly assume that $`n4`$, but the final solution is valid for any $`n2`$.
$$\mathrm{\Omega }^{\alpha +\beta }=\mathrm{\Omega }^\alpha ,$$
(B.21)
$$\epsilon ^{\alpha +\beta }+\eta ^{\alpha +\beta }=\epsilon ^\alpha +\eta ^\alpha ,$$
(B.22)
$$\epsilon ^\alpha \eta ^\alpha =2\epsilon ^\beta (\beta b^\alpha ),$$
(B.23)
$$\eta ^{\alpha +\beta }\epsilon ^{\alpha +\beta }=2\epsilon ^\beta (\beta b^\alpha ).$$
(B.24)
These tell us that
$$\mathrm{\Omega }^{\alpha +\beta }=\mathrm{\Omega }^\alpha ,\epsilon ^{\alpha +\beta }=\epsilon ^\alpha ,\eta ^{\alpha +\beta }=\eta ^\alpha .$$
(B.25)
In conclusion, there exist some constants $`\epsilon `$, $`\eta `$, $`\mathrm{\Omega }`$ that
$$\epsilon ^\alpha =\epsilon ,\eta ^\alpha =\eta ,\mathrm{\Omega }^\alpha =\mathrm{\Omega },\alpha \mathrm{\Phi }.$$
(B.26)
In addition, we can compute from (B.12) that in the above case $`2\beta b^\alpha =(\eta \epsilon )`$, and by substituting this back into (B.23) we obtain
$$(\epsilon +1)(\eta \epsilon )=0.$$
(B.27)
At the same time we know from (B.11) that $`\epsilon `$ must be equal to $`1`$ or $`1`$.
The first solution of (B.27) is $`\epsilon =1=\eta `$. In this case we can determine $`b_k^\alpha `$ from (B.12) in terms of the arbitrary constant $`\mathrm{\Omega }`$ as
$$b_k^{\lambda _i\lambda _j}=\delta _{ki}+\mathrm{\Omega }.$$
(B.28)
We can then also calculate $`\beta r^\alpha `$ from the above equations, and thereby find from (B.2) that $`C_\alpha =H_\alpha `$ must hold. This is precisely the result stated in case I of Proposition 2. We have obtained it as a consequence of considering a subset of all cases of (B.1), but it can checked to satisfy this equation in all remaining cases (for $`\alpha =(\lambda _i\lambda _j)`$, $`\beta =(\lambda _l\lambda _i)`$ etc.) as well.
The other solution of (B.27) is $`\epsilon =1`$, but then we still have to determine $`\eta `$. For this we consider $`\alpha =(\lambda _i\lambda _j)`$, $`\beta =(\lambda _j\lambda _l)`$ and calculate that
$$b_k^\alpha =\frac{1}{2}\left(\left(\eta 1\right)\delta _{ki}+\left(\eta +1\right)\delta _{kj}\right)+\mathrm{\Omega },$$
(B.29)
$$b_k^\beta =\frac{1}{2}\left(\left(\eta 1\right)\delta _{kj}+\left(\eta +1\right)\delta _{kl}\right)+\mathrm{\Omega }.$$
(B.30)
We then look at (B.1) for the $`(\alpha ,\beta )`$ and $`(\beta ,\alpha )`$ pairs of roots and add these two equations, which gives
$$0=(\alpha \widehat{r}^\beta +\alpha b^\beta )b_k^\alpha +(\beta \widehat{r}^\alpha +\beta b^\alpha )b_k^\beta .$$
(B.31)
Since $`b_k^\alpha `$ and $`b_k^\beta `$ are linearly independent $`n`$-component vectors for any $`\eta `$, we obtain
$$\alpha \widehat{r}^\beta +\alpha b^\beta =0,\beta \widehat{r}^\alpha +\beta b^\alpha =0.$$
(B.32)
By subtracting these equations and taking into account that by (B.2) now
$$\alpha \widehat{r}^\beta \beta \widehat{r}^\alpha =\frac{1}{2}\left(\beta K_\alpha \alpha K_\beta \right)=1,$$
(B.33)
we find that $`\eta =1`$. So we have completely determined $`b_k^\alpha `$ again, and it is easy to confirm that the final formula agrees with case II of Proposition 2. Thus the proof is complete.
## C Proof of Proposition 5
In this appendix we verify the statement of Proposition 5.
By combining eq. (18) and Proposition 4, the constant r-matrix that we wish to calculate can be written in the form
$$\stackrel{~}{r}^{}=\left(\phi (q)\phi (q)\right)\rho (q)\left(\phi (q)\phi (q)\right)^1$$
(C.1)
with
$$\rho (q)=\left(\chi (q)\chi (q)\right)\left(\stackrel{~}{r}(q)+\underset{k}{}A_k(q)H_k\right)\left(\chi (q)\chi (q)\right)^1.$$
(C.2)
The formulas in (48), (50) together with (10) and (11) result in
$`\rho ={\displaystyle \underset{kl}{}}{\displaystyle \frac{1}{F_kF_l}}(e_{kl}e_{ll})(e_{lk}e_{kk})+{\displaystyle \underset{kl}{}}{\displaystyle \frac{F_kF_l}{F_kF_l}}(e_{kl}e_{ll})(e_{lk}e_{kk})`$
$`+{\displaystyle \underset{kl}{}}F_ke_{kk}e_{kl}{\displaystyle \underset{kl}{}}F_le_{lk}e_{ll}.`$ (C.3)
Therefore, to prove Proposition 5 it is enough to verify that
$$\left(\phi (q)\phi (q)\right)\rho (q)=\stackrel{~}{r}^{}\left(\phi (q)\phi (q)\right)$$
(C.4)
holds for $`\rho `$ in (C.3) and $`\stackrel{~}{r}^{}`$ in (61). We obtain in a straightforward manner that
$$\left(\phi (q)\phi (q)\right)\rho (q)=\underset{a,b,c,d=1}{\overset{n}{}}\left(B_{abcd}+\stackrel{~}{B}_{abcd}\right)e_{ab}e_{cd},$$
(C.5)
where
$$B_{abcd}=\frac{(\phi _{ad}\phi _{ab})(\phi _{cd}\phi _{cb})}{F_dF_b},\text{if}bd,$$
(C.6)
$$\stackrel{~}{B}_{abcd}=\frac{F_dF_b}{F_dF_b}(\phi _{ad}\phi _{ab})(\phi _{cd}\phi _{cb})+F_d\phi _{ad}\phi _{cd}F_b\phi _{ab}\phi _{cb},\text{if}bd,$$
(C.7)
and $`B_{abcd}=\stackrel{~}{B}_{abcd}=0`$ if $`b=d`$. From (49) and (61), the right hand side of (C.4) is found to be
$$\stackrel{~}{r}^{}\left(\phi (q)\phi (q)\right)=\underset{a,b,c,d=1}{\overset{n}{}}\left(D_{abcd}+\stackrel{~}{D}_{abcd}\right)e_{ab}e_{cd}$$
(C.8)
with
$$D_{abcd}=\underset{(a,x,c,y)S}{}\phi _{xb}\phi _{yd}\underset{(c,y,a,x)S}{}\phi _{xb}\phi _{yd},$$
(C.9)
$$\stackrel{~}{D}_{abcd}=\underset{(a1,x,c1,y)S}{}\phi _{xb}\phi _{yd}\underset{(c1,y,a1,x)S}{}\phi _{xb}\phi _{yd},$$
(C.10)
where the set $`S`$ is defined in Proposition 5 and by an empty sum we mean zero.
We now observe that $`\stackrel{~}{D}_{abcd}=0=\stackrel{~}{B}_{abcd}`$ if $`a=1`$ or $`c=1`$, and
$$\stackrel{~}{D}_{a,b,c,d}=D_{a1,b,c1,d},\stackrel{~}{B}_{a,b,c,d}=B_{a1,b,c1,d},\text{if}2a,cn.$$
(C.11)
These properties are obvious for $`\stackrel{~}{D}`$, while for $`\stackrel{~}{B}`$ they follow from the formula (49). In particular, the second equality in (C.11) is checked by inserting into (C.6) the identity
$$\phi _{a1,d}\phi _{a1,b}=F_b\phi _{ab}F_d\phi _{ad},2an,$$
(C.12)
which is consequence of (49). We conclude that it is sufficient to show that $`B_{abcd}=D_{abcd}`$.
Let us examine the expressions of $`B_{abcd}`$ and $`D_{abcd}`$. First, we notice that for all indices
$$B_{abcd}=B_{cbad},D_{abcd}=D_{cbad},$$
(C.13)
and
$$B_{abcd}=0=D_{abcd}\text{if}a=n\text{or}c=n\text{or}b=d.$$
(C.14)
Hence it is enough to show that $`B_{abcd}=D_{abcd}`$ for such indices that $`ac<n`$ and $`bd`$. We now introduce the notation
$$F_P:=\underset{tP}{}F_tP\{1,\mathrm{},n\},$$
(C.15)
and also put $`F_P:=1`$ if $`P=\mathrm{}`$, for which $`|P|=0`$. We then rewrite $`B_{abcd}`$ as
$$B_{abcd}=(F_dF_b)\left(\underset{\begin{array}{c}PI_b^nI_d^n\\ |P|=n1a\end{array}}{}F_P\right)\left(\underset{\begin{array}{c}PI_b^nI_d^n\\ |P|=n1c\end{array}}{}F_P\right),$$
(C.16)
where $`I_k^n`$ is defined in (47). This is derived from (C.6) by using that as a result of (49)
$$\phi _{al}\phi _{ak}=(F_kF_l)\underset{\begin{array}{c}PI_k^nI_l^n\\ |P|=n1a\end{array}}{}F_P.$$
(C.17)
Next, by inserting (49) into (C.9) and using that $`ac`$, we get the expression
$`D_{abcd}=(F_dF_b){\displaystyle \underset{\begin{array}{c}x+y=a+c+1\\ 1xa<yn\end{array}}{}}\left(\right({\displaystyle \underset{\begin{array}{c}PI_b^nI_d^n\\ |P|=n1x\end{array}}{}}F_P\left)\right({\displaystyle \underset{\begin{array}{c}PI_b^nI_d^n\\ |P|=ny\end{array}}{}}F_P)`$ (C.24)
$`\left({\displaystyle \underset{\begin{array}{c}PI_b^nI_d^n\\ |P|=n1y\end{array}}{}}F_P\right)\left({\displaystyle \underset{\begin{array}{c}PI_b^nI_d^n\\ |P|=nx\end{array}}{}}F_P\right)).`$ (C.29)
The $`x=a`$, $`y=(c+1)`$ term in the first line of the right hand side of (C.29) clearly equals the right hand side of (C.16). The proof is completed by a close inspection of the ranges of the summation indices, which shows that all the remaining terms cancel pairwise between the two lines of (C.29) for any $`ac(n1)`$.
|
warning/0005/math0005190.html
|
ar5iv
|
text
|
# Differential operators on equivariant vector bundles over symmetric spaces
## 1 Equivariant vector bundles
Let $`X`$ denote a manifold with a smooth action of a Lie group $`G`$. An equivariant vector bundle $`E`$ over $`X`$ is a smooth vector bundle $`\pi :EX`$ together with a smooth action of $`G`$ on $`E`$ such that $`\pi (gv)=g\pi (v)`$ for $`vE`$ and such that all maps of the fibres $`g:E_xE_{gx}`$, $`gG`$, $`xX`$, are linear. An example is given by $`E=X\times V`$, where $`(\sigma ,V)`$ is a finite dimensional representation of $`G`$ and $`G`$ acts on $`E`$ by $`g(x,v)=(gx,\sigma (g)v)`$.
Now assume $`X`$ to be a homogeneous space, i.e. $`X=G/H`$ for a closed subgroup $`H`$ of $`G`$. Given an equivariant bundle $`E`$ over $`X`$ we get a representation of $`H`$ on the fibre over $`eH`$, where $`e`$ denotes the neutral element of $`G`$. Conversely given a representation $`\tau `$ of $`H`$ on a finite dimensional vector space $`V`$ we let $`H`$ act from the right on $`G\times V`$ by $`(g,v)h=(gh,\tau (h)^1v)`$ and define $`E_\tau =(G\times V)/H`$. This is a vector bundle over $`X`$. This construction gives an equivalence of categories between the category of equivariant vector bundles over $`X`$ and the category of finite dimensional representations of $`H`$. If $`\tau `$ splits as a direct sum $`\tau =\tau _1\tau _2`$ then $`E_\tau E_{\tau _1}E_{\tau _2}`$ and vice versa.
Now let $`G`$ denote a connected semisimple Lie group with finite center and let $`K`$ denote a maximal compact subgroup. The quotient $`X=G/K`$ is diffeomorphic to $`^n`$ for some $`n`$. Using the Killing form of $`G`$ one defines a Riemannian metric on $`X`$ such that the group $`G`$ acts by isometries. This is the most general symmetric space without compact factors. See for further details.
## 2 Differential operators
Let $`\widehat{K}`$ denote the set of isomorphism classes of irreducible unitary representations of the group $`K`$. Since $`K`$ is compact, every $`\tau \widehat{K}`$ is finite dimensional. We do not distinguish between a class in $`\widehat{K}`$ and a representative. For $`(\tau ,V_\tau )\widehat{K}`$ let $`E_\tau G/K`$ denote the vector bundle as in section 1. The group $`G`$ acts on the space of smooth sections $`\mathrm{\Gamma }^{\mathrm{}}(E_\tau )`$ of the bundle $`E_\tau `$ by $`g.s(x)=gs(g^1x)`$, $`xX`$, $`gG`$, $`s\mathrm{\Gamma }^{\mathrm{}}(E_\tau )`$. Hence $`G`$ also acts on differential operators by conjugation. Let $`𝒟_\tau `$ denote the algebra of $`G`$-invariant differential operators on $`E_\tau `$, i.e. those operators $`D`$ that satisfy $`D(g.s)=g.D(s)`$ for any $`s\mathrm{\Gamma }^{\mathrm{}}(E_\tau )`$.
Let $`C_\tau ^{\mathrm{}}(G)`$ denote the space of all infinitely often differentiable maps from $`G`$ to $`V`$ with $`f(kx)=\tau (k)f(x)`$ for $`kK`$ and $`xG`$. The group $`G`$ acts on the space $`C_\tau ^{\mathrm{}}(G)`$ by translations from the right. The map $`sf_s`$ with $`f_s(x)=xf_s(x^1)`$ gives a $`G`$-isomorphism of $`\mathrm{\Gamma }^{\mathrm{}}(E_\tau )`$ to $`C_\tau ^{\mathrm{}}(G)`$. We conclude that the algebra $`𝒟_\tau `$ acts on $`C_\tau ^{\mathrm{}}(G)`$.
Let the compact group $`K`$ act on the space $`C^{\mathrm{}}(G)`$ of smooth functions on $`G`$ by
$$L_k(x)=f(k^1x),kK,xG.$$
Then there is a decomposition into $`K`$-isotypes:
$$C^{\mathrm{}}(G)=\underset{\tau \widehat{K}}{}C^{\mathrm{}}(G)(\tau ).$$
The sum here means that finite sums on the right hand side are dense in the Fréchet space $`C^{\mathrm{}}(G)`$. Further $`C^{\mathrm{}}(G)(\tau )`$ is the space of functions in $`C^{\mathrm{}}(G)`$ that transform under $`L`$ according to $`\tau `$.
For $`(\tau ,V)\widehat{K}`$ fix any nonzero $`v^{}`$ in the dual space $`V^{}`$ and for $`fC_\tau ^{\mathrm{}}(G)`$ let $`Bf(x)=v^{}(f(x))`$. Then $`BfC^{\mathrm{}}(G)(\tau )`$ and the map $`B`$ is a $`G`$-isomorphism. We have
$$\mathrm{\Gamma }^{\mathrm{}}(E_\tau )C_\tau ^{\mathrm{}}(G)C^{\mathrm{}}(G)(\tau ).$$
Let $`𝔤`$ denote the Lie algebra of $`G`$. The universal enveloping algebra $`U(𝔤)`$ may be viewed as the algebra of all right invariant differential operators on $`G`$, i.e. the algebra of all differential operators $`D`$ on $`G`$ such that $`R_gD=DR_g`$ for any $`gG`$. Here for a function $`\phi C^{\mathrm{}}(G)`$ and $`gG`$ we have $`R_g\phi (x)=\phi (xg)`$ for all $`xG`$. Let $`U(𝔤)^K`$ denote the subalgebra of all differential operators which are $`K`$-invariant on the left side, i.e. which satisfy $`L_kD=DL_k`$ for any $`kK`$. The algebra $`U(𝔤)^K`$ leaves stable the decomposition of $`C^{\mathrm{}}(G)`$ and thus acts on $`C^{\mathrm{}}(G)(\tau )`$. We therefore get an algebra homomorphism
$$\phi _\tau :U(𝔤)^K𝒟_\tau .$$
###### Proposition 2.1
The homomorphism $`\phi _\tau `$ is surjective and the intersection of all kernels $`\mathrm{ker}(\phi _\tau )`$ for varying $`\tau `$ is zero.
Proof: Let the algebra $`U(𝔤)\mathrm{End}(V)`$ acts on the space $`C^{\mathrm{}}(G,V)`$ of smooth functions on $`G`$ with values in $`V`$. The group $`K`$ acts on $`U(𝔤)`$ via the adjoint representation and on $`\mathrm{End}(V)`$ by means of conjugation via $`\tau `$. Then the algebra of $`K`$-invariants, $`(U(𝔤)\mathrm{End}(V))^K`$ acts on $`C_\tau ^{\mathrm{}}`$. The annihilator $`I`$ of $`C_\tau ^{\mathrm{}}`$ in $`U(𝔤)\mathrm{End}(V)`$ is generated by elements of the form $`XYT+XT\tau (Y)`$ with $`XU(𝔤)`$, $`Y𝔨`$ and $`T\mathrm{End}(V)`$. Since $`K`$ is reductive we have
$$𝒟_\tau (U(𝔤)\mathrm{End}(V))/I)^K(U(𝔤)\mathrm{End}(V))^K/I^K$$
Since $`\tau `$ is irreducible and $`K`$ is connected it follows $`\tau (U(k))=\mathrm{End}(V)`$. This implies that any element of $`𝒟_\tau `$ can be written in the form $`Z1`$ for some $`ZU(𝔤)`$. It follows that $`Z`$ must be in $`U(𝔤)^K`$, which implies the surjectivity of $`\phi _\tau `$.
For the second assertion assume $`X`$ is in the intersection of all kernels. Then, since $`C^{\mathrm{}}(G)=_\tau C^{\mathrm{}}(G)(\tau )`$, we get that $`Xf=0`$ for every $`fC^{\mathrm{}}(G)`$, this gives $`X=0`$.
$`\frac{}{}`$ Q.E.D.
###### Corollary 2.4
$`𝒟_\tau `$ is finitely generated as $``$-algebra.
Proof: $`U(𝔤)`$ has a natural filtration by order. The associated graded version equals the symmetric algebra $`S(𝔤)`$ over $`𝔤`$. Since the adjoint action of $`K`$ preserves the filtration, we have a filtration on $`U(𝔤)^K`$ with graded version $`S(𝔤)^K`$. The latter is finitely generated by invariant theory, hence the former is, too.
$`\frac{}{}`$ Q.E.D.
Examples.
For $`\tau =1`$ the trivial representation, the algebra $`𝒟_\tau `$ is the algebra of $`G`$-invariant differential operators on $`G/K`$. In this case $`𝒟_\tau `$ is isomorphic to the polynomial ring in $`r`$ generators, where $`r`$ is the rank of the symmetric space $`G/K`$.
Let $`G=SO(n,1)^+`$ the group of motions on the hyperbolic space $`H_n`$. For $`\tau =^p(\mathrm{Ad})`$ the space $`\mathrm{\Gamma }^{\mathrm{}}(E_\tau )`$ is just the space of $`p`$-differential forms on $`H_n`$. For $`p=0`$ or $`p=n`$ the algebra $`𝒟_\tau `$ is the polynomial ring in one variable, generated by $`\delta d`$ and $`d\delta `$ respectively, where $`d`$ is the exterior differential and $`\delta `$ its formal adjoint. For $`n`$ odd and $`p=(n\pm 1)/2`$ the algebra $`𝒟_\tau `$ is generated by $`d\delta `$ and $`d`$ resp. $`\delta d`$ and $`d`$, where $``$ is the Hodge operator, with the generating relations
$$d\delta d=0=dd\delta ,\delta dd=d\delta d=0.$$
In all other cases $`𝒟_\tau `$ is generated by $`d\delta `$ and $`\delta d`$ qith the relations $`d\delta \delta d=0=\delta dd\delta `$. Summarizing we get the structure of $`𝒟_\tau `$ as:
$$𝒟_\tau \{\begin{array}{cc}[x,y]/xy& \mathrm{for}1pn1,\hfill \\ [x]& 𝔪forp=0,n,\hfill \end{array}$$
## 3 Integral operators
A smooth function $`\mathrm{\Phi }:G\mathrm{End}(V)`$ which satisfies $`\mathrm{\Phi }(kxl)=\tau (k)\mathrm{\Phi }(x)\tau (l)`$ for $`xG`$, $`k,lK`$ is called $`\tau `$-sherical. The algebra $`𝒟_\tau `$ acts on the set of $`\tau `$-spherical functions. Compactly supported $`\tau `$-spherical functions form an algebra under convolution:
$$\mathrm{\Phi }\mathrm{\Psi }(x)=_G\mathrm{\Phi }(y)\mathrm{\Psi }(y^1x)𝑑y,$$
where $`dy`$ denotes a Haar measure on $`G`$. This algebra is denoted by $`𝒜_\tau `$. The algebra $`𝒜_\tau `$ acts on $`C_\tau ^{\mathrm{}}(G)`$ by $`L_\mathrm{\Phi }f=\mathrm{\Phi }f`$ for $`\mathrm{\Phi }𝒜_\tau `$ and $`fC_\tau ^{\mathrm{}}`$. The algebra $`𝒜_\tau `$ contains an approximate identity.
Let $`D_1,\mathrm{},D_n`$ be a set of generators of the algebra $`𝒟_\tau `$. For $`z^n`$ let
$$\mathrm{Eig}(z)=\{fC^{\mathrm{}}(G,\mathrm{End}(V))|f\mathrm{is}\tau \mathrm{spherical}\mathrm{and}D_jf=z_jf\}.$$
###### Proposition 3.5
For any $`z^n`$ we have $`dim\mathrm{Eig}(z)1`$.
Proof: By Lemma 1 in any $`f\mathrm{Eig}(z)`$ is analytic. So let $`f\mathrm{Eig}(z)`$ and $`H𝔤`$. If $`H`$ is small enough then
$$f(\mathrm{exp}(H))=\underset{n0}{}\frac{1}{n!}\left(\frac{}{t}\right)^nf(\mathrm{exp}(tH))|_{t=0}=\underset{n0}{}\frac{1}{n!}H^nf(0).$$
We get
$$_K\tau (k)f(\mathrm{exp}(H))\tau (k^1)𝑑k=\underset{n0}{}_K(\mathrm{Ad}(k)H)^nf(e)𝑑k.$$
Now $`_K(\mathrm{Ad}(k)H)^n𝑑k`$ is a right invariant differential operator mapping $`C_\tau (G)`$ to itself, hence defines an element of $`𝒟_\tau `$. So the values $`_K(\mathrm{Ad}(k)H)^nf(e)𝑑k`$ only depend on $`z`$ and not on $`f`$. We conclude that the function $`\mathrm{tr}f(x)`$ is determined by $`z`$ up to scalar. But the map $`𝒜_\tau C^{\mathrm{}}(G)`$, $`f\mathrm{tr}(f)`$ is injective .
$`\frac{}{}`$ Q.E.D.
###### Theorem 3.10
Suppose $`fC_\tau ^{\mathrm{}}(G)`$ is an eigenform for any $`D𝒟_\tau `$. Then $`f`$ is an eigenform for every $`T𝒜_\tau `$ with an eigenvalue only depending on $`T`$ and the eigenvalues on $`𝒟_\tau `$.
Proof: For $`\phi C_\tau ^{\mathrm{}}(G)`$ and $`wV`$ define $`\phi _wC^{\mathrm{}}(G,\mathrm{End}(V))`$ by
$$\phi _w(x)v=v,w\phi (x),$$
where $`.,.`$ denotes the scalar product on $`V`$. Fot $`psi:G\mathrm{End}(V)`$ let
$$M\psi (x)=_K\psi (xk)\tau (k^1)𝑑k.$$
Now let $`f`$ be as in the theorem and assume $`f(e)0`$. (Otherwiese replace $`f`$ by $`R_gf`$ for a suitable $`gG`$.) Fix some $`wV`$ such that $`\mathrm{tr}f_w(e)=1`$. Now $`M(f_w)`$ lies in $`\mathrm{Eig}(z)`$ for some $`z`$. Let $`\mathrm{\Phi }`$ be $`\tau `$-spherical and compactly supported. Since $`L_\mathrm{\Phi }M(f_w)=M((L_\mathrm{\Phi }f)_w)`$ we see that there is a $`\lambda `$ such that
$$M((\lambda fL_\mathrm{\Phi }f)_w)=0,$$
and $`\lambda `$ does not depend on $`f`$ or $`w`$. The claim now follows by the proposition.
$`\frac{}{}`$ Q.E.D.
Let $`𝔞`$ denote the Lie algebra of the maximal $``$-split torus $`A`$ of $`G`$. Let $`G=KNA`$ be a corresponding Iwasawa decomposition. For $`\lambda 𝔞_{}^{}`$ the complex dual space of $`𝔞`$ and any $`vV`$ let
$$p_{\lambda ,v}(kna)=\tau (k)e^{\lambda (\mathrm{log}(a))}v.$$
This defines $`p_{\lambda ,v}C_\tau ^{\mathrm{}}(G)`$. Let $`P=MAN`$ be the minimal parabolic given by $`A`$ and $`N`$. Then $`M`$ is a closed subgroup of $`K`$.
###### Lemma 3.13
For every $`T𝒜_\tau 𝒟_\tau `$ and every $`\lambda 𝔞_{}^{}`$ there is a $`S_\lambda (T)`$ in $`\mathrm{End}(V)`$ such that for all $`vV`$
$$T(p_{\lambda ,v})=p_{\lambda ,S_\lambda (T)v}.$$
We have $`S_\lambda (T)\tau (m)=\tau (m)S_\lambda (T)`$ for all $`mM`$
Proof: The lemma is clear by group invariance and the fact that $`A`$ normalizes $`N`$.
$`\frac{}{}`$ Q.E.D.
For a simultaneous eigenform $`fC_\tau ^{\mathrm{}}(G)`$ of $`𝒟_\tau `$ let $`\chi _f`$ denote the eigencharacter $`\chi _f:𝒜_\tau 𝒟_\tau `$ defined by
$$Tf=\chi _f(T)f.$$
###### Theorem 3.14
Let $`f`$ denote a bounded simultaneous eigenform. Then there is a $`vV`$ and a $`\lambda 𝔞_{}^{}`$ such that $`p_{\lambda ,v}`$ is an eigenform and
$$\chi _f=\chi _{p_{\lambda ,v}}.$$
Proof: The character $`\chi _f`$ is determined by its values on $`𝒜_\tau `$. Let $`p`$ denote the trivial seminorm on $`G`$, i.e. $`p(g)=1`$ for all $`gG`$. Set
$$\mathrm{\Phi }_p=_G\mathrm{\Phi }(y)p(y)𝑑y,$$
where $`.`$ is the norm on $`\mathrm{End}(V)`$. Assume $`f(e)=1`$ and $`f(x)M`$, $`xG`$. Then we get for $`\mathrm{\Phi }𝔸_\tau `$:
$`|\chi _f(L_\mathrm{\Phi })|`$ $`=`$ $`{\displaystyle _G}\mathrm{\Phi }(y)f(y^1)𝑑y`$
$``$ $`M{\displaystyle _G}\mathrm{\Phi }(y)𝑑y`$
$``$ $`M\mathrm{\Phi }_p.`$
Thus $`\chi _f`$ is a $`p`$-continuous representation of $`𝒜_\tau `$. The claim now follows from the theorem of Glover , p.40.
$`\frac{}{}`$ Q.E.D.
We now give the computation of the eigencharacters of $`𝒜_\tau `$. Let $`\rho =\frac{1}{2}_{\alpha >0}m_\alpha \alpha 𝔞^{}`$ be the usual modular shift, i.e. the sum runs over all positive roots and $`m_\alpha `$ is the dimension of the root space to the root $`\alpha `$.
###### Theorem 3.15
Let $`S_\lambda (L_\mathrm{\Phi })`$ denote the endomorphism of Lemma 3.13. Then with
$$g_\mathrm{\Phi }(a)=a^\rho _N\mathrm{\Phi }(na)𝑑n,$$
(Abel transform), we get
$$S_\lambda (L_\mathrm{\Phi })=_Aa^{\rho \lambda }g_\mathrm{\Phi }(a)𝑑a,$$
(Fourier transform on $`A`$). Moreover, $`g_\mathrm{\Phi }(a)`$ is in the center of $`\tau |_M`$ and $`g_{\mathrm{\Phi }\mathrm{\Psi }}=g_\mathrm{\Phi }g_\mathrm{\Psi }`$ with $`A`$-convolution on the right hand side. The map $`\mathrm{\Phi }g_\mathrm{\Phi }`$ is injective.
Proof: A calculation using the integral formula of the Iwasawa decomposition gives the first claim. The injectivity is proved in , p.35.
$`\frac{}{}`$ Q.E.D.
## 4 Surjectivity of differential operators
Let $`\theta `$ denote the Cartan involution fixing $`K`$ pointwise. Let $`\mathrm{\Gamma }_1,\mathrm{},\mathrm{\Gamma }_r`$ denote a complete system of nonconjugate $`\theta `$-stable Cartan subgroups of $`G`$ and let $`A_i=\mathrm{\Gamma }_i\mathrm{exp}(𝔭)`$, where $`𝔤=𝔨𝔭`$ denotes the polar decomposition of $`𝔤`$. Let $`A`$ denote one of the $`A_i`$. Let $`L=MA`$ denote the centralizer of $`A`$ in $`G`$. Let $`\tau _M`$ denote the restriction f $`\tau `$ to $`K_M=KM`$. Let $`C^{\mathrm{}}(MA,\tau )`$ denote the set of $`\tau _M`$-spherical functions on $`MA`$. For $`gC^{\mathrm{}}(MA,\tau _M)`$ let $`g^\mathrm{\#}(kman)=\tau (k)g(ma)`$. The set of these functions $`g^\mathrm{\#}`$ is stable under $`𝒟_\tau `$ and we get a homomorphism
$$\gamma :𝒟_\tau 𝒟_{\tau _M}(MA)=𝒟_{\tau _M}U(𝔞).$$
Now every $`T𝒟_{\tau _M}`$ operates on the finite dimensional space of rapidly decreasing cusp forms $`{}_{}{}^{0}𝒞(M,\tau _M)`$ as defined in . So $`\gamma (D)`$ defines an element $`{}_{}{}^{0}\gamma (D)|End(^0𝒞(M,\tau _M))U(𝔞)`$. Let $`\mathrm{det}(^0\gamma (D))U(𝔞)`$ denote the determinant. Call $`D`$ regular if $`\mathrm{det}(^0\gamma (D))0`$ for all $`A=A_i`$.
###### Theorem 4.16
If $`D𝒟_\tau `$ is regular, then $`D`$ is surjective as a map from $`C_\tau ^{\mathrm{}}`$ to $`C_\tau ^{\mathrm{}}`$.
Proof: We are going to formulate a vector bundle version of Holmgren’s uniqueness theorem. Let $`X`$ denote a real analytuc manifold and $`E`$ an analytic vector bundle over $`X`$. Let $`P`$ be a differential operator on $`E`$ with analytic coefficients. For a susbset $`A`$ of $`X`$ let $`N(A)`$ denote the set of normal vectors to $`A`$ in $`T^{}X`$, (see , chap. 8). Consider the determinant of the principal symbol $`\sigma _P`$ as a map:
$$\mathrm{det}\sigma _p:T^{}X.$$
###### Proposition 4.21
(Holmgren’s principle) Let $`u\mathrm{\Gamma }_c^{\mathrm{}}(E)^{}`$ be a distribution with $`Pu=0`$. Then we have
$$\mathrm{det}\sigma _P(N(\mathrm{supp}(u)))=0.$$
Proof: Since the assertion is local in nature the proof for the trivial bundle (, Theorem 8.6.5) carries over to the present case.
$`\frac{}{}`$ Q.E.D.
Let $`P`$ be as above. A point $`xX`$ is called a singular point of $`P`$ if
$$\mathrm{det}\sigma _P(T_x^{}X)=0.$$
Now consider the case when $`X`$ is a symmetric space $`G/K`$ and let $`E=E_\tau `$ be a homogeneous bundle. Assume that $`P`$ is invariant, i.e. $`P𝒟_\tau `$. Let $``$ denote the canonical homogeneous connection on $`E`$. For $`\lambda 𝔞_{}^{}`$ let $`p_\lambda (ank)=e^{\lambda (\mathrm{log}(a))}`$.
###### Lemma 4.26
There is a section $`\gamma _P`$ of the bundle $`S(𝔞)\mathrm{End}(E)`$ such that
$$P(p_\lambda s)=\gamma _P(\lambda )p_\lambda s,$$
for all parallel sections $`s`$ of $`E`$. For this section $`\gamma _P`$ and any parallel $`s`$ we have
$$\sigma _P(dp_\lambda )s=\gamma _{P,m}s,$$
where $`m=\mathrm{deg}(P)`$ and $`\gamma _{P,m}`$ is the principal part of $`\gamma _P`$ with respect to the gradation of $`S(𝔞)`$.
Proof: The first part is well known, the second follows from a calculation in Iwasawa coordinates.
$`\frac{}{}`$ Q.E.D.
A differential operator $`P`$ on $`E`$ is called $`D`$-convex, if for any Weyl group stable compact convex subset $`𝔰`$ of $`𝔞`$ and anny section $`s`$ of $`E`$ with compact support and $`\mathrm{supp}(Ps)K\mathrm{exp}(𝔰)K`$ we already have $`\mathrm{supp}(s)K\mathrm{exp}(𝔰)K`$.
###### Theorem 4.23
Any $`P𝒟_\tau `$ with $`\mathrm{det}\sigma _P0`$ is $`D`$-convex.
Proof: Let $`P`$ be as in the theorem. Since $`\mathrm{det}\sigma _P`$ is $`G`$-invariant, it follows from $`\mathrm{det}\sigma _P0`$ that $`P`$ has no singular points. Let for $`xX`$, $`x=k\mathrm{exp}(H)K`$, $`H𝔞`$,
$`\delta (x)`$ $`=`$ $`inf\{t>0|{\displaystyle \frac{H}{t}}𝔰\},`$
$`\delta (A)`$ $`=`$ $`sup\{\delta (x)|xA\},\mathrm{for}AX.`$
Assume $`\delta (\mathrm{supp}Ps)=1`$ and $`\delta (\mathrm{supp}s)>>1`$. Let $`x_0\mathrm{supp}s`$ with $`\delta (x_0)=\alpha `$ and $`x_0=\mathrm{exp}(H)`$, $`H`$ in the positive Weyl chamber of $`𝔞`$. Let $`\lambda 𝔞_+^{}`$. With Kostants convexity theorem we get
$$p_\lambda (x)p_\lambda (x_0),$$
for all $`x\mathrm{supp}s`$. So $`dp_\lambda (x_0)N(\mathrm{supp}s)`$. Hence $`\mathrm{det}\sigma _P(dp_\lambda (x_0))=0`$ for all $`\lambda >>0`$, hence for all $`\lambda `$. By group invariance this gives that $`x_0`$ is a singular point for $`P`$, a contradiction. This proves 4.23.
$`\frac{}{}`$ Q.E.D.
Now regular operators $`P`$ satisfy the condition of 4.23. They further admit fundamental solutions . Now 4.16 follows as in .
$`\frac{}{}`$ Q.E.D.
School of Mathematical Sciences
University of Exeter
Laver Building, North Park Road
Exeter, EX4 4QE
Devon, UK
|
warning/0005/hep-th0005001.html
|
ar5iv
|
text
|
# Adding Matter to Poincare Invariant Branes
## I Introduction
The cosmological constant problem has evaded solution since its inception. If only, it has become more severe: the triumph of Quantum Field Theories as the correct description of the fundamental interactions came at the cost of additive contributions to the cosmological constant at big and disparate scales. It is natural to expect that the cosmological constant gets contributions of order $`M_P^4`$, where $`M_P`$ is the Planck scale, from short distance gravitational dynamics; $`M_W^4`$, where $`M_W`$ is the mass of the $`W`$-boson, from the phase transition associated with electroweak symmetry breaking; $`\mathrm{\Lambda }_{\mathrm{QCD}}^4`$, where $`\mathrm{\Lambda }_{\mathrm{QCD}}`$ is the scale of Quantum Chromo Dynamics, from the chiral symmetry breaking phase transition; etc. A proper solution to the problem has to explain how all such contributions are cancelled to absurdly high precision.
An intriguing solution has been proposed in which spacetime is five dimensional, but the observable universe is constrained to a four-dimensional hypersurface, a “3-brane.” The authors exhibit solutions to the field equations which give a flat, Poincare invariant, brane regardless of the value of the cosmological constant. The geometry of space includes naked singularities which are four dimensional hypersurfaces on either side of the brane on which spacetime ends. The significance of and consistency of theories with these singularities remains unclear. It has been suggested that the singularities hide a fine tuning equivalent to that required to set the cosmological constant to zero. In Ref. the gravitational action is modified by including a Gauss-Bonnet term, of second order in the curvature tensor, and it is found that the singularity can be smoothed out but only at the price of a fine tuning. It has also been shown that in some cases the field equations admit solutions that correspond to an Einstein-de Sitter universe on the brane. It has been suggested that this type of models may be derived from string theory.
The solution is however incomplete. There is no matter in the toy model of . It is necessary to incorporate matter if the solution is to be relevant to cosmology. This is a non-trivial issue: the standard paradigm, namely the big-bang cosmology based on a Friedman-Robertson-Walker (FRW) metric, is observationally very successful, so one should aim at reproducing, or at least approximating this paradigm once matter is added.
We investigate inclusion of matter into these models. We discover new, exact solutions to the brane models with matter. They are static and therefore incompatible with observational cosmology. Since the solutions are not unique, it is possible that other solutions exist which appropriately describe the expansion of the universe. Therefore we look for time dependent solutions by linearizing the field equations about our new solutions. As will be seen, generically the small perturbations correspond to propagating modes or to non-expanding universes (static solutions). Perturbations about special, non-generic backgrounds yield expanding universes that seem, however, inconsistent with observational cosmology.
The paper is organized as follows. In section II we briefly review the basic equations and the models of . In section III we present our new solutions with non-zero matter on the brane. In section IV we perform a small perturbation analysis about these new solutions. We present our conclusions in section V.
The cosmology of brane models has been investigated in a number of papers. A general formulation was given in Ref. . The work in Refs. is concerned with the cosmology of brane models of the Randall-Sundrum type. In addition, Randall-Sundrum models to which scalars are added have been of interest. A method to generate solutions to the non-linear field equations in classes of Randall-Sundrum models with scalars was given in Ref. . However, there has been little, and only very recent, work on the cosmology of automatically Poincare invariant branes.
## II Preliminaries
We denote the coordinates of spacetime by $`x^A`$, $`A=0,\mathrm{},4`$, and often use $`t=x^0`$ and $`y=x^4`$. The 3-brane is located at $`y=0`$. The class of spherically symmetric metrics we study is parameterized by three functions of $`t`$ and $`y`$ only
$$ds^2=G_{AB}dx^Adx^B=n^2(t,y)dt^2a^2(t,y)d\stackrel{}{x}^2b^2(t,y)dy^2.$$
(1)
Fixing $`y=0`$ we see that the metric gives a flat FRW cosmology on the brane with scale factor $`R(t^{})=a(t(t^{}),0)`$ where $`dt^{}=n(t,0)dt`$. We will denote by $`g_{\mu \nu }`$, with $`\mu ,\nu =0,\mathrm{},3`$, the induced metric on the brane.
The models of Refs. include a scalar field $`\varphi `$. The action is
$$S=d^5x\sqrt{G}\left[R+\frac{4}{3}(\varphi )^2\mathrm{\Lambda }e^{\overline{a}\varphi }\right]+d^4x\sqrt{g}\left[Ve^{\overline{b}\varphi }\right].$$
(2)
We have adopted the notation of Ref. , save for the constants $`\overline{a}`$ and $`\overline{b}`$ which we have adorned with a bar to distinguish them from the metric components $`a^2(t,y)`$ and $`b^2(t,y)`$. $`R`$ denotes the Ricci scalar.
The peculiar normalization of the scalar field is adopted from string theory: when $`\varphi `$ is a string theory dilaton its couplings are fixed. In particular, in this normalization, $`\overline{b}=2/3`$ at lowest order. We are not interested solely in this particular set of string theory inspired parameters, so we keep the values unspecified.
The constants $`\mathrm{\Lambda }`$ and $`V`$ represent the cosmological constant in the bulk (5-dimensional space) and on the brane, respectively. For our analysis we set $`\mathrm{\Lambda }=0`$. As seen in Ref. this simplifies the analysis without compromising the essential features of the model. Moreover, one could imagine that if the model is embedded in a supersymmetric setting, $`\mathrm{\Lambda }`$ could naturally vanish. The cosmological constant problem is associated with standard model fields’ contributions to $`V`$, $`VM_P^4+M_W^4+\mathrm{\Lambda }_{\mathrm{QCD}}^4+\mathrm{}`$
Einstein’s equations are
$$R^{AB}\frac{1}{2}G^{AB}R=\kappa ^2T^{AB}.$$
(3)
Here $`R^{AB}`$ and $`R`$ are the Ricci tensor and scalar. The gravitational constant is $`\kappa ^2`$ and from now on we work in units of $`\kappa ^2=1`$. $`T^{AB}`$ is the stress-energy tensor, which has two components:
$$T^{AB}=\stackrel{~}{T}^{AB}+\frac{S^{AB}}{b}\delta (y),$$
(4)
where $`\stackrel{~}{T}^{AB}`$ is derived as usual by varying the action with respect to the metric, and $`S^{AB}`$ is a contribution from a perfect fluid of density $`\rho `$ and pressure $`p`$ on the brane,
$$S_B^A=\text{diag}(\rho ,p,p,p,0).$$
(5)
Alternatively one may take the matter fluid to couple to the scalar, thus
$$S_B^A=\text{diag}(e^{\overline{b}\varphi }\rho ,e^{\overline{b}\varphi }p,e^{\overline{b}\varphi }p,e^{\overline{b}\varphi }p,0).$$
(6)
These two ways of writing $`S_B^A`$ correspond to distinct physical models. However the distinction turns out to be irrelevant for the exact solutions that we present in Sec. III. The field equation for the scalar $`\varphi `$ is
$$\frac{8}{3}^2\varphi \frac{\sqrt{g}}{\sqrt{G}}\delta (y)\overline{b}Ve^{\overline{b}\varphi }=0.$$
(7)
For the particular metric (1) Einstein’s equations and the $`\varphi `$ field equation are
$`3\left[\left(({\displaystyle \frac{\dot{a}}{a}})^2+{\displaystyle \frac{\dot{a}\dot{b}}{ab}}\right)+{\displaystyle \frac{n^2}{b^2}}\left({\displaystyle \frac{a^{\prime \prime }}{a}}({\displaystyle \frac{a^{}}{a}})^2+{\displaystyle \frac{a^{}b^{}}{ab}}\right)\right]`$ $`=`$ $`{\displaystyle \frac{2}{3}}n^2\left({\displaystyle \frac{\dot{\varphi }^2}{n^2}}+{\displaystyle \frac{\varphi _{}^{}{}_{}{}^{2}}{b^2}}\right)+\delta (y){\displaystyle \frac{n^2}{b}}({\displaystyle \frac{1}{2}}Ve^{\overline{b}\varphi }+\rho ),`$ (8)
$`3\left({\displaystyle \frac{\dot{a}n^{}}{an}}+{\displaystyle \frac{a^{}\dot{b}}{ab}}{\displaystyle \frac{\dot{a}^{}}{a}}\right)`$ $`=`$ $`{\displaystyle \frac{4}{3}}\dot{\varphi }\varphi ^{},`$ (9)
$`{\displaystyle \frac{a^2}{n^2}}\left({\displaystyle \frac{\dot{a}^2}{a^2}}2{\displaystyle \frac{\ddot{a}}{a}}+2{\displaystyle \frac{\dot{a}\dot{n}}{an}}2{\displaystyle \frac{\dot{a}\dot{b}}{ab}}{\displaystyle \frac{\ddot{b}}{b}}+{\displaystyle \frac{\dot{n}\dot{b}}{nb}}\right)+`$ (10)
$`{\displaystyle \frac{a^2}{b^2}}\left(({\displaystyle \frac{a^{}}{a}})^2+2{\displaystyle \frac{a^{\prime \prime }}{a}}+2{\displaystyle \frac{a^{}n^{}}{an}}2{\displaystyle \frac{a^{}b^{}}{ab}}+{\displaystyle \frac{n^{\prime \prime }}{n}}{\displaystyle \frac{n^{}b^{}}{nb}}\right)`$ $`=`$ $`{\displaystyle \frac{2}{3}}a^2\left({\displaystyle \frac{\dot{\varphi }^2}{n^2}}{\displaystyle \frac{\varphi _{}^{}{}_{}{}^{2}}{b^2}}\right)\delta (y){\displaystyle \frac{a^2}{b}}({\displaystyle \frac{1}{2}}Ve^{\overline{b}\varphi }p),`$ (11)
$`3\left[{\displaystyle \frac{b^2}{n^2}}\left({\displaystyle \frac{\dot{a}^2}{a^2}}{\displaystyle \frac{\ddot{a}}{a}}+{\displaystyle \frac{\dot{a}\dot{n}}{an}}\right)+\left(\left({\displaystyle \frac{a^{}}{a}}\right)^2+{\displaystyle \frac{a^{}n^{}}{an}}\right)\right]`$ $`=`$ $`{\displaystyle \frac{2}{3}}b^2\left({\displaystyle \frac{\dot{\varphi }^2}{n^2}}+{\displaystyle \frac{\varphi _{}^{}{}_{}{}^{2}}{b^2}}\right),`$ (12)
$`{\displaystyle \frac{8}{3}}{\displaystyle \frac{1}{b^2}}[(\ddot{\varphi }{\displaystyle \frac{b^2}{n^2}}\varphi ^{\prime \prime })\varphi ^{}(3{\displaystyle \frac{a^{}}{a}}{\displaystyle \frac{b^{}}{b}}+{\displaystyle \frac{n^{}}{n}})+`$ (13)
$`\dot{\varphi }{\displaystyle \frac{b^2}{n^2}}(3{\displaystyle \frac{\dot{a}}{a}}+{\displaystyle \frac{\dot{b}}{b}}{\displaystyle \frac{\dot{n}}{n}})]`$ $`=`$ $`\delta (y){\displaystyle \frac{1}{b}}V\overline{b}e^{\overline{b}\varphi }.`$ (14)
Here a dot is a shorthand for $`/t`$ and a prime for $`/y`$. The first four equations correspond to the 00, 04, 11 and 44 components of Einstein’s equations.
Conservation of the stress-energy tensor would be automatic were it derived from a local action. However, since a fluid component has been added on the brane, the equation $`T_{;B}^{AB}=0`$ contains additional non-trivial information. Conservation of energy, $`T_{;B}^{0B}=0`$, gives
$`{\displaystyle \frac{8}{3}}{\displaystyle \frac{1}{b^2}}\dot{\varphi }[(\ddot{\varphi }{\displaystyle \frac{b^2}{n^2}}\varphi ^{\prime \prime })\varphi ^{}(3{\displaystyle \frac{a^{}}{a}}{\displaystyle \frac{b^{}}{b}}+{\displaystyle \frac{n^{}}{n}})+`$ (15)
$`\dot{\varphi }{\displaystyle \frac{b^2}{n^2}}(3{\displaystyle \frac{\dot{a}}{a}}+{\displaystyle \frac{\dot{b}}{b}}{\displaystyle \frac{\dot{n}}{n}})]`$ $`=`$ $`{\displaystyle \frac{\delta (y)}{b}}\left[\dot{\varphi }V\overline{b}e^{\overline{b}\varphi }+2\left(\dot{\rho }+3{\displaystyle \frac{\dot{a}}{a}}(\rho +p)\right)\right]`$ (16)
while $`T_{;B}^{4B}=0`$ yields
$`{\displaystyle \frac{8}{3}}{\displaystyle \frac{1}{b^2}}\varphi ^{}[(\ddot{\varphi }{\displaystyle \frac{b^2}{n^2}}\varphi ^{\prime \prime })\varphi ^{}(3{\displaystyle \frac{a^{}}{a}}{\displaystyle \frac{b^{}}{b}}+{\displaystyle \frac{n^{}}{n}})+`$ (17)
$`\dot{\varphi }{\displaystyle \frac{b^2}{n^2}}(3{\displaystyle \frac{\dot{a}}{a}}+{\displaystyle \frac{\dot{b}}{b}}{\displaystyle \frac{\dot{n}}{n}})]`$ $`=`$ $`\delta (y){\displaystyle \frac{1}{b}}\left({\displaystyle \frac{n^{}}{n}}(Ve^{\overline{b}\varphi }+2\rho )+3{\displaystyle \frac{a^{}}{a}}(Ve^{\overline{b}\varphi }2p)\right).`$ (18)
Using the field equation for the scalar, Eq. (14), in the conservation of energy equation gives
$$\dot{\rho }+3\frac{\dot{a}}{a}(\rho +p)=0$$
(19)
on the brane. The second conservation equation, $`T_{;B}^{4B}=0`$, is always satisfied as a consequence of the $`\varphi `$ equation of motion and Einstein’s equations. Only the brane part of $`T_{;B}^{4B}=0`$ does not follow immediately from Eq. (14). The identity for $`y=0`$ can be confirmed by taking the $`y`$ derivative of Eq. (12) and then using Eqs. (8) and (11).
It is stated in Ref. that for brane geometries of the form given by Eq. (1) the equation of conservation of transverse momentum, $`T^{4A}{}_{;A}{}^{}=0`$, reduces, on the brane, to
$$\frac{n^{}}{n}\rho =3\frac{a^{}}{a}p.$$
(20)
It is understood here that when discontinuous quantities are evaluated on the brane, like $`n^{}`$ and $`a^{}`$, they are given by their average, i.e.,
$$n^{}(y=0)=\frac{1}{2}(n^{}(y=0+)+n^{}(y=0)).$$
(21)
The equation is seldom considered any further because, for $`Z_2`$ symmetric brane-spaces, that is for metrics with $`yy`$ symmetry, which are overwhelmingly most common, the averages both vanish separately, $`a^{}=0=n^{}`$, and the equation is trivially satisfied.
However the solutions we consider are not $`Z_2`$ symmetric. This is also true of the solutions in Refs. but there Eq. (20) is still trivially satisfied since the matter density and pressure both vanish. It is easy to see that our solutions below, Eqs. (34)–(36), do not satisfy Eq. (20). The reason is, in fact, that Eq. (20) does not apply to the case in which there are bulk scalars. The conservation of $`y`$-momentum is automatically satisfied in the bulk because the action is translation invariant in the bulk. So at issue here is only the conservation equation on the brane. Retaining only the bulk terms involving second derivatives, the conservation Eq. (17) gives
$$\frac{4}{3}\varphi ^{}\varphi ^{\prime \prime }+b\delta (y)\left[\frac{n^{}}{n}\rho +3\frac{a^{}}{a}p\frac{1}{2}Ve^{\overline{b}\varphi }\left(\frac{n^{}}{n}+3\frac{a^{}}{a}\right)\right]=0.$$
(22)
Following Ref. we interpret the discontinuous derivatives on the brane as averages, and this will remain implicit in what follows. Integrating gives a jump equation,
$$\frac{2}{3b}\mathrm{\Delta }\left(\varphi ^{}\right)^2=\frac{n^{}}{n}\rho 3\frac{a^{}}{a}p+\frac{1}{2}Ve^{\overline{b}\varphi }\left(\frac{n^{}}{n}+3\frac{a^{}}{a}\right)$$
(23)
We have verified that our solutions, Eqs. (34)–(36), and also the perturbations, Eqs. (59)–(61) and (63)–(65) satisfy this equation. In fact one can prove this without reference to the explicit form of the solution.
In Ref. it is advocated that the correct form of the conservation of transverse momentum equation is
$$T^{4\mu }{}_{,\mu }{}^{}=0,$$
(24)
where the index $`\mu `$ runs from $`0`$ to $`3`$ only. However, the first term in Eq. (22), involving the all important second derivative term $`\varphi ^{\prime \prime }\varphi ^{}`$, arises from the derivative $`T^{44}_{,4}`$, which is ommitted from Eq. (24).
Let us now describe the models studied in Refs. . They take $`\rho =p=0`$. The solutions all have $`n(t,y)=a(t,y)`$ and $`b(t,y)=1`$, and are static, $`\dot{a}=\dot{\varphi }=0`$. For example, case I studied in Ref. has $`\mathrm{\Lambda }=0`$ and solutions
$`n=a`$ $`=`$ $`\{\begin{array}{cc}(1y/y_+)^{\gamma _+},\hfill & \text{for }y>0\hfill \\ (1y/y_{})^\gamma _{},\hfill & \text{for }y<0\hfill \end{array}`$ (25)
$`\varphi `$ $`=`$ $`\{\begin{array}{cc}\phi _+\mathrm{log}(1y/y_+)+c,\hfill & y>0\hfill \\ \phi _{}\mathrm{log}(1y/y_{})+c,\hfill & y<0\hfill \end{array}.`$ (26)
The constants $`\gamma _\pm =1/4`$ and $`|\phi _\pm |=3/4`$ are fixed by the field equations in the bulk. Case I has, in particular, $`\phi _\pm =3/4`$. The constants $`y_\pm `$ are determined by “jump” conditions, that is, by requiring that the second derivatives of fields in the field equations correctly reproduce the $`\delta `$-function terms from the brane. The constant $`c`$ is an irrelevant constant shift of the scalar field.
## III Solutions with Matter
The models of Refs. provide a solution to the cosmological constant problem which is deficient in several ways: (1) There are naked singularities. Whether these are problematic remains an open question; see, for example, Refs. . (2) There is a massless scalar which interacts with all matter with a universal, gravity-like coupling strength. This is ruled out unless the coupling is made sufficiently weak. It can be arranged, however, by choosing the parameter $`\overline{b}`$ small enough. (3) It describes a static cosmology, in conflict with observation (see, however, Ref. ). (4) It does not include matter density (and pressure) on the brane.
Here we address the last two of these problems. One hopes the two are connected: when matter is included in the model the universe will evolve in time. Of course, not only should the universe evolve, but the rate of expansion should be adequate.
However, even after the introduction of matter the model admits static solutions. This is somewhat surprising, particularly if it is contrasted with the FRW cosmology which admits a static solution only if the matter density is precisely balanced by a cosmological term giving the Einstein universe. Since we are ultimately interested in time dependent solutions, we will look at time dependent small fluctuations about these solutions in the next section.
We look for solutions to the field equations in the bulk with the ansatz
$`a`$ $`=`$ $`y^\alpha ,`$ (27)
$`n`$ $`=`$ $`y^\nu ,`$ (28)
$`b`$ $`=`$ $`1,`$ (29)
$`\varphi `$ $`=`$ $`\phi \mathrm{log}y.`$ (30)
The $`\varphi `$ field equation gives
$$3\alpha +\nu =1.$$
(31)
All Einstein field equations give then
$$2\alpha ^2\alpha +\frac{2}{9}\phi ^2=0.$$
(32)
For definiteness we consider solutions akin to case I of Ref. , with
$$\phi =3\sqrt{\alpha /2\alpha ^2},$$
(33)
where the upper and lower signs correspond to the regions $`y>0`$ and $`y<0`$, respectively.
The full solution is found by shifting $`y`$ by $`y_+`$($`y_{}`$) on $`y>0`$($`y<0`$), and pasting these using the jump equations. We have
$`a`$ $`=`$ $`A\left(1{\displaystyle \frac{y}{y_\pm }}\right)^{\alpha _\pm },`$ (34)
$`n`$ $`=`$ $`N\left(1{\displaystyle \frac{y}{y_\pm }}\right)^{\nu _\pm },`$ (35)
$`\varphi `$ $`=`$ $`\phi _\pm \mathrm{log}(1y/y_\pm )+c.`$ (36)
Here $`A`$ and $`N`$ are arbitrary constants that can be set to unity by a coordinate rescaling. The jump equations are
$`{\displaystyle \frac{\mathrm{\Delta }a^{}}{a}}`$ $`=`$ $`{\displaystyle \frac{b}{3}}\left({\displaystyle \frac{1}{2}}Ve^{\overline{b}\varphi }+\rho \right),`$ (37)
$`2{\displaystyle \frac{\mathrm{\Delta }a^{}}{a}}+{\displaystyle \frac{\mathrm{\Delta }n^{}}{n}}`$ $`=`$ $`b\left({\displaystyle \frac{1}{2}}Ve^{\overline{b}\varphi }+p\right),`$ (38)
$`\mathrm{\Delta }\varphi ^{}`$ $`=`$ $`{\displaystyle \frac{3}{8}}b\overline{b}Ve^{\overline{b}\varphi },`$ (39)
where $`\mathrm{\Delta }a^{}=a^{}(y=0+)a^{}(y=0)`$, etc. These give three equations for five unknowns,
$`{\displaystyle \frac{\alpha _+}{y_+}}{\displaystyle \frac{\alpha _{}}{y_{}}}`$ $`=`$ $`{\displaystyle \frac{1}{3}}\left({\displaystyle \frac{1}{2}}Ve^{\overline{b}c}+\rho \right),`$ (40)
$`{\displaystyle \frac{1}{y_+}}{\displaystyle \frac{1}{y_{}}}`$ $`=`$ $`p+{\displaystyle \frac{1}{3}}\rho +{\displaystyle \frac{2}{3}}Ve^{\overline{b}c},`$ (41)
$`{\displaystyle \frac{\sqrt{\alpha _+/2\alpha _+^2}}{y_+}}+{\displaystyle \frac{\sqrt{\alpha _{}/2\alpha _{}^2}}{y_{}}}`$ $`=`$ $`{\displaystyle \frac{1}{8}}\overline{b}Ve^{\overline{b}c}.`$ (42)
These can always be solved for three unknowns (say $`\alpha _+`$ and $`y_\pm `$) in terms of $`\rho `$, $`p`$ and two other unknowns (say $`\alpha _{}`$ and $`c`$). The reason not all unknowns are determined is twofold. First, gauge (diffeomorphism) invariance allows us to make unphysical changes to our solutions. This will be explained below in detail, but for now it suffices to know that one may fix the gauge freedom by setting, say, $`c=0`$. And secondly, even for vanishing $`\rho `$ and $`p`$ one can find a class of solutions parameterized by one parameter. These are new solutions to the model considered in , which could not be discovered there because it was assumed that $`n=a`$. Indeed, imposing this one has $`\nu =\alpha `$ which together with Eqs. (31) and (33) imply
$$\alpha =\nu =\frac{1}{4}$$
(43)
and
$$\phi =\frac{3}{4}$$
(44)
as found in case I of .
We have searched for time dependent solutions. We found one special solution, valid only for $`\rho =0`$ and provided $`\overline{b}=2/3`$. It is given by
$`ds^2`$ $`=`$ $`(1y/y_\pm )^2dt^2(1y/y_\pm )^2t^2d\stackrel{}{x}^2t^2dy^2,`$ (45)
$`\varphi `$ $`=`$ $`{\displaystyle \frac{3}{2}}\mathrm{ln}[t(1y/y_\pm )]+c,`$ (46)
$`{\displaystyle \frac{1}{y_+}}{\displaystyle \frac{1}{y_{}}}`$ $`=`$ $`{\displaystyle \frac{1}{6}}Ve^{2c/3}.`$ (47)
Perplexingly, the special value $`\overline{b}=2/3`$ corresponds to the string theory value of this parameter. Other non-static solutions were given in Ref. . In this paper we will not consider these type of non-static matter-free solutions further.
## IV Small perturbations about Large Matter Density
Armed with the new solutions with static matter density, we proceed to investigate the time dependence of small matter perturbations. Let us denote the static solution of the previous section by $`n_0`$, $`a_0`$, $`b_0`$, $`\varphi _0`$, $`\rho _0`$ and $`p_0`$. We look for solutions to the field equations, Eqs. (8)–(14), of the form
$`n`$ $`=`$ $`n_0(1+\delta n),`$ (48)
$`a`$ $`=`$ $`a_0(1+\delta a),`$ (49)
$`b`$ $`=`$ $`b_0(1+\delta b),`$ (50)
$`\varphi `$ $`=`$ $`\varphi _0+\delta \varphi ,`$ (51)
$`\rho `$ $`=`$ $`\rho _0+\delta \rho ,`$ (52)
$`p`$ $`=`$ $`p_0+\delta p.`$ (53)
We count orders of the perturbative expansion parametrically in $`\delta \rho `$ and $`\delta p`$. That is, we re-scale $`\delta \rho ϵ\delta \rho `$, count powers of $`ϵ`$ and set $`ϵ=1`$ at the end of the calculation. In particular this implies that we make no assumption as to the relative importance of temporal or spatial derivatives.
To derive the linearized equations in the bulk, we use again this parameterization and the explicit form of the zeroth order solutions, Eq. (27). The 00, 04, 11 and 44 components of Einstein’s equations and the $`\varphi `$ field equation give
$`\delta a^{\prime \prime }+4{\displaystyle \frac{\alpha }{y}}\delta a^{}{\displaystyle \frac{\alpha }{y}}\delta b^{}+{\displaystyle \frac{4}{9}}{\displaystyle \frac{\phi }{y}}\delta \varphi ^{}`$ $`=`$ $`0,`$ (54)
$`{\displaystyle \frac{}{t}}\left({\displaystyle \frac{\nu }{y}}\delta a+{\displaystyle \frac{\alpha }{y}}\delta b{\displaystyle \frac{\alpha }{y}}\delta a\delta a^{}{\displaystyle \frac{4}{9}}{\displaystyle \frac{\phi }{y}}\delta \varphi \right)`$ $`=`$ $`0,`$ (55)
$`2\delta a^{\prime \prime }+{\displaystyle \frac{2}{y}}\delta a^{}+2{\displaystyle \frac{\alpha +\nu }{y}}\delta n^{}+\delta n^{\prime \prime }{\displaystyle \frac{2\alpha +\nu }{y}}\delta b^{}+{\displaystyle \frac{4}{3}}{\displaystyle \frac{\phi }{y}}\delta \varphi ^{}`$ $`=`$ $`{\displaystyle \frac{1}{n_0^2}}\left(2\delta \ddot{a}+\delta \ddot{b}\right),`$ (56)
$`{\displaystyle \frac{2\alpha +\nu }{y}}\delta a^{}+{\displaystyle \frac{\alpha }{y}}\delta n^{}{\displaystyle \frac{4}{9}}{\displaystyle \frac{\phi }{y}}\delta \varphi ^{}`$ $`=`$ $`{\displaystyle \frac{1}{n_0^2}}\delta \ddot{a},`$ (57)
$`{\displaystyle \frac{1}{y}}{\displaystyle \frac{}{y}}\left[y\delta \varphi ^{}+\phi (3\delta a+\delta n\delta b)\right]`$ $`=`$ $`{\displaystyle \frac{1}{n_0^2}}\delta \ddot{\varphi }.`$ (58)
The solution to these equations gives $`\delta \varphi `$, $`\delta b`$ and $`\delta n^{}`$ in terms of $`\delta a`$:
$`\delta \varphi `$ $`=`$ $`{\displaystyle \frac{1}{\alpha }}(\phi \delta aF),`$ (59)
$`\delta b`$ $`=`$ $`{\displaystyle \frac{1}{\alpha }}\left[(y\delta a)^{}{\displaystyle \frac{4}{9}}{\displaystyle \frac{\phi }{\alpha }}F\xi \right],`$ (60)
$`\delta n^{}`$ $`=`$ $`{\displaystyle \frac{1}{\alpha }}\left[{\displaystyle \frac{y}{n_0^2}}\delta \ddot{a}{\displaystyle \frac{}{y}}\left((3\alpha 1)\delta a+{\displaystyle \frac{4}{9}}{\displaystyle \frac{\phi }{\alpha }}F\right)\right].`$ (61)
Here $`\xi `$ is an arbitrary constant and $`F`$ is a function of $`t`$ and $`y`$ satisfying
$$\frac{1}{y}\frac{}{y}\left(y\frac{F}{y}\right)\frac{1}{n_0^2}\frac{^2F}{t^2}=0.$$
(62)
We connect the bulk solutions for $`y>0`$ and $`y<0`$ demanding continuity of the fields at the brane, $`y=0`$, and using the jump equations (37) for the discontinuous derivatives at $`y=0`$. The latter give jump conditions for the perturbations:
$`\mathrm{\Delta }\delta a^{}`$ $`=`$ $`{\displaystyle \frac{1}{6}}Ve^{\overline{b}\varphi _0}(\delta b+\overline{b}\delta \varphi ){\displaystyle \frac{1}{3}}\delta \rho {\displaystyle \frac{1}{3}}\rho _0\delta b,`$ (63)
$`\mathrm{\Delta }\delta n^{}`$ $`=`$ $`{\displaystyle \frac{1}{6}}Ve^{\overline{b}\varphi _0}(\delta b+\overline{b}\delta \varphi )+\delta p+{\displaystyle \frac{2}{3}}\delta \rho +(p_0+{\displaystyle \frac{2}{3}}\rho _0)\delta b,`$ (64)
$`\mathrm{\Delta }\delta \varphi ^{}`$ $`=`$ $`{\displaystyle \frac{3}{8}}\overline{b}Ve^{\overline{b}\varphi _0}(\delta b+\overline{b}\delta \varphi ).`$ (65)
It must be observed that both $`\xi `$ and $`F`$ can have different values on either side of the brane and, moreover, that $`F`$ may be discontinuous at $`y=0`$. In addition, conservation of energy gives, on the brane,
$$\delta \rho +3(\rho _0+p_0)\delta a=\delta \rho _0,$$
(66)
where $`\delta \rho _0`$ is a constant of integration.
There is a gauge freedom, that is, reparameterization invariance consistent with the form of our metric. Starting from the metric
$$ds^2=n^2(t^{},y^{})dt^2a^2(t^{},y^{})d\stackrel{}{x}^2b^2(t^{},y^{})dy^2,$$
(67)
we look for infinitesimal transformations
$`t^{}`$ $`=`$ $`t+T(t,y)`$ (68)
$`y^{}`$ $`=`$ $`y+Y(t,y)`$ (69)
that leave the form of the metric invariant. Here $`T`$ and $`Y`$ are infinitesimal. The only constraint on these functions comes from the absence of off-diagonal terms in the metric:
$$n^2T^{}b^2\dot{Y}=0.$$
(70)
Under the gauge transformation the metric variations are
$`\delta n`$ $`=`$ $`{\displaystyle \frac{n_0^{}}{n_0}}Y+\dot{T},`$ (71)
$`\delta a`$ $`=`$ $`{\displaystyle \frac{a_0^{}}{a_0}}Y,`$ (72)
$`\delta b`$ $`=`$ $`Y^{},`$ (73)
$`\delta \varphi `$ $`=`$ $`\varphi _0^{}Y.`$ (74)
For simplicity we have indicated the variation about a static solution with $`b_0=1`$ and $`\dot{a}_0=\dot{n}_0=\dot{\varphi }_0=0`$. It is instructive to check that our solutions of the field equations for the perturbations are invariant under these transformations, that is, that the perturbations (71)-(74) satisfy the field equations automatically. Indeed, the solution for $`\delta \varphi `$, Eq. (59), is satisfied provided one takes $`F=0`$. Then, the solution for $`\delta b`$, Eq. (60), requires $`\xi =0`$. Finally the solution for $`\delta n^{}`$, Eq. (61), is satisfied provided
$`\dot{T}^{}{\displaystyle \frac{1}{n_0^2}}\ddot{Y}=0,`$
which is a consequence of the condition (70).
One may fix the gauge by imposing, for example,
$$\delta b(y,t)=0\text{and}\delta \varphi (y=0,t)=0.$$
(75)
There is some residual gauge freedom: there are further transformations with $`Y=0`$ and $`T=T(t)`$. These are uninteresting time reparameterizations. Were we to impose a gauge condition on any of the fields that have a discontinuous derivative on the brane, we could only require that it vansihes for either $`y>0`$ or $`y<0`$, but not both, since the reparametrization functions are smooth.
We are ready to present our solution to the linearized field equations. The solutions to the bulk equations express $`\delta \varphi `$, $`\delta b`$ and $`\delta n^{}`$ in terms of $`\delta a`$ and $`F`$. Once $`F`$ is determined in terms of $`\delta a`$ our task is to determine $`\delta a`$ only. We will proceed as follows. First we use the gauge conditions and our solutions in the bulk, Eqs. (59)–(61), to eliminate the function $`F`$ everywhere, and to relate $`\delta a^{}`$ to $`\delta a`$ on the brane. Then we consider the jump equations which impose further restrictions on $`\delta a`$ and its derivatives on the brane. These constraints on the brane are used, finally, to determine $`\delta a`$ in the bulk. The cosmlogy depends only on the fields on the brane. Therefore, we concentrate on determining as fully as possible the fields on the brane.
We eliminate the function $`F`$ by fixing the gauge as in Eq. (75) and using our solution of the field equations in the bulk, Eq. (60), thus
$$F=\frac{9}{4}\frac{\alpha }{\phi }\left[(y\delta a)^{}\xi \right].$$
(76)
We can now fix $`\delta a^{}`$ on the brane in terms of $`\delta a`$. Consider the constraints from requiring continuity of $`\delta \varphi `$ at $`y=0`$. We have in addition the gauge choice in Eq. (75), $`\delta \varphi (y=0,t)=0`$, so we obtain
$`(4\alpha _+1)\delta a`$ $`=`$ $`(y_+\delta a_+^{}+\xi _+),`$ (77)
$`(4\alpha _{}1)\delta a`$ $`=`$ $`(y_{}\delta a_{}^{}+\xi _{}).`$ (78)
Here and below we denote the limiting values of fields as $`y0\pm `$ with a corresponding subscript, eg, $`\delta a_\pm ^{}lim_{y0\pm }\delta a^{}`$.
The first jump equation, Eq. (63), gives a constraint between the constants of integration. From the jump equations for the exact solution, Eqs. (40) and (41), one has
$$\rho _0+p_0=\frac{4\alpha _+1}{y_+}+\frac{4\alpha _{}1}{y_{}}.$$
(79)
Using this and Eqs. (77) and (78) in Eq. (63) we obtain
$$\frac{1}{3}\delta \rho _0=\frac{\xi _+}{y_+}+\frac{\xi _{}}{y_{}}.$$
(80)
The remaining two jump equations, (64) and (65), involve $`\delta a_\pm ^{\prime \prime }`$ in addition to $`\delta a_\pm ^{}`$, $`\delta \ddot{a}`$ and $`\delta a`$. They can be solved to give, on the brane, $`\delta a_\pm ^{\prime \prime }`$ in terms of $`\delta \ddot{a}`$ and $`\delta a`$. However, the resulting expresions are long and not terribly illuminating so we refrain form presenting them here.
The jump equations do not fix the time behavior of $`\delta a`$ on the brane. We now describe a procedure that determines the time behavior and the bulk dependence of $`\delta a`$. Note that since $`F`$ is given in terms of $`\delta a`$ through Eq. (76), one must now impose that $`\delta a`$ satisfy the wave-like equation (62). The solution to Eq. (62) is straightforward,
$$F_\pm (y,t)=𝑑\omega e^{i\omega t}\stackrel{~}{F}_\pm (\omega )\frac{J_0(\frac{\omega }{3\alpha _\pm }|yy_\pm |^{3\alpha _\pm })}{J_0(\frac{\omega }{3\alpha _\pm }|y_\pm |^{3\alpha _\pm })},$$
(81)
where $`J_0`$ is a Bessel function. There is a second solution involving the Neuman function. However, we have dismissed it since it is arbitrarily large as $`|yy_\pm |0`$, which is outside the validity of perturbation theory. The Fourier coefficient $`\stackrel{~}{F}_\pm (\omega )`$ is determined by the function on the brane, $`F(0,t)`$. Using the conditions on derivatives of $`\delta a`$ on the brane one may express this solely in terms of $`\delta a`$:
$$\frac{9}{2}\frac{\alpha _\pm (12\alpha _\pm )}{\phi _\pm }\delta a(0,t)=𝑑\omega e^{i\omega t}\stackrel{~}{F}_\pm (\omega ).$$
(82)
Thus, knowledge of $`\delta a`$ on the brane determines all fields everywhere. The time-dependence of $`\delta a`$, however, is not arbitrary. Consistency of Eq. (81) and the jump equations gives an integro-differential equation for $`\delta a`$ on the brane. To see this, differentiate Eq. (81) and evaluate at $`y=0`$:
$$y_\pm \delta a_\pm ^{\prime \prime }+2\delta a_\pm ^{}=𝑑\omega e^{i\omega t}(4\alpha _\pm 2)\delta \stackrel{~}{a}(\omega )\omega |y_\pm |^{3\alpha _\pm 1}\frac{J_1(\frac{\omega }{3\alpha _\pm }|y_\pm |^{3\alpha _\pm })}{J_0(\frac{\omega }{3\alpha _\pm }|y_\pm |^{3\alpha _\pm })},$$
(83)
where $`(12\alpha _\pm )\delta \stackrel{~}{a}(\omega )=2\phi _\pm /9\alpha _\pm \stackrel{~}{F}_\pm (\omega )`$. Now the left hand side of this equation can be expressed in terms of $`\delta \ddot{a}`$ and $`\delta a`$ through the jump equations. One obtains
$$y_\pm \delta a_\pm ^{\prime \prime }+2\delta a_\pm ^{}=A_\pm \delta \ddot{a}+B_\pm \delta a+C_\pm ,$$
(84)
where $`A_\pm `$ and $`B_\pm `$ are constants given in terms of the parameters of the background solution, and $`C_\pm `$ also depends on these parameters but is, in addition, linear in the small constants $`\xi _\pm `$. We have ommitted the complicated and long expressions for these. Thus, we obtain an equation that determines the time dependence of $`\delta a`$
$$A_\pm \delta \ddot{a}+B_\pm \delta a+C_\pm =𝑑\omega e^{i\omega t}(4\alpha _\pm 2)\delta \stackrel{~}{a}(\omega )\omega |y_\pm |^{3\alpha _\pm 1}\frac{J_1(\frac{\omega }{3\alpha _\pm }|y_\pm |^{3\alpha _\pm })}{J_0(\frac{\omega }{3\alpha _\pm }|y_\pm |^{3\alpha _\pm })}.$$
(85)
A particular solution to this eqution is easily found by Fourier transform:
$$\delta a(0,t)=C_\pm /B_\pm .$$
(86)
To this solution one may add arbitrary linear combinations of solutions to the associated homogeneous equation (obtained by setting $`C_\pm =0`$). These solutions have time dependence of the form of simple exponentials, $`\mathrm{exp}i\omega _0t`$, with characterisitic frequencies $`\omega _0`$ that are solutions to
$$\omega _0^2A_\pm +B_\pm (24\alpha _\pm )\omega _0|y_\pm |^{3\alpha _\pm 1}\frac{J_1(\frac{\omega _0}{3\alpha _\pm }|y_\pm |^{3\alpha _\pm })}{J_0(\frac{\omega _0}{3\alpha _\pm }|y_\pm |^{3\alpha _\pm })}=0.$$
(87)
There is an infinite number of solutions to this equation. The function $`J_0`$ has an infinite number of simple zeros. Between two succesive zeros the last term takes on any value. There may be additional complex solutions when $`A_\pm B_\pm <0`$.
The physical interpretation of our solution is straightforward. The infinite set of oscilatory modes simply correspond to field excitations in the bulk. Neither their amplitudes nor their frequencies depend on the additional density perturbation $`\delta \rho _0`$. This is to be expected because even in the absence of additional matter there can be propagating gravitational and scalar field waves. It is only the particular solution that actually depends on the added matter. Ignoring the possible excitation of propagating modes, we see that the new solution is static.
There are some caveats. This conclusion does not hold if the parameters of the background are such that $`B_\pm =0`$ or if there are complex solutions $`\omega _0`$. In the former case, if in addition $`C_\pm =0`$, there are further solutions of the form
$$\delta at.$$
(88)
This is of particular interest beacuse the perturbations may now be time dependent. In fact, such time dependent perturbations are expected since they must correspond to the linearization of the exact solutions of Ref. for which $`a\mathrm{exp}(f(y)+\sqrt{\overline{\mathrm{\Lambda }}}t)1+f(y)+\sqrt{\overline{\mathrm{\Lambda }}}t`$, where $`\overline{\mathrm{\Lambda }}`$ is an arbitrary constant. We have verified that $`B_\pm =0`$ and $`A_\pm 0`$ for parameters corresponding to the background of Ref. ($`\alpha _\pm =1/4`$, $`\phi _\pm =3/4`$, $`\rho _0=\delta \rho _0=0`$). There are also solutions of the form $`\delta at^2`$. If there are modes of complex $`\omega _0`$, the solutions are not simply oscilatory, but will in addition have exponential dependence. None of these solutions are suitable for observational cosmology.
## V Conclusions
The proposed solution to the cosmological constant problem of Refs. is incomplete in that it does not include matter on the brane, i.e., in our universe. We have found new exact solutions including matter on the brane. These solutions are static and therefore describe an unacceptable cosmology.
However, the equations admit other solutions even under the same assumptions on the symmetry of the metric. This had been recognized in Ref. which found curved brane solutions to the model dubbed case I in Ref. . Here we have studied new solutions obtained as small perturbations about our new static solutions with matter. The small perturbations that are not simply propagating waves are generally static. In special cases the non-propagating, small perturbations grow linearly with time, but we have identified these as the De-Sitter-like solutions of Ref. . On this basis it is tempting to rule out these as viable cosmologies.
In standard FRW cosmology the evolution of the scale factor is completely determined once the equation of state is fixed. Once spherical symmetry is chosen and matter is specified, Einstein’s equations determine the cosmology. However, this is not the case in the peculiar brane cosmologies of Refs. . To understand what is happening one could consider this as an initial value problem. In standard cosmology if the metric and matter content were specified at an initial time (in a fixed gauge), one could evolve forward using the field equations. Thus one would recover the standard picture. Clearly this is not the case of the brane models. What else must be specified and why? This is an interesting question that we hope to explore further. Our guess is that the naked singularities introduce additional information that has been implicitly specified. In the absence of a new general principle that specifies these additional data one would have to give up the notion of causality (at least on a global scale). This may be the price one must pay in order to solve the cosmological constant problem.
Acknowledgments We would like to thank Csaba Csáki and Ira Rothstein for discussions. This work is supported by the Department of Energy under contract No. DOE-FG03-97ER40546.
|
warning/0005/hep-ph0005020.html
|
ar5iv
|
text
|
# I Introduction
## I Introduction
The potential of static heavy quarks illuminates the most important features of QCD dynamics: the asymptotic freedom and confinement. Trying to study subtle electroweak phenomena in the heavy quark sector of Standard Model, we need quite an accurate quantitative understanding of effects caused by the strong interactions. In addition to the perturbative calculations for hard contributions, at present there are three general approaches to get a systematic description of how the heavy quarks are bound into the hadrons and what are the relations between the measured properties of such the hadrons and the characteristics of heavy quarks as relevant to the electroweak interactions and QCD. These approaches are the Operator Product Expansion (OPE) in the inverse powers of heavy quark mass, the Sum Rules (SR) of QCD and the Potential Models (PM) for the systems containing the heavy quarks by exploring various approximations of Bethe-Salpeter equation with the static potential treated in the framework of effective theory with a power counting in terms of powers of the inverse heavy quark mass. The first method is usually exploited in the inclusive estimates, while the second and third techniques are the frameworks of exclusive calculations. The important challenge is a consistency of evaluations obtained in such the ways, that requires the comparative analysis of calculations. A wide variety of systems and processes for the analysis provides a more complete qualitative and quantitative understanding of heavy quark dynamics.
In the leading order of perturbative QCD at short distances and with a linear confining term in the infrared region, the potential of static heavy quarks was considered in the Cornell model , incorporating the simple superposition of both asymptotic limits (the effective coulomb and string-like interactions). The observed heavy quarkonia posed in the intermediate distances, where both terms are important for the determination of mass spectra. So, the phenomenological approximations of potential (logarithmic one and power law ), taking into account the regularities of such the spectra, were quite successful , while the quantities more sensitive to the global properties of potential are the wave functions at the origin as related to the leptonic constants and production rates. So, the potentials consistent with the asymptotic freedom to one and two loops as well as the linear confinement were proposed by Richardson , Buchmüller and Tye , respectively. Technically, using a given scheme of regularization, say, $`\overline{\mathrm{MS}}`$, one has to calculate the perturbative expansion for the potential of static quarks. This potential can be written down as the coulomb one with the running coupling constant in the so-called V scheme. Thus, the perturbative calculations provide us with the matching of $`\overline{\mathrm{MS}}`$ scheme with V-one. The $`n`$ loop running of $`\alpha _s^{\overline{\mathrm{MS}}}`$ requires the $`n1`$ loop matching to $`\alpha _\mathrm{V}`$. Note, that initial two coefficients of corresponding $`\beta `$ functions are scheme and gauge independent, while others generally depend. With the dynamical fields integrated out, the V scheme is defined in terms of the action depending on the static sources (the distance $`r`$), so that its $`\beta `$ function is gauge invariant. The motivation by Buchmüller and Tye was to write down the $`\beta `$ function of $`\alpha _\mathrm{V}`$ consistent with two known asymptotic regimes at short and long distances. They proposed the function, which results in the effective charge determined by two parameters, only: the perturbative parameter is the scale in the running of coupling constant at large virtualities and the nonperturbative parameter is the string tension. The necessary inputs are the coefficients of $`\beta `$ function. Two loop results and the one loop matching condition were available to the moment. Recently, the progress in calculations has provided us with the two loop matching of V and $`\overline{\mathrm{MS}}`$ schemes , that can be combined with the three loop running of $`\alpha _s^{\overline{\mathrm{MS}}}`$. So, the modification of Buchmüller–Tye (BT) potential of static quarks as dictated by the current status of perturbative calculations is of great interest. Moreover, to the moment two peculiar questions become open. First, the asymptotic perturbative expansion of BT $`\beta `$ function to the third order results in the three loop coefficient, which is wrong even in its sign. Second, the elaborated $`\mathrm{\Lambda }_{\overline{\mathrm{MS}}}`$ parameter by BT is in a deep contradiction with the measured value . To clarify the situation, we are tending to derive the static quark potential consistent with the state of the art.
Thus, our motivation is to combine high order multi-loop calculations of the perturbative static potential with the string tension ansatz. So, we improve the perturbative input for the potential model in order to remove the contradiction between the modern high energy data on the QCD coupling constant and the description of heavy quark potential in the framework of one-loop Buchmüller–Tye model, which accepts extremely high value of coupling constant evolved to the $`Z`$ mass scale. In other words, if we accept the current normalization of coupling constant and introduce its value into the Buchmüller–Tye approach to the one-loop potential, then we get the contradiction of such the potential with the heavy quarkonium mass spectra, certainly, since we find about 200 MeV smaller splitting between the $`1S`$ and $`2S`$ levels in comparison with experimental 580 MeV. This discrepancy cannot be removed by the modification of nonperturbative part in the potential with no contradiction with the data on the slope of Regge trajectories. Therefore, the modification of perturbative input for the model of static potential in QCD is meaningful in this sense even when the nonperturbative contribution is conserved in the old string tension form. So, the significant improvement of perturbative $`\beta `$ function for the charge in the coulomb potential is combined with the consequent evolution from high virtualities to low ones with taking into account the influence of nonperturbative term on the evolution, that becomes essential numerically below the scale of 4 GeV.
We have to emphasize that to the moment of paper by Buchmüller and Tye a theory for the heavy nonrelativistic $`Q\overline{Q}`$ pair did not exist. So, the phenomenological derivation of the static potential including perturbative short-distance and non-perturbative long-distance elements made by BT was all one could do. At present, at least for very heavy quarks, such a theory does exist in the form of pNRQCD and vNRQCD , and we address the comparison of static potential model developed in this work with these sound theoretical approaches in QCD to the physics of heavy quarkonium.
Another aspect of this work is devoted to the heavy quark masses. After the potential is given, the heavy quark masses incorporated in the corresponding Schrödinger equation determine the heavy quarkonium spectra with no ambiguity<sup>§</sup><sup>§</sup>§We deal with the so-called spin-averaged spectra, since the consideration of spin-dependent splitting involves some additional parameters beyond the static potential.. These masses involved in the potential model are denoted by $`m_Q^\mathrm{V}`$. Such the mass should be distinguished from the pole mass which is a purely perturbative concept defined unambiguously at each order of perturbation theory through the pole of the perturbative heavy quark propagator. Thus, we need to test the consistency of estimates for the masses in the QCD potential of static quarks and in SR, say.
In Section 2 we generalize the BT approach to three loops and derive the static potential of heavy quarks. The numerical values of potential parameters and their consistency with the relevant quantities are considered. The implications for the heavy quark masses, spectra of heavy quarkonia and leptonic constants are discussed in Section 3. The obtained results are summarized in Conclusion.
## II QCD and potential of static quarks
In this section, first, we discuss two regimes for the QCD forces between the static heavy quarks: the asymptotic freedom and confinement. Second, we formulate how they can be combined in a unified $`\beta `$ function obeyed both limits of small and large QCD couplings.
### A Perturbative results at short distances
The static potential is defined in a manifestly gauge invariant way by means of the vacuum expectation value of a Wilson loop ,
$`V(r)`$ $`=`$ $`\underset{T\mathrm{}}{lim}{\displaystyle \frac{1}{iT}}\mathrm{ln}𝒲_\mathrm{\Gamma },`$ (1)
$`𝒲_\mathrm{\Gamma }`$ $`=`$ $`\stackrel{~}{\mathrm{tr}}𝒫\mathrm{exp}\left(ig{\displaystyle _\mathrm{\Gamma }}𝑑x_\mu A^\mu \right).`$ (2)
Here, $`\mathrm{\Gamma }`$ is taken as a rectangular loop with time extension $`T`$ and spatial extension $`r`$. The gauge fields $`A_\mu `$ are path-ordered along the loop, while the color trace is normalized according to $`\stackrel{~}{\mathrm{tr}}(..)=\mathrm{tr}(..)/\mathrm{tr1}1`$.
Generally, one introduces the V scheme of QCD coupling constant by the definition of QCD potential of static quarks in momentum space as follows:
$$V(𝐪^2)=C_F\frac{4\pi \alpha _\mathrm{V}(𝐪^2)}{𝐪^2},$$
(3)
while $`\alpha _\mathrm{V}`$ can be matched with $`\alpha _{\overline{\mathrm{MS}}}`$
$`\alpha _\mathrm{V}(𝐪^2)`$ $`=`$ $`\alpha _{\overline{\mathrm{MS}}}(\mu ^2){\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\stackrel{~}{a}_n(\mu ^2/𝐪^\mathrm{𝟐})\left({\displaystyle \frac{\alpha _{\overline{\mathrm{MS}}}(\mu ^2)}{4\pi }}\right)^n`$ (4)
$`=`$ $`\alpha _{\overline{\mathrm{MS}}}(𝐪^2){\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}a_n\left({\displaystyle \frac{\alpha _{\overline{\mathrm{MS}}}(𝐪^2)}{4\pi }}\right)^n.`$ (5)
At present, our knowledge of this expansionOn a possible peculiar behaviour in the expansion see ref.. is restricted by
$$a_0=\stackrel{~}{a}_0=1,a_1=\frac{31}{9}C_A\frac{20}{9}T_Fn_f,\stackrel{~}{a}_1=a_1+\beta _0\mathrm{ln}\frac{\mu ^2}{𝐪^\mathrm{𝟐}},$$
(6)
which is the well-known one-loop result, and the recent two-loop calculations , which gave
$`a_2`$ $`=`$ $`\left({\displaystyle \frac{4343}{162}}+4\pi ^2{\displaystyle \frac{\pi ^4}{4}}+{\displaystyle \frac{22}{3}}\zeta (3)\right)C_A^2\left({\displaystyle \frac{1798}{81}}+{\displaystyle \frac{56}{3}}\zeta (3)\right)C_AT_Fn_f`$ (8)
$`\left({\displaystyle \frac{55}{3}}16\zeta (3)\right)C_FT_Fn_f+\left({\displaystyle \frac{20}{9}}T_Fn_f\right)^2,`$
$`\stackrel{~}{a}_2`$ $`=`$ $`a_2+\beta _0^2\mathrm{ln}^2{\displaystyle \frac{\mu ^2}{𝐪^\mathrm{𝟐}}}+(\beta _1+2\beta _0a_1)\mathrm{ln}{\displaystyle \frac{\mu ^2}{𝐪^\mathrm{𝟐}}}.`$ (9)
We have used here the ordinary notations for the SU($`N_c`$) gauge group: $`C_A=N_c`$, $`C_F=\frac{N_c^21}{2N_c}`$, $`T_F=\frac{1}{2}`$. The number of active flavors is denoted by $`n_f`$.
After the introduction of $`𝔞=\frac{\alpha }{4\pi }`$, the $`\beta `$ function is actually defined by
$$\frac{d𝔞(\mu ^2)}{d\mathrm{ln}\mu ^2}=\beta (𝔞)=\underset{n=0}{\overset{\mathrm{}}{}}\beta _n𝔞^{n+2}(\mu ^2),$$
(10)
so that $`\beta _{0,1}^\mathrm{V}=\beta _{0,1}^{\overline{\mathrm{MS}}}`$ and
$`\beta _2^\mathrm{V}`$ $`=`$ $`\beta _2^{\overline{\mathrm{MS}}}a_1\beta _1^{\overline{\mathrm{MS}}}+(a_2a_1^2)\beta _0^{\overline{\mathrm{MS}}}`$ (11)
$`=`$ $`\left({\displaystyle \frac{618+242\zeta (3)}{9}}+{\displaystyle \frac{11(16\pi ^2\pi ^4)}{12}}\right)C_A^3`$ (15)
$`\left({\displaystyle \frac{445+704\zeta (3)}{9}}+{\displaystyle \frac{16\pi ^2\pi ^4}{3}}\right)C_A^2T_Fn_f`$
$`+{\displaystyle \frac{2+224\zeta (3)}{9}}C_A(T_Fn_f)^2{\displaystyle \frac{686528\zeta (3)}{9}}C_AC_FT_Fn_f`$
$`+2C_F^2T_Fn_f+{\displaystyle \frac{184192\zeta (3)}{9}}C_F(T_Fn_f)^2.`$
The coefficients of $`\beta `$ function, calculated in the $`\overline{\mathrm{MS}}`$ scheme are given by
$`\beta _0^{\overline{\mathrm{MS}}}`$ $`=`$ $`{\displaystyle \frac{11}{3}}C_A{\displaystyle \frac{4}{3}}T_Fn_f,`$ (16)
$`\beta _1^{\overline{\mathrm{MS}}}`$ $`=`$ $`{\displaystyle \frac{34}{3}}C_A^24C_FT_Fn_f{\displaystyle \frac{20}{3}}C_AT_Fn_f,`$ (17)
$`\beta _2^{\overline{\mathrm{MS}}}`$ $`=`$ $`{\displaystyle \frac{2857}{54}}C_A^3+2C_F^2T_Fn_f{\displaystyle \frac{205}{9}}C_AC_FT_Fn_f{\displaystyle \frac{1415}{27}}C_A^2T_Fn_f`$ (19)
$`+{\displaystyle \frac{44}{9}}C_F(T_Fn_f)^2+{\displaystyle \frac{158}{27}}C_A(T_Fn_f)^2.`$
The Fourier transform results in the position-space potential
$`V(r)`$ $`=`$ $`C_F{\displaystyle \frac{\alpha _{\overline{\mathrm{MS}}}(\mu ^2)}{r}}(1+{\displaystyle \frac{\alpha _{\overline{\mathrm{MS}}}(\mu ^2)}{4\pi }}(2\beta _0\mathrm{ln}(\mu r^{})+a_1)`$ (22)
$`+\left({\displaystyle \frac{\alpha _{\overline{\mathrm{MS}}}(\mu ^2)}{4\pi }}\right)^2(\beta _0^2(4\mathrm{ln}^2(\mu r^{})+{\displaystyle \frac{\pi ^2}{3}})`$
$`+2(\beta _1+2\beta _0a_1)\mathrm{ln}(\mu r^{})+a_2)+\mathrm{})`$
with $`r^{}r\mathrm{exp}(\gamma _E)`$. Defining the new running coupling constant, depending on the distance,
$$V(r)=C_F\frac{\overline{\alpha }_\mathrm{V}(1/r^2)}{r}.$$
(23)
we can calculate its $`\beta `$-function from (22), so that
$$\overline{\beta }_2^\mathrm{V}=\beta _2^\mathrm{V}+\frac{\pi ^2}{3}\beta _0^3,$$
(24)
and the minor coefficients $`\overline{\beta }_{0,1}^\mathrm{V}`$ are equal to the scheme-independent values given above.
To normalize the couplings, we use (5) at $`𝐪^2=m_Z^2`$.
### B Confining term
The nonperturbative behaviour of QCD forces between the static heavy quarks at long distances $`r`$ is usually represented by the linear potential (see discussion in ref.)
$$V^{\mathrm{conf}}(r)=kr,$$
(25)
which corresponds to the square-law limit for the Wilson loop.
We can represent this potential in terms of constant chromoelectric field between the sources posed in the fundamental representation of SU($`N_c`$). So, in the Fock-Schwinger gauge of fixed point
$$x_\mu A^\mu (x)=0,$$
we can represent the gluon field by means of strength tensor
$$A_\mu (x)\frac{1}{2}x^\nu G_{\nu \mu }(0),$$
so that for the static quarks separated by the distance $`𝐫`$
$$\overline{Q}_i(0)G_{m0}^a(0)Q_j(0)=\frac{𝐫_m}{r}ET_{ij}^a,$$
where the heavy quark fields are normalized to unit. Then, the confining potential is written down as
$$V^{\mathrm{conf}}(r)=\frac{1}{2}g_sC_FEr.$$
Supposing, that the same strength of the field is responsible for the forming of gluon condensate, and introducing the colored sources $`n_i`$, which have to be averaged in the vacuum, we can easily find
$$G_{\mu \nu }^2=4G_{m0}^a(0)G_{m0}^a(0)=4C_FE^2\overline{n}n,$$
where we have supposed the relation
$$\overline{n}T^aT^bn=\overline{n}T^an\overline{n}T^bn,$$
(26)
which ensures that the sources conserve the massless of the gluon, and, hence, the gauge invarianceThe mass term generated by the sources should be equal to $`A_\nu ^aA_\nu ^b[\overline{n}T^aT^bn+\overline{n}T^an\overline{n}T^bn]`$, so that the averaging of sources yields zero, if we suppose (26).. Further, it is evident that
$$\overline{n}T^aT^bn=C_F\frac{\delta ^{ab}}{N_c^21}\overline{n}n.$$
Then, we conclude that the relation between the strength $`E`$ and the string tension depends on the normalization of vacuum sources $`n_i`$. We put
$$\overline{n}_in_j=n_l\delta _{ij},$$
where $`n_l`$ denotes the number of light stochastic flavors, which is the free parameter of such the representation. Of course, the value of $`n_l`$ should be finite even in the case of pure gluodynamics with no light quarks in the infrared region. Moreover, the light quark loops could cause the breaking of gluon string, i.e. the strong decays of higher excitations. We assume that $`n_l`$ is basically determined by the gluon dynamics (i.e. the number of colors), and it slightly correlates with the number of quark flavors. After a simple consideration of potential strength between two colored sources in the fundamental and adjoint representations, i.e. the color factors in front of single gluon coulomb potential, we assume that in the pure gluodynamics the number of stochastic sources substituting for the vacuum gluons can be accepted in the form<sup>\**</sup><sup>\**</sup>\**This assumption corresponds to the definition of vacuum properties in QCD in terms of notations under the consideration, which is in agreement with the value of gluon condensate and Regge trajectories slope.
$$n_l=\frac{1}{N_c}\frac{C_A}{C_F}=\frac{3}{4}=\frac{1}{4}\stackrel{~}{n}_l,$$
where the factor $`1/N_c`$ normalizes the source to unit, and $`C_A/C_F`$ is the appropriate ratio of color charges. To the moment, the shift of $`n_l`$ in QCD with light quarks is not explicitly fixed, while the lattice calculations shown that the dependence of string tension on the number of light quarks is weak . Finally, we find for the linear term of the potential
$$k=\frac{\pi }{\sqrt{C_FN_c\stackrel{~}{n}_l}}C_F\sqrt{\frac{\alpha _s}{\pi }G_{\mu \nu }^2}=\frac{\pi }{2\sqrt{N_c}}C_F\sqrt{\frac{\alpha _s}{\pi }G_{\mu \nu }^2}.$$
(27)
The $`k`$ term is usually represented through a parameter $`\alpha _P^{}`$ as
$$k=\frac{1}{2\pi \alpha _P^{}}.$$
Buchmüller and Tye put $`\alpha _P^{}=1.04`$ GeV<sup>-2</sup>, which we use throughout of this paper. This value of tension, that is related with a slope of Regge trajectories, can be compared with the estimate following from (27). At $`\frac{\alpha _s}{\pi }G_{\mu \nu }^2=(1.6\pm 0.1)10^2`$ GeV<sup>4</sup> we have found
$$\alpha _P^{}=1.04\pm 0.03\mathrm{GeV}^2,$$
which is in a good agreement with the fixed value<sup>††</sup><sup>††</sup>††The ambiguity in the choice of $`n_l`$ can change the appropriate value of gluon condensate..
The form of (25) corresponds to the limit, when at low virtualities $`𝐪^20`$ the coupling $`\alpha _\mathrm{V}`$ tends to
$$\alpha _\mathrm{V}(𝐪^\mathrm{𝟐})\frac{K}{𝐪^\mathrm{𝟐}},$$
so that
$$\frac{d\alpha _\mathrm{V}(𝐪^\mathrm{𝟐})}{\mathrm{ln}𝐪^\mathrm{𝟐}}\alpha _\mathrm{V}(𝐪^\mathrm{𝟐}),$$
(28)
which gives the confinement asymptotics for the $`\beta _\mathrm{V}`$ function.
A special comment should be done on the role of linear term in the potential. Considering the power corrections, which can be attributed to various sources such as the renormalon, topological effects caused by monopoles and vorteces, deviations from the operator product expansion, the authors of argued that this term responsible for the quark confinement can contribute at short distances, too. This conclusion is essentially different from the point of view based on the notion about a low energy phase transition leading to the condensation of gluons and quarks. This condensation provides the forming of chromoelectric string between the static quarks. Thus, at short distances (or high virtualities $`q^2`$) one could expect the decomposition of condensates, that means the scale of confinement (or the string tension) should disappear from the physical quantities at large $`q^2`$. In contrast, the nonperturbative scale can contribute as the factor in front of power corrections $`1/q^2`$ even at $`q^2\mathrm{}`$. So, in several indications of linear term contribution at small distances were considered. We repeat the items relevant to the question on the static potential here.
First, the lattice simulation does not show any change in the slope of the full $`Q\overline{Q}`$ potential as the distances are changed from the largest to the smallest ones where the coulombic part becomes dominant. Hence, no rapid energy jump, characteristic for the phase transition, is found on the lattice. An explicit subtraction of the perturbative corrections at small distances from the potential in the lattice gluodynamics was performed in . This procedure gives an essential nonzero linear term at very small distances.
Second, there are the lattice measurements of the fine splitting in the heavy quarkonium levels as a function of the heavy quark mass. The approach by Voloshin and Leutwyler predicts a particular pattern of such the dependence. Indeed, the multipole expansion of heavy quarkonium interaction with the external gluon field leads to the dominant contribution by the second order of chromoelectric dipole. Therefore, the quark distance squared appears as the leading term in the perturbation due to soft gluons at short distances. These predictions are very different from the evaluations based on the static quark potential with the linear term, such as the potential by Buchmüller and Tye . The numerical results from the lattice simulations favor the linear correction to the potential at short distances.
Third, an interesting manifestation of short strings might be the power corrections to current correlation functions $`\mathrm{\Pi }_j(q^2)`$. Calculations of a relevant coefficient in front of the $`1/q^2`$ terms involve the model assumptions. So, in it was suggested to simulate this power correction by a tachyonic gluon mass. The tachyonic mass can imitate the stringy piece in the potential at short distances. Rather unexpectedly, the use of the tachyonic gluon mass ($`m_g^2=0.5\text{ GeV}^2`$) explains well the behavior of $`\mathrm{\Pi }_j(q^2)`$ in various channels. This fact implies again we see the confirmation of short distance linear term in the potential.
Thus, we do not involve any additional assumptions on the possible scale and properties of quark-gluon condensate decomposition at short distances in the description of static potential in QCD.
### C Unified $`\beta `$ function and potential
Buchmüller and Tye supposed the following procedure for the reconstruction of $`\beta `$ function in the whole region of charge variation by the known limits of asymptotic freedom to a given order in $`\alpha _s`$ and confinement regime. So, in the framework of asymptotic perturbative theory (PT) to one loop, the $`\beta _{\mathrm{PT}}`$ is transformed to the Richardson one,
$$\frac{1}{\beta _{\mathrm{PT}}(𝔞)}=\frac{1}{\beta _0𝔞^2}\frac{1}{\beta _{\mathrm{Rich}}(𝔞)}=\frac{1}{\beta _0𝔞^2\left(1\mathrm{exp}\left[\frac{1}{\beta _0𝔞}\right]\right)}.$$
(29)
The Richardson function has the essential peculiarity at $`𝔞0`$, so that the expansion is the asymptotic series in $`𝔞`$. At $`𝔞\mathrm{}`$ the $`\beta `$ function tends to the confinement limit represented in (28).
To the two loop accuracy, following in the same way results in the $`\beta `$ function by Buchmüller–Tye,
$$\frac{1}{\beta _{\mathrm{PT}}(𝔞)}=\frac{1}{\beta _0𝔞^2}+\frac{\beta _1}{\beta _0^2𝔞}\frac{1}{\beta _{\mathrm{BT}}(𝔞)}=\frac{1}{\beta _0𝔞^2\left(1\mathrm{exp}\left[\frac{1}{\beta _0𝔞}\right]\right)}+\frac{\beta _1}{\beta _0^2𝔞}\mathrm{exp}[l𝔞].$$
(30)
The exponential factor in the second term contributes to the next order in $`𝔞`$ at small $`𝔞`$, so that the perturbative limit is restored. However, we can easily find that third coefficient of $`\beta _{\mathrm{BT}}`$ function is equal to
$$\beta _{2,\mathrm{BT}}=\frac{\beta _1}{\beta _0}(\beta _1l\beta _0),$$
and it is negative at the chosen value of $`l=24`$ , which is in contradiction with the recent result , shown in (15).
To incorporate the three loop results on the perturbative $`\beta `$ function, we introduce
$`{\displaystyle \frac{1}{\beta _{\mathrm{PT}}(𝔞)}}`$ $`=`$ $`{\displaystyle \frac{1}{\beta _0𝔞^2}}+{\displaystyle \frac{\beta _1+\left(\beta _2^\mathrm{V}\frac{\beta _1^2}{\beta _0}\right)𝔞}{\beta _0^2𝔞}}`$ (31)
$`{\displaystyle \frac{1}{\beta (𝔞)}}`$ $`=`$ $`{\displaystyle \frac{1}{\beta _0𝔞^2\left(1\mathrm{exp}\left[\frac{1}{\beta _0𝔞}\right]\right)}}+{\displaystyle \frac{\beta _1+\left(\beta _2^\mathrm{V}\frac{\beta _1^2}{\beta _0}\right)𝔞}{\beta _0^2𝔞}}\mathrm{exp}\left[{\displaystyle \frac{l^2𝔞^2}{2}}\right],`$ (32)
where again the exponential factor in the second term contributes to the next order in $`𝔞0`$. In the perturbative limit the usual solution
$`𝔞(\mu ^2)={\displaystyle \frac{1}{\beta _0\mathrm{ln}\frac{\mu ^2}{\mathrm{\Lambda }^2}}}`$ $`\left[1{\displaystyle \frac{\beta _1}{\beta _0^2}}{\displaystyle \frac{1}{\mathrm{ln}\frac{\mu ^2}{\mathrm{\Lambda }^2}}}\mathrm{ln}\mathrm{ln}{\displaystyle \frac{\mu ^2}{\mathrm{\Lambda }^2}}+{\displaystyle \frac{\beta _1^2}{\beta _0^4}}{\displaystyle \frac{1}{\mathrm{ln}^2\frac{\mu ^2}{\mathrm{\Lambda }^2}}}\left(\mathrm{ln}^2\mathrm{ln}{\displaystyle \frac{\mu ^2}{\mathrm{\Lambda }^2}}\mathrm{ln}\mathrm{ln}{\displaystyle \frac{\mu ^2}{\mathrm{\Lambda }^2}}1+{\displaystyle \frac{\beta _2^\mathrm{V}\beta _0}{\beta _1^2}}\right)\right],`$ (33)
is valid. Using the asymptotic limits of (28) and (33), one can get the equations for any $`\beta `$ function, satisfying these boundary conditions, as follows:
$`\mathrm{ln}{\displaystyle \frac{\mu ^2}{\mathrm{\Lambda }^2}}`$ $`=`$ $`{\displaystyle \frac{1}{\beta _0𝔞(\mu ^2)}}+{\displaystyle \frac{\beta _1}{\beta _0^2}}\mathrm{ln}\beta _0𝔞(\mu ^2)+{\displaystyle _0^{𝔞(\mu ^2)}}𝑑x\left[{\displaystyle \frac{1}{\beta _0x^2}}{\displaystyle \frac{\beta _1}{\beta _0^2x}}+{\displaystyle \frac{1}{\beta (x)}}\right],`$ (34)
$`\mathrm{ln}{\displaystyle \frac{K}{\mu ^2}}`$ $`=`$ $`\mathrm{ln}𝔞(\mu ^2)+{\displaystyle _{𝔞(\mu ^2)}^{\mathrm{}}}𝑑x\left[{\displaystyle \frac{1}{x}}+{\displaystyle \frac{1}{\beta (x)}}\right].`$ (35)
In general, at a given $`\beta `$ function, Eqs.(34) and (35) determine the connection between the scale $`\mathrm{\Lambda }`$ and the parameter of linear potential $`K`$,
$$k=2\pi C_FK.$$
Supposing (32) we can easily integrate out (34) to get the implicit solution of charge dependence on the scale
$`\mathrm{ln}{\displaystyle \frac{\mu ^2}{\mathrm{\Lambda }^2}}=`$ $`\mathrm{ln}\left[\mathrm{exp}\left({\displaystyle \frac{1}{\beta _0𝔞(\mu ^2)}}\right)1\right]+{\displaystyle \frac{\beta _1}{\beta _0^2}}\left[\mathrm{ln}{\displaystyle \frac{\beta _0\sqrt{2}}{l}}{\displaystyle \frac{1}{2}}\left(\gamma _E+\mathrm{E}_1\left[{\displaystyle \frac{l^2𝔞^2(\mu ^2)}{2}}\right]\right)\right]+`$ (37)
$`{\displaystyle \frac{\beta _2^\mathrm{V}\beta _0\beta _1^2}{\beta _0^3}}{\displaystyle \frac{\sqrt{\frac{\pi }{2}}}{l}}\mathrm{Erf}\left[{\displaystyle \frac{l𝔞(\mu ^2)}{\sqrt{2}}}\right],`$
where E$`{}_{1}{}^{}[x]=_x^{\mathrm{}}dtt^1\mathrm{exp}[t]`$ is the exponential integral, and Erf$`[x]=\frac{2}{\sqrt{\pi }}_0^x𝑑t\mathrm{exp}[t^2]`$ is the error function.
Eq.(37) can be inverted by the iteration procedure as it was explored in the derivation of (33). So, the approximate solution of (37) has the following form:
$$𝔞(\mu ^2)=\frac{1}{\beta _0\mathrm{ln}\left(1+\eta (\mu ^2)\frac{\mu ^2}{\mathrm{\Lambda }^2}\right)},$$
(38)
where
$`\eta (\mu ^2)`$ $`=`$ $`\left({\displaystyle \frac{l}{\beta _0\sqrt{2}}}\right)^{\frac{\beta _1}{\beta _0^2}}\mathrm{exp}\left[{\displaystyle \frac{\beta _1}{2\beta _0^2}}\left(\gamma _E+\mathrm{E}_1\left[{\displaystyle \frac{l^2𝔞_1^2(\mu ^2)}{2}}\right]\right){\displaystyle \frac{\beta _2^\mathrm{V}\beta _0\beta _1^2}{\beta _0^3}}{\displaystyle \frac{\sqrt{\frac{\pi }{2}}}{l}}\mathrm{Erf}\left[{\displaystyle \frac{l𝔞_1(\mu ^2)}{\sqrt{2}}}\right]\right],`$ (39)
while $`𝔞_1`$ is obtained in two iterations
$`𝔞_1(\mu ^2)`$ $`=`$ $`{\displaystyle \frac{1}{\beta _0\mathrm{ln}\left(1+\eta _1(\mu ^2)\frac{\mu ^2}{\mathrm{\Lambda }^2}\right)}},`$ (40)
$`\eta _1(\mu ^2)`$ $`=`$ $`\left({\displaystyle \frac{l}{\beta _0\sqrt{2}}}\right)^{\frac{\beta _1}{\beta _0^2}}\mathrm{exp}\left[{\displaystyle \frac{\beta _1}{2\beta _0^2}}\left(\gamma _E+\mathrm{E}_1\left[{\displaystyle \frac{l^2𝔞_0^2(\mu ^2)}{2}}\right]\right)\right],`$ (41)
$`𝔞_0(\mu ^2)`$ $`=`$ $`{\displaystyle \frac{1}{\beta _0\mathrm{ln}\left(1+\frac{\mu ^2}{\mathrm{\Lambda }^2}\right)}}.`$ (42)
Taking the limit of $`\mu ^20`$ we find the relation
$$\mathrm{ln}4\pi ^2C_F\alpha _P^{}\mathrm{\Lambda }^2=\mathrm{ln}\beta _0+\frac{\beta _1}{2\beta _0^2}\left(\gamma _E+\frac{l^2}{2\beta _0^2}\right)\frac{\beta _2^\mathrm{V}\beta _0\beta _1^2}{\beta _0^3}\frac{\sqrt{\frac{\pi }{2}}}{l},$$
(43)
which completely fixes the $`\beta `$ function and charge in terms of scale $`\mathrm{\Lambda }`$ and the slope $`\alpha _P^{}`$, since we have expressed the parameter $`l`$ in terms of above quantities.
Remember, that at $`\mu ^2\mathrm{}`$ the perturbative expression (33) becomes valid as the limit of effective charge (38).
To the moment we are ready to discuss the numerical values of parameters.
### D Setting the scales
As we have already mentioned the slope of Regge trajectories, determining the linear part of potential, is fixed as
$$\alpha _P^{}=1.04\mathrm{GeV}^2.$$
We use also the measured value of QCD coupling constant and pose
$$\alpha _s^{\overline{\mathrm{MS}}}(m_Z^2)=0.123,$$
as the basic input of the potential.
At the given choice of normalization value for the QCD coupling constant we get the scale $`\mathrm{\Lambda }_{n_f=5}^{\overline{\mathrm{MS}}}273`$ MeV, which certainly differs from the world average value resulted in the analysis of PDG , where $`\mathrm{\Lambda }_{n_f=5}^{\overline{\mathrm{MS}}}208_{23}^{+25}`$ MeV, that corresponds to the coupling constant $`\alpha _s^{\overline{\mathrm{MS}}}(m_Z^2)=0.1181\pm 0.002`$ . However, this average value including various data is generally determined by the most precise measurements: the data on the hadronic events in the peak of $`Z`$ boson at LEP (the hadronic width), the decays of $`\tau `$ lepton, the data on the deep inelastic scattering (DIS) for leptons off nucleons and the lattice simulations for the systems of heavy quarkonia. In this set of estimates, the high energy measurements at LEP for $`Z`$ and at HERA for the evolution of nucleon structure functions give the average values $`\alpha _s^{\overline{\mathrm{MS}}}(m_Z^2)=0.123\pm 0.004`$ and $`\alpha _s^{\overline{\mathrm{MS}}}(m_Z^2)=0.122\pm 0.004`$, respectively, while the evolution of structure functions at low virtualities, where an ambiguity in the description of nonperturbative effects and contributions of higher twists are essential, as well as the energy-dependent sum rules for the structure functions at low energies significantly displace down the common average value for the coupling constant extracted from the DIS data. Thus, we argue that the methodical uncertainty for such the averaging is underestimated, since the low-energy data have got some little calculated sources of theoretical uncertainties. The analysis of data on the decays of $`\tau `$ lepton resulting in $`\alpha _s^{\overline{\mathrm{MS}}}(m_Z^2)=0.121\pm 0.003`$, is based on the sum rules, where the control of nonperturbative corrections is much better than in DIS, though there are some theoretical problems on the formulation of sum rules in the region of physical states in contrast with the classic variant of sum rules in the deep euclidean region. Finally, the lattice simulations investigate the splitting between the states of heavy quarkonia, i.e. they operate with the low-energy data and rely on the approximation with the zero number of light quarks $`n_f=0`$ or $`n_f=2`$ under the extrapolation to both the real number of $`n_f=3`$ and the region of high virtualities due to the evolution. A high accuracy of such lattice estimates is announced. As we have seen the spectroscopic characteristics for the systems of heavy quarks need an extremely careful interpretation, since the evolution of potential parameters from the region of bound states to the high virtualities is affected by the nonperturbative factors. Thus, we see that the normalization value of QCD coupling constant accepted above agrees with the direct high-energy measurements, while the data obtained at low energies allow the agreement, if we take into account their systematic uncertainties, which are not well estimated.
Note that the decrease of normalization value to $`\alpha _s^{\overline{\mathrm{MS}}}(m_Z^2)=0.120`$, for example, leads to the discrepancy with the data on the splitting of heavy quarkonium masses between the levels of $`1S`$ and $`2S`$ states, which is very sensitive to the normalization of QCD coupling constant, so that instead of $`M(2S)M(1S)580`$ MeV we get the value which is less by about $`100`$ MeV. In this respect, the variation of other dimensional parameter, the Regge trajectory slope, from the accepted value of $`\alpha _P^{}=1.04`$ GeV<sup>-2</sup> to $`\alpha _P^{}=0.87`$ GeV<sup>-2</sup> leads to unessential change in both the splitting and the corresponding value for the scale in the coupling constant evolved to low virtualities.
Then, we evaluate
$$\alpha _\mathrm{V}(m_Z^2)0.1306,$$
and put it as the normalization point for $`𝔞(m_Z^2)=\alpha _\mathrm{V}(m_Z^2)/(4\pi )`$. Further, we find the following values of $`\mathrm{\Lambda }`$ for the effective charge, depending on the number of active flavors:
$`\mathrm{\Lambda }_{n_f=3}`$ $`=`$ $`643.48\mathrm{MeV},l=56,`$ (44)
$`\mathrm{\Lambda }_{n_f=4}`$ $`=`$ $`495.24\mathrm{MeV},l=37.876,`$ (45)
$`\mathrm{\Lambda }_{n_f=5}`$ $`=`$ $`369.99\mathrm{MeV},l=23.8967,`$ (46)
where we set the threshold values for the switching the number of flavors to be equal to $`m_5=4.6`$ GeV and $`m_4=1.5`$ GeV. After such the fixing the momentum space dependence of the charge, we perform the Fourier transform to get
$$V(r)=kr\frac{8C_F}{r}u(r),$$
(47)
with
$$u(r)=_0^{\mathrm{}}\frac{dq}{q}\left(𝔞(q^2)\frac{K}{q^2}\right)\mathrm{sin}(qr),$$
which is calculated numerically at $`r>0.01`$ fm and represented in the MATHEMATICA file in the format of notebook at the site http://www.ihep.su/kiselev/Potential.nb.
Note, that at short distances the potential behaviour is purely perturbative, so that at $`r<0.01`$ fm we put
$$V(r)=C_F\frac{\overline{\alpha }_\mathrm{V}(1/r^2)}{r},$$
(48)
where the running $`\overline{\alpha }_\mathrm{V}(1/r^2)`$ is given by eq.(33) with the appropriate value of $`\overline{\beta }_2^\mathrm{V}`$ at $`n_f=5`$, and with the matching with the potential (47) at $`r_s=0.01`$ fm, where we have found
$$\overline{\alpha }_\mathrm{V}(1/r_s^2)=0.22213,$$
which implies $`\mathrm{\Lambda }_{n_f=5}^{\overline{\mathrm{V}}}=617.42`$ MeV.
Thus, we have completely determined the model for the potential of static heavy quarks in QCD. In Fig. 1 we present it versus the distance between the quarks. As we can see the potential is very close to what was obtained in the Cornell model in the phenomenological manner by fitting the mass spectra of heavy quarkonia.
The visual deviation between the QCD potential derived and the Cornell model at long distances is caused by a numerical difference in the choice of string tension: we adopt the value given by Buchmüller and Tye, while in the Cornell model the tension is slightly greater than that of we have used. A more essential point is the deviation between the potentials at short distances (see Fig. 2), because of clear physical reason, the running of coupling constant in QCD in contrast to the constant effective value in the Cornell model.
To compare, we show the differences between the $`\beta `$ functions (29), (30) and (32) in Fig. 3 at the fixed values of $`l`$ and $`n_f=3`$. Wee see that the asymptotic perturbative expansion of $`\beta `$ at $`𝔞0`$ dominates at $`𝔞<𝔞_0`$, where $`𝔞_00.03`$ corresponding to $`\alpha _{\mathrm{V},0}0.37`$. This value of coupling $`\alpha _{\mathrm{V},0}`$ coincides with the effective coulomb constant used in the Cornell model. At larger values of coupling the contributions related with the confinement regime are essential.
Two comments are to the point. First, the resulting potential is obtained by the perturbative normalization to the measured value of $`\alpha _s^{\overline{\mathrm{MS}}}(m_Z^2)`$ as combined with the three-loop evolution to the lower virtualities. Second, the running of coupling constant is modified (numerically the deviation from the perturbative regime begins at $`\mu <34`$ GeV) to reach the confinement limit at $`\mu 0`$, so that the perturbative connection between the scales $`\mathrm{\Lambda }`$ and $`\mathrm{\Lambda }^{\overline{\mathrm{MS}}}`$ is broken at virtualities under touch by the charmed and bottom quarks, that was the reason for the error in the assignment of $`\mathrm{\Lambda }^{\overline{\mathrm{MS}}}`$ by Buchmüller and Tye.
## III Heavy quark masses and leptonic constants
Considering the characteristics of heavy quark bound states we should emphasize a significant necessity to certainly separate two distinct theoretical problems. The first problem is the calculation of heavy quark potential, where the leading approximation is the static limit of $`m_Q\mathrm{}`$ in the operator product expansion over the powers of inverse heavy quark mass. we have considered this problem in the previous section. The other problem is the calculation of bound state masses. In the heavy quarkonium the kinetic energy of quark motion is comparable with the potential energy. So, the leading approximation for the effective lagrangian in the operator product expansion over the inverse heavy quark mass is the sum of nonrelativistic kinetic term and the static potential, which give the dominant contribution in the Schrödinger equation for the bound states. Corrections are relativistic terms in the kinetic energy and perturbations of the static potential in the form of operators suppressed by the inverse powers of heavy quark mass, as well as nonpotential retardation effects. The magnitude of such the corrections can be restricted numerically, that leads to a systematic uncertainty in the calculations of mass spectra for the heavy quarkonia in the framework of potential approach with the static potential.
### A Masses
The determination of potential provides us with the extraction of heavy quark masses in the static approximation by comparison of heavy quarkonium mass-spectra with the calculated ones. The predicted charmonium and bottomonium masses are presented in Tables<sup>‡‡</sup><sup>‡‡</sup>‡‡We suppose that the $`\psi (3770)`$-state is a mixture of $`3S`$ and $`3D`$ levels with unessential shift of $`3D`$-mass. I and II at the following values of heavy quark masses in the potential approach
$$m_c^\mathrm{V}=1.468\mathrm{GeV},m_b^\mathrm{V}=4.873\mathrm{GeV},$$
(49)
with no taking into account relativistic corrections, which can be sizable for the charmonium (say, $`\mathrm{\Delta }M(\overline{c}c)40`$ MeV). To the moment, the only measured splitting of $`nS`$-levels is that of $`\eta _c`$ and $`J/\psi `$, which allows us to evaluate the so-called spin-averaged mass
$$\overline{M}(1S)=(3M_{J/\psi }+M_{\eta _c})/4.$$
Supposing the simple relation : $`\overline{M}(ns)=M_V(nS)\frac{1}{4n}(M_{J/\psi }M_{\eta _c})`$, we estimate also the expected values for the excited states with an accuracy better than 10 MeV, we think. For the $`P`$-wave levels we explore the masses
$$\overline{M}(P)=M_1+\frac{1}{3}(M_2M_0)+\frac{2}{9}(M_2M_1+2(M_0M_1)),$$
where $`M_J`$ denotes the mass of state with the total spin $`J`$ and the sum of quark spins $`S=1`$, and we have supposed the spin-dependent forces in the form
$$V_{SD}=A(𝐋𝐒)+B(𝐋𝐒)^2\frac{1}{3}B𝐋^2𝐒^2,$$
where the third term corresponds to the third term in the above expression for $`\overline{M}(P)`$
and it results in the $`L`$-dependent shift of levels.
We have supposed also
$$M_\mathrm{{\rm Y}}M_{\eta _b}\frac{\alpha _s(m_b)}{\alpha _s(m_c)}\frac{m_c^2}{m_b^2}\frac{|R_{\overline{b}b}(0)|^2}{|R_{\overline{c}c}(0)|^2}(M_{J/\psi }M_{\eta _c})56\mathrm{MeV}.$$
We have found that the sizes of quarkonia are the same as they were predicted by Buchmüller and Tye, while the masses of states are slightly different since we have used the other prescription for the input values of ground state masses:
$$M_{\overline{c}c}(1S)=3.068\mathrm{GeV},M_{\overline{b}b}(1S)=9.446\mathrm{GeV}.$$
Then, we predict the masses of $`\overline{b}c`$ quarkonium<sup>\**</sup><sup>\**</sup>\**The experimental error in the ground state mass is still large, $`\delta M=\pm 0.39`$ GeV ., as shown in Table III. The calculated values of masses agree with those of estimated in the Buchmüller–Tye and Martin potentials . The wave functions at the origin are related with the production rates of heavy quarkonia. These parameters are close to what was predicted in the BT potential, but slightly smaller because of both the change in the charmed quark mass and the asymptotic behaviour at $`r0`$.
To the moment we have fixed the potential masses of heavy quarks (49) as independent of scale. To compare with the masses evaluated in the framework of QCD sum rules, we note that in the sum rules for the heavy quarkonia one usually explores the NRQCD with the perturbative potential (22) explicitly dependent of the normalization point $`\mu `$ (referred as $`\mu _{\mathrm{soft}}`$ in the SR). We have checked that at short distances and high $`\mu _{\mathrm{soft}}`$ the perturbative potential (22) and that of present paper coincide with each other, while a deviation appears at $`r1/\mu _{\mathrm{soft}}`$. However, at the distances characteristic for the ground states of heavy quarkonia: $`r_{\overline{b}b(1S)}0.22`$ fm and $`r_{\overline{c}c(1S)}0.42`$ fm, the shape of the potential can be approximated by the perturbative term at $`\mu _{\mathrm{soft}}=1.52.0`$ GeV (see Figs. 4 and 5) with the additive shift of energy scale $`\delta V(\mu _{\mathrm{soft}})`$, which is defined by the expression
$$\delta V(\mu _{\mathrm{soft}})=[V(r)V_{\mathrm{pert}}(r;\mu _{\mathrm{soft}})]|{}_{_{r={\scriptscriptstyle \frac{1}{\mu _{\mathrm{soft}}}}\zeta }}{}^{},$$
(50)
where the parameter $`\zeta `$ has been put in the region of $`\zeta =12`$, where the energy shift $`\delta V`$ has got a little variation about 30-40 MeV, which is, on the first hand, a characteristic uncertainty of potential approach, and on the other hand, it points to a similar form of perturbative potential with the calculated model potential in the region of distance variation. The dependence of energy shift is represented in Fig. 6.
So, if we redefine the heavy quark masses<sup>\*†</sup><sup>\*†</sup>\*†This redefinition is the indication of perturbative renormalon (see review in ). Indeed, there are two sources for the deviation $`\delta V`$. The first is the linear confining term in the potential of static quarks. However, it is a small fraction of $`\delta V`$. The second source is the infrared singularity in the perturbative running coupling. One can easily find that subtracting the singular term of the form $`\frac{1}{\mu _{\mathrm{soft}}\mathrm{\Lambda }}`$ from $`\delta V`$ results in a small value slowly depending on $`\mu _{\mathrm{soft}}`$. In the effective theory for the nonrelativistic heavy quarks, the subtraction that connects the pole mass and the threshold mass can be calculated explicitly (see and references therein). by
$$m^{\mathrm{pole}}(\mu )_{b,c}=m_{b,c}^\mathrm{V}+\frac{1}{2}\delta V(\mu ),$$
the solution of Schrödinger equation with the perturbative potential and $`m^{\mathrm{pole}}(\mu )`$ results in the quarkonia masses close to the experimental values. Thus, we have matched the values of potential masses $`m^\mathrm{V}`$ in the QCD potential with the perturbative pole masses standing in the two-loop calculations. We stress that the dependence on the soft scale in both the energy shift $`\delta V(\mu )`$ and the pole mass $`m^{\mathrm{pole}}(\mu )`$ does not reflect a nonzero anomalous dimension, since these quantities are renormalization group invariants. This scale dependence is due to the truncation of perturbative expansion, wherein the coefficients in front of powers of coupling constant can contain the factorial growth (the renormalon), so that even at the scale close to the charmed quark mass the infrared singularity in the running coupling constant of QCD provides the significant custodial scale dependence.
Numerically, we estimate the running masses $`\overline{m}(\overline{m})`$ in the $`\overline{\mathrm{MS}}`$ scheme using the two and three-loop relations with the pole mass derived in and adjusting the scale $`\mu _{\mathrm{soft}}`$ to be equal to $`\overline{m}`$. So, in two loops we get
$$\overline{m}_c(\overline{m}_c)_{2\mathrm{loops}}=1.40\pm 0.09\mathrm{GeV},\overline{m}_b(\overline{m}_b)_{2\mathrm{loops}}=4.20\pm 0.06\mathrm{GeV},$$
while the three-loop approximation , which is consistent with the three-loop evolution of coupling constant, results in slightly smaller masses, especially, for the charmed quark, where the uncertainty of estimate increases because of stronger sensitivity of quantities involved to the scale variation,
$$\overline{m}_c(\overline{m}_c)_{3\mathrm{loops}}=1.17\pm 0.10\mathrm{GeV},\overline{m}_b(\overline{m}_b)_{3\mathrm{loops}}=4.15\pm 0.06\mathrm{GeV},$$
which are in agreement with the various estimates in the sum rules on $`m_b`$ and $`m_c`$ <sup>\*‡</sup><sup>\*‡</sup>\*‡Note, there is the difference between the usually quoted values of $`\overline{m}(\overline{m})`$ and $`\overline{m}(m^{\mathrm{pole}})`$..
The uncertainty of estimates is determined by the deviations in the calculations of heavy quarkonium masses $`\frac{1}{2}\delta M=20`$ MeV (as shown in Tables I and II) and the error in the extraction of $`\delta V`$ mentioned above. The uncertainty in the running mass of charmed quark is slightly larger than in the bottom mass, since, in addition, its value is more sensitive to a small variation of scale, pole mass and energy shift.
Note, that the calculations in the framework of sum rules were performed for the $`b`$-quark mass in both the full QCD and the effective theory of nonrelativistic heavy quarks NRQCD . The mass extraction of Ref. has been carried out in the nonrelativistic effective theory at next-to-leading order (NLO), whereas Refs. carried out NNLO analyses in the same framework. The calculations in the nonrelativistic effective theory are the calculations in the framework of first principles in QCD, where the results of full QCD are determined in a systematic expansion in $`\alpha _s`$ and the velocity. In Ref. the analysis has also been carried out in the nonrelativistic situation, but no systematic expansion in $`\alpha _s`$ and the velocity has been carried out. That the results for the $`\overline{\mathrm{MS}}`$ mass obtained in agree with the other analyses is not understood and requires further examination (see the conclusions of ).
Recently, the charmed quark mass was evaluated from the NRQCD sum rules in , so that the result on the running mass is in a good agreement with the value given above, too. There is a recent sum rule extraction of the bottom $`\overline{\mathrm{MS}}`$ mass, where the charmed quark mass effects are also included. The estimate of potential approach under consideration is in a good agreement with this recent SR result.
In the dependence of ‘pole’ mass on the scale $`\mu _{\mathrm{soft}}`$ was explicitly calculated in the N<sup>2</sup>LO. The uncertainty of mass extraction from the sum rules for bottomonium was given by $`0.1`$ GeV for the running $`\overline{\mathrm{MS}}`$ mass and $`0.06`$ GeV for the low-energy running mass (‘kinetic’ mass). The result on the $`b`$-quark pole mass depends on both the scale of calculations and the order in $`\alpha _s`$ of perturbative QCD. To compare the results in the sum rules with those of given in the present paper we fix the order in $`\alpha _s`$ by the two-loop corrections. Then we have found that, say, at $`\mu _{\mathrm{soft}}=2.5`$ GeV the results of estimates in the perturbative potential approach and in the framework of sum rules are the same within the uncertainty mentioned. So, putting the above value as the matching point we show the sum rule results in the form of energy shift in Fig. 6. For the sake of representability in Fig. 6 we show the $`\mu `$-dependent ‘pole’ mass extracted in with the uncertainty of $`\delta m=80`$ MeV, which is characteristic inherent error for the short-distance masses in the analysis of . Despite of various choice for the normalization of QCD coupling constant (in $`\alpha _s^{\overline{\mathrm{MS}}}(m_Z^2)=0.118`$), we see a good agreement between the $`\mu `$-dependencies of both the energy shift in the perturbative potential with respect to the static potential of QCD and the variation of perturbative ‘pole’ mass of $`b`$-quark in the sum rules of QCD. As for the one-loop matching of perturbative potential, we mention only that the corresponding sum rules in the NLO give the value of energy shift close to zero at $`\mu _{\mathrm{soft}}>2`$ GeV within the uncertainty of the method, and this estimate is consistent with the result of potential approach as shown in Fig. 6. Thus, the energy shift of perturbative potential with the two-loop matching of $`\mathrm{V}`$ and $`\overline{\mathrm{MS}}`$ schemes indicates the form of QCD potential in agreement with the corresponding soft scale dependence of perturbative pole mass in sum rules of QCD for the bottomonium.
To the moment we can compare the obtained $`\mu `$-dependence of ‘pole’ mass with the relation between the running $`\overline{\mathrm{MS}}`$-mass of heavy quark and the pole mass derived in , where we find
$$m^{\mathrm{pole}}=\overline{m}(\mu )\left(1+c_1(\mu )\frac{\alpha _s^{\overline{\mathrm{MS}}}(\mu ^2)}{4\pi }+c_2(\mu )\left(\frac{\alpha _s^{\overline{\mathrm{MS}}}(\mu ^2)}{4\pi }\right)^2\right),$$
(51)
with
$`c_1(\mu )`$ $`=`$ $`C_F(4+3L),`$ (52)
$`c_2(\mu )`$ $`=`$ $`C_FC_A\left({\displaystyle \frac{1111}{24}}8\zeta (2)4I_3(1)+{\displaystyle \frac{185}{6}}L+{\displaystyle \frac{11}{2}}L^2\right)`$ (55)
$`C_FT_Fn_f\left({\displaystyle \frac{71}{6}}+8\zeta (2)+{\displaystyle \frac{26}{3}}L+2L^2\right)`$
$`+C_F^2\left({\displaystyle \frac{121}{8}}+30\zeta (2)+8I_3(1)+{\displaystyle \frac{27}{2}}L+{\displaystyle \frac{9}{2}}L^2\right)12C_FT_F(12\zeta (2)),`$
where $`I_3(1)=\frac{3}{2}\zeta (3)6\zeta (2)\mathrm{ln}2`$, and $`L=2\mathrm{ln}(\mu /m^{\mathrm{pole}})`$. At $`\mu =m^{\mathrm{pole}}`$, the result of is reproduced. We check that the logs in the definitions of $`c_{1,2}`$ can be removed by the expression of running values $`\overline{m}(\mu )`$ and $`\alpha _s^{\overline{\mathrm{MS}}}(\mu )`$ in terms of $`\overline{m}(m^{\mathrm{pole}})`$ and $`\alpha _s^{\overline{\mathrm{MS}}}(m^{\mathrm{pole}})`$ in (51). Nevertheless, we find that the explicit $`\mu `$-dependence in (51) repeats the form of renormalon contribution as we see it in the perturbative potential, where the similar effect takes place because of both the truncation of perturbative series and the infrared pole in the running coupling constant of QCD. Following (51), we show the value of difference $`2(m_b^{\mathrm{pole}}(\mu )m_b)`$ in Fig. 6 at $`\overline{m}(\overline{m})=4.3`$ GeV. We see that, first, the results of QCD sum rules in agree with the values expected from (51), and second, the $`\mu `$-dependent shift of pole mass approximately coincides with the shift of perturbative potential with respect to the static QCD potential free off renormalon ambiguity caused by infrared singularity of perturbative coupling constant at finite energy scale. This fact implies the cancellation of infrared uncertainties. Thus, we can define the unambiguous pole mass by
$$\widehat{m}^{\mathrm{pole}}=m^{\mathrm{pole}}(\mu )\frac{1}{2}\delta V(\mu ),$$
(56)
where we use the pole mass of (51). The basis for the validity of (56) was observed in , where in the context of perturbative bottom mass extractions, the cancellation of the leading renormalon at $`u=1/2`$ of the Borel plane in the total static perturbative energy of a heavy $`Q\overline{Q}`$ pair was shown.
We find that for the bottom quark the defined mass is given by the value of mass extracted from the potential approach
$$\widehat{m}_b^{\mathrm{pole}}m_b^\mathrm{V},$$
with the accuracy about $`80`$ MeV.
### B Heavy quark masses and pNRQCD
In this section we discuss the modern development in the theory of heavy quarkonium $`QQ^{}`$ on the basis of effective theory called pNRQCD , naturally incorporating the potential interactions between the heavy quarks and external ultrasoft fields in QCD, and compare the pNRQCD results with the values of heavy quark masses obtained above in the QCD potential of static quarks.
First, pNRQCD argues that in the heavy quarkonium the nonrelativistic motion of heavy quarks inside the bound state allows us to introduce three actual physical scales: the heavy quark mass $`m`$, the soft scale of heavy quark momentum inside the hadron $`mv`$ and ultrasoft scale of energy $`mv^2`$, which are distinctly separated by a small parameter $`v`$ being the velocity of heavy quark. After the matching with full QCD at a hard scale $`\mu _{\mathrm{hard}}m`$, in NRQCD the hard fields are integrated out, that results in the perturbative Wilson coefficients of OPE in the effective theory, and we deal with the heavy quarks interacting with the gluons at virtualities $`\mu _{\mathrm{fact},\mathrm{soft}}`$ about $`mv`$. In order to consider the heavy quark fields at lower $`\mu `$ up to $`mv^2`$ we should introduce the effective lagrangian of pNRQCD, where the soft fields are integrated out, and we deal with the potential interaction of heavy quarks and the ultrasoft external gluon fields in the framework of multipole expansion. The matching of pNRQCD with NRQCD takes place at a scale $`\mu _{\mathrm{fact}}mv`$. Recently, the effective theory of vNRQCD was formulated in , using the velocity renormalization group to match the vNRQCD operators with the full QCD at a scale about $`m`$ with the single-step evolution to a soft scale, which can be either $`mv`$ or $`mv^2`$. The current status of vNRQCD provides us with the one-loop matching of heavy quark potential to order $`v^2`$, i.e. up to spin-dependent $`1/m^2`$ terms, which are beyond the current consideration. Therefore, we concentrate our discussion on pNRQCD.
The pNRQCD lagrangian has the following form:
$`_{\mathrm{pNRQCD}}=\mathrm{Tr}\{\mathrm{S}^{}(i_0{\displaystyle \frac{𝐏^2}{4m}}{\displaystyle \frac{𝐩^2}{m}}+{\displaystyle \frac{𝐩^4}{4m^3}}V_s(r){\displaystyle \frac{V_s^{(1)}}{m}}{\displaystyle \frac{V_s^{(2)}}{m^2}}+\mathrm{})\mathrm{S}`$ (57)
$`+\mathrm{O}^{}(iD_0{\displaystyle \frac{𝐏^2}{4m}}{\displaystyle \frac{𝐩^2}{m}}+{\displaystyle \frac{𝐩^4}{4m^3}}V_o(r){\displaystyle \frac{V_o^{(1)}}{m}}{\displaystyle \frac{V_o^{(2)}}{m^2}}+\mathrm{})\mathrm{O}\}`$ (58)
$`+gV_A(r)\mathrm{Tr}\left\{\mathrm{O}^{}𝐫𝐄\mathrm{S}+\mathrm{S}^{}𝐫𝐄\mathrm{O}\right\}+g{\displaystyle \frac{V_B(r)}{2}}\mathrm{Tr}\left\{\mathrm{O}^{}𝐫𝐄\mathrm{O}+\mathrm{O}^{}\mathrm{O}𝐫𝐄\right\}`$ (59)
$`{\displaystyle \frac{1}{4}}F_{\mu \nu }^aF^{\mu \nu a},`$ (60)
where $`𝐏`$ is the momentum associated to the centre-of-mass coordinate. In Eq. (58) the $`1/m`$ corrections to $`V_A`$, $`V_B`$ and to pure gluonic operators as well as the higher order terms in the multipole expansion are not displayed. To the leading order the singlet and octet operators $`S`$, $`O`$ are represented by the appropriate products of nonrelativistic heavy quark and antiquark spinors. The matching of $`S`$ and $`O`$ operators with the NRQCD spinors was done in up to three loops for both the potentials and the normalization factors in OPE. In this lagrangian the singlet and octet potentials $`V_s(r)`$ and $`V_o(r)`$ are treated as the corresponding Wilson coefficients in front of bilinear forms in $`S`$ and $`O`$ to the leading order in $`1/m`$. In ref. the authors shown that this definition of static quark potential is consistent with the definition in terms of Wilson loop (1).
The other result of pNRQCD is the cancellation of renormalon ambiguity in the sum of heavy quark pole masses and the potential up to two loops, which is a confirmation of general consideration in QCD, that was first derived in .
A new feature appears by the consideration of three-loop leading log matching of $`\mathrm{V}`$ and $`\overline{\mathrm{MS}}`$ schemes. So, for the distance-dependent running coupling the result reads off
$`\alpha _\mathrm{V}(1/r^2,\mu )`$ $`=`$ $`\alpha _{\overline{\mathrm{MS}}}(1/r^2)\{1+(a_1+2\gamma _E\beta _0){\displaystyle \frac{\alpha _{\overline{\mathrm{MS}}}(1/r^2)}{4\pi }}`$ (63)
$`+\left[\gamma _E\left(4a_1\beta _0+2\beta _1\right)+\left({\displaystyle \frac{\pi ^2}{3}}+4\gamma _E^2\right)\beta _0^2+a_2\right]{\displaystyle \frac{\alpha _{\overline{\mathrm{MS}}}^2(1/r^2)}{16\pi ^2}}`$
$`+{\displaystyle \frac{C_A^3}{12}}{\displaystyle \frac{\alpha _{\overline{\mathrm{MS}}}^3(1/r^2)}{\pi }}\mathrm{ln}r\mu \},`$
where the two-loop contribution was taken from and it is coincides with (22), of course. However, the three-loop term leads to the explicit dependence on the scale in the perturbative pNRQCD calculations, which has to be expected from the general note on the infrared singularity observed by Appelquist, Dine and Muzinich , that was rederived in pNRQCD by supplementing a certain infrared subtraction. This dependence was considered in for two cases, when the scales of confinement $`\mathrm{\Lambda }_{QCD}`$ and binding energy $`mv^2`$ have the arrangements a) $`\mathrm{\Lambda }_{QCD}mv^2`$ or b) $`mv^2\mathrm{\Lambda }_{QCD}`$. If a), the singlet potential of static quarks suffers from the nonperturbative effects, and it can be treated only after introduction of some model dependent terms coming from the ultrasoft gluons, which form the gluon sea in the heavy quarkonium, so that the sea has its excitations, and the characteristic excitation energy of gluelumps should replace the scale $`\mu `$, that results in the scale-independent nonperturbative potential<sup>*</sup><sup>*</sup>*Possible non-potential terms are discussed in .. If b), the potential is purely perturbative. However, calculating the physical quantities such as the masses of bound states, we have to take into account the contributions coming from the perturbative ultrasoft gluons with the virtualities less than $`\mu `$, which can produce a $`\mu `$-dependent shift of energy, that should be cancelled with the $`\mu `$-dependence in the potential (63) and, probably, in the heavy quark masses. In both cases, the perturbative calculations of singlet potentialWe do not concern for the octet potential of static quarks in the present consideration, though some qualitative conclusions could be straightforwardly generalized from the siglet state to the octet one. explicitly indicate the necessity to take into account the gluon degrees of freedom inside the heavy quarkonium. As was noted in , apparently, this feature is characteristic for the nonabelian theory (see the factor of $`C_A`$ in front of log term in (63)).
To our opinion, this dependence of potential on the ultrasoft gluon fields (the infrared singularity in terms of Appelquist, Dine and Muzinich) inside the heavy quarkonium naturally indicates the formation of gluon string between the heavy quarks at long distances. Indeed, expression (45) was derived under the following arrangement of scales: $`rmv`$, $`mv^2<\mu <mv`$. So, if we put
$$\mu =\frac{u}{r}\sigma r,$$
with
$$v<u<1,\mathrm{and}\sigma \frac{u}{m^2v^2},$$
then, perturbatively expanding in the small parameter $`\sigma r`$, we get the linear correction to the potential in pNRQCD, so that
$$\mathrm{\Delta }V_{\mathrm{pNRQCD}}=\mathrm{\Delta }kr,$$
where
$$\mathrm{\Delta }k=\frac{C_FC_A^3}{12\pi }\alpha _{\overline{\mathrm{MS}}}^4\frac{\sigma }{u}\frac{C_A^3}{12\pi }\alpha _{\overline{\mathrm{MS}}}^3\sigma ,$$
so that we have dropped the scale dependence of strong coupling constant, since it is beyond the accuracy under study, and we have substituted the coulomb relation for the quark velocity inside the bound state $`uC_F\alpha _{\overline{\mathrm{MS}}}`$ to the given order. Numerically, for the charmed quarks this perturbative contribution could be of the order of $`\mathrm{\Delta }k0.1`$ GeV<sup>2</sup>. Thus, we can motivate the relation between the nonperturbative string and the three-loop scale dependent term in the pNRQCD potential.
Indifferently of the arrangement for the confinement and binding energy scales, the introduction of such the string should remove the explicit dependence of full potential on the scale. This has been done above by introduction of unified $`\beta `$ function of coupling in the $`\mathrm{V}`$ scheme. This solution of the problem qualitatively agrees with the consideration in pNRQCD, since, first, in the perturbative regime the contribution of log term is negligibly small as we see for the linear confining term of potential at short distances, and, second, at long distances the nonperturbative confining term is essential, where the string tension is the natural physical scale. In the static potential of QCD given above we do not consider possible ‘nontrivial’ excitations with the broken string geometry, where the break point moves on the string with the speed of light. Such the excitations would correspond to hybrid states with the gluelumps. Thus, we find that the QCD potential of static quarks in the form offered in the present paper has no conflicts with the current status of pNRQCD.
However, to our opinion the problem can be more deep. The static potential, introduced by the Wilson loop, is renormalization group invariant, and it does not contain any separation between the potential gluons and the ultrasoft gluons forming the sea, since it gives the total energy of dynamical fields. In contrast, the pNRQCD introduced the singlet potential as the Wilson coefficient in front of four quark operator, so that it intrinsically operates with the separation of potential and sea, as well as the nonrelativistic quarks, which act as sources, so that some gluons with virtualities greater than $`\mu `$ are considered as emitted, while others with virtualities less than $`\mu `$ are included into the origin of sources, and the gluons with virtualities about $`mv`$ mediate the potential interaction. Generally, this separation of heavy quarks, potential gluons and sea gluons in the operator product expansion can involve nonzero anomalous dimensions for the singlet pNRQCD-potential, say. This fact does not contradict with the OPE basis, but it reflects the point that the static potential of Wilson loop generally differs from the pNRQCD-potential. In addition, the ultrasoft gluon sea introduced in pNRQCD in terms of multipole interaction with local external chromoelectric and chromomagnetic fields is not a local object, indeed.
A note should be done on the linear confining term of potential. In a model of infrared behaviour was used, so that at long distances between the heavy quarks the ultrasoft correction was derived in the form of constant energy shift $`\delta V_0`$ and quadratic term $`\sigma _2r^2`$. The corresponding conclusion was drawn to stress that the linear term could appear in a more complicated case of infrared behaviour. We show in the previous sections how this confinement regime can be reached.
Recently, several papers were devoted to the calculations of ground states in the heavy quarkonia in the way, combining the pNRQCD potential with the nonperturbative corrections to the binding energy as they produced by the multipole expansion of QCD in the form of pNRQCD explicitly shown in (58). Ref. does not strictly estimate the gluon condensate effects in the multipole expansion, and it presents purely perturbative results. It follows the perturbative ground state mass technique as a mass definition that leads to the cancellation of the $`u=1/2`$ renormalon that was considered in the approach of upsilon expansion introduced by Hoang et al. in . So, in the perturbative mass of $`B_c`$ meson was calculated on the base of perturbative expansion for the static potential with the leading approximation in the form of coulomb wave functions. As we see above the perturbative potential suffers from the renormalon ambiguity. In order to remove this dependence on the choice of scale $`\mu `$ in the potential, the authors of calculated the masses of $`J/\mathrm{\Psi }`$ and $`\mathrm{{\rm Y}}`$ in the same technique at the same point $`\mu `$ and inverted the problem on the heavy quark masses by equalizing the perturbative masses of ground states in the charmonium and bottomonium to the measured values. This procedure leads to the $`\mu `$-dependent pole masses of heavy quarks as expressed by the series in $`\alpha _s(\mu )`$. We expect that such the procedure could cancel the renormalon with the accuracy about 50 MeV in the mass of hadron. As a results, the perturbative mass of $`B_c`$ has quite a stable value
$$M_{\mathrm{pert}}(B_c)=6326_9^{+29}\mathrm{MeV},$$
(64)
in the range of $`1.2<\mu <2.0`$ GeV, which should be compared with the results in Table III and the range of $`\mu `$ described above in the study of matching the perturbative potential with the full QCD one. The authors of did not present the $`\mu `$-dependent heavy quark masses. Nevertheless, due to the almost coinciding estimates of $`B_c`$ mass in (64) and Table III, we expect that this dependence should be given by the form of $`\delta V(\mu )`$.
In refs. the same technique for the perturbative contribution with the account for both the gluon condensate corrections in the multipole expansion of QCD and a small $`\alpha _s^5\mathrm{log}\alpha _s`$ term, was used to extract the heavy quark masses. The authors determined the ‘pole’ mass, which is scale dependent, indeed, by putting $`\mu =C_F\alpha _sm_Q`$ in the potential. As we understand, they introduced the mass suffered from the renormalon and got
$$m_b=5022\pm 58\mathrm{MeV},$$
which is greater than we determine in the current presentation. The reason is quite evident. It is the energy shift $`\delta V(\mu )`$. The running $`\overline{\mathrm{MS}}`$ mass quoted in is about 260 MeV greater than we find in the same order in $`\alpha _s`$ for the relation between the pole and running masses. The difference becomes unessential by using the three-loop matching of the masses in , however, the same correction will also decrease the value obtained in the spectroscopy with the full QCD potential. Thus, to our opinion the values of heavy quark masses given in should be kept with a large care.
Finally, in the dependence of potential on the finite heavy quark masses was considered. This dependence is due to the smooth variation of number of active flavors in the expressions for the coefficients of perturbative $`\beta `$ function as well as in the matching coefficients of $`\alpha _\mathrm{V}`$. As we have described above we use the step-like change of active flavor number, which infers implicit model-dependence, which is practically unavoidable in the case under study.
As for the lattice simulations in QCD for the relevant problem, the review can be found in ref. . We emphasize only that the lattice potential of static quarks is close to what is given by the Cornell model. A modern review of phenomenological potential models can be found in lectures . The finite mass effects in the nonrelativistic bound states was recently considered at nest-to-leading order in and . A next-to-next-to-leading order analysis of light quark mass effects in the heavy nonrelativistic $`Q\overline{Q}`$ systems was given in . Some applications of pNRQCD to the heavy quarkonia were done in .
### C Leptonic constants
In the static approximation for the heavy quarks the calculation of leptonic constants for the heavy quarkonia with the two-loop accuracy involves the matching of leptonic currents in NRQCD with the currents of full QCD,
$$J_\nu ^{QCD}=\overline{Q}\gamma _\nu Q,𝒥_\nu ^{NRQCD}=\chi ^{}\sigma _\nu ^{}\varphi ,$$
with the relativistic quark fields $`Q`$ and their nonrelativistic two-component limits of antiquark $`\chi `$ and quark $`\varphi `$, $`\sigma _\nu ^{}=\sigma _\nu v_\nu (\sigma v)`$, and $`v`$ is the four-velocity of heavy quarkonium, so that
$$J_\nu ^{QCD}=𝒦(\mu _{\mathrm{hard}};\mu _{\mathrm{fact}})𝒥_\nu ^{NRQCD},$$
where the scale $`\mu _{\mathrm{hard}}`$ determines the normalization point for the matching of NRQCD with full QCD, while $`\mu _{\mathrm{fact}}`$ refers to the point of perturbative calculations in NRQCD. Using the matching of potential for the static quarks in QCD with the two-loop perturbative potential, we argue that the most appropriate choice of scale relevant to the charmonium and bottomonium is
$$\mu _{\mathrm{fact}}=\mu _{\mathrm{soft}}=1.32\mathrm{GeV}.$$
(65)
For the heavy quarkonium composed by quarks of the same flavour the Wilson coefficient $`𝒦`$ is known up to the two-loop accuracy
$$𝒦(\mu _{\mathrm{hard}};\mu _{\mathrm{fact}})=1\frac{8}{3}\frac{\alpha _s^{\overline{\mathrm{MS}}}(\mu _{\mathrm{hard}})}{\pi }+\left(\frac{\alpha _s^{\overline{\mathrm{MS}}}(\mu _{\mathrm{hard}})}{\pi }\right)^2c_2(\mu _{\mathrm{hard}};\mu _{\mathrm{fact}}),$$
(66)
and $`c_2`$ is explicitly given in . The additional problem is the convergency of (66) at the fixed choice of scales. So, putting $`\mu _{\mathrm{hard}}=(12)m_b`$ and (65) we find a good convergency of QCD corrections for the bottomonium and estimate its leptonic constant defined by
$$0|J_\nu ^{QCD}|\mathrm{{\rm Y}},\lambda =ϵ_\nu ^\lambda f_\mathrm{{\rm Y}}M_\mathrm{{\rm Y}},$$
where $`\lambda `$ denotes the polarization of vector state $`ϵ_\nu `$, so that
$$f_\mathrm{{\rm Y}}=685\pm 30\mathrm{MeV},$$
while the experimental value is equal to $`f_\mathrm{{\rm Y}}^{\mathrm{exp}}=690\pm 13\mathrm{MeV}`$ .
As we can see in Fig. 7 the variation of hard scale in broad limits leads to existence of stable point, where the result is slowly sensitive to such the variation. The stability occurs at $`\mu _{\mathrm{soft}}2.6`$ GeV, where the perturbative potential is still close to the potential of static quarks at the distances characteristic for the $`1S`$-level of $`\overline{b}b`$.
The estimate of leptonic constant for the charmonium $`J/\psi `$ is more sensitive to the choice of factorization scale. Indeed, the size of this system, $`r_{\overline{c}c(1S)}0.42`$ fm, makes more strict constraints on $`\mu _{\mathrm{fact}}1.31.5`$ GeV, since at higher scales the perturbative potential significantly deviates from the potential of static quarks in QCD in the region of bound $`\overline{c}c`$ states, while at lower scales the perturbative potential in two loops does not match the QCD potential in all of the form.
Another problem is the energy shift $`\delta V(\mu )=1.01.2`$ GeV, which essentially renormalizes the pole mass of charmed quark: $`m_c^{\mathrm{pole}}=1.9682.068`$ GeV. This shift does not perturb the mass of the ground state, but it is significant for the value of wave function at the origin. So, following the well-adjusted scaling relation for the leptonic constants , we put $`P(\mu )=\kappa \mathrm{\Psi }(0)m_c^{\mathrm{pole}}(\mu )/m_c`$ and use it in the calculations of the leptonic constantSolving the Schrödinger equation with the shifted masses and potential, we check that this mass dependence of wave function is valid with the accuracy better than 6%, so we put $`\kappa =0.95`$.. We get
$$f_{J/\psi }=400\pm 35\mathrm{MeV},$$
to compare with the experimental value $`f_{J/\psi }^{\mathrm{exp}}=409\pm 15\mathrm{MeV}`$.
In Fig. 8 we see that again the stability point can be reached in the variation of $`\mu _{\mathrm{hadr}}`$ at reasonable value of $`\mu _{\mathrm{soft}}1.35`$ GeV. However, the stability takes place in the narrow region of $`\mu _{\mathrm{hadr}}`$ close to the charm quark mass.
At present, the matching condition for the heavy quarkonium composed by the quarks of different flavors, $`\overline{b}c`$, is known to one loop, only . So, for the pseudoscalar state we have
$$𝒦(\mu _{\mathrm{hard}};\mu _{\mathrm{fact}})=1\frac{\alpha _s^{\overline{\mathrm{MS}}}(\mu _{\mathrm{hard}})}{\pi }\left(2\frac{m_bm_c}{m_b+m_c}\mathrm{ln}\frac{m_b}{m_c}\right),$$
(67)
which is independent of the factorization scale. The matching of perturbative potential to the one-loop accuracy with the QCD potential of static quarks at $`r0.30.4`$ fm relevant to the ground state of $`B_c`$ meson , is rather questionable, since the deviation in the forms of potentials is quite sizable. In addition we have to pose $`\mu _{\mathrm{fact}}=\mu _{\mathrm{hard}}`$, because we cannot distinguish these scales, while the nonzero anomalous dimension to two loops is not taken into account. Nevertheless, we can put $`\mu _{\mathrm{hard}}=1.31.8`$ GeV and neglect $`\delta V`$, which is beyond the actual control in the one-loop accuracy. Indeed, as we see in Fig. 6 the one loop value of energy shift for the matching of perturbative and QCD potentials is quite small at the large virtualities about 2 GeV, and it can be neglected, while at smaller virtualities the form of perturbative potential is close to that of given by QCD only in the short range of distances $`r=0.10.25`$ fm, hence, the results on the matching are not reliable for the extracting the heavy quark masses from the parameters of bound states. So, we estimate
$$f_{B_c}=400\pm 45\mathrm{MeV},$$
to compare with the estimates in the SR, where $`f_{B_c}^{\mathrm{SR}}=400\pm 25\mathrm{MeV}`$ .
Finally, we present the ratios of leptonic constants for the excited $`nS`$-levels of $`\overline{b}b`$ and $`\overline{c}c`$ in Table IV in comparison with the experimental data. We see that the predictions are in a good agreement with the measured values. For completeness, we predict also the constant of $`2S`$-level in the $`\overline{b}c`$ system
$$f_{B_c(2S)}=280\pm 50\mathrm{MeV},$$
which agrees with the scaling relation .
Thus, we have analyzed the estimates following from the potential of static quarks in QCD for the masses of quarks and heavy quarkonia as well as for the leptonic constants, and found both the good agreement with the experimental data available and the consistency with the QCD sum rules.
## IV Conclusion
We have derived the potential of static heavy quarks in QCD on the base of known limits at short and long distances: the asymptotic freedom to the three loop accuracy and the confinement regime. The inputs of potential are the coefficients of perturbative $`\beta `$ function, the matching of $`\overline{\mathrm{MS}}`$ scheme with the V scheme of potential, the normalization of running coupling constant of QCD at $`\mu ^2=m_Z^2`$ and the slope of Regge trajectories, determining the linear term in the potential. Thus, the approach by Buchmüller and Tye has been modified in accordance with the current status of perturbative calculations.
In the static limit the two-loop improvement of coulomb potential results in the significant correction to the $`\beta `$ function for the effective charge, $`\mathrm{\Delta }\beta /\beta 10\%`$ as shown in Fig. 3. This correction is important for the determination of critical values of charge, i.e. the value in the intermediate region between the perturbative and nonperturbative regimes. Moreover, the two-loop matching condition and the three-loop running of coupling constant normalized by the data at the high energy of $`m_Z`$ determine the region of energetic scale for changing the regimes mentioned above. This scale strongly correlates with the data on the mass spectra of heavy quarkonia. So, it is connected with the splitting of masses between the $`1S`$ and $`2S`$ levels. We stress that the consistent consideration of two-loop improvement gives the appropriate value of effective coulomb coupling constant as it was fitted in the Cornell model of potential. This is achieved in the present paper in contrast to the one-loop consideration by Buchmüller and Tye, who found the value of $`\mathrm{\Lambda }_{QCD}`$ inconsistent with the current normalization at high energies. So, the two-loop improvement gives the correct normalization of effective coulomb exchange at the distances characteristic for the average separation between the heavy quarks inside the heavy quarkonium and determines the deviations at short distances $`r<0.08`$ fm (see Fig. 2), that is important in the calculations of leptonic constants related with the wave functions at the origin.
Other corrections to the potential of heavy quarks are connected with the finite mass effects and cannot be treated in the framework of static approximation. For example, the spin-dependent forces, relativistic corrections and specific non-abelian potential terms<sup>§</sup><sup>§</sup>§They have the form of $`\alpha _s^2/r^2`$ with the factor given by the inverse heavy quark masses. in the heavy quarkonium should be taken in the analysis of mass spectra. A magnitude of leading nonstatic corrections can be evaluated by the characteristic shifts of levels due to the hyper-fine splitting of $`S`$-wave levels in the heavy quarkoniaThe splitting is about 100 MeV or less.. So, we conservatively evaluate the uncertainty of heavy quark mass analysis $`\delta m80`$ MeV.
Thus, the non-abelian term of potential $`\alpha _s^2/r^2`$, say, has the factors in the form of $`1/m_Q`$, and it is equal to zero in the static limit $`m_Q\mathrm{}`$, while the uncertainty in the heavy quark masses due to the omission of such the terms is estimated in the paragraph above. Formally, if we consider the perturbation theory for the calculation of bound state levels in the heavy quarkonium with the coulomb functions taken as the leading approximation, which is not a scope of our consideration, then the mentioned non-abelian potential contributes in the same order in $`\alpha _s`$ as the two-loop corrections to the matching of perturbative static potential $`\alpha _s^4`$, since the averaging of $`1/r^2`$ results in $`\alpha _s^2m_Q^2`$ factor. However, the two-loop effects are important for the consistent consideration of static potential and the high energy normalization, i.e. these corrections are significant in the running of effective charge in the potential from the high energies to the scale relevant to the heavy quark bound states even in the static limit, while the nonstatic contributions can be consistently neglected in the numerical analysis. We see that our consideration is consistent in the static approximation, which we have addressed in the present paper.
The matching of two loop perturbative potential with the QCD potential of static quarks has been performed to get estimates of heavy quark masses, which can be compared with the results of QCD sum rules. A good agreement between two approaches has been found.
The recent determinations of heavy quark masses in Refs. were done in the framework of QCD sum rules, which is a systematic approach, indeed. It is based on the separation of short-distance region from the nonperturbative effects at some values of parameters defining the scheme of calculations in the sum rules. In this approach the nonperturbative terms are given in the form of quark-gluon condensates contributing with corresponding short-distance Wilson coefficients, so that as was shown in , a numerical contribution of gluon condensate term in the sum rules is negligibly small in comparison with the perturbative part. However, it would be incorrect to think that these explicit contributions suppressed in some region of parameters are the only terms caused by the nonperturbative infrared dynamics of QCD. Indeed, neglecting the condensate terms, we find that the perturbative correlators suffer from the renormalon ambiguity, which implies that the perturbative expansion in series of $`\alpha _s`$ is asymptotic, and the summation of series depends on a method used. The physical reason for such the divergency and ambiguity is the infrared singularity in the QCD coupling constant. This singularity is regularized by introducing the threshold mass parameters free of renormalon. Such the approach is independent of any assumptions on the gluon condensate, since, generally the pole mass renormalon and the gluon condensates are different issues.
So, the perturbative pole mass used in the QCD sum rules is not well defined quantity, and some relevant quantities are introduced in Refs. . These quantities are constructed from the perturbative pole mass of heavy quark with specific infrared subtractions, which are treated independently of the quark-gluon condensates. These constructions are author-dependent. Though the authors of subtracted masses gave some physical motivations, which are more or less strict, but justified. These infrared subtractions imply the introduction of infrared regulators.
In the present paper the unified $`\beta `$ function for the effective charge in the potential is considered, and its definition supposes the infrared stability. Thus, we see that the analysis of heavy quark masses in both the QCD sum rules and potential approach involves the consideration of relevant effects caused by the infrared dynamics of QCD, though the explicit constructive procedures are certainly different, but they have similar inherent uncertainties.
The calculated mass spectra of heavy quarkonia and the leptonic constants of vector $`nS`$-levels are in agreement with the measured values. The characteristics of $`B_c`$ meson have been predicted.
The authors are grateful to prof. A.K.Likhoded and A.L.Kataev for stimulating discussions and A.A.Pivovarov for clarifying the results of QCD sum rules on the $`b`$-quark mass. We thank Dr. Antonio Vairo for valuable remarks, references and explanations concerning the approach of pNRQCD, and for discussions.
This work is in part supported by the Russian Foundation for Basic Research, grants 01-02-99315, 01-02-16585 and 00-15-96645. The work of A.I.Onishchenko was supported, in part, by International Center of Fundamental Physics in Moscow, International Science Foundation, and INTAS-RFBR-95I1300 grants.
|
warning/0005/cond-mat0005026.html
|
ar5iv
|
text
|
# A Rigorous Derivation of the Gross-Pitaevskii Energy Functional for a Two-dimensional Bose Gas
## 1 Introduction
Motivated by recent experimental realizations of Bose-Einstein condensation the theory of dilute, inhomogeneous Bose gases is currently a subject of intensive studies. Most of this work is based on the assumption that the ground state properties are well described by the Gross-Pitaevskii (GP) energy functional (see the review article ). A rigorous derivation of this functional from the basic many-body Hamiltonian in an appropriate limit is not a simple matter, however, and has only been achieved recently for bosons with a short range, repulsive interaction in three spatial dimensions .
The present paper is concerned with the justification of the GP functional in two spatial dimensions. Several new issues arise. One is the form of the nonlinear interaction term in the energy functional for the GP wave function $`\mathrm{\Phi }`$. In three dimensions this term is $`4\pi a|\mathrm{\Phi }|^4`$, where $`a`$ is the scattering length of the interaction potential. The rationale is the well known formula for the energy density of a homogeneous Bose gas, which, for dilute gases with particle density $`\rho `$, is $`4\pi a\rho ^2`$. This fact has been ‘known’ since the early 50’s but a rigorous proof is fairly recent . In two dimensions the corresponding formula is $`4\pi \rho ^2|\mathrm{ln}(\rho a^2)|^1`$ as proved in by extension of the method of . The formula was first stated by Schick ; other early references to this formula are . It would seem natural to consider $`4\pi |\mathrm{\Phi }|^4|\mathrm{ln}(|\mathrm{\Phi }|^2a^2)|^1`$ as the interaction term in the GP functional, and this has indeed been suggested in . Such a term, however, is unnecessarily complicated for the purpose of leading order calculations. In fact, since the logarithm varies only slowly it turns out that one can use the same form as in the three dimensional case, but with an appropriate dimensionless coupling constant $`g`$ replacing the scattering length, and still retain an exact theory (to leading order in $`\rho `$).
It is often assumed that a justification of the GP functional depends on the existence of Bose Einstein condensation. Several remarks can be made about this: 1. We neither assume nor prove the existence of BE condensation, but we do demonstrate a kind of condensation over a distance that is fixed (i.e., non-thermodynamic) but whose length goes to infinity as the density goes to zero; 2. BE condensation does not exist in two dimensions when the temperature is positive, but it can, and most likely does, exist in the ground state; 3. In any event, when the density is low and the temperature is zero it appears to be likely that the system can be described for many purposes in terms of only a few macroscopic order parameters such as the density and phase – at least this is true for the dependence of the ground state energy and density upon an external potential.
The functional we shall consider is
$$^{\mathrm{GP}}[\mathrm{\Phi }]=\left(|\mathrm{\Phi }(𝐱)|^2+V(𝐱)|\mathrm{\Phi }(𝐱)|^2+4\pi g|\mathrm{\Phi }(𝐱)|^4\right)\mathrm{d}^2𝐱,$$
(1.1)
where $`V`$ is the external confining potential and all integrals are over $`^2`$.
The choice of $`g`$ is an issue on which there has not been unanimous opinion in the recent papers on this subject. We shall prove that a right choice is $`g=|\mathrm{ln}(\overline{\rho }a^2)|^1`$ where $`\overline{\rho }`$ is a mean density that will be defined more precisely below. This mean density depends on the particle number $`N`$, which implies that the scaling properties of the GP functional are quite different in two and three dimensions. In the three-dimensional case the natural parameter is $`Na/a_{\mathrm{osc}}`$, with $`a_{\mathrm{osc}}`$ being the length scale defined by the external confining potential. If $`a/a_{\mathrm{osc}}`$ is scaled like $`1/N`$ as $`N\mathrm{}`$ this parameter is fixed and the gradient term $`|\mathrm{\Phi }|^2`$ in the GP functional is of the same order as the other terms. In two dimensions the corresponding parameter is $`N|\mathrm{ln}(\overline{\rho }a^2)|^1`$. For a quadratic external potential $`\overline{\rho }`$ behaves like $`N^{1/2}/a_{\mathrm{osc}}^2`$ and hence the parameter can only be kept fixed if $`a/a_{\mathrm{osc}}`$ decreases exponentially with $`N`$. A slower decrease means that the parameter tends to infinity. This corresponds to the so-called Thomas Fermi (TF) limit where the gradient term has been dropped altogether and the functional is
$$^{\mathrm{TF}}[\rho ]=\left(V(𝐱)\rho (𝐱)+4\pi g\rho (𝐱)^2\right)\mathrm{d}^2𝐱,$$
(1.2)
defined for nonnegative functions $`\rho `$. Our main result, stated in Theorems 1.3 and 1.4 below, is that minimization of (1.2) reproduces correctly the ground state energy and density of the many-body Hamiltonian in the limit when $`N\mathrm{}`$, $`\overline{\rho }a^20`$, but $`N|\mathrm{ln}(\overline{\rho }a^2)|^1\mathrm{}`$. Only in the exceptional situation that $`N|\mathrm{ln}(\overline{\rho }a^2)|^1`$ stays bounded is there need for the full GP functional (1.1), cf. Theorems 1.1 and 1.2.
We shall now describe the setting more precisely. The starting point is the Hamiltonian for $`N`$ identical bosons in an external potential $`V`$ and with pair interaction $`v`$,
$$H^{(N)}=\underset{i=1}{\overset{N}{}}\left(_i^2+V(𝐱_i)\right)+\underset{i<j}{}v(𝐱_i𝐱_j),$$
(1.3)
acting on the totally symmetric wave functions in $`^NL^2(^2)`$. Units have been chosen so that $`\mathrm{}=2m=1`$, where $`m`$ is the particle mass. We assume that $`v`$ is nonnegative and spherically symmetric with a finite scattering length $`a`$. (For the definition of scattering length in two dimensions see the appendix.) The external potential should be continuous and tend to $`\mathrm{}`$ as $`|𝐱|\mathrm{}`$. It is then possible and convenient to shift the energy scale so that $`\mathrm{min}_𝐱V(𝐱)=0`$. For the TF limit theorem we shall require some additional properties of $`V`$ to be specified later.
The ground state energy $`\omega `$ of the one-particle operator $`^2+V`$ is a natural energy unit and gives rise to the length unit $`a_{\mathrm{osc}}\omega ^{1/2}`$. In the sequel we shall be considering a limit where $`a/a_{\mathrm{osc}}`$ tends to zero while $`N\mathrm{}`$. Experimentally $`a/a_{\mathrm{osc}}`$ can be changed in two ways: One can either vary $`a_{\mathrm{osc}}`$ or $`a`$. The first alternative is usually simpler in practice but very recently a direct tuning of the scattering length itself has also been shown to be feasible . Mathematically, both alternatives are equivalent, of course. The first corresponds to writing $`V(𝐱)=a_{\mathrm{osc}}^2\widehat{V}(𝐱/a_{\mathrm{osc}})`$ and keeping $`\widehat{V}`$ and $`v`$ fixed. The second corresponds to writing the interaction potential as $`v(𝐱)=a^2\widehat{v}(𝐱/a)`$, where $`\widehat{v}`$ has unit scattering length, and keeping $`V`$ and $`\widehat{v}`$ fixed. This is equivalent to the first, since for given $`\widehat{V}`$ and $`\widehat{v}`$ the ground state energy of (1.3), measured in units of $`\omega `$, depends only on $`N`$ and $`a/a_{\mathrm{osc}}`$. In the dilute limit when $`a`$ is much smaller than the mean particle distance, the energy becomes independent of $`\widehat{v}`$.
We shall measure all energies in terms of $`\omega `$ and lengths in terms of $`a_{\mathrm{osc}}`$ and regard $`\widehat{V}`$ and $`\widehat{v}`$ as fixed. The notation $`E^{\mathrm{QM}}(N,a)`$ for the ground state energy of (1.3) is then justified.
The quantum mechanical particle density is defined by
$$\rho _{N,a}^{\mathrm{QM}}(𝐱)=N|\mathrm{\Psi }^{(N)}(𝐱,𝐱_2,\mathrm{},𝐱_N)|^2\mathrm{d}^2𝐱_2\mathrm{}\mathrm{d}^2𝐱_N,$$
(1.4)
where $`\mathrm{\Psi }^{(N)}`$ is a ground state for (1.3).
The GP functional (1.1) has an obvious domain of definition (cf. Eq. (2.1) in ). The infimum of $`^{\mathrm{GP}}[\mathrm{\Phi }]`$ under the condition $`|\mathrm{\Phi }|^2=N`$ will be denoted by $`E^{\mathrm{GP}}(N,g)`$. The infimum is obtained for a unique, positive function, denoted $`\mathrm{\Phi }_{N,g}^{\mathrm{GP}}`$, and the GP density is defined as $`\rho _{N,g}^{\mathrm{GP}}(𝐱)=\mathrm{\Phi }_{N,g}^{\mathrm{GP}}(𝐱)^2`$.
The ground state energy of the TF functional (1.2) with the subsidiary condition $`\rho =N`$ is denoted $`E^{\mathrm{TF}}(N,g)`$. The corresponding minimizer can be written explicitly; it is
$$\rho _{N,g}^{\mathrm{TF}}(𝐱)=\frac{1}{8\pi g}[\mu ^{\mathrm{TF}}V(𝐱)]_+,$$
(1.5)
where $`[t]_+\mathrm{max}\{t,0\}`$ and $`\mu ^{\mathrm{TF}}`$ is chosen so that the normalization condition $`\rho _{N,g}^{\mathrm{TF}}=N`$ holds.
We now define the mean density $`\overline{\rho }`$ as the average of the TF density $`\rho _{N,1}^{\mathrm{TF}}`$ at coupling constant $`g=1`$, weighted with $`N^1\rho _{N,1}^{\mathrm{TF}}`$, i.e.,
$$\overline{\rho }=\frac{1}{N}\rho _{N,1}^{\mathrm{TF}}(𝐱)^2\mathrm{d}^2𝐱.$$
(1.6)
It is clear that $`\overline{\rho }`$ depends on $`N`$ and when we wish to emphasize this we write $`\overline{\rho }_N`$. The definition (1.6) has the advantage that $`\overline{\rho }`$ is easily computed; for instance, if $`V(𝐱)|𝐱|^s`$ for some $`s>0`$, then $`\overline{\rho }_NN^{s/(s+2)}`$. It may appear more natural to define $`\overline{\rho }`$ self-consistently as $`\overline{\rho }=\frac{1}{N}\rho _{N,g}^{\mathrm{TF}}(𝐱)^2\mathrm{d}^2𝐱`$ with $`g=|\mathrm{ln}(\overline{\rho }a^2)|^1`$, which amounts to solving a nonlinear equation for $`\overline{\rho }`$. Also, the TF density could be replaced by the GP density. However, since $`\overline{\rho }`$ will only appear under a logarithm such sophisticated definitions are not needed for the leading order result we are after. The simple formula (1.6) is adequate for our purpose, but it should be kept in mind that the self-consistent definition may be relevant in computations beyond the leading order.
With this notation we can now state the two dimensional analogue of Theorem I.1 in .
###### Theorem 1.1 (GP limit for the energy).
If, for $`N\mathrm{}`$, $`a^2\overline{\rho }_N0`$ with $`N/|\mathrm{ln}(a^2\overline{\rho }_N)|`$ fixed, then
$$\underset{N\mathrm{}}{lim}\frac{E^{\mathrm{QM}}(N,a)}{E^{\mathrm{GP}}(N,1/|\mathrm{ln}(a^2\overline{\rho }_N)|)}=1.$$
(1.7)
The corresponding theorem for the density, c.f. Theorem I.2 in , is
###### Theorem 1.2 (GP limit for the density).
If, for $`N\mathrm{}`$, $`a^2\overline{\rho }_N0`$ with $`\gamma N/|\mathrm{ln}(a^2\overline{\rho }_N)|`$ fixed, then
$$\underset{N\mathrm{}}{lim}\frac{1}{N}\rho _{N,a}^{\mathrm{QM}}(𝐱)=\rho _{1,\gamma }^{\mathrm{GP}}(𝐱)$$
(1.8)
in the sense of weak convergence in $`L^1(^2)`$.
These theorems, however, are not particularly useful in the two dimensional case, because the hypothesis that $`N/|\mathrm{ln}(a^2\overline{\rho }_N)|`$ stays bounded requires an exponential decrease of $`a`$ with $`N`$. As remarked above, the TF limit, where $`N/|\mathrm{ln}(a^2\overline{\rho }_N)|\mathrm{}`$, is much more relevant. Our treatment of this limit requires that $`V`$ is asymptotically homogeneous and sufficiently regular in a sense made precise below. This condition can be relaxed, but it seems adequate for most practical applications and simplifies things considerably.
###### Definition 1.1.
We say that $`V`$ is asymptotically homogeneous of order $`s>0`$ if there is a function $`W`$ with $`W(𝐱)0`$ for $`𝐱\mathrm{𝟎}`$ such that
$$\frac{\lambda ^sV(\lambda 𝐱)W(𝐱)}{1+|W(𝐱)|}0\mathrm{as}\lambda \mathrm{}$$
(1.9)
and the convergence is uniform in $`𝐱`$.
The function $`W`$ is clearly uniquely determined and homogeneous of order $`s`$, i.e., $`W(\lambda 𝐱)=\lambda ^sW(𝐱)`$ for all $`\lambda 0`$.
###### Theorem 1.3 (TF limit for the energy).
Suppose $`V`$ is asymptotically homogeneous of order $`s>0`$ and its scaling limit $`W`$ is locally Hölder continuous, i.e., $`|W(𝐱)W(𝐲)|(\mathrm{const}.)|𝐱𝐲|^\alpha `$ for $`|𝐱|,|𝐲|=1`$ for some fixed $`\alpha >0`$. If, for $`N\mathrm{}`$, $`a^2\overline{\rho }_N0`$ but $`N/|\mathrm{ln}(a^2\overline{\rho }_N)|\mathrm{}`$, then
$$\underset{N\mathrm{}}{lim}\frac{E^{\mathrm{QM}}(N,a)}{E^{\mathrm{TF}}(N,1/|\mathrm{ln}(a^2\overline{\rho }_N)|)}=1.$$
(1.10)
To state the corresponding theorem for the density we need the minimizer of (1.2) with $`g=1`$, $`V`$ replaced by $`W`$, and normalization $`\rho =1`$. We shall denote this minimizer by $`\stackrel{~}{\rho }_{1,1}^{\mathrm{TF}}`$; an explicit formula is
$$\stackrel{~}{\rho }_{1,1}^{\mathrm{TF}}(𝐱)=\frac{1}{8\pi }[\stackrel{~}{\mu }^{\mathrm{TF}}W(𝐱)]_+,$$
(1.11)
where $`\stackrel{~}{\mu }^{\mathrm{TF}}`$ is determined by the normalization condition.
###### Theorem 1.4 (TF limit for the density).
Let $`V`$ satisfy the same hypothesis as in Theorem 1.3. If, for $`N\mathrm{}`$, $`a^2\overline{\rho }_N0`$ but $`\gamma =N/|\mathrm{ln}(a^2\overline{\rho }_N)|\mathrm{}`$, then
$$\underset{N\mathrm{}}{lim}\frac{\gamma ^{2/(s+2)}}{N}\rho _{N,a}^{\mathrm{QM}}(\gamma ^{1/(s+2)}𝐱)=\stackrel{~}{\rho }_{1,1}^{\mathrm{TF}}(𝐱)$$
(1.12)
in the sense of weak convergence in $`L^1(^2)`$.
###### Remark 1.1.
For large $`N`$, $`\overline{\rho }_N`$ behaves like $`(\mathrm{const}.)N^{s/(s+2)}`$. Moreover, prefactors are unimportant in the limit $`N\mathrm{}`$, because $`\overline{\rho }_N`$ stands under a logarithm. Hence Theorems 1.3 and 1.4 could also be stated with $`N^{s/(s+2)}`$ in place of $`\overline{\rho }_N`$.
The proofs of these theorems follow from upper and lower bounds on the ground state energy $`E^{\mathrm{QM}}(N,a)`$ that are derived in Sections 3 and 4. For these bounds some properties of the minimizers of the functionals (1.1) and (1.2), discussed in the following section, are needed.
## 2 GP and TF theory
In this section we consider the functionals (1.1) and (1.2) with an arbitrary positive coupling constant $`g`$. Existence and uniqueness of minimizers is shown in the same way as in Theorem II.1 in . The GP energy $`E^{\mathrm{GP}}(N,g)`$ has the simple scaling property $`E^{\mathrm{GP}}(N,g)=NE^{\mathrm{GP}}(1,Ng)`$. Likewise, $`N^{1/2}\mathrm{\Phi }_{N,g}^{\mathrm{GP}}\varphi _\gamma ^{\mathrm{GP}}`$ depends only on
$$\gamma Ng$$
(2.1)
and satisfies the normalization condition $`|\varphi _\gamma ^{\mathrm{GP}}|^2=1`$. The variational equation (GP equation) for the GP minimization problem, written in terms of $`\varphi _\gamma ^{\mathrm{GP}}`$, is
$$\mathrm{\Delta }\varphi _\gamma ^{\mathrm{GP}}+V\varphi _\gamma ^{\mathrm{GP}}+8\pi \gamma (\varphi _\gamma ^{\mathrm{GP}})^3=\mu ^{\mathrm{GP}}(\gamma )\varphi _\gamma ^{\mathrm{GP}},$$
(2.2)
where the Lagrange multiplier (chemical potential) $`\mu ^{\mathrm{GP}}(\gamma )`$ is determined by the subsidiary normalization condition. Multiplying (2.2) with $`\varphi _\gamma ^{\mathrm{GP}}`$ and integrating we obtain
$$\mu ^{\mathrm{GP}}(\gamma )=E^{\mathrm{GP}}(1,\gamma )+4\pi \gamma \varphi _\gamma ^{\mathrm{GP}}(𝐱)^4\mathrm{d}^2𝐱.$$
(2.3)
For the upper bound on the quantum mechanical energy in the next section we shall need a bound on the absolute value of the minimizer $`\varphi _\gamma ^{\mathrm{GP}}`$.
###### Lemma 2.1 (Upper bound for the GP minimizer).
$$\varphi _\gamma ^{\mathrm{GP}}_{\mathrm{}}^2\frac{\mu ^{\mathrm{GP}}(\gamma )}{8\pi \gamma }$$
(2.4)
###### Proof.
$`\varphi _\gamma ^{\mathrm{GP}}`$ is a continuous and positive function that satisfies the variational equation
$$\mathrm{\Delta }\varphi _\gamma ^{\mathrm{GP}}+U\varphi _\gamma ^{\mathrm{GP}}=\mu ^{\mathrm{GP}}\varphi _\gamma ^{\mathrm{GP}}$$
(2.5)
with $`U=V+8\pi \gamma (\varphi _\gamma ^{\mathrm{GP}})^2`$. Let $`=\{𝐱|\varphi _\gamma ^{\mathrm{GP}}(𝐱)^2>\mu ^{\mathrm{GP}}/(8\pi \gamma )\}`$. Since $`V0`$ we see that $`\mathrm{\Delta }\varphi _\gamma ^{\mathrm{GP}}0`$ on $``$, i.e., $`\varphi _\gamma ^{\mathrm{GP}}`$ is subharmonic on $``$. Hence $`\varphi _\gamma ^{\mathrm{GP}}`$ achieves its maximum on the boundary of $``$, where $`\varphi _\gamma ^{\mathrm{GP}}(𝐱)^2=\mu ^{\mathrm{GP}}/(8\pi \gamma )`$, so $``$ is empty. ∎
The ground state energy $`E^{\mathrm{TF}}(N,g)`$ of the TF functional (1.2) scales in the same way as $`E^{\mathrm{GP}}(N,g)`$, i.e., $`E^{\mathrm{TF}}(N,g)=NE^{\mathrm{TF}}(1,Ng)`$, and the corresponding minimizer $`\rho _{N,g}^{\mathrm{TF}}`$ is equal to $`N\rho _{1,Ng}^{\mathrm{TF}}`$. For short, we shall denote $`\rho _{1,\gamma }^{\mathrm{TF}}`$ by $`\rho _\gamma ^{\mathrm{TF}}`$. By (1.5) we have
$$\rho _\gamma ^{\mathrm{TF}}(𝐱)=\frac{1}{8\pi \gamma }[\mu ^{\mathrm{TF}}(\gamma )V(𝐱)]_+,$$
(2.6)
with the chemical potential $`\mu ^{\mathrm{TF}}(\gamma )`$ determined by the normalization condition $`\rho _\gamma ^{\mathrm{TF}}=1`$. In the same way as in (2.3) we have
$$\mu ^{\mathrm{TF}}(\gamma )=E^{\mathrm{TF}}(1,\gamma )+4\pi \gamma \rho _\gamma ^{\mathrm{TF}}(𝐱)^2\mathrm{d}^2𝐱.$$
(2.7)
The chemical potential can also be computed from a variational principle:
###### Lemma 2.2 (Variational principle for $`\mu ^{\mathrm{TF}}`$).
$$\mu ^{\mathrm{TF}}(\gamma )=\underset{\rho 0,{\scriptscriptstyle \rho }=1}{inf}V\rho +8\pi \gamma \rho _{\mathrm{}}$$
(2.8)
###### Proof.
Obviously, the infimum is achieved for a multiple of a characteristic function for some measurable set $`^2`$. If $`||`$ denotes the Lebesgue measure of $``$, then
$`\underset{{\scriptscriptstyle \rho }=1}{inf}{\displaystyle V\rho }+8\pi \gamma \rho _{\mathrm{}}`$
$`=\underset{}{inf}\left({\displaystyle _{}}V+8\pi \gamma \right){\displaystyle \frac{1}{||}}`$ (2.9)
$`=\underset{}{inf}\left({\displaystyle _{}}\left(V\mu ^{\mathrm{TF}}(\gamma )\right)+8\pi \gamma +\mu ^{\mathrm{TF}}(\gamma )||\right){\displaystyle \frac{1}{||}}.`$ (2.10)
Now $`_{}(V\mu ^{\mathrm{TF}}(\gamma ))8\pi \gamma `$, with equality for
$$\left\{𝐱\right|V(𝐱)<\mu ^{\mathrm{TF}}(\gamma )\}\left\{𝐱\right|V(𝐱)\mu ^{\mathrm{TF}}(\gamma )\}.$$
(2.11)
###### Corollary 2.1 (Properties of $`\mu ^{\mathrm{TF}}(\gamma )`$).
$`\mu ^{\mathrm{TF}}(\gamma )`$ is a concave and monotonously increasing function of $`\gamma `$ with $`\mu ^{\mathrm{TF}}(0)=0`$. Hence $`\mu ^{\mathrm{TF}}(\gamma )/\gamma `$ is decreasing in $`\gamma `$. Moreover, $`\mu ^{\mathrm{TF}}(\gamma )\mathrm{}`$ and $`\mu ^{\mathrm{TF}}(\gamma )/\gamma 0`$ as $`\gamma \mathrm{}`$.
###### Proof.
Immediate consequences of Lemma 2.2, using that $`\mathrm{min}_𝐱V(𝐱)=0`$ and $`lim_{|𝐱|\mathrm{}}V(𝐱)=\mathrm{}`$. ∎
Note that since $`E^{\mathrm{TF}}(1,\gamma )\frac{1}{2}\mu ^{\mathrm{TF}}(\gamma )`$ we also see that $`E^{\mathrm{TF}}(1,\gamma )\mathrm{}`$ with $`\gamma `$. In this limit the GP energy converges to the TF energy, provided the external potential satisfies a mild regularity and growth condition:
###### Lemma 2.3 (TF limit of the GP energy).
Suppose for some constants $`\alpha >0`$, $`L_1`$ and $`L_2`$
$$|V(𝐱)V(𝐲)|L_1|𝐱𝐲|^\alpha e^{L_2|𝐱𝐲|}(1+V(𝐱)).$$
(2.12)
Then
$$\underset{\gamma \mathrm{}}{lim}\frac{E^{\mathrm{GP}}(1,\gamma )}{E^{\mathrm{TF}}(1,\gamma )}=1.$$
(2.13)
###### Proof.
It is clear that $`E^{\mathrm{TF}}(1,\gamma )E^{\mathrm{GP}}(1,\gamma )`$. For the other direction, we use $`(j_ϵ\rho _\gamma ^{\mathrm{TF}})^{1/2}`$ as a test function for $`^{\mathrm{GP}}`$, where
$$j_ϵ(𝐱)=\frac{1}{2\pi ϵ^2}\mathrm{exp}\left(\frac{1}{ϵ}|𝐱|\right).$$
(2.14)
Note that $`j_ϵ=1`$ and $`|j_ϵ|=ϵ^1j_ϵ`$. Therefore
$`E^{\mathrm{GP}}(1,\gamma )`$ $``$ $`{\displaystyle \left(\frac{1}{4j_ϵ\rho _\gamma ^{\mathrm{TF}}}|j_ϵ\rho _\gamma ^{\mathrm{TF}}|^2+V(j_ϵ\rho _\gamma ^{\mathrm{TF}})+4\pi \gamma (j_ϵ\rho _\gamma ^{\mathrm{TF}})^2\right)}`$ (2.15)
$``$ $`{\displaystyle \frac{1}{4ϵ^2}}+{\displaystyle \left((j_ϵV)\rho _\gamma ^{\mathrm{TF}}+4\pi \gamma (\rho _\gamma ^{\mathrm{TF}})^2\right)},`$
where we have used convexity for the last term. Moreover,
$`{\displaystyle (j_ϵVV)\rho _\gamma ^{\mathrm{TF}}}`$ $`=`$ $`{\displaystyle \mathrm{d}^2𝐱\mathrm{d}^2𝐲j_ϵ(𝐱𝐲)\left(V(𝐱)V(𝐲)\right)\rho _\gamma ^{\mathrm{TF}}(𝐱)}`$ (2.16)
$``$ $`{\displaystyle \frac{L_1}{2\pi ϵ^2}}{\displaystyle \mathrm{d}^2𝐱\mathrm{d}^2𝐲|𝐱𝐲|^\alpha e^{(ϵ^1+L_2)|𝐱𝐲|}(1+V(𝐱))\rho _\gamma ^{\mathrm{TF}}(𝐱)}`$
$``$ $`(\mathrm{const}.)ϵ^\alpha (1+E^{\mathrm{TF}}(1,\gamma )),`$
as long as $`ϵ<L_2^1`$. So we have
$$E^{\mathrm{GP}}(1,\gamma )(1+(\mathrm{const}.)ϵ^\alpha )E^{\mathrm{TF}}(1,\gamma )+\frac{1}{4ϵ^2}+(\mathrm{const}.)ϵ^\alpha .$$
(2.17)
Optimizing over $`ϵ`$ gives as a final result
$$E^{\mathrm{GP}}(1,\gamma )E^{\mathrm{TF}}(1,\gamma )(1+(\mathrm{const}.)E^{\mathrm{TF}}(1,\gamma )^{\alpha /(\alpha +2)}).$$
(2.18)
Condition (2.12) is in particular fulfilled if $`V`$ is homogeneous of some order $`s>0`$ and locally Hölder continuous. In this case,
$$E^{\mathrm{TF}}(1,\gamma )=\gamma ^{s/(s+2)}E^{\mathrm{TF}}(1,1)$$
(2.19)
and
$$\gamma ^{2/(s+2)}\rho _\gamma ^{\mathrm{TF}}(\gamma ^{1/(s+2)}𝐱)=\rho _{1,1}^{\mathrm{TF}}(𝐱).$$
(2.20)
By (2.7) we also have
$$\mu ^{\mathrm{TF}}(\gamma )=\gamma ^{s/(s+2)}\mu ^{\mathrm{TF}}(1).$$
(2.21)
If $`V`$ is asymptotically homogeneous with a locally Hölder continuous limiting function $`W`$, we can prove corresponding formulas for the limit $`\gamma \mathrm{}`$. This is the content of the next theorem, where we have included results on the GP $``$ TF limit as well:
###### Theorem 2.1 (Scaling limits).
Suppose $`V`$ satisfies the condition of Theorem 1.3. Let $`\stackrel{~}{E}^{\mathrm{TF}}(1,1)`$ be the minimum of the TF functional (1.2) with $`g=1`$ and $`N=1`$ and $`V`$ replaced by $`W`$, and let $`\stackrel{~}{\rho }_{1,1}^{\mathrm{TF}}`$ be the corresponding minimizer. Then
1. $`lim_\gamma \mathrm{}E^{\mathrm{GP}}(1,\gamma )/\gamma ^{s/(s+2)}=lim_\gamma \mathrm{}E^{\mathrm{TF}}(1,\gamma )/\gamma ^{s/(s+2)}=\stackrel{~}{E}^{\mathrm{TF}}(1,1)`$.
2. $`lim_\gamma \mathrm{}\gamma ^{2/(s+2)}\rho _{1,\gamma }^{\mathrm{GP}}(\gamma ^{1/(s+2)}𝐱)=\stackrel{~}{\rho }_{1,1}^{\mathrm{TF}}(𝐱)`$, strongly in $`L^2(^2)`$.
3. $`lim_\gamma \mathrm{}\gamma ^{2/(s+2)}\rho _\gamma ^{\mathrm{TF}}(\gamma ^{1/(s+2)}𝐱)=\stackrel{~}{\rho }_{1,1}^{\mathrm{TF}}(𝐱)`$, uniformly in $`𝐱`$.
###### Proof.
With the demanded properties of $`V`$, (2.13) holds. Using this and (1.9) one easily verifies (i). Moreover, $`\gamma ^{2/(s+2)}\rho _{1,\gamma }^{\mathrm{GP}}(\gamma ^{1/(s+2)}𝐱)`$ is a minimizing sequence for the functional in question, so we can conclude as in Theorem II.2 in that it converges to $`\stackrel{~}{\rho }_{1,1}^{\mathrm{TF}}(𝐱)`$ strongly in $`L^2`$, proving (ii). (Remark: In Eq. (2.10) in there is a misprint, instead of $`\rho _{1,Na}^{\mathrm{GP}}`$ one should have $`\stackrel{~}{\rho }_{1,Na}^{\mathrm{GP}}`$ on the left side.) To see (iii) let us define
$$\widehat{\rho }_\gamma (𝐱)=\gamma ^{2/(s+2)}\rho _\gamma ^{\mathrm{TF}}\left(\gamma ^{1/(s+2)}𝐱\right).$$
(2.22)
We can write
$$\widehat{\rho }_\gamma (𝐱)=\frac{1}{8\pi }\left[\gamma ^{s/(s+2)}\mu ^{\mathrm{TF}}(\gamma )W(𝐱)ϵ(\gamma ,𝐱)\right]_+$$
(2.23)
with
$$ϵ(\gamma ,𝐱)=\gamma ^{s/(s+2)}V(\gamma ^{1/(s+2)}𝐱)W(𝐱).$$
(2.24)
By assumption, $`|ϵ(\gamma ,𝐱)|<\delta (\gamma )(1+W(𝐱))`$ for some $`\delta (\gamma )`$ with $`lim_\gamma \mathrm{}\delta (\gamma )=0`$. Because $`\widehat{\rho }_\gamma =1`$ for all $`\gamma `$, we see from Eq. (2.23) that $`\mu ^{\mathrm{TF}}(\gamma )\gamma ^{s/(s+2)}`$ converges to some $`c`$ as $`\gamma \mathrm{}`$. Moreover, we can conclude that the support of $`\widehat{\rho }_\gamma `$ is for large $`\gamma `$ contained in some bounded set $``$ independent of $`\gamma `$. Therefore
$$1=\underset{\gamma \mathrm{}}{lim}\widehat{\rho }_\gamma =(8\pi )^1[cW(𝐱)]_+$$
(2.25)
by dominated convergence, so $`c`$ is equal to the $`\stackrel{~}{\mu }^{\mathrm{TF}}`$ of Eq. (1.11). Now
$$\widehat{\rho }_\gamma (𝐱)=\frac{1}{8\pi }\left[\stackrel{~}{\mu }^{\mathrm{TF}}W(𝐱)\overline{ϵ}(\gamma ,𝐱)\right]_+$$
(2.26)
with
$$\overline{ϵ}(\gamma ,𝐱)=ϵ(\gamma ,𝐱)+\stackrel{~}{\mu }^{\mathrm{TF}}\gamma ^{s/(s+2)}\mu ^{\mathrm{TF}}(\gamma ).$$
(2.27)
Again $`|\overline{ϵ}(\gamma ,𝐱)|<\overline{\delta }(\gamma )(1+W(𝐱))`$ for some $`\overline{\delta }(\gamma )`$ with $`lim_\gamma \mathrm{}\overline{\delta }(\gamma )=0`$. By Eqs. (1.11) and (2.26) we thus have
$$\widehat{\rho }_\gamma \stackrel{~}{\rho }_{1,1}^{\mathrm{TF}}_{\mathrm{}}<C\overline{\delta }(\gamma ).$$
(2.28)
with $`C=(8\pi )^1sup_𝐱(1+W(𝐱))<\mathrm{}`$. ∎
The mean density for the TF theory is defined by
$$\overline{\rho }_\gamma N\rho _\gamma ^{\mathrm{TF}}(𝐱)^2\mathrm{d}^2𝐱.$$
(2.29)
For $`\gamma =N`$, i.e., $`g=1`$ this is the same as (1.6). It satisfies
###### Lemma 2.4 (Bounds on $`\overline{\rho }_\gamma `$).
For some constant $`C>0`$
$$N\frac{\mu ^{\mathrm{TF}}(\gamma )}{8\pi \gamma }\overline{\rho }_\gamma CN\frac{\mu ^{\mathrm{TF}}(\gamma )}{\gamma }.$$
(2.30)
###### Proof.
The upper bound is trivial. Because $`\widehat{\rho }_\gamma `$, defined in (2.22), converges uniformly to $`\stackrel{~}{\rho }_{1,1}^{\mathrm{TF}}`$ and $`\mu ^{\mathrm{TF}}(\gamma )\gamma ^{s/(s+2)}\stackrel{~}{\mu }^{\mathrm{TF}}`$ as $`\gamma \mathrm{}`$, we have the lower bound
$$\frac{\gamma \overline{\rho }_\gamma }{N\mu ^{\mathrm{TF}}(\gamma )}8\pi \gamma ^{s/(s+2)}\mu ^{\mathrm{TF}}(\gamma )^1\left((\stackrel{~}{\rho }_{1,1}^{\mathrm{TF}})^22\stackrel{~}{\rho }_{1,1}^{\mathrm{TF}}\widehat{\rho }_{\mathrm{}}\right)>C$$
(2.31)
for some $`C>0`$. ∎
###### Remark 2.1.
With $`V`$ asymptotically homogeneous of order $`s`$, $`\mu ^{\mathrm{TF}}(\gamma )\gamma ^{s/(s+2)}`$ converges as $`\gamma \mathrm{}`$, i.e. $`\mu ^{\mathrm{TF}}(\gamma )\gamma ^{s/(s+2)}`$ for large $`\gamma `$. So the mean TF density for coupling constant $`g=1`$, defined in (1.6), has the asymptotic behavior $`\overline{\rho }N^{s/(s+2)}`$.
## 3 Upper bound to the QM energy
As in the three dimensional case, cf. Eqs. (3.29) and (3.27) in , one has the upper bound
$$\frac{E^{\mathrm{QM}}(N,a)}{N}\frac{|\varphi _\gamma ^{\mathrm{GP}}|^2+V(\varphi _\gamma ^{\mathrm{GP}})^2}{1N\varphi _\gamma ^{\mathrm{GP}}_{\mathrm{}}^2I}+\frac{NJ(\varphi _\gamma ^{\mathrm{GP}})^4+\frac{2}{3}N^2(\varphi _\gamma ^{\mathrm{GP}}_{\mathrm{}}^2K)^2}{(1N\varphi _\gamma ^{\mathrm{GP}}_{\mathrm{}}^2I)^2},$$
(3.1)
where we have implicitly used that $`\mathrm{\Delta }\varphi _\gamma ^{\mathrm{GP}}+V\varphi _\gamma ^{\mathrm{GP}}0`$, which is justified by Lemma 2.1. The coefficients $`I`$, $`J`$ and $`K`$ are given by Eqs. (2.4)–(2.10) in . They depend on the scattering length and a parameter $`b`$. We choose $`\gamma =N/|\mathrm{ln}(a^2\overline{\rho })|`$ and $`b=\overline{\rho }^{1/2}`$. (Recall that $`\overline{\rho }`$ is short for $`\overline{\rho }_N`$.) With this choice we have (as long as $`a^2\overline{\rho }<1`$)
$$J=\frac{4\pi }{|\mathrm{ln}(a^2\overline{\rho })|},$$
(3.2)
and the error terms
$$N\varphi _\gamma ^{\mathrm{GP}}_{\mathrm{}}^2I(\mathrm{const}.)\frac{\mu ^{\mathrm{GP}}(\gamma )}{\overline{\rho }}(1+O(|\mathrm{ln}(a^2\overline{\rho })|^1))$$
(3.3)
and
$$K^2N^2\varphi _\gamma ^{\mathrm{GP}}_{\mathrm{}}^4(\mathrm{const}.)E^{\mathrm{GP}}(1,\gamma )\frac{\mu ^{\mathrm{GP}}(\gamma )}{\overline{\rho }}(1+O(|\mathrm{ln}(a^2\overline{\rho })|^1)),$$
(3.4)
where we have used Lemma 2.1. So we have the upper bound
$$\frac{E^{\mathrm{QM}}(N,a)}{E^{\mathrm{GP}}(N,1/|\mathrm{ln}(a^2\overline{\rho })|)}1+O(\mu ^{\mathrm{GP}}(\gamma )/\overline{\rho })+O((|\mathrm{ln}(a^2\overline{\rho })|^1)).$$
(3.5)
Now if $`\gamma `$ is fixed as $`N\mathrm{}`$
$$\frac{\mu ^{\mathrm{GP}}(\gamma )}{\overline{\rho }}\frac{1}{|\mathrm{ln}(a^2\overline{\rho })|}\frac{1}{N}.$$
(3.6)
If $`\gamma \mathrm{}`$ with $`N`$ we have instead, assuming that the external potential is asymptotically homogeneous of order $`s`$,
$$\frac{\mu ^{\mathrm{GP}}(\gamma )}{\overline{\rho }}\frac{\mu ^{\mathrm{TF}}(\gamma )}{\mu ^{\mathrm{TF}}(N)}\left(\frac{\gamma }{N}\right)^{s/(s+2)},$$
(3.7)
so in any case
$$\frac{E^{\mathrm{QM}}(N,a)}{E^{\mathrm{GP}}(N,1/|\mathrm{ln}(a^2\overline{\rho })|)}1+O\left(|\mathrm{ln}(a^2\overline{\rho })|^{s/(s+2)}\right)$$
(3.8)
holds as $`N\mathrm{}`$ and $`a^2\overline{\rho }0`$.
## 4 Lower bound to the QM energy
Compared to the treatment of the 3D problem in the new issue here is the TF case, i.e., $`\gamma =N/|\mathrm{ln}(a^2\overline{\rho })|\mathrm{}`$, and we discuss this case first. The GP limit with $`\gamma `$ fixed can be treated in complete analogy with the 3D case, cf. Remark 4.1 below.
We introduce again the rescaled $`\widehat{\rho }_\gamma `$ as in (2.22) and also
$$\widehat{v}(𝐱)=\gamma ^{2/(s+2)}v\left(\gamma ^{1/(s+2)}𝐱\right).$$
(4.1)
Note that the scattering length of $`\widehat{v}`$ is $`\widehat{a}=a\gamma ^{1/(s+2)}`$. Using $`V\mu ^{\mathrm{TF}}(\gamma )8\pi \gamma \rho _\gamma ^{\mathrm{TF}}`$ and (2.7) we see that
$`E^{\mathrm{QM}}(N,a)`$ $``$ $`E^{\mathrm{TF}}(N,\gamma /N)+4\pi N\gamma ^{s/(s+2)}{\displaystyle \widehat{\rho }_\gamma ^2}+\gamma ^{2/(s+2)}Q`$ (4.2)
$`8\pi N\gamma ^{s/(s+2)}\widehat{\rho }_\gamma \stackrel{~}{\rho }_{1,1}^{\mathrm{TF}}_{\mathrm{}},`$
with
$$Q=\underset{{\scriptscriptstyle |\mathrm{\Psi }|^2}=1}{inf}\underset{i}{}\left(|_i\mathrm{\Psi }|^2+\underset{j<i}{}\widehat{v}(𝐱_i𝐱_j)|\mathrm{\Psi }|^28\pi \gamma \stackrel{~}{\rho }_{1,1}^{\mathrm{TF}}(𝐱_i)|\mathrm{\Psi }|^2\right).$$
(4.3)
Dividing space into boxes $`\alpha `$ of side length $`L`$ with Neumann boundary conditions we get
$$Q\underset{\alpha }{}E^{\mathrm{hom}}(n_\alpha ,L)8\pi \gamma \rho _{\alpha ,\mathrm{max}}n_\alpha ,$$
(4.4)
where $`\rho _{\alpha ,\mathrm{max}}`$ denotes the maximal value of $`\stackrel{~}{\rho }_{1,1}^{\mathrm{TF}}`$ in the box $`\alpha `$, and $`E^{\mathrm{hom}}(n,L)`$ is the energy of a homogeneous gas of $`n`$ bosons in a box of side length $`L`$ and Neumann boundary conditions. We can forget about the boxes where $`\rho _{\alpha ,\mathrm{max}}=0`$, because the energy of particles in these boxes is positive.
We now want to use the lower bound on $`E^{\mathrm{hom}}`$ given in , namely
$$E^{\mathrm{hom}}(n,L)4\pi \frac{n^2}{L^2}\frac{1}{|\mathrm{ln}(\widehat{a}^2n/L^2)|}\left(1C|\mathrm{ln}(\widehat{a}^2n/L^2)|^{1/5}\right).$$
(4.5)
This bound holds for $`n>(\mathrm{const}.)|\mathrm{ln}(\widehat{a}^2n/L^2)|^{1/5}`$ and small enough $`\widehat{a}^2n/L^2`$. Now if the minimum in (4.4) is taken in some box $`\alpha `$ for some value $`n_\alpha `$, we have
$$E^{\mathrm{hom}}(n_\alpha +1,L)E^{\mathrm{hom}}(n_\alpha ,L)8\pi \gamma \rho _{\alpha ,\mathrm{max}}.$$
(4.6)
By a computation analogous to the upper bound (see ) one shows that
$`E^{\mathrm{hom}}(n+1,L)E^{\mathrm{hom}}(n,L)`$
$`8\pi {\displaystyle \frac{n}{L^2}}{\displaystyle \frac{1}{|\mathrm{ln}(\widehat{a}^2n/L^2)|}}\left(1+O\left(|\mathrm{ln}(\widehat{a}^2n/L^2)|^1\right)\right).`$ (4.7)
Using Lemma 2.4 and the asymptotics of $`\mu ^{\mathrm{TF}}`$ (Remark 2.1) we see that
$$\frac{\widehat{a}^2n}{L^2}\frac{\widehat{a}^2N}{L^2}=N^{s/(s+2)}\left(\frac{N}{\gamma }\right)^{2/(s+2)}\frac{a^2}{L^2}a^2\overline{\rho }\frac{C}{L^2}\left(\frac{N}{\gamma }\right)^{2/(s+2)},$$
(4.8)
for some constant $`C`$, so (4.7) reads
$`E^{\mathrm{hom}}(n+1,L)E^{\mathrm{hom}}(n,L)`$
$`8\pi {\displaystyle \frac{n}{L^2}}{\displaystyle \frac{1}{|\mathrm{ln}(a^2\overline{\rho })|}}\left(1+O\left({\displaystyle \frac{1+|\mathrm{ln}((\gamma /N)^{2/(s+2)}L^2/C)|}{|\mathrm{ln}(a^2\overline{\rho })|}}\right)\right).`$ (4.9)
So if $`L`$ is fixed, our minimizing $`n_\alpha `$ is at least $`\rho _{\alpha ,\mathrm{max}}L^2N`$. If $`N`$ is large enough and $`a^2\overline{\rho }`$ is small enough, we can thus use (4.5) in (4.4) to get
$$Q\underset{\alpha }{}4\pi \left(\frac{n_\alpha ^2}{L^2}\frac{1}{|\mathrm{ln}\left(\frac{\widehat{a}^2n_\alpha }{L^2}\right)|}\left(1\frac{C}{|\mathrm{ln}\left(\frac{\widehat{a}^2N}{L^2}\right)|^{1/5}}\right)2\frac{N\rho _{\alpha ,\mathrm{max}}}{|\mathrm{ln}(a^2\overline{\rho })|}\right).$$
(4.10)
###### Lemma 4.1.
For $`0<x,b<1`$ we have
$$\frac{x^2}{|\mathrm{ln}x|}2\frac{b}{|\mathrm{ln}b|}x\frac{b^2}{|\mathrm{ln}b|}\left(1+\frac{1}{(2|\mathrm{ln}b|)^2}\right).$$
(4.11)
###### Proof.
Since $`\mathrm{ln}x\frac{1}{de}x^d`$ for all $`d>0`$ we have
$$\frac{x^2}{b^2}\frac{|\mathrm{ln}b|}{|\mathrm{ln}x|}2\frac{x}{b}\frac{|\mathrm{ln}b|}{b^2}edx^{2+d}\frac{2x}{b}c(d)(b^ded|\mathrm{ln}b|)^{1/(1+d)}$$
(4.12)
with
$$c(d)=2^{(2+d)/(1+d)}\left(\frac{1}{(2+d)^{(2+d)/(1+d)}}\frac{1}{(2+d)^{1/(1+d)}}\right)1\frac{1}{4}d^2.$$
(4.13)
Choosing $`d=1/|\mathrm{ln}b|`$ gives the desired result. ∎
Note that the Lemma above implies for $`k1`$
$$\frac{x^2}{|\mathrm{ln}x|}2\frac{b}{|\mathrm{ln}b|}xk\frac{b^2}{|\mathrm{ln}b|}\left(1+\frac{1}{(2|\mathrm{ln}b|)^2}\right)k^2.$$
(4.14)
Applying this with $`x=\widehat{a}^2n_\alpha /L^2`$ and $`b=N\widehat{a}^2\rho _{\alpha ,\mathrm{max}}`$ we get the bound
$`Q`$ $`4\pi N\gamma {\displaystyle \underset{\alpha }{}}\rho _{\alpha ,\mathrm{max}}^2L^2`$
$`\times \left[\left(1+{\displaystyle \frac{1}{4|\mathrm{ln}(\widehat{a}^2N\rho _{\alpha ,\mathrm{max}})|^2}}\right){\displaystyle \frac{|\mathrm{ln}(\widehat{a}^2N\rho _{\alpha ,\mathrm{max}})|}{|\mathrm{ln}(a^2\overline{\rho })|}}\left(1{\displaystyle \frac{C}{|\mathrm{ln}\left(\frac{\widehat{a}^2N}{L^2}\right)|^{1/5}}}\right)^1\right]`$
for (4.10). To estimate the error terms, note that as in (4.8)
$$\widehat{a}^2Na^2\overline{\rho }\left(\frac{N}{\gamma }\right)^{2/(s+2)},$$
(4.16)
so $`|\mathrm{ln}(\widehat{a}^2N)|=|\mathrm{ln}(a^2\overline{\rho })|+O(\mathrm{ln}|\mathrm{ln}(a^2\overline{\rho })|)`$ for small $`a^2\overline{\rho }`$. Using $`\widehat{\rho }_\gamma \stackrel{~}{\rho }_{1,1}^{\mathrm{TF}}_{\mathrm{}}0`$ (Theorem 2.1 (iii)) and $`\widehat{\rho }_\gamma ^2(\stackrel{~}{\rho }_{1,1}^{\mathrm{TF}})^2`$ as $`\gamma \mathrm{}`$ (which follows from the uniform convergence and boundedness of the supports) we get
$$\underset{N\mathrm{}}{lim\; inf}\frac{E^{\mathrm{QM}}(N,a)}{E^{\mathrm{TF}}(N,1/|\mathrm{ln}(a^2\overline{\rho })|)}1(\mathrm{const}.)(\underset{\alpha }{}\rho _{\alpha ,\mathrm{max}}^2L^2(\stackrel{~}{\rho }_{1,1}^{\mathrm{TF}})^2).$$
(4.17)
Since this holds for all choices of the boxes $`\alpha `$ with arbitrary small side length $`L`$, and by the assumptions on $`V`$ $`\stackrel{~}{\rho }_{1,1}^{\mathrm{TF}}`$ is continuous and has compact support, we can conclude
$$\underset{N\mathrm{}}{lim\; inf}\frac{E^{\mathrm{QM}}(N,a)}{E^{\mathrm{TF}}(N,1/|\mathrm{ln}(a^2\overline{\rho })|)}1$$
(4.18)
in the limit $`N\mathrm{}`$, $`a^2\overline{\rho }0`$ and $`N/|\mathrm{ln}(a^2\overline{\rho })|\mathrm{}`$.
###### Remark 4.1 (The GP case).
In the derivation of the lower bound we have assumed that $`\gamma \mathrm{}`$ with $`N`$, i.e. $`N|\mathrm{ln}(a^2\overline{\rho })|`$, which seems natural because otherwise the scattering length would have to decrease exponentially with $`N`$. However, for fixed $`\gamma `$ one can use the methods of (with slight modifications: One uses the 2D bounds on the homogeneous gas and Lemma 4.1) to compute a lower bound in terms of the GP energy. The result is
$$\underset{N\mathrm{}}{lim\; inf}\frac{E^{\mathrm{QM}}(N,a)}{E^{\mathrm{GP}}(N,1/|\mathrm{ln}(a^2\overline{\rho })|)}1$$
(4.19)
in the limit $`N\mathrm{}`$, $`a^2\overline{\rho }0`$ with $`\gamma =N/|\mathrm{ln}(a^2\overline{\rho })|`$ fixed.
## 5 The limit theorems
We have now all the estimates needed for Theorems 1.11.4. The upper bound (3.8) and the lower bound (4.19) prove Theorem 1.1. The energy limit Theorem 1.3 for the TF case follows from (3.8), Theorem 2.1 (i) and (4.18).
The convergence of the energies implies the convergence of the densities in the usual way by variation of the external potential. Replacing $`V(𝐱)`$ by $`V(𝐱)+\delta \gamma ^{s/(s+2)}Y(\gamma ^{1/(s+2)}𝐱)`$ for some positive $`YC_0^{\mathrm{}}`$ and redoing the upper and lower bounds we see that Theorem 1.3 and Theorem 2.1 (i) hold with $`W`$ replaced by $`W+\delta Y`$. Differentiating with respect to $`\delta `$ at $`\delta =0`$ yields
$$\underset{N\mathrm{}}{lim}\frac{\gamma ^{2/(s+2)}}{N}\rho _{N,a}^{\mathrm{QM}}(\gamma ^{1/(s+2)}𝐱)=\stackrel{~}{\rho }_{1,1}^{\mathrm{TF}}(𝐱)$$
(5.1)
in the sense of distributions. Since the functions all have norm 1, we can conclude that there is even weak $`L^1`$-convergence.
###### Remark 5.1 (The 3D case).
In the analogues of Theorems 1.1 and 1.2 were shown for the three-dimensional Bose gas. Using the methods developed here one can extend these results to analogues of Theorems 1.3 and 1.4. In 3D the coupling constant is $`g=a`$, so $`\gamma =Na`$. Moreover, the relevant mean 3D density is $`\overline{\rho }_\gamma N(Na)^{3/(s+3)}`$.
## Appendix A Appendix: Scattering length in two dimensions
Due to the logarithmic behavior of the Green function of the two dimensional Laplacian the definition of the scattering length is slightly more delicate in two dimensions than in three. For a nonnegative potential $`v(𝐱)`$, depending only on $`|𝐱|`$ and with finite range $`R_0`$, it is naturally defined by the following variational principle:
###### Theorem A.1.
Let $`R>R_0`$ and consider the functional
$$_R[\varphi ]=_{|𝐱|R}\left\{|\varphi (𝐱)|^2+\frac{1}{2}v(𝐱)|\varphi (𝐱)|^2\right\}\mathrm{d}^2𝐱.$$
(A.1)
Then, in the subclass of functions such that $`(|\varphi |^2+|\varphi |^2)<\mathrm{}`$ and $`\varphi (𝐱)=1`$ for $`|𝐱|=R`$, there is a unique function $`\varphi _0`$ that minimizes $`_R[\varphi ]`$. This function is nonnegative and rotationally symmetric, and satisfies the equation
$$\mathrm{\Delta }\varphi _0(𝐱)+\frac{1}{2}v(𝐱)\varphi _0(𝐱)=0$$
(A.2)
for $`|𝐱|R`$ in the sense of distributions, with boundary condition $`\varphi _0(𝐱)=1`$ for $`|𝐱|=R`$.
For $`R_0<|𝐱|<R`$
$$\varphi _0(𝐱)=\mathrm{ln}(|𝐱|/a)/\mathrm{ln}(R/a)$$
(A.3)
for a unique number $`a`$ called the scattering length.
For the proof see , where generalizations to other dimensions and potentials with a negative part are also discussed. Note that the factor $`\frac{1}{2}`$ in (A.1) and (A.2) is due to the reduced mass of the two body problem.
If $`v`$ has infinite range it is easy to extend the definition of the scattering length for nonnegative $`v`$ under the assumption that $`_{|𝐱|R_1}^{\mathrm{}}v(𝐱)\mathrm{d}^2𝐱<\mathrm{}`$ for some $`R_1`$. In fact, one may then simply cut off the potential at some point $`R_0>R_1`$ (i.e., set $`v(𝐱)=0`$ for $`|𝐱|>R_0`$) and consider the limit of the scattering lengths of the cut off potentials as $`R_0\mathrm{}`$. See for details.
|
warning/0005/astro-ph0005413.html
|
ar5iv
|
text
|
# 1 Magnetic field vs. spin period diagram. Solid diamonds are observed data summarized in table 1. Initial spin period is chosen to be 100 s. The initial magnetic fields 𝐵₀=5×10¹² G and 𝐵₀=1×10¹² G are used, and the luminosity is varied from the Eddington limited luminosity 𝐿₃₈=1 to 𝐿₃₆=1. The upper heavy solid line is the spin-up line (equilibrium period line) given in equation (11) and the lower solid line is the death line defined as 𝐵₁₂/𝑃² = 0.2. Curve 1 and curve 3 are calculated from equation (6) with luminosity 𝐿₃₈=1 and 𝐿₃₆=1 respectively. Curve 2 and curve 4 are our numerical solutions for equations (1)-(5) with luminosity 𝐿₃₈=1. and 𝐿₃₆=1. respectively. In all model evolution curves, we have used M = 1𝑀_⊙ and R = 106cm.
The Evolution of Magnetic Field and Spin Period of Accreting Neutron Stars
K.S. Cheng<sup>1</sup> and C.M. Zhang<sup>1,2,3</sup>
1. Department of Physics,
The University of Hong Kong,
Pokfulam Road, Hong Kong, P.R. China
2. Department of Physics,
Hebei University of Technology,
Tianjin-300130, P.R. China
3.Instituto de Física Teórica
Universidade Estadual Paulista
Rua Pamplona 145
01405-900 São Paulo
Brazil
Received:
Accepted:
ABSTRACT
Based on the accretion induced magnetic a field decay model, in which a frozen field and an incompressible fluid are assumed, we obtain the following results. (1) An analytic relation between the magnetic field and spin period, if the fastness parameter of the accretion disk is neglected. The evolutionary tracks of accreting neutron stars in the P-B diagram in our model are different from the equilibrium period lines when the influence of the fastness parameter is taken into account. (2) The theoretical minimum spin period of an accreting neutron star is $`\mathrm{max}(1.1\mathrm{ms}\left(\frac{\mathrm{\Delta }M}{M_{}}\right)^1R_6^{5/14}I_{45}\left(\frac{M}{M_{}}\right)^{1/2},1.1\mathrm{ms}\left(\frac{M}{M_{}}\right)^{1/2}R_6^{17/14})`$, independent of the accretion rate (X-ray luminosity) but dependent on the total accretion mass $`\mathrm{\Delta }M`$. However, the minimum magnetic field depends on the accretion rate. (3) The magnetic field strength decreases faster with time than the period.
Subject headings: stars: neutron – stars: rotation – stars: magnetic field –X-rays: stars–stars: binaries: close
1. Introduction
The magnetic field of a neutron star has long been a complex issue and one which is not yet to be solved(Bhattachaya & van den Heuvel 1991; Chanmugam 1992; Phinney & Kulkarni 1994). On the evolution of the magnetic field of a neutron star, there is not yet a commonly accepted model (for a general review cf. Bhattacharya and Srinivasan 1995). The currently popular idea seems to ascribe the field decay of a neutron star in an X-ray binary to the period during which accretion occurs. There is evidence that the magnetic fields of X-ray neutron stars and recycled pulsars are correlated with the duration of the mass accretion phase, or the total amount of matter accreted(Taam & van den Heuvel 1986; van den Heuvel et al 1986). In fact, Taam and van den Heuvel have already discovered a possible inverse correlation between the magnetic field and the estimated total mass of accreted matter for the binary X-ray sources. Later, Shibazaki et al. (1989) presented an assumed formula relating the decay of magnetic field with accretion mass,which seems to reproduce the observed field-period relations of the recycled pulsars quite well.
Theoretically, for explaining the accretion induced field decay, some suggestions and models have been proposed(e.g. Bisnovatyi-Kogan & Komberg 1974; Romani 1990; Ruderman 1991a,b,c; Ding et al. 1993; Zhang et al. 1994; Urpin & Geppert 1995; Geppert et al 1996; Urpin & Konenkov 1997; Zhang et al 1997; Cheng & Dai 1997; Zhang 1998; Ruderman, Zhu & Chen 1998). Recently, van den Heuvel & Bitzaraki(1995a, 1995b), from the statistical analysis of 24 binary radio pulsars with nearly circular orbits and low mass companions, discovered a clear correlation between spin period and orbital period, as well as between the magnetic field and orbital period. These relations strongly suggest that an increase in the amount of accreted mass leads to a decay of the magnetic field, and a ’bottom’ field strength of about 10<sup>8</sup> G is also implied. White & Zhang(1997) discovered that the spin periods of LMXBs, implied by killo-hertz X-ray QPO, constitute a homogenouse group with spin period of about 2 milliseconds, which has little correlation with X-ray luminosity. The above two recent observational statistics seem to place some constraints on the construction of a theoretical model of accretion induced magnetic field decay.
In this paper, we study the evolution of a magnetic field and spin period of accreting neutron stars according to an accretion induced magnetic field decay model(Zhang et al 1997; Cheng & Zhang 1998), based on the idea of van den Heuvel & Bitzaraki(1995a, 1995b) and Romani(1990) for the mechanism of accretion induced field decay. In this model, the accretion matter starts to be channeled onto the two polar caps by the strong magnetic field near the Alfven radius. Part of the accreted matter flowing towards the equator pushes the field lines aside and thus dilutes the polar field strength. The bottom field should be reached when the polar cap extends over the entire stellar surface, which corresponds to the Alfven radius matching the star’s radius, and gives a stellar magnetic field of about 10<sup>8</sup> G. Since the spin period of an accreting X-ray neutron star depends sensitively on the magnetic field together with the influence of the fastness parameter, some interesting results can be obtained, which are consistent with the recent observational data on low mass X-ray neutron stars by White & Zhang (1997) and on the millisecond pulsars by van den Heuvel & Bitzaraki(1995a, 1995b).
2. Models
Under the assumption that the magnetic field lines of the neutron star are frozen in the entire crust which has homogenous average mass density, we (Cheng & Zhang 1998) have obtained an analytical expression for the field evolution as follows,
$$B=\frac{B_f}{\{1C\mathrm{exp}(\frac{\mathrm{\Delta }M}{M_{cr}})\}^{7/4}}$$
(1)
where $`C=1x_0^2`$ and $`x_0=(\frac{B_f}{B_0})^{2/7}`$, $`B_0`$ is the initial magnetic field strength, $`M_{cr}`$ is the crustal mass, $`\mathrm{\Delta }M=\dot{M}t`$ and $`B_f`$ is the magnetic field defined by the Alfven radius matching the radius of neutron star, i.e., $`R_A(B_f)=R`$, which gives,
$$B_f=4.3\times 10^8(\frac{\dot{M}}{\dot{M}_{Ed}})^{1/2}(\frac{M}{M_{}})^{1/4}R_6^{5/4}G,$$
(2)
where $`\dot{M}_{Ed}=10^{18}R_6\mathrm{g}\mathrm{s}^1`$ is the Eddington accretion rate and $`R_6`$ is the radius of the neutron star in units of $`10^6`$ cm.
Three consistent observational conclusions can be obtained, (1) the field decay is inversely related to the accreted mass, (2) the bottom field strength is about $`10^8`$ Gauss, and (3) the bottom field strength is proportionally related to the X-ray luminosity. On the basis of the above solution, we will study the spin period evolution of accreting neutron stars.
To acquire the magnetic field versus period relation, we use the formula for the variation of the rotation due to accretion given by Ghosh & Lamb (1979, hereafter GL)
$$\dot{P}=5.0\times 10^5[(\frac{M}{M_{}})^{3/7}R_6^{12/7}I_{45}^1]B_{12}^{2/7}(PL_{37}^{3/7})^2n(\omega _s)\mathrm{s}\mathrm{yr}^1,$$
(3)
where $`B_{12}`$ is the surface field in units of $`10^{12}`$ G, $`I_{45}`$ is the moment of inertia in units of $`10^{45}\mathrm{g}\mathrm{cm}^2`$, $`L_{37}`$ is the X-ray brightness in units of $`10^{37}\mathrm{erg}\mathrm{s}^1`$ and $`n(\omega _s)`$ is a dimensionless function that depends primarily on the fastness parameter,
$$\omega _s=1.35[(\frac{M}{M_{}})^{2/7}R_6^{15/7}]B_{12}^{6/7}P^1L_{37}^{3/7}.$$
(4)
For a star rotating slowly in the same sense as the disk flow $`(\omega _s1)`$, GL found that $`n(\omega _s)1.4`$. They also found that the dimensionless function $`n(\omega _s)`$ decreases with increasing $`\omega _s`$ and becoming negative for $`\omega _s>\omega _c`$. A simple expression for $`n(\omega _s)`$ that agrees approximately with numerical results over the whole range of $`\omega _s`$ is,
$$n(\omega _s)=1.4\times \left(\frac{1\omega _s/\omega _c}{1\omega _s}\right).$$
(5)
GL found $`\omega _c0.35`$ from their model, but stressed that the actual value of this critical fastness parameter was relatively uncertain. Subsequent work(Ghosh & Lamb 1991) indicates that $`\omega _c`$ is unlikely to be less than 0.2, but it could be as large as 0.9. In the following subsections, we study the influence of the field decay on the evolution of the spin period of the neutron star in an X-ray binary.
2.1 Spin Evolution with a Constant Fastness Parameter
If the variation of the fastness parameter can be ignored, i.e. $`n(\omega _s)`$ =1, we can solve equation (3) analytically. First we can rewrite equation (3) as $`\frac{dP}{P^2}=5.0\times 10^5[(\frac{M}{M_{}})^{3/7}R_6^{12/7}I_{45}^1]B_{f12}^{2/7}L_{37}^{6/7}z^{2/7}dt\mathrm{s}\mathrm{yr}^1`$, where $`B_{f12}=B_f/10^{12}G`$ and $`z=B/B_f`$. From equation (1), we obtain $`dt=\frac{4M_{cr}}{7\dot{M}z}\frac{dz}{z^{4/7}1}`$, which gives $`\frac{dP}{P^2}=10^4\frac{M_{cr}}{\dot{M}}(\frac{M}{M_{}})^{3/7}R_6^{12/7}I_{45}^1`$ $`B_{f12}^{2/7}L_{37}^{6/7}\frac{dy}{y^21}\mathrm{s}\mathrm{yr}^1`$, where $`y=z^{2/7}`$. Using equation (2) to eliminate $`B_f`$, assuming mass and radius as constants and the initial period condition $`P_0=\mathrm{}`$ , we obtain the field-period relation (B-P) in the following analytic form,
$$P=\frac{1.5\mathrm{ms}}{\mathrm{atanh}(\mathrm{x})\mathrm{atanh}(\mathrm{x}_0)}(M/M_{})^{1/2}R_6^{5/14}I_{45}(M_{cr}/0.1M_{})^1$$
(6)
where $`x=y^1=(\frac{B_f}{B})^{2/7},x_0=(\frac{B_f}{B_0})^{2/7}`$ and $`\mathrm{atanh}(\mathrm{x})=[\mathrm{ln}(1+\mathrm{x})\mathrm{ln}(1\mathrm{x})]/2`$. There are two interesting limits for equations (1) and (6). For $`\mathrm{\Delta }MM_{cr}`$, equations (1) and (6) can be approximated as
$$BB_f\frac{M_{cr}}{\mathrm{\Delta }M}t^1$$
(7)
and
$$P\frac{1.5\mathrm{ms}}{x}\left(\frac{M}{M_{}}\right)^{1/2}R_6^{5/14}I_{45}\left(\frac{M_{cr}}{0.1M_{}}\right)^1t^{2/7}.$$
(8)
For $`\mathrm{\Delta }MM_{cr}`$, equations (1) and (6) can be approximated as
$$BB_f\left(1+\frac{7}{4}\mathrm{exp}\left(\frac{\mathrm{\Delta }M}{M_{cr}}\right)\right)$$
(9)
and
$$P1.1\mathrm{ms}\left(\frac{\mathrm{\Delta }M}{M_{}}\right)^1R_6^{5/14}I_{45}\left(\frac{M}{M_{}}\right)^{1/2}.$$
(10)
We should note that the period of the neutron star cannot be shorter than the equilibrium spin-up line which represents the minimum period (Bhattacharya & van den Heuvel 1991)
$$P_{eq}=2.4\mathrm{ms}B_9\left(\frac{\dot{M}}{\dot{M}_{Ed}}\right)^{3/7}\left(\frac{M}{M_{}}\right)^{5/7}R_6^{16/7},$$
(11)
to which such a spin-up may proceed at the Eddington accretion rate. If we substitute the minimum B field in equation (11), we obtain the minimum equilibrium period,
$$P_{min}^{eq}=1.1\mathrm{ms}\left(\frac{M}{M_{}}\right)^{1/2}R_6^{17/14}.$$
(12)
However, we want to point out that equation (12) is valid only in cases where the neutron star has accreted a sufficient amount of matter to spin-up to that period. For a given amount of accreted matter, the minimum period of the neutron star is
$$P_{min}=\mathrm{max}(1.1\left(\frac{\mathrm{\Delta }M}{M_{}}\right)^1R_6^{5/14}I_{45}\left(\frac{M}{M_{}}\right)^{1/2},1.1\left(\frac{M}{M_{}}\right)^{1/2}R_6^{17/14})\mathrm{ms},$$
(13)
an expression which is independent of the accretion rate (X-ray luminosity) but depends on the total accretion mass $`\mathrm{\Delta }M`$.
We plot the field-period relation in figure 1, which shows the evolutionary track curves of the magnetic field and spin period (curve 1 and curve 3). Solid diamonds in figure 1 are observed data summarized in table 1. Initially, a small amount of mass is transferred, and the neutron star is spun-up from the death valley (Chen & Ruderman 1993) where it has a long period, which causes a modest field decay and produces systems such as PSR0655+64 and PSR1913+16. The binaries with longer-lived accretion phases, e.g.LMXB, will accept sufficient mass from their companions, and yield a substantial field decay as in the case of millisecond pulsars such as PSR1953+10 and PSR1620+21. Our analytic model B-P curves, which are obtained by assuming $`n(\omega _s)`$ = 1, can go beyond the spin-up line (the equilibrium period line). This results from the fact that the influence of the fastness parameter is neglected, therefore the spin-up torque still exists even when the spin of the neutron star is faster than the Keplerian angular velocity of the accretion disk at the inner radius. ¿From the analytical P-B relation, we find that the minimum period can be obtained if the magnetic field arrives the bottom field strength(cf. equation 6). The maximum mass accreted in LMXB from the companion could be $`1.0M_{}`$ (van den Heuvel & Bitzaraki 1995a,1995b), the mass of the neutron star after accretion could reach $`2.4M_{}`$ with a radius $`R_61`$ for realistic equations of state (cf. Table I of Cheng & Dai 1997) and I$`{}_{45}{}^{}`$ M. So the minimum spin period given by equation (10) is about $`P_{min}1.7`$ ms. Further, the more interesting thing is that, unlike the bottom field, the minimum period is independent of the accretion rate(X-ray luminosity) (cf. equations (10) and (12)). Our expression seems to be supported by the recent work of White & Zhang(1997). They find that the luminosities in the 10 samples of QPO LMXBs vary by two order of magnitude from $`L_{36}=1`$ to $`L_{38}=1`$, but the spin periods of the sources diffuse into a narrow region from 2.76 ms to 3.8 ms.
2.2 Spin Evolution with a Non-constant Fastness Parameter
Numerical solutions (curve 2 and curve 4) for equations (1), (3), (4) and (5) where the critical fastness parameter is set at 0.9 are plotted in figure 1. However, we find that the influence of the fastness parameter has little effect in the low magnetic field region. The main effect of the fastness parameter is to force the evolution curves back/below the equilibrium period line. The evolutionary curves in the P-B diagram, however, are not sensitive to the fastness parameter near the bottom field and/or near the minimum period. ¿From figure 1 these curves show that the field decay time scale is longer than the spin-up time scale because the field reaches the bottom value first, then the neutron star evolves towards the minimum spin period horizontally in the P-B diagram. The numerical solution are consistent with our analytic expressions in equations (7) and (8).
However, the fastness parameter effect is important when the evolutionary track is close to the equilibrium period line, which means that the spin angular velocity of the star matches the Keplerian angular velocity at the inner edge of the accretion disk, and the accretion torque produced by magnetic lines immersed in the accretion disk tends to produce the negative(spin-down) torque. This effect ensures that the evolution track cannot go beyond the equilibrium period line. It is interesting to note that none of the model B-P curves go along the equilibrium period line, which may seem confusing. However, in fact, the equilibrium period line only represents the final position of the evolution track in the P-B diagram. Therefore it is not surprising that the real evolutionary track deviates from the equilibrium period line. Physically, the accretion induced field decay arises from the contraction of the corotation radius of the magnetosphere during the accretion spin-up phase, and the evolution track has little chance to meet the equilibrium line if the field decay time scale is shorter than the spin-up time scale at the early stage of the accretion phase. Some X-ray sources in HMXB such as Her X-1 and Vela X-1 should be very close to the equilibrium period position if the field decay really exists in the accretion phase.
3.Conclusion
We have presented a simple model for the evolution of the magnetic field and spin period of accreting neutron stars. Analytic formulae for evolution trajectories in B-P are derived. The theoretical minimum period of the neutron star does not depend on the accretion rate but instead depends on the total amount of accreted matter, and the stellar parameters including the moment of inertia, stellar mass and radius, which depend on the equations of state. Our model results seem to be supported by the observed data (White & Zhang 1997). However, in this paper we have ignored the fact that the stellar parameters, i.e. $`M`$, $`R`$ and $`I`$, are all time dependent. The exact evolution curves must take this factor into account. On the other hand, the minimum period should not depend on the details of the evolution trajectories instead it only depends on the final values of the stellar parameters which are equations of state dependent. Observing the minimum period in LMXB may provide useful constraints on the equation of state for high density matter. For $`\mathrm{\Delta }M<1M_{}`$, the minimum period is longer than 1.7ms for a wide range of realistic equations of state (Wiringa, Fiks & Fabrocini 1988).
We thank A.Potekhin for his useful comments and P.K.MacKeown for a critical reading of our manuscript. This research program is partially supported by a RGC grant of the Hong Kong Government, a Croucher Foundation Senior Research Fellowship and a NSF grant of PRC.
REFERENCES
Bailes, M. et al 1994, ApJL,425, 41
Bell, J.F. et al 1997, MNRAS,286, 463
Bhattacharya, D. & van den Heuvel, E.P.J. 1991, Phys. Rep., 203,1
Bhattacharya, D. & Srinivasan, G. 1995, in X-ray Binaries, eds. W.H.G. Lewin, J. van Paradijs & E.P.J. van den Heuvel,(Cambridge University Press).
Bisnovatyi-Kogan, G. & Komberg, B.1974, Soviet Astron,18,217
Camilo, F. et al 1996, APJ, 469, 819
Chen, K. & Ruderman, M.A., 1993, ApJ, 402,264
Cheng, K.S.& Dai, Z.G.1997, ApJ, 476, L39
Cheng, K.S.& Zhang, C.M. 1998, A&A, 337, 441
Chanmugam, G. 1992, Ann. Rev. Astron. Astrophys., 30,143
Corbet, et al 1995,APJ,443,786
Ding, K.Y.,Cheng, K.S.& Chau,H.F.1993,ApJ,408,167
Ghosh, P. & Lamb, F.K., 1979, ApJ, 234, 296
Ghosh, P. & Lamb, F.K., 1991 in ”Neutron Stars : Theory and Observation” eds J. Ventura & D. Pines (Kluwer:Netherlands), p.363
Lorimer, D.R., et al 1996, MNRAS, 283, 1383
Makishima, K. 1992, in ”The structure and evolution of neutron stars”, eds D. Pines, R. Tamagaki & S. Tsuruta (Addison-Wesley), p.86
Phinney, E.S. & Kulkarni, S.R. 1994, Ann. Rev. Astron. Astrophys., 32, 591.
Romani, G.W. 1990,Nature,347,741
Ruderman, M. 1991a,ApJ,366,261
Ruderman, M. 1991b,ApJ,382,576
Ruderman, M. 1991c,ApJ,382,587
Ruderman, M., Zhu, T. & Chen, K. 1998,ApJ,492, 267
Shibazaki, N., Murakami, T., Shaham, J. & Nomoto, K. 1989, Nature,342,656
Taam, R.E. & van den Heuvel, E.P.J. 1986, ApJ, 305, 235
Taylor, J.H., Manchester, R.N. & Lyne, A.G. 1993, APJS,88,529
Urpin, V. & Geppert, U. 1995, MNRAS,275,1117
Urpin, V. & Konenkov, D. 1997,MNRAS,284,741
van den Heuvel, E.P.J., van Paradijs, J.A. & Taam, R.E.1986,Nature,322,153
van den Heuvel, E.P.J. ,Bizaraki, O .1995a,A&A,297,L41
van den Heuvel, E.P.J. ,Bizaraki, O. 1995b,in The Lives of the Neutron Stars, eds. M.A. Alpar, $`\ddot{\mathrm{U}}`$. Kilzilo$`\stackrel{ˇ}{\mathrm{g}}`$lu and J. van Paradijs, (Kluwer Academic Publishers, Dordrecht).
White, N. & Zhang, W. 1997, APJ, 490, L87
Wiringa, R.B., Fiks, V. & Fabrocini, A., 1988, Phys. Rev. C38, 1010
Zhang, C.M., Wu, X.J. & Yang, G.C.1994,A&A,283,889
Zhang, C.M., 1998, A&A, 330, 195
Zhang, C.M., Cheng, K.S., & Zhang, L. 1997, in ”21st Centuary Chinese Astronomy Conference”, eds. K.S.Cheng & S.Chan, (World Scientific, Singapore), p301
|
warning/0005/hep-th0005241.html
|
ar5iv
|
text
|
# 1 Space dependence of the external fields.
hep-th/0005241
D0-Branes As Confined Quarks <sup>1</sup><sup>1</sup>1Talk presented at Isfahan String Workshop 2000, May 13-14, IRAN.
Amir H. Fatollahi
Institute for Advanced Studies in Basic Sciences (IASBS),
P.O.Box 45195-159, Zanjan, IRAN
and
Institute for Studies in Theoretical Physics and Mathematics (IPM),
P.O.Box 19395-5531, Tehran, IRAN
fath@theory.ipm.ac.ir
## Abstract
The possibility of using the quantum mechanics of D0-branes for the bound-states of quarks and QCD strings is investigated. Issues such as the inter D0-branes potential, the whiteness of the D0-branes bound-states and the large-$`N`$ limit of D0-branes effective theory are studied. A possible role of the non-commutativity of relative distances of D0-branes in a study of ordinary QCD is discussed.
D0-branes are defined as particles which strings end on them . The question is: Can one use D0-branes dynamics as the effective theory of bound-states of quarks and QCD strings (QCD electric fluxes)? We study the following issues to approach this question:
* Inter D0-branes potential to compare with phenomenological one and that of electric-flux picture.
* Whiteness of D0-branes bound-states under $`SU(N)`$ electric field.
* Large-$`N`$ behaviour of D0-branes bound-states to compare with QCD baryonic states at large-$`N`$.
Also a possible role of involving non-commutativity like the same one in relative distances of D0-branes in a study of ordinary QCD is discussed. Discussions here are mostly coming from the results appeared in .
Dynamics of $`N`$ D0-branes is given by the matrix quantum mechanics resulted from dimensional reduction of $`U(N)`$ gauge theory to 0+1 dimension, by replacements $`A_iX_i`$ : <sup>2</sup><sup>2</sup>2Here we ignore supersymmetry. Also we work in arbitrary dimensions $`d`$.
$`S={\displaystyle }dtm_0\mathrm{Tr}({\displaystyle \frac{1}{2}}D_tX_i^2+{\displaystyle \frac{[X_i,X_j]^2}{4(2\pi \alpha ^{})^2}}),`$
$`i,j=1,\mathrm{},d,D_t=_ti[a_0,],`$
with $`\frac{1}{2\pi \alpha ^{}}`$=string tension ($`l_s=\sqrt{\alpha ^{}}`$ and $`g_s`$=string coupling). $`X`$’s are in the algebra by the usual expansion $`X_i=x_{i(a)}T_{(a)}`$, $`(a)=1,\mathrm{},N^2`$. <sup>3</sup><sup>3</sup>3To avoid confusion, we put group indices always in ( ).
The action is invariant under gauge transformations:
$`\stackrel{}{X}`$ $``$ $`\stackrel{}{X^{}}=U\stackrel{}{X}U^{},`$
$`a_0`$ $``$ $`a_0^{}=Ua_0U^{}+iU_tU^{},`$ (2)
with $`U`$ as arbitrary unitary matrix, and consequently one finds:
$`D_t\stackrel{}{X}`$ $``$ $`D_t^{}\stackrel{}{X^{}}=U(D_t\stackrel{}{X})U^{},`$
$`D_tD_t\stackrel{}{X}`$ $``$ $`D_t^{}D_t^{}\stackrel{}{X^{}}=U(D_tD_t\stackrel{}{X})U^{}.`$ (3)
D0-branes are presented “classically” by diagonal matrices and the action takes the form of $`N`$ free particles for them:
$`S={\displaystyle 𝑑t\underset{(a)=1}{\overset{N}{}}\frac{1}{2}m_0\dot{\stackrel{}{x}}_{(a)}^2}.`$ (4)
The action is non-relativistic, but can be used for covariant formulation by Light-Cone Frame (LCF) interpretation with the following identifications :
$`m_0=p^+,t=x^+,X_i=\mathrm{transverse}\mathrm{directions}.`$ (5)
By the scalings
$`tg_s^{1/3}t,a_0g_s^{1/3}a_0,Xg_s^{1/3}X,`$ (6)
one finds the relevant energy and size scales as:
$`Eg_s^{1/3}/l_s,l_{d+2}=g_s^{1/3}l_s.`$ (7)
The length $`l_{d+2}`$ should be identified as the fundamental length scale of the covariant $`d+2`$ dimensional theory which is expected that its LCF formulation is presented by the action (1). So we take (for $`d=2`$) $`l_{d+2}`$ as the inverse of 4 dimensional QCD mass scale, denoted by $`\mathrm{\Lambda }_{QCD}`$.
1. D0-Branes Potential: The effective potential between D0-branes comes from the effect of quantum fluctuations around a classical configuration, presented here by diagonal matrices. This work is equivalent with integrating over oscillations of strings stretched between D0-branes. One-loop effective action is given by ($`2\pi \alpha ^{}=1`$):
$`({\displaystyle 𝑑t})V(X_\mu ^{cl})={\displaystyle \frac{1}{2}}\mathrm{Tr}\mathrm{log}\left(P_\lambda ^2\delta _{\mu \nu }2iF_{\mu \nu }\right)\mathrm{Tr}\mathrm{log}\left(P_\lambda ^2\right),`$ (8)
with
$`P_\mu [X_\mu ^{cl},],F_{\mu \nu }[f_{\mu \nu },],f_{\mu \nu }[X_\mu ^{cl},X_\nu ^{cl}],`$
$`\mu ,\nu =0,1,\mathrm{},d,X_0=i_t+a_0,`$
$`P_\lambda ^2=_t^2+{\displaystyle \underset{i=1}{\overset{d}{}}}P_i^2,\mathrm{for}a_0^{cl}=0.`$ (9)
For two static D0-branes at distance $`r`$ we may take:
$`X_1^{cl}={\displaystyle \frac{1}{2}}\left(\begin{array}{cc}r& 0\\ 0& r\end{array}\right),X_0^{cl}=i_t\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),`$ (10)
$`a_0^{cl}=X_{i>1}^{cl}=0.`$
So one finds
$`P_1={\displaystyle \frac{r}{2}}\mathrm{\Sigma }_3,P_0=i_t1_4,P_{i>1}=0,`$ (11)
with $`\mathrm{\Sigma }_3=[\sigma _3,]`$ which has 0, 0, $`\pm 2`$ as eigenvalues. Also we find the operator $`P_\lambda ^2=_t^21_4+\frac{r^2}{4}\mathrm{\Sigma }_3^2`$ as a harmonic oscillator operator with frequency $`\omega r/\alpha ^{}`$. One-loop is a good approximation for $`\omega m_0\dot{r}^2`$ or $`rg_sl_s\dot{r}^2`$ which for $`g_s0`$ ($`m_0l_s^1`$) is satisfied for large separations and low velocities.
One-loop effective action can be calculated easily to find:
$`V(r)`$ $`=`$ $`({\displaystyle \frac{d1}{2}})\mathrm{Tr}\mathrm{log}\left(P_\lambda ^2\right)`$ (12)
$`=`$ $`\mathrm{\hspace{0.33em}2}({\displaystyle \frac{d1}{2}}){\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{ds}{s}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑k_0\mathrm{e}^{s(k_0^2+r^2)}`$
$`+\mathrm{traces}\mathrm{independent}\mathrm{of}r,`$
and after integrations one obtains
$`V(r)`$ $`=`$ $`\mathrm{\hspace{0.33em}2}({\displaystyle \frac{d1}{2}}){\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{ds}{s}}({\displaystyle \frac{\pi }{s}})^{\frac{1}{2}}\mathrm{e}^{sr^2}`$ (13)
$`=`$ $`\mathrm{\hspace{0.33em}4}\pi ({\displaystyle \frac{d1}{2}})|r|`$
$`\mathrm{}(\mathrm{independent}\mathrm{of}r).`$
which is the linear potential with phenomenology interests . So the effective theory for the relative dynamics of two D0-branes is given by:
$`S={\displaystyle 𝑑t\left(\frac{1}{2}\frac{m_0}{2}\dot{\stackrel{}{r}}^24\pi (\frac{d1}{2})\frac{|\stackrel{}{r}|}{2\pi \alpha ^{}}\right)}`$ (14)
and one finds the energy scale as $`E\alpha ^{2/3}m_0^{1/3}g_s^{1/3}/l_s`$, as pointed in eq. (7). By assuming the dynamics in LCF with the longitudinal momentum $`m_0`$, we have $`M^2p^+p^{}m_0Eg_s^{2/3}l_s^2l_{d+2}^2`$, by eq. (7).
This potential is also true for every pair of D0-branes, and one can write the effective theory for $`N`$ D0-branes as:
$`S={\displaystyle 𝑑t\left(\frac{1}{2}m_0\underset{(a)=1}{\overset{N}{}}\dot{\stackrel{}{r}}_{(a)}^24\pi (\frac{d1}{2})\underset{(a)>(b)=1}{\overset{N}{}}\frac{|\stackrel{}{r}_{(a)}\stackrel{}{r}_{(b)}|}{2\pi \alpha ^{}}\right)}.`$ (15)
In a recent work by taking the linear potential in transverse directions of LCF between the quarks of a baryonic state, the structure functions are obtained with a good agreement with observed ones.
One can relate the parameter $`1/\alpha ^{}`$ in the front of the potential to gauge theory parameters. To do so one needs a string theoretic description of the gauge theory in LCF, and the natural guess for this is “Light-Cone–lattice gauge theory” (LClgt) . In LClgt one assumes time direction and one of the spatial directions to be continuous to define LC variables $`x^\pm t\pm z`$. Other spatial directions play the role of transverse directions of LCF which are assumed to be lattices. As usual in LCF, time is $`x^+`$ and continuous and so we have a Hamiltonian formulation of lgt . The linear potential in LClgt, related to string tension is known to be:
$`V(r){\displaystyle \frac{g_{_{YM}}^2}{a^2}}|\stackrel{}{r}|,`$ (16)
with $`a`$ as the lattice spacing parameter in the transverse directions. Via this one finds the relation
$`{\displaystyle \frac{1}{\alpha ^{}}}{\displaystyle \frac{g_{_{YM}}^2}{a^2}},`$ (17)
for the parameters.
2. Whiteness: To find the charge and colour of D0-branes bound-states we need to know their dynamics in YM backgrounds. In the case of electromagnetism there is a simple relation:
$`m_0\ddot{\stackrel{}{x}}=q(\stackrel{}{E}_{ext.}+\stackrel{}{v}\times \stackrel{}{B}_{ext.}).`$ (18)
The concept of gauge invariance here is understood as the invariance of the equations of motion under the gauge symmetry transformations. In the case of chromodynamics in r.h.s. matrices in adjoint representation are placed and so they transform like:
$`\stackrel{}{E}\stackrel{}{E}^{}=U\stackrel{}{E}U^{},\stackrel{}{B}\stackrel{}{B}^{}=U\stackrel{}{B}U^{}.`$ (19)
So we need to replace the l.h.s. with matrices with the correct behaviour under gauge transformations. Now we have good candidate for non-commutative coordinates: D0-branes coordinates. One may write for “matrix” coordinates
$`m_0\ddot{\stackrel{}{X}}=q(\stackrel{}{E}_{ext.}+\dot{\stackrel{}{X}}\times \stackrel{}{B}_{ext.}),`$ (20)
but yet the l.h.s. does not have correct behaviour under gauge transformations! Here the world-line gauge symmetry eq. (S0.Ex2) helps us, to write the generalized “Lorentz” equation as
$`m_0D_tD_t\stackrel{}{X}=q(\stackrel{}{E}_{ext.}+D_t\stackrel{}{X}\times \stackrel{}{B}_{ext.}),`$ (21)
and now by eq. (3) both sides have an equal behaviour under gauge transformations. The space dependence of the fields is a subtle point, because the coordinates themselves change under transformation on the gauge fields $`A_{(a)}`$’s. Resolving the space dependence may be done by assuming the D0-branes bound-states very small and then taking the space dependence of external fields just for the centre-of mass (c.m.) (see fig. below).
The coordinates and momenta of c.m. are given by the trace of matrices, as:
$`\stackrel{}{x}_{cm}{\displaystyle \frac{1}{N}}\mathrm{Tr}\stackrel{}{X},\stackrel{}{p}_{cm}\mathrm{Tr}\stackrel{}{P},`$ (22)
and so are invariant under transformations $`XUXU^{}`$. So by taking the space dependence just for c.m. we have:
$`\stackrel{}{E}_{ext}=\stackrel{}{E}_{ext}(x_{cm}),\stackrel{}{B}_{ext}=\stackrel{}{B}_{ext}(x_{cm}).`$ (23)
To specify the charge or colour of an extended object (e.g. a bound-state), we study the dynamics in absence of magnetic field ($`\stackrel{}{B}=0`$) and in uniform electric field ($`\stackrel{}{E}(x)=\stackrel{}{E}_0`$). In our case the c.m. dynamics decouples from non-Abelian parts due to the trace nature of $`U(1)`$ and $`SU(N)`$ parts. So we have:
$`m_0\ddot{\stackrel{}{x}}_{c.m.}=q\stackrel{}{E}_{(1)0ext.},`$ (24)
which (1) is for $`U(1)`$ part of $`U(N)`$. So the dynamics of c.m. will not be affected by the non-Abelian part: the c.m. is white. It means that each D0-brane sees the net effect of other D0-branes as the white-complement of its colour: the field fluxes extracted from one D0-brane to other ones are as the same of one flux between a colour and an anti-colour. The linear potential of previous part is consistent with flux-string picture. The number of D0-branes in the bound-state is equal to the same of baryons, $`N`$.
3. Large-$`N`$: Baryons show special properties at large-$`N`$ limit of gauge theories :
* Their mass grows linearly by $`N`$.
* Their size is not dependent on $`N`$. So their density goes to infinity at large-$`N`$.
* Baryon-baryon force grows in proportion to $`N`$.
These properties are mainly extracted from the study of quantum mechanics of $`N`$ quarks and their bound-states as an $`N`$-body problem. The problem is approached by approximations (e.g. Hartree) for general potentials which have two characters: 1) they are attractive, and 2) their strength decreases by 1/$`N`$ at large-$`N`$. Here we check the same behaviours for our problem, by reminding LCF interpretations. The effective theory for $`N`$ D0-branes is obtained to be
$`S={\displaystyle 𝑑t\left(\frac{1}{2}m_0\underset{(a)=1}{\overset{N}{}}\dot{\stackrel{}{r}}_{(a)}^24\pi (\frac{d1}{2})\underset{(a)>(b)=1}{\overset{N}{}}\frac{|\stackrel{}{r}_{(a)}\stackrel{}{r}_{(b)}|}{2\pi \alpha ^{}}\right)}.`$ (25)
with the relation $`1/\alpha ^{}g_{_{YM}}^2/a^2`$. Also we have the replacement at large-$`N`$:
$`g_{_{YM}}{}_{}{}^{^{largeN}}{\displaystyle \frac{g_{_{YM}}}{\sqrt{N}}},`$ (26)
and so the action is read
$`S={\displaystyle 𝑑t\left(\frac{1}{2}m_0\underset{(a)=1}{\overset{N}{}}\dot{\stackrel{}{r}}_{(a)}^24\pi (\frac{d1}{2})\frac{g_{_{YM}}^2}{a^2}\frac{1}{N}\underset{(a),(b)=1}{\overset{N}{}}|\stackrel{}{r}_{(a)}\stackrel{}{r}_{(b)}|\right)}.`$ (27)
The associated Hamiltonian of this action is the same used before except for the potential term, which is Coulomb one there. Here we just check the mass: The kinetic term of c.m. ($`\frac{\stackrel{}{P}^2}{Nm_0}`$) grows with $`N`$, and the net potential for each D0-brane takes a factor $`\frac{1}{2}N(N1)`$ due to pair interactions. So the potential term grows with $`\frac{1}{2}N(N1)g_{_{YM}}^2/NN`$. The energy grows as $`EN`$ at large-$`N`$. In LCF the energy is $`P^{}`$. Also the total longitudinal momentum of this bound-state is $`P^+=Np^+`$ with $`p^+=m_0`$. So the invariant mass $`M`$ is read
$`M^2=2P^+P^{}\stackrel{}{P}^2N^2MN.`$ (28)
Space-Time Considerations: Non-Commutativity
Relative coordinates of D0-branes are matrices and so non-commutative. If the correspondence between the dynamics of D0-branes and confined quarks has a root in Nature, the question will be about possible justification of this non-commutativity. In the following 3 comments are in order:
1. Special Relativity Idea: In the way to find a consistent theory for the propagation of electromagnetic fields, special relativity learns to us that space and time should be treated as a 4-vector $`X_\mu `$ under boost transformations, such as the gauge field 4-vector, $`A_\mu `$.
Also the idea of supersymmetry (SUSY) can be considered as a natural continuation of the special relativity program: Adding spin half sector to the coordinates of space-time as the representative of the fermions of the Nature. This idea leads one to the super-space formulation of SUSY theories. Also it is the same way which one introduces fermions to the bosonic string theory.
Now, what may be modified if in some regions of space and time there exists non-Abelian (non-commutative) gauge fields? In the present Nature non-Abelian gauge fields can not make spatially long coherent states; they are confined or too heavy. But the picture may be changed inside a hadron. In fact recent developments of string theories sound this change and it is understood that non-commutative coordinates and non-Abelian gauge fields are two sides of one coin. As we discussed, the interaction between D-branes is the result of path-integrations over fluctuations of the non-commutative parts of coordinates. It means that in this picture “non-commutative” fluctuations of space-time are the source of “non-Abelian” interactions. One may summarize this discussion as in the table below:
$`\mathrm{Field}`$ $``$ $`\mathrm{Space}\mathrm{Time}`$
$`A_\mu (\mathrm{Photons})`$ $``$ $`X_\mu (4\mathrm{Vector}):\mathrm{Maxwell}`$
$`\psi (\mathrm{Fermions})`$ $``$ $`\theta ,\overline{\theta }(\mathrm{Super}\mathrm{Coordinates}):\mathrm{SUSY}`$
$`A_\mu ^{(a)}(\mathrm{Gluons})`$ $``$ $`X_\mu ^{(a)}(\mathrm{Matrix}\mathrm{Coordinates}):\mathrm{QCD}`$
As it has been mentioned previously, the non-commutativity of D0-branes coordinates just come back to their relative distances and the c.m.’s of different bound-states of D0-branes presented by the trace of the position matrices, are commutative objects. We know that QCD fields are zero outside of hadrons, so the non-commutativity should be restricted to relative distances of hadron constituents (see fig.).
2. Recent Example: Pure U(1) gauge theory on ordinary space has free photons. On non-commutative space the theory has interacting photons and the structure of the theory becomes very similar to the same of non-Abelian gauge theory, summarized in the table below :
$`\mathrm{Commutative}\mathrm{Space}`$ $``$ $`\mathrm{Non}\mathrm{Commutative}\mathrm{Space}`$
$`[X,X]=0`$ $``$ $`[X,X]=i\theta `$
$`\mathrm{free}\mathrm{photons}`$ $``$ $`\mathrm{interacting}\mathrm{photons}`$
$`F=A`$ $``$ $`F=A+g\{A,A\}`$
$`F^{}=F`$ $``$ $`F^{}=UFU^1.`$ (29)
The lesson of this example is that one may cover the aspects of non-Abelian gauge theories by changing the structure of space-time. It means that by assuming non-commutativity between the coordinates of space-time one can get a theory with properties similar to non-Abelian theories.
For our special case the question will be about “Is the structure of space-time suggested by D0-brane–quark correspondence appropriate to cover the non-Abelian structure of $`U(N)`$ gauge theories or QCD?”
3. Lattice Continuum Limit: Firstly let us have a look to the procedure of taking the continuum limit of lattice gauge theories <sup>4</sup><sup>4</sup>4This discussion is borrowed from .. Consider the correlation length $`\xi `$ between two plaquettes:
$`\xi =a\xi _{latt}`$ (30)
which is expressed in terms of dimensionless parameter $`\xi _{latt}`$, and lattice spacing parameter $`a`$, appearing just as a scale factor. $`\xi `$ is the physical quantity and in the continuum limit it should remain constant, providing:
$`\xi _{latt}_{_{_{a0}}}\mathrm{}.`$ (31)
The correlation function has the behaviour:
$`G(r)___r\mathrm{}e^{r/\xi },`$ (32)
and by setting $`r=na`$, the $`\xi _{latt}`$ does not depends on $`a`$, becoming a function of coupling constant $`g`$:
$`\xi _{latt}=\underset{n\mathrm{}}{lim}{\displaystyle \frac{n}{\mathrm{ln}G(n)}}=f(g),`$ (33)
which we have for it at a critical value $`g_c`$:
$`f(g_c)=\mathrm{}.`$ (34)
One can find the $`g`$-dependence of any physical quantity by the function $`f(g)`$. Assume $`Q`$ is a physical quantity with dimension \[length\]<sup>d</sup>; putting $`Q=a^dQ_{latt}`$ we have:
$`Q\xi ^d=Q_{latt}\xi _{latt}^d.`$ (35)
Assuming $`Q\xi ^d`$ is finite in the continuum limit we have:
$`Q_{latt}=C[f(g)]^d,`$ (36)
with $`C`$ as a constant.
One can find the behaviour of the function $`f(g)`$ at the strong and weak coupling limit. At strong coupling limit lattice gauge theory gives the string tension $`K`$ with dimension \[length\]<sup>-2</sup>, so:
$`f(g)={\displaystyle \frac{1}{(a^2K)^2}}={\displaystyle \frac{C}{\mathrm{ln}^2(\kappa g)}},g1,`$ (37)
At weak coupling we have the perturbative result as:
$`{\displaystyle \frac{g^2}{4\pi }}={\displaystyle \frac{1}{\gamma _0\mathrm{ln}(M^2/\mathrm{\Lambda }^2)}},\gamma _0={\displaystyle \frac{33}{12\pi }},`$ (38)
where $`M`$ is a mass scale. So we have:
$`f(g)=C^{}\mathrm{e}^{[1/(2\gamma _0g^2)]}=C^{}\mathrm{e}^{[8\pi ^2/(11g^2)]},g1.`$ (39)
These two behaviours are plotted in the figure, and one can see that the continuum limit ($`f(g)=\mathrm{}`$) is just gained at $`g_c=0`$. Based on this, for every finite value of the coupling constant lattice formulation does not reach to continuum limit.
On the other hand, we know that the natural framework of formulation theories on discrete space-time is Non-Commutative Geometry, with the known examples two-point world or the standard model of particles on two-sheet world .
|
Subsets and Splits
No community queries yet
The top public SQL queries from the community will appear here once available.