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# ∗-Ideals and Formal Morita Equivalence of ∗-Algebras
## 1 Introduction
The concept of formal deformation quantization, first introduced in , has been used successfully to construct and classify the quantum observable algebras of physical systems in terms of classical data, generally given by a symplectic or Poisson manifold (see for existence and classifications and for recent reviews). One is naturally led to investigate representations of these algebras, as different representations can be used to describe different physical situations. It is then important to decide what kind of representations (or category of modules) of the observable algebras have physical meaning and should be considered. It turns out that the usual -representation theory of $`C^{}`$-algebras provides a guide for the formulation of representations of quantum observable algebras arising in deformation quantization. In particular, one can always consider star-products compatible with the natural -involution on functions (given by complex conjugation), and we note that the underlying ring of real formal power series $`[[\mathrm{}]]`$ has a natural order structure. Combining these two facts, one can define positive $`[[\mathrm{}]]`$-linear functionals on the star-product algebras and -representations on pre-Hilbert spaces over $`[[\mathrm{}]]`$ analogous to the $`C^{}`$-algebra case. For example, every positive linear functional gives rise to a GNS construction of a -representation (see for a detailed exposition of these ideas). These notions have yielded physically interesting representations and, in fact, most known operator representations are of this type, see e.g. .
With this motivation, we started in a more systematic study of -algebras over $`𝖢=𝖱(\mathrm{i})`$, where $`𝖱`$ is an arbitrary (commutative, unital) ordered ring and $`𝖢`$ is its quadratic extension by $`\mathrm{i}`$, with $`\mathrm{i}^2=1`$, and their representations on pre-Hilbert spaces over $`𝖢`$. We note that not only star-product algebras over $`[[\mathrm{}]]`$ but also arbitrary -algebras over $``$ (such as -algebras of unbounded operators, see e.g. ) fit into this framework. In this setting, one can consider a purely algebraic version of Rieffel induction and define the notion of formal Morita equivalence , analogous to $`C^{}`$-algebras . In the present paper, we continue the investigation of -representations of -algebras and formal Morita equivalence.
First, we clarify the relationship between formal Morita equivalence, Ara’s notion of Morita -equivalence of rings with involution and strong Morita equivalence of $`C^{}`$-algebras . For non-degenerate and idempotent -algebras over $`𝖢`$, we show that formal Morita equivalence implies Morita -equivalence. In the case where the -algebras are given by Pedersen ideals of $`C^{}`$-algebras , we show that, in fact, formal Morita equivalence and Morita -equivalence are equivalent. It then follows immediately from Ara’s recent algebraic characterization of strong Morita equivalence of $`C^{}`$-algebras in terms of their Pedersen ideals \[1, Thm. 2.4\] that two $`C^{}`$-algebras are strongly Morita equivalent if and only if their Pedersen ideals are formally Morita equivalent.
Second, we discuss -ideals in -algebras over $`𝖢`$, their behavior under formal Morita equivalence and how they can be used to describe particular properties of star-product algebras which general -algebras over $`𝖢`$ may lack. In a -algebra over $`𝖢`$, we define an ideal to be ‘closed’ if it occurs as the kernel of some -representation on a pre -Hilbert space over $`𝖢`$. These ideals form a lattice, which we show is invariant under formal Morita equivalence. In the case of $`C^{}`$-algebras, an ideal is closed in this algebraic sense if and only if it is norm closed and the aforementioned invariance reduces to the so-called Rieffel correspondence theorem (see \[26, Prop. 3.24\]). We note that -algebras having a “large” number of positive linear functionals possess nice features which do not hold in general. For instance, they admit faithful representations on pre-Hilbert spaces (see \[9, Prop. 2.8\]). We observe that not only $`C^{}`$-algebras but also star-product algebras arising in deformation quantization have this property \[10, Prop. 5.3\]. Admitting a faithful representation is equivalent to the triviality of the minimal closed ideal (defined as the intersection of all closed -ideals), which we also show is a formal Morita invariant. We use this invariant to describe a large class of examples of -algebras having equivalent categories of (strongly non-degenerate) -representations on pre-Hilbert spaces but which are not formally Morita equivalent, generalizing \[9, Cor. 5.20\]. (It was shown in \[9, Thm. 5.10\] that formal Morita equivalence implies the equivalence of categories of strongly non-degenerate -representations.)
Finally, we apply these general algebraic notions to the particular case of the -algebra $`C^{\mathrm{}}(M)`$ of complex-valued smooth functions on a manifold $`M`$. We show that closed -ideals in this example are given by the ideals of functions vanishing on closed subsets of $`M`$, and hence agree with the closed ideals with respect the (weak) $`C^0`$-topology (\[35, Thm. 3\]). We also show that $`C^{\mathrm{}}(M)`$ is formally Morita equivalent to the -algebra of smooth sections of the endomorphism bundle of any Hermitian vector bundle $`E`$ over $`M`$. These facts are not surprising since their analogues in the continuous ($`C^{}`$-algebraic) category (in the sense of strong Morita equivalence) are well-known. They are, however, conceptually important since they show that -algebras in the smooth category can be treated independently through general algebraic methods (we note that $`M`$ is *not* assumed to be compact and hence the -algebras in questions are *not* pre-$`C^{}`$-algebras). We point out that this discussion is the ‘classical’ starting point in the investigation of these ideas in the framework of formal deformation quantization, see also for the case of a trivial bundle.
The paper is organized as follows. In Section 2, we recall the notions of ordered rings, -algebras, and their -representations, including the GNS construction (details can be found in ). In Section 3, we describe how the notions of formal Morita equivalence, Morita -equivalence and strong Morita equivalence are related. We define closed -ideals in Section 4 and discuss basic properties of the minimal closed -ideal and its relation to GNS representations and positive functionals. In Section 5, we focus on algebraic Rieffel induction and formal Morita equivalence of -algebras and prove that the lattice of closed -ideals is a formal Morita invariant. Section 6 contains the description of closed -ideals of $`C^{\mathrm{}}(M)`$. As a -algebra with non-trivial minimal ideal we consider the sections of the Grassmann bundle over $`M`$. We also show in this section that for any Hermitian vector bundle $`E`$ over a manifold $`M`$, $`C^{\mathrm{}}(M)`$ and $`\mathrm{\Gamma }^{\mathrm{}}(\mathrm{𝖤𝗇𝖽}(E))`$ are formally Morita equivalent, as well as $`C_0^{\mathrm{}}(M)`$ and $`\mathrm{\Gamma }_0^{\mathrm{}}(\mathrm{𝖤𝗇𝖽}(E))`$.
## 2 Preliminaries
In this section, we just recall some basic facts on ordered rings, pre-Hilbert spaces and -algebras to establish our notation (see for details).
Let $`𝖱`$ be an associative commutative ring with $`10`$ which is *ordered*, i.e. equipped with a distinguished subset $`𝖯`$, the positive elements, such that $`𝖱`$ is the disjoint union $`𝖯\{0\}𝖯`$ and one has $`𝖯+𝖯,𝖯𝖯𝖯`$. We write $`a>b`$ if $`ab𝖯`$ etc. In particular $`1>0`$, and $`𝖱`$ has characteristic zero and no zero divisors. Next we consider the canonical quadratic ring extension $`𝖢=𝖱(\mathrm{i})`$ where $`\mathrm{i}^2=1`$. Then $`𝖢`$ becomes an associative commutative ring with $`10`$ which has again characteristic zero and no zero divisors. Complex conjugation $`z\overline{z}`$ is defined as usual for $`z𝖢`$ and the embedding $`𝖱𝖢`$ determines the real elements of $`𝖢`$, i.e. those with $`z=\overline{z}`$. Finally, note that $`\overline{z}z0`$ and $`\overline{z}z=0`$ if and only if $`z=0`$. Examples of ordered rings are of course $``$, $``$, and the formal power series $`[[\lambda ]]`$, which is the important example for formal deformation quantization.
A *pre-Hilbert space* over $`𝖢`$ is a $`𝖢`$-module $``$ with positive-definite Hermitian product $`,:\times 𝖢`$ which is linear in the second argument (physicists’ convention). We immediately have the Cauchy-Schwarz inequality $`\varphi ,\psi \psi ,\varphi \varphi ,\varphi \psi ,\psi `$ for all $`\varphi ,\psi `$. For a pre-Hilbert space $``$, we say that $`A\mathrm{𝖤𝗇𝖽}_𝖢()`$ has a (necessarily unique) adjoint $`A^{}\mathrm{𝖤𝗇𝖽}_𝖢()`$ if for all $`\varphi ,\psi `$ the equality $`A^{}\varphi ,\psi =\varphi ,A\psi `$ holds. One defines $`𝔅():=\{A\mathrm{𝖤𝗇𝖽}_𝖢|A^{}`$ exists $`\}`$ and easily finds that $`𝔅()`$ is an associative algebra over $`𝖢`$ with unit element $`\mathrm{𝗂𝖽}_{}`$ and a -involution $`AA^{}`$. More generally, one defines a *-algebra* $`𝒜`$ to be an associative algebra over $`𝖢`$ equipped with a -involution, i.e. an antilinear involutive anti-automorphism. Hermitian, normal, isometric and unitary elements in $`𝒜`$ are defined as usual. An *approximate identity* consists of a filtration $`𝒜=_{\alpha I}𝒜_\alpha `$ by $`𝖢`$-submodules $`\{𝒜_\alpha \}_{\alpha I}`$ indexed by a directed set $`I`$ and elements $`\{E_\alpha \}_{\alpha I}`$ such that $`E_\alpha ^{}=E_\alpha `$, $`E_\alpha E_\beta =E_\alpha =E_\beta E_\alpha `$ for $`\alpha <\beta `$ and $`E_\alpha A=A=AE_\alpha `$ for all $`A𝒜_\alpha `$ and $`\alpha I`$. A linear functional $`\omega :𝒜𝖢`$ is called *positive* if $`\omega (A^{}A)0`$ for all $`A𝒜`$. Then one has the Cauchy-Schwarz inequality $`\omega (A^{}B)\overline{\omega (A^{}B)}\omega (A^{}A)\omega (B^{}B)`$ as well as the ‘almost reality’ property $`\omega (A^{}B)=\overline{\omega (B^{}A)}`$, which implies reality $`\omega (A^{})=\overline{\omega (A)}`$ if e.g. $`𝒜`$ has an approximate identity. Using positive functionals, we are able to define positive algebra elements as well: by $`𝒜^{++}`$ we denote the *algebraically positive elements* of the form $`A=a_1A_1^{}A_1+\mathrm{}+a_nA_n^{}A_n`$ where $`a_i>0`$ and $`A_i𝒜`$. By $`𝒜^+`$ we denote the *positive* elements of $`𝒜`$, where $`A𝒜^+`$ if $`A`$ is Hermitian and $`\omega (A)0`$ for all positive linear functionals $`\omega :𝒜𝖢`$. We note that if $`𝒜`$ is a $`C^{}`$-algebra, then these notions of positivity agree with the standard ones. Moreover, in this case $`𝒜^+=𝒜^{++}`$. These notions also yield the expected results for $`C^{\mathrm{}}(M)`$, see \[9, App. B\].
A *-representation* of a -algebra $`𝒜`$ on a pre-Hilbert space $``$ is a -homomorphism $`\pi :𝒜𝔅()`$. Such a -representation is called *non-degenerate* if $`\pi (A)\varphi =0`$ for all $`A𝒜`$ implies $`\varphi =0`$ and *strongly non-degenerate* if the $`𝖢`$-span of all $`\pi (A)\varphi `$ with $`A𝒜`$ and $`\varphi `$ coincides with $``$. In general, strong non-degeneracy implies non-degeneracy and if $`𝒜`$ has a unit both notions agree and are equivalent to the condition $`\pi (\mathrm{𝟣})=\mathrm{𝗂𝖽}_{}`$.
If $`\omega :𝒜𝖢`$ is a positive linear functional, then $`𝒥_\omega :=\{A𝒜|\omega (A^{}A)=0\}`$ is a left ideal in $`𝒜`$, the so-called Gel’fand ideal and hence $`_\omega :=𝒜/𝒥_\omega `$ is a $`𝒜`$-left module. Moreover, $`_\omega `$ becomes a pre-Hilbert space via $`\psi _B,\psi _C:=\omega (B^{}C)`$ such that the left action $`\pi _\omega (A)\psi _B:=\psi _{AB}`$ is a -representation, the so-called GNS representation. Here $`\psi _B_\omega `$ denotes the equivalence class of $`B𝒜`$. This GNS construction will play a crucial role in this paper (see for an extensive treatment).
Let us recall the notion of algebraic Rieffel induction, as discussed in . Let $`𝒜`$ be a -algebra over $`𝖢`$ and let $`𝔛`$ be a right $`𝒜`$-module (all modules over -algebras considered here will be assumed to have a compatible $`𝖢`$-module structure). An *$`𝒜`$-valued inner product* is a map $`,_𝒜:𝔛\times 𝔛𝒜`$, which is $`𝖢`$-linear and $`𝒜`$-right linear in the second argument and satisfies $`x,y_𝒜=y,x_𝒜^{}`$, for all $`x,y𝔛`$. We call this inner product *positive semi-definite* if $`x,x_𝒜𝒜^+`$ and *positive definite* if in addition $`x,x_𝒜=0`$ implies $`x=0`$. We say it is *full* if the $`𝖢`$-span of all $`x,y_𝒜`$ is $`𝒜`$. A (full, positive semi-definite) $`𝒜`$-valued inner product on a left $`𝒜`$-module is defined analogously, but with linearity properties (with respect to $`𝖢`$ and $`𝒜`$) in the first argument. If $``$ is another -algebra over $`𝖢`$, and $`{}_{}{}^{}𝔛_{𝒜}^{}`$ is a $`(`$-$`𝒜)`$-bimodule with an $`𝒜`$-valued inner product, we call the $``$-action *compatible* if $`Bx,y_𝒜=x,B^{}y_𝒜`$ for all $`B`$ and $`x,y{}_{}{}^{}𝔛_{𝒜}^{}`$.
Let $`{}_{}{}^{}𝔛_{𝒜}^{}`$ be a $`(`$-$`𝒜)`$-bimodule with positive semi-definite $`𝒜`$-valued inner product and compatible $``$-action. If $`\pi _𝒜`$ is a -representation of $`𝒜`$ on $``$, then one can consider the $`𝖢`$-module $`\stackrel{~}{𝔎}:={}_{}{}^{}𝔛_{𝒜}^{}_𝒜`$, where the $`𝒜`$-balanced tensor product is defined using $`\pi _𝒜`$. Then $`\stackrel{~}{𝔎}`$ is a $``$-left module and by $`x\psi ,y\varphi _{\stackrel{~}{𝔎}}:=\psi ,\pi _𝒜(x,y_𝒜)\varphi _{}`$ one obtains a Hermitian inner product on $`\stackrel{~}{𝔎}`$. We say that the bimodule $`{}_{}{}^{}𝔛_{𝒜}^{}`$ satisfies the property P if this Hermitian product is positive semi-definite for all -representations $`(\pi _𝒜,)`$ of $`𝒜`$, see \[9, Sect. 3 and 4\] for various sufficient conditions guaranteeing P. In this case one can quotient out the length zero vectors $`\stackrel{~}{𝔎}^{}`$ to obtain a pre-Hilbert space $`𝔎=\stackrel{~}{𝔎}/\stackrel{~}{𝔎}^{}`$. It is easy to check that the left action of $``$ passes to the quotient and yields a -representation $`\pi _{}`$ on $`𝔎`$. The whole procedure, called *algebraic Rieffel induction*, is functorial and the functor is denoted by $`_𝔛`$. Given a $`(`$-$`𝒜)`$-bimodule with a positive semi-definite $``$-valued inner product and compatible $`𝒜`$-action, we can consider the corresponding conjugated $`(𝒜`$-$`)`$-bimodule, denoted by $`\overline{𝔛}`$, which is just $`𝔛`$ as additive group, but with $`𝒜`$ and $``$ acting by adjoints and $`𝖢`$ by complex conjugates. We then say that $`{}_{}{}^{}𝔛_{𝒜}^{}`$ satisfies property Q if $`\overline{𝔛}`$ satisfies P as above (see \[9, Sect. 5\]).
Let us now consider formal Morita equivalence, as in . Suppose $`{}_{}{}^{}𝔛_{𝒜}^{}`$ is a $`(`$-$`𝒜)`$-bimodule, equipped with full positive semi-definite $`𝒜`$-valued and $``$-valued inner products, compatible $`𝒜`$ and $``$ actions and so that properties P and Q are satisfied. We also assume that the inner products are *compatible*, in the sense that $`xy,z_𝒜={}_{}{}^{}x,yz`$ for all $`x,y,z𝔛`$. We call this object an *equivalence bimodule* and say that two -algebras $`𝒜`$ and $``$ are *formally Morita equivalent* if there exists a $`(`$-$`𝒜)`$-equivalence bimodule. We recall that if the actions of $`𝒜`$ and $``$ on $`{}_{}{}^{}𝔛_{𝒜}^{}`$ are in addition strongly non-degenerate, then $`{}_{}{}^{}𝔛_{𝒜}^{}`$ is called a *non-degenerate equivalence bimodule* and one can show that in this case that $`𝒜`$ and $``$ have equivalent categories of (strongly non-degenerate) -representations, with the functors $`_𝔛`$ and $`_{\overline{𝔛}}`$ defining this equivalence (\[9, Thm. 5.10\]). As for notation, we will write $`{}_{}{}^{}\text{-}\mathrm{𝗋𝖾𝗉}(𝒜)`$ for the category of -representations of $`𝒜`$ (usually considered with isometric interwiners as morphisms) and $`{}_{}{}^{}\text{-}\mathrm{𝖱𝖾𝗉}(𝒜)`$ for the category of strongly non-denerate -representations.
A -algebra over $`𝖢`$ has *sufficiently many positive linear functionals* if for every non-zero Hermitian element $`H`$ one finds a positive linear functional $`\omega `$ with $`\omega (H)0`$. We recall the following important fact: a -algebra $`𝒜`$ *with approximate identity* has sufficiently many positive linear functionals if and only if $`𝒜`$ has a faithful -representation. As a consequence, these algebras present some nice properties, resembling $`C^{}`$-algebras. For instance, they are torsion-free ($`zA=0`$ implies $`A=0`$ or $`z=0`$, for $`z𝖢`$), $`A^{}A=0`$ implies $`A=0`$ and $`H^n=0`$ for a normal element $`H`$ implies $`H=0`$ (\[9, Prop. 2.8\]). We observe that the assumption on the existence of an approximate identity can be weakened in particular cases. Indeed, if $`𝒜`$ is a $`C^{}`$-algebra or a star-product algebra, then it always has sufficiently many positive linear functionals but it does not admit an algebraic approximate identity in general. Nevertheless, these algebras always have faithful representations and all the aforementioned nice properties. We observe that in both cases, however, a topological approximate identity exists (with respect to the norm topology in the case of $`C^{}`$-algebras and with respect to the $`\lambda `$-adic topology for star-product algebras, see \[9, Sect. 8\]).
On the other hand, the existence of sufficiently many positive linear functionals, without any further requirement, is generally not enough to guarantee the existence of faithful representations and desirable algebraic properties. For instance, consider $`𝖢=𝖱(\mathrm{i})`$ with its natural $`𝖢`$-module structure. Define the -algebra $`𝒜`$ to be this $`𝖢`$-module, with the natural -involution but with zero multiplication. Then *all* linear functionals on $`𝒜`$ are positive and $`𝒜`$ has sufficiently many positive linear functionals. Note that $`𝒜`$ does not admit faithful representations (since we have $`A0`$ with $`A^{}A=0`$). For an example of a -algebra *without* sufficiently many positive linear functionals, let $`^{}(𝖢^n)`$ be the Grassmann algebra over $`𝖢^n`$, equipped with a -involution given by $`e_i^{}:=e_i`$, where $`e_1,\mathrm{},e_n`$ is the canonical basis of $`𝖢^n`$. Denote by $`^+(𝖢^n)`$ the ideal of elements with positive degree and write $`^{}(𝖢^n)=\mathrm{𝖢𝟣}^+(𝖢^n)`$. Then every positive linear functional $`\omega :^{}(𝖢^n)𝖢`$ has to satisfy $`\omega (\mathrm{𝟣})0`$ and $`\omega |_{^+(𝖢^n)}=0`$. Hence, up to normalization, there is only one positive functional. In particular, every (Hermitian) element in $`^+(𝖢^n)`$ is positive but $`^{}(𝖢^n)`$ does not have sufficiently many positive linear functionals for $`n1`$. Note that this -algebra is nevertheless torsion-free. For an example with torsion, set $`𝖱=`$ and $`𝖢=\mathrm{i}`$ and consider $`_2=\{0,\mathrm{𝟣}\}`$. Then $`_2`$ becomes a unital -algebra over $`𝖢`$ by $`\mathrm{i}\mathrm{𝟣}=\mathrm{𝟣}`$ and $`\mathrm{𝟣}^{}=\mathrm{𝟣}`$. Note that $`_2`$ is not torsion-free, since for instance $`2\mathrm{𝟣}=0`$, and there are no $`𝖢`$-linear functionals $`\omega :_2𝖢`$ except for the zero functional.
Thus the general category of -algebras over $`𝖢`$ seems too wide for the interpretation of these algebras as ‘observable algebras’. We will discuss in Sections 4 and 5 how the ideal strucuture of these algebras help characterizing the more ‘observable algebra’-like objects.
## 3 Formal Morita equivalence and $`C^{}`$-algebras
The aim of this section is to clarify the relationship between the notions of formal Morita equivalence , Morita -equivalence and strong Morita equivalence . We start by briefly recalling some of the definitions.
Let $`𝒜`$ be a $`C^{}`$-algebra. A *right pre-Hilbert $`𝒜`$-module* is a right $`𝒜`$-module $`𝔛`$, equipped with a positive definite $`𝒜`$-valued inner product $`,_𝒜`$ We call $`𝔛_𝒜`$ a *right Hilbert $`𝒜`$-module* if $`𝔛`$ is in addition complete under the norm $`x_𝒜=x,x_𝒜^{1/2}`$. A *left Hilbert $`𝒜`$-module* is defined analogously. A Hilbert $`𝒜`$-module is called *(topologically) full* if the ideal $`𝔛,𝔛_𝒜`$ is dense in $`𝒜`$. If $`𝒜`$ and $``$ are $`C^{}`$-algebras, then an *imprimitivity bimodule* is a $`(`$-$`𝒜)`$-bimodule $`{}_{}{}^{}𝔛_{𝒜}^{}`$ which is a full right $`𝒜`$-Hilbert module and a full left $``$-Hilbert module so that for all $`x,y,z{}_{}{}^{}𝔛_{𝒜}^{}`$ and $`A𝒜`$, $`B`$, we have $`Bx,y_𝒜=x,B^{}y_𝒜`$, $`{}_{}{}^{}xA,y={}_{}{}^{}x,yA^{}`$ and $`xy,z_𝒜={}_{}{}^{}x,yz`$. We finally say that two $`C^{}`$-algebras $`𝒜`$ and $``$ are *strongly Morita equivalent* if there is a $`(`$-$`𝒜)`$-imprimitivity bimodule. An equivalent definition of strong Morita equivalence can be given as follows: $`𝒜`$ and $``$ are strongly Morita equivalent if there exists a full right $`𝒜`$-Hilbert module $`𝔛_𝒜`$ so that $`𝒦(𝔛_𝒜)`$, where $`𝒦(𝔛_𝒜)`$ is the completion of the -algebra $`𝔉(𝔛_𝒜)`$ of “finite-rank operators” on $`𝔛_𝒜`$, that is, the linear span of operators of the form $`\mathrm{\Theta }_{x,y}`$, where $`\mathrm{\Theta }_{x,y}z=xy,z_𝒜`$, for $`x,y,z𝔛_𝒜`$.
To state Ara’s result on the algebraic characterization of strong Morita equivalence, we will need the following definitions. Recall that a ring $`R`$ is called *non-degenerate* if $`rR=0`$ or $`Rr=0`$ implies that $`r=0`$ for all $`rR`$, and *idempotent* if elements of the form $`r_1r_2`$ span $`R`$. We also consider the category $`\mathrm{𝖬𝗈𝖽}`$-$`R`$ of right $`R`$-modules $``$ satisfying $`R=`$ (in our terminology, this means that the action of $`R`$ on $``$ is strongly non-degenerate) and $`xR=0`$ implies $`x=0`$ for all $`x`$, i.e. the $`R`$-action on $``$ is non-degenerate. Let now $`R`$ and $`S`$ be two non-degenerate and idempotent rings with involution. Then $`R`$ and $`S`$ are called *Morita -equivalent* if there exists a right $`R`$-module $`\mathrm{𝖬𝗈𝖽}`$-$`R`$ equipped with a full $`R`$-valued inner product $`,`$ so that $`x,=0`$ implies $`x=0`$ for $`x,y`$ and $`S𝔉(_R)`$. Here $`𝔉(_R)`$ again denotes the -algebra of “finite rank” operators on $`_R`$ and is defined as before for $`C^{}`$-algebras. We remark that if $`R`$ is a $`K`$-algebra, for some fixed unital and commutative ring $`K`$, and $``$ has a compatible $`K`$-module structure, then it follows from the non-degeneracy of $`R`$ that $`,`$ is automatically $`K`$-linear in the second argument.
With this notion of Morita -equivalence, one can give a purely algebraic characterization of strong Morita equivalence of $`C^{}`$-algebras. Recall that if $`𝒜`$ is a $`C^{}`$-algebra, then its *Pedersen ideal* $`𝒫_𝒜`$ is defined as the minimal dense ideal of $`𝒜`$, see \[25, Sect. 5.6\]. One has then the following result, due to Ara (\[1, Thm. 2.4\]): *Two $`C^{}`$-algebras $`𝒜`$ and $``$ are strongly Morita equivalent if and only if $`𝒫_𝒜`$ and $`𝒫_{}`$ are Morita -equivalent*. We will now discuss how this result is related to the notion of formal Morita equivalence. We start with the following observation.
###### Lemma 3.1
Let $`𝒜`$ and $``$ be non-degenerate, idempotent -algebras over $`𝖢`$. Then if $`𝒜`$ and $``$ are formally Morita equivalent, they are also Morita -equivalent.
Proof: Note that it suffices to show that there exists an equivalence bimodule $`{}_{}{}^{}𝔛_{𝒜}^{}`$ so that $`x,𝔛_𝒜=0`$ implies $`x=0`$, for $`x𝔛`$, $`(𝔛_𝒜,,_𝒜)\mathrm{𝖬𝗈𝖽}`$-$`𝒜`$ (i.e., the $`𝒜`$-action on $`𝔛`$ is strongly non-degenerate and $`x𝒜=0`$ implies $`x=0`$ for $`x𝔛`$) and $`𝔉(𝔛_𝒜)`$.
Let $`{}_{}{}^{}𝔛_{𝒜}^{}`$ be an arbitrary equivalence bimodule. We start by noticing that we can choose it so that the $``$ and $`𝒜`$ actions are strongly non-degenerate. Indeed, it follows from the compatibility of $`,_𝒜`$ and $`{}_{}{}^{},`$ and their fullness properties that $`𝔛𝒜=𝔛`$. If we denote $`\widehat{𝔛}=𝔛𝒜=𝔛`$, then $`\widehat{𝔛}`$ has a natural $`(`$-$`𝒜)`$-bimodule structure and also compatible $`𝒜`$ and $``$ valued inner products, defined by simply restricting $`{}_{}{}^{},`$ and $`,_𝒜`$ to $`\widehat{𝔛}`$. Note that given $`A𝒜`$, by idempotency we can write $`A=_iA_i^1A_i^2A_i^3`$, for $`A_i^j𝒜`$, and by fullness of $`,_𝒜`$, we can write $`A_i^2=_jx_j^i,y_j^i_𝒜`$. So $`A=_iA_i^1_jx_j^i,y_j^i_𝒜A_i^3=_{i,j}x_j^iA_{i}^{1}{}_{}{}^{},y_j^iA_i^3_𝒜`$, showing that $`,_𝒜`$ restricted to $`\widehat{𝔛}`$ is still full, and the same holds for $`,_{}`$. Therefore $`\widehat{𝔛}`$ is a $`(`$-$`𝒜)`$-equivalence bimodule with strongly non-degenerate $`𝒜`$ and $``$ actions.
We also observe that if $`{}_{}{}^{}𝔛_{𝒜}^{}`$ is an (arbitrary) equivalence bimodule, then $``$ and $`𝒜`$ act on $`𝔛`$ faithfully (i.e., $`Bx=0`$ for all $`x𝔛`$ implies $`B=0`$, for all $`B`$, and the same for $`𝒜`$). To see that, suppose $`Bx=0`$ for all $`x`$. Then given an arbitrary $`B^{}`$, it follows that we can write $`B^{}=_i{}_{}{}^{}x_i,y_i`$ and hence $`BB^{}=B_i{}_{}{}^{}x_i,y_i=_i{}_{}{}^{}Bx_i,y_i=0`$. Therefore $`B=0`$ by non-degeneracy of $``$ and the same argument holds for $`𝒜`$. It then follows (see the proof of \[9, Prop. 5.16\]) that $`\{x{}_{}{}^{}𝔛_{𝒜}^{}|x,y_𝒜=0,y{}_{}{}^{}𝔛_{𝒜}^{}\}=\{x{}_{}{}^{}𝔛_{𝒜}^{}|{}_{}{}^{}x,y=0,y{}_{}{}^{}𝔛_{𝒜}^{}\}=N`$ and $`𝔛/N`$ is still a $`(`$-$`𝒜)`$-equivalence bimodule, with the property that $`x,𝔛_𝒜=0`$ implies that $`x=0`$, the same holding for $`,_𝒜`$. Moreover, if the actions of $`𝒜`$ and $``$ were strongly non-degenerate on $`{}_{}{}^{}𝔛_{𝒜}^{}`$, this still holds for the quotient $`𝔛/N`$. So, we can always find an equivalence bimodule $`{}_{}{}^{}𝔛_{𝒜}^{}`$ with strongly non-degenerate $`𝒜`$ and $``$ actions and inner products satisfying $`x,𝔛_𝒜=0x=0`$ and $`{}_{}{}^{}x,𝔛=0x=0`$. We observe that, in this case (see proof of \[9, Prop. 6.1\]), $`𝔉(𝔛_𝒜)`$.
Given $`{}_{}{}^{}𝔛_{𝒜}^{}`$ an equivalence bimodule with these properties, we note that if $`x𝒜=0`$, then in particular $`xy,z_𝒜=0`$ for all $`y,z𝔛`$. But then $`{}_{}{}^{}x,yz=0`$ for all $`z`$, and hence $`{}_{}{}^{}x,y=0`$. But since $`y`$ is also arbitrary, it follows that $`x=0`$. Therefore, $`(𝔛_𝒜,,_𝒜)\mathrm{𝖬𝗈𝖽}`$-$`𝒜`$ and clearly establishes a Morita -equivalence between $`𝒜`$ and $`𝔉(𝔛_𝒜)`$. $`\mathrm{}`$
We observe that -algebras with approximate identities are non-degenerate and idempotent. We will now show the converse of the previous lemma in the case of $`C^{}`$-algebras and Pedersen ideals (recalling that both types of algebras are also non-degenerate and idempotent). We start with some general lemmas.
###### Lemma 3.2
Let $`𝒜`$ be a $`C^{}`$-algebra and $`𝒫_𝒜`$ its Pedersen ideal. Then $`A𝒫_𝒜`$ is positive in $`𝒜`$ if and only if $`A𝒫_{𝒜}^{}{}_{}{}^{+}`$ (positivity in the purely algebraic sense of -algebras over $`𝖢`$).
Proof: Let $`A𝒫_𝒜`$ be positive in $`𝒜`$ and let $`\omega :𝒫_𝒜`$ be an arbitrary positive linear functional (in the algebraic sense, that is, $`\omega (A^{}A)0`$ for all $`A𝒫_𝒜`$). We must show that $`\omega (A)0`$. Note that by the general properties of Pedersen ideals (see \[25, Thm. 5.6.2\]), if $`A𝒫_𝒜`$, then $`C^{}(A)𝒫_𝒜`$, where $`C^{}(A)`$ denotes the $`C^{}`$-algebra generated by $`A`$. Also observe that if $`A𝒜^+`$, then $`AC^{}(A)^+`$, and hence $`A=A_1^{}A_1`$, for some $`A_1C^{}(A)𝒫_𝒜`$. Therefore, $`\omega (A)=\omega (A_1^{}A_1)0`$. For the converse, note that if $`\omega `$ is a positive linear functional in $`𝒜`$, then $`\omega |_{𝒫_𝒜}`$ is positive in $`𝒫_𝒜`$ and hence if $`A𝒫_𝒜^+`$, we immediately have $`\omega (A)0`$. Hence $`A𝒜^+`$. $`\mathrm{}`$
We now observe a general property of -representations of Pedersen ideals on pre-Hilbert spaces which should be well-known.
###### Lemma 3.3
Let $`𝒜`$ be a $`C^{}`$-algebra, $`𝒫_𝒜`$ its Pedersen ideal and suppose $`\pi :𝒫_𝒜𝔅()`$ is a -representation on a complex pre-Hilbert space $``$. Then $`\pi `$ extends to a -representation $`\pi ^{\mathrm{cl}}:𝒫_𝒜𝔅(^{\mathrm{cl}})`$, where $`^{\mathrm{cl}}`$ is the closure of $``$. Moreover, $`\pi ^{\mathrm{cl}}(A)A`$ for all $`A𝒫_𝒜`$ and hence $`\pi ^{\mathrm{cl}}`$ also extends to a -representation of $`𝒜`$ on $`^{\mathrm{cl}}`$.
###### Remark 3.4
As pointed out to us by P. Ara, the above lemma follows from a more general fact. If $`𝒜`$ is any complex -algebra with positive definite involution (\[15, pp. 338\]), let $`𝒜_b`$ be the -subalgebra of bounded elements of $`𝒜`$ (\[15, pp. 338\]). Then any -representation of $`𝒜_b`$ on a complex pre-Hilbert space $``$ naturally extends to a -representation of $`\overline{𝒜_b}`$ (closure of $`𝒜_b`$ with respect to its natural seminorm, as defined in \[15, pp. 342\]) on $`^{\mathrm{cl}}`$.
We can now discuss general positivity properties of modules over Pedersen ideals.
###### Proposition 3.5
Let $`𝒜`$ be a $`C^{}`$-algebra, $`𝒫_𝒜`$ its Pedersen ideal and $``$ an arbitrary -algebra over $``$. Suppose $`{}_{}{}^{}𝔛_{𝒫_𝒜}^{}`$ is a bimodule, equipped with a positive semi-definite $`𝒫_𝒜`$-valued inner product $`,_{𝒫_𝒜}`$ and compatible $``$-action. Then $`𝔛`$ satisfies property P.
Proof: We must show that given a complex pre-Hilbert space $``$ and a -representation $`\pi :𝒫_𝒜𝔅()`$, the formula $`x_1\varphi _1,x_2\varphi _2_{\stackrel{~}{𝔎}}=\varphi _1,\pi (x_1,x_2_{𝒫_𝒜})\varphi _2_{}`$, for $`x_1,x_2𝔛`$ and $`\varphi _1,\varphi _2`$ defines a positive semi-definite inner product on $`\stackrel{~}{𝔎}=𝔛_{𝒫_𝒜}`$. Note that since $`\pi `$ extends to a -representation $`\pi ^{\mathrm{cl}}:𝒫_𝒜𝔅(^{\mathrm{cl}})`$ on the Hilbert space completion of $``$, there is a natural isometric embedding $`\iota :𝔛_{𝒫_𝒜}𝔛_{𝒫_𝒜}^{\mathrm{cl}}`$, where the $`𝒫_𝒜`$-balanced tensor products are defined using $`\pi `$ and $`\pi ^{\mathrm{cl}}`$, respectively. So for $`\varphi _1,\mathrm{},\varphi _n^{\mathrm{cl}}`$ and $`x_1,\mathrm{},x_n𝔛`$, we have
$$\underset{i}{}x_i\varphi _i,\underset{i}{}x_i\varphi _i_{\stackrel{~}{𝔎}}=\underset{i,j}{}\varphi _i,\pi \left(x_i,x_j_𝒜\right)\varphi _j_{}=\underset{i,j}{}\varphi _i,\pi ^{\mathrm{cl}}\left(x_i,x_j_𝒜\right)\varphi _j_{^{\mathrm{cl}}}.$$
(3.1)
Therefore, to show the positivity of $`,_{\stackrel{~}{𝔎}}`$, it suffices to show that the right hand side of (3.1) is non-negative. This can be done by the usual procedure (see \[18, Ch. IV, Sect. 2.2\]), using the fact that any -representation of $`𝒫_𝒜`$ on a Hilbert space can be decomposed into the sum of a trivial representation and a non-degenerate one, and the non-degenerate part itself can be decomposed into the sum (in the topological sense) of cyclic (also in the topological sense) -representations (see \[18, Ch. I, Prop. 1.5.2\]). $`\mathrm{}`$
Let $`𝒜`$ and $``$ be $`C^{}`$-algebras, let $`𝒫_𝒜`$ and $`𝒫_{}`$ be the corresponding Pedersen ideals and suppose there exists $`𝔛\mathrm{𝖬𝗈𝖽}`$-$`𝒫_𝒜`$, equipped with a pairing $`,`$ so that $`(𝔛,,)`$ establishes a Morita -equivalence between $`𝒫_𝒜`$ and $`𝔉(𝔛_{𝒫_𝒜})𝒫_{}`$. In \[1, Thm. 2.4\], it is shown that, in this case, one automatically has $`x,x𝒜^+`$ and $`\mathrm{\Theta }_{x,x}^+`$ (recall that for $`C^{}`$-algebras, the usual notion of positivity coincides with positivity as defined for arbitrary -algebras over $`𝖢`$). Then, by Lemma 3.2, it follows that $`x,x𝒫_𝒜^+`$ and $`\mathrm{\Theta }_{x,x}𝒫_{}^+`$. It is then easy to check (see e.g. \[9, Prop. 6.2\]) that $`𝔛`$ is a $`(𝒫_{}`$-$`𝒫_𝒜)`$-bimodule satisfying all the properties of an equivalence bimodule, except possibly for P and Q. But these properties follow directly from Proposition 3.5. Hence we have
###### Corollary 3.6
If $`𝒫_𝒜`$ and $`𝒫_{}`$ are Morita -equivalent, then they are automatically formally Morita equivalent.
It then follows from Lemma 3.1, Proposition 3.5 and \[1, Thm. 2.4\] that
###### Theorem 3.7
Two $`C^{}`$-algebras $`𝒜`$ and $``$ are strongly Morita equivalent if and only if their corresponding Pedersen ideals, $`𝒫_𝒜`$ and $`𝒫_{}`$, are formally Morita equivalent. In particular, if $`𝒜`$ and $``$ are unital, then they are strongly Morita equivalent if and only if they are formally Morita equivalent.
It follows from \[9, Thm. 5.10\] that if two $`C^{}`$-algebras $`𝒜`$ and $``$ are strongly Morita equivalent, then -$`\mathrm{𝖱𝖾𝗉}(𝒫_𝒜)`$ and -$`\mathrm{𝖱𝖾𝗉}(𝒫_{})`$ are equivalent categories. We recall from that strong Morita equivalent $`C^{}`$-algebras have equivalent categories of Hermitian modules. We conclude this section discussing how these two results are related.
For a $`C^{}`$-algebra $`𝒜`$, let $`\mathrm{𝗁𝖾𝗋𝗆𝗈𝖽}(𝒜)`$ be the category of -representations of $`𝒜`$ on Hilbert spaces, with isometric intertwiners as morphisms. The category of *non-degenerate* -representations is denoted by $`\mathrm{𝖧𝖾𝗋𝗆𝗈𝖽}(𝒜)`$ and objects in this category are called *Hermitian modules* . In order to understand how these categories are related to the categories of purely algebraic representations -$`\mathrm{𝗋𝖾𝗉}(𝒫_𝒜)`$ and -$`\mathrm{𝖱𝖾𝗉}(𝒫_𝒜)`$, we introduce two other categories. By $`{}_{}{}^{}\text{-}\mathrm{𝗋𝖾𝗉𝖣𝖲}(𝒜)`$ we denote the category of -representations of $`𝒜`$ on Hilbert spaces with a dense $`𝒫_𝒜`$-invariant subspace, while for $`{}_{}{}^{}\text{-}\mathrm{𝖱𝖾𝗉𝖣𝖲}(𝒜)`$ we require in addition that the induced -representation of $`𝒫_𝒜`$ restricted to the dense subspace be strongly non-degenerate (in the algebraic sense). For the morphisms we take isometric intertwiners preserving the dense invariant subspaces. Note that objects in $`{}_{}{}^{}\text{-}\mathrm{𝖱𝖾𝗉𝖣𝖲}(𝒜)`$ are in particular Hermitian modules over $`𝒜`$ (with a distinguished dense subspace). One can check that the process of extending the -representations, discussed in Lemma 3.3, gives rise to a functor from -$`\mathrm{𝖱𝖾𝗉}(𝒫_𝒜)`$ to $`{}_{}{}^{}\text{-}\mathrm{𝖱𝖾𝗉𝖣𝖲}(𝒜)`$ (-$`\mathrm{𝗋𝖾𝗉}(𝒫_𝒜)`$ to $`{}_{}{}^{}\text{-}\mathrm{𝗋𝖾𝗉𝖣𝖲}(𝒜)`$) for which the natural restriction functor is an inverse. It then follows that -$`\mathrm{𝖱𝖾𝗉}(𝒫_𝒜)`$ and $`{}_{}{}^{}\text{-}\mathrm{𝖱𝖾𝗉𝖣𝖲}(𝒜)`$ (-$`\mathrm{𝗋𝖾𝗉}(𝒫_𝒜)`$ and $`{}_{}{}^{}\text{-}\mathrm{𝗋𝖾𝗉𝖣𝖲}(𝒜)`$) are equivalent categories. We point out that $`{}_{}{}^{}\text{-}\mathrm{𝖱𝖾𝗉𝖣𝖲}(𝒜)`$ and $`\mathrm{𝖧𝖾𝗋𝗆𝗈𝖽}(𝒜)`$ ($`{}_{}{}^{}\text{-}\mathrm{𝗋𝖾𝗉𝖣𝖲}(𝒜)`$ and $`\mathrm{𝗁𝖾𝗋𝗆𝗈𝖽}(𝒜)`$) are not equivalent categories in general. We finally remark that the categories $`{}_{}{}^{}\text{-}\mathrm{𝗋𝖾𝗉𝖣𝖲}(𝒜)`$ and $`{}_{}{}^{}\text{-}\mathrm{𝖱𝖾𝗉𝖣𝖲}(𝒜)`$ are interesting in their own right since, in many applications, -representations of $`𝒜`$ come automatically with a dense invariant subspace, as in GNS representations. As another example, certain unbounded operators associated to -representations, such as generators of symmetries, distinguish dense subspaces by their domains of definition.
## 4 Closed -Ideals and the Minimal Ideal
After having clarified the connection between formal Morita equivalence and $`C^{}`$-algebras, we consider in this section general -algebras over $`𝖢`$ and start the discussion about (closed) -ideals.
If $`𝒜`$ is a -ideal, i.e. an ideal closed under the -involution, then $`𝒜/`$ is again a -algebra. If $``$ is a left (or right) ideal which is closed under the -involution then it is automatically a two-sided ideal and hence a -ideal.
###### Definition 4.1
A -ideal $``$ of $`𝒜`$ is called closed if it is the kernel of a -representation.
Clearly, the kernel $`\mathrm{ker}\pi `$ of a -representation $`\pi `$ is a -ideal, and unitary equivalent -representations have the same kernel. But $`\mathrm{ker}\pi `$ does not determine this equivalence class since e.g. $`\mathrm{ker}\pi =\mathrm{ker}(\pi \pi )`$. We observe that if $`𝒜`$ is a $`C^{}`$-algebra, then a -ideal $``$ in $`𝒜`$ is closed in this algebraic sense if and only if it is norm closed. Indeed, recall that a -ideal in a $`C^{}`$-algebra is norm closed if and only if it is the kernel of some -representation of $`𝒜`$ on a Hilbert space (\[18, Ch. I, Thm. 1.3.10\]). But since any representation of a $`C^{}`$-algebra on a pre-Hilbert space extends to a representation on its completion (see Lemma 3.3), the conclusion follows.
###### Lemma 4.2
Let $`𝒜`$ be a -algebra over $`𝖢`$ and $``$ a -ideal. Then $`𝒜/`$ has a faithful -representation if and only if $``$ is closed. If $`\{_\alpha \}_{\alpha \mathrm{\Lambda }}`$ are closed -ideals then $`=_{\alpha \mathrm{\Lambda }}_\alpha `$ is again a closed -ideal.
Proof: For the first part use the projection $`\rho :𝒜𝒜/`$ to pull-back -representations. For the second observe that if $`_\alpha =\mathrm{ker}\pi _\alpha `$ then $`=\mathrm{ker}(_{\alpha \mathrm{\Lambda }}\pi _\alpha )`$. $`\mathrm{}`$
Since formal Morita equivalence deals with strongly non-degenerate -representations, it would be nice if these particular representations were sufficient to define all closed -ideals. If $`𝒜`$ is idempotent, then this is indeed the case.
###### Lemma 4.3
Let $`𝒜`$ be an idempotent -algebra over $`𝖢`$ and $`\pi :𝒜𝔅()`$ a -representation. Let
$$_{\mathrm{snd}}:=\left\{\varphi \right|\varphi _i,A_i𝒜,i=1,\mathrm{},n:\varphi =_i\pi (A_i)\varphi _i\}.$$
(4.1)
Then $`_{\mathrm{snd}}`$ is $`\pi `$-invariant and the restriction $`\pi _{\mathrm{snd}}:=\pi |_{_{\mathrm{snd}}}`$ is strongly non-degenerate with $`\mathrm{ker}\pi _{\mathrm{snd}}=\mathrm{ker}\pi `$. Moreover, every GNS representation is strongly non-degenerate.
Proof: The invariance of $`_{\mathrm{snd}}`$ is obvious and the fact that $`\pi _{\mathrm{snd}}`$ is strongly non-degenerate follows from the idempotency of $`𝒜`$. Finally let $`\pi _{\mathrm{snd}}(A)=0`$. Then in particular $`\pi (A)\pi (A^{})\varphi =0`$ for all $`\varphi `$ whence $`\pi (A)=0`$ follows and thus $`\mathrm{ker}\pi _{\mathrm{snd}}\mathrm{ker}\pi `$. The other inclusion is trivial. The strong non-degeneracy of a GNS representation follows directly from the definition. $`\mathrm{}`$
Using the closed -ideals, we define the following closure operation. If $`𝒥𝒜`$ is an arbitrary subset then $`𝒥^{\mathrm{cl}}`$ is the smallest closed -ideal containing $`𝒥`$, i.e.
$$𝒥^{\mathrm{cl}}:=\underset{𝒥,\text{ closed }\text{}\text{-ideal}}{},$$
(4.2)
where $``$ runs over the closed -ideals of $`𝒜`$. Clearly $`𝒥𝒥^{\mathrm{cl}}`$ and $`(𝒥^{\mathrm{cl}})^{\mathrm{cl}}=𝒥^{\mathrm{cl}}`$ for any subset of $`𝒜`$. This implies that for $`𝒥`$ we have $`𝒥^{\mathrm{cl}}^{\mathrm{cl}}`$ and $`(_{\alpha \mathrm{\Lambda }}_\alpha )^{\mathrm{cl}}=(_{\alpha \mathrm{\Lambda }}_\alpha ^{\mathrm{cl}})^{\mathrm{cl}}`$ for an arbitrary index set $`\mathrm{\Lambda }`$ and $`_\alpha 𝒜`$. Finally note that $`_{\alpha \mathrm{\Lambda }}_\alpha ^{\mathrm{cl}}=\left(_{\alpha \mathrm{\Lambda }}_\alpha ^{\mathrm{cl}}\right)^{\mathrm{cl}}`$. We can define the following operations. For two subsets $`,𝒥𝒜`$, we set
$$𝒥:=(𝒥)^{\mathrm{cl}}\text{ and }𝒥:=^{\mathrm{cl}}𝒥^{\mathrm{cl}}.$$
(4.3)
One can check that these two operations define a lattice structure on the set of closed -ideals of $`𝒜`$. In this case, the operation ‘$``$’, which is defined by $`𝒥`$ if $`𝒥=`$, coincides with the set-theoretic ‘$``$’. We then have a lattice structure on the set of all closed -ideals in $`𝒜`$, that we denote by $`𝔏(𝒜)`$. Note that all -ideals, closed or not, also have a lattice structure, with the ‘closure map’ defined by taking the smallest -ideal containing the corresponding subset. However, we will be mainly interested in the closed -ideals.
For a $`C^{}`$-algebra $`𝒜`$, $`𝔏(𝒜)`$ coincides with its usual lattice of norm closed -ideals. If $`𝒫_𝒜`$ is its Pedersen ideal, note that we have a natural map $`F:𝔏(𝒫_𝒜)𝔏(𝒜)`$ defined as follows. Let $`𝔏(𝒫_𝒜)`$. Then $`=\mathrm{ker}\pi `$, for some -representation of $`𝒫_𝒜`$ on a pre-Hilbert space $``$. As we observed in Lemma 3.3, $`\pi `$ extends to a -representation $`\pi ^{\mathrm{cl}}`$ of $`𝒜`$ on $`^{\mathrm{cl}}`$. We then define $`F`$ as the map $`=\mathrm{ker}\pi F()=\mathrm{ker}\pi ^{\mathrm{cl}}`$.
###### Proposition 4.4
$`F:𝔏(𝒫_𝒜)𝔏(𝒜)`$ is a lattice isomorphism.
Proof: Let us first observe that $`F`$ is a well-defined map. Note that if $`\pi `$ is a -representation of $`𝒫_𝒜`$, then $`\mathrm{ker}\pi `$ is always norm dense in $`\mathrm{ker}\pi ^{\mathrm{cl}}`$. To see that, let $`\{e_\lambda \}_{\lambda \mathrm{\Lambda }}`$ be an approximate identity in $`𝒜`$, with $`e_\lambda 𝒫_𝒜`$. This is possible since $`𝒫_𝒜`$ is dense in $`𝒜`$ (\[12, Thm. I.7.2\]). Then given $`A\mathrm{ker}\pi ^{\mathrm{cl}}`$, we have $`e_\lambda A\mathrm{ker}\pi ^{\mathrm{cl}}𝒫_𝒜=\mathrm{ker}\pi `$ and $`e_\lambda AA`$. Now let $`𝔏(𝒫_𝒜)`$ and suppose $`\pi `$ and $`\rho `$ are two different -representations of $`𝒫_𝒜`$ so that $`=\mathrm{ker}\pi =\mathrm{ker}\rho `$. We must show that $`\mathrm{ker}\pi ^{\mathrm{cl}}=\mathrm{ker}\rho ^{\mathrm{cl}}`$. Note that if $`A\mathrm{ker}\pi ^{\mathrm{cl}}`$, then we can find $`A_\lambda \mathrm{ker}\pi =\mathrm{ker}\rho `$ with $`A_\lambda A`$. So $`\rho (A_\lambda )=0`$, for all $`\lambda `$, and hence $`\rho (A)=0`$. This shows that $`\mathrm{ker}\pi ^{\mathrm{cl}}\mathrm{ker}\rho ^{\mathrm{cl}}`$ and by changing the roles of $`\pi `$ and $`\rho `$ we conclude that $`\mathrm{ker}\pi ^{\mathrm{cl}}=\mathrm{ker}\rho ^{\mathrm{cl}}`$. This proves well-definedness. It is not hard to check that $`F`$ is a bijection which preserves the order structure of the lattices. This guarantees that $`F`$ is a lattice isomorphism. $`\mathrm{}`$
Recall that for a $`C^{}`$-algebra, any (topologically) non-degenerate -representation can be decomposed into a direct sum (also in the topological sense) of GNS representations. Thus for $`C^{}`$-algebras, any closed -ideal can be obtained as the kernel of a (topological) sum of GNS representations. We observe that for a general -algebra over $`𝖢`$, it is not true that an arbitrary (strongly non-degenerate) -representation decomposes into an algebraic sum of GNS representations. In fact, even if $`𝖢`$ is an algebraically closed field (which is complete in its canonical topology) and we allow topological sums, this decomposition does not hold in general. This follows since otherwise every Hilbert space over $`𝖢`$ would have a Hilbert basis (this notion makes sense in this situation), which is *not* the case (see e.g. \[8, Thm. 7\]). Nevertheless, it is still true that arbitrary closed -ideals on a -algebra arise as kernel of sums of GNS representations, as we now observe.
Let $`\omega :𝒜𝖢`$ be a positive linear functional. For $`B𝒜`$, we consider the positive linear functional $`\omega _B(A):=\omega (B^{}AB)=\psi _B,\pi _\omega (A)\psi _B`$. Then $`A\mathrm{ker}\pi _\omega `$ if and only if $`\omega _B(A)=0`$ for all $`B𝒜`$. Using a polarization argument and the fact that $`𝔅(_\omega )`$ is torsion-free (\[9, Prop. 2.8\]), we see that this is equivalent to $`\omega _B(A^{}A)=0`$ for all $`B𝒜`$. Hence we have
$$\mathrm{ker}\pi _\omega =\underset{B𝒜}{}𝒥_{\omega _B}=\underset{B𝒜}{}\mathrm{ker}\omega _B.$$
(4.4)
Note that neither $`𝒥_{\omega _B}`$ nor $`\mathrm{ker}\omega _B`$ is a -ideal in general but only the above intersection. Now consider an arbitrary -representation $`\pi `$ on $``$. For every $`\phi `$, one defines the expectation value $`\omega _\phi (A):=\phi ,\pi (A)\phi `$, which is clearly a positive linear functional of $`𝒜`$. Denoting the Gel’fand ideal by $`𝒥_\phi `$ and the corresponding GNS representation by $`\pi _\phi `$, it follows that $`A\mathrm{ker}\pi _\phi `$ if and only if $`\pi (A)\pi (B)\phi =0`$ for all $`B𝒜`$. But this implies that $`A𝒥_\phi `$ by taking $`B=A^{}`$. Moreover, $`A𝒥_\phi `$ if and only if $`\pi (A)\phi =0`$. So $`A\mathrm{ker}\omega _\phi `$. Thus we have the inclusions
$$\mathrm{ker}\pi _\phi 𝒥_\phi \mathrm{ker}\omega _\phi .$$
(4.5)
###### Lemma 4.5
Let $`𝒜`$ be a -algebra over $`𝖢`$ and $`\pi :𝒜𝔅()`$ a -representation. Consider $`\pi _M:=_\phi \pi _\phi `$. Then one has
$$\mathrm{ker}\pi =\mathrm{ker}\pi _M=\underset{\phi }{}\mathrm{ker}\pi _\phi =\underset{\phi }{}𝒥_\phi =\underset{\phi }{}\mathrm{ker}\omega _\phi .$$
(4.6)
Proof: It remains to show $`_\phi \mathrm{ker}\omega _\phi =\mathrm{ker}\pi `$. But this follows by polarization and the fact that $`𝔅()`$ is torsion-free. $`\mathrm{}`$
###### Lemma 4.6
Let $`𝒜`$ be a -algebra over $`𝖢`$ with an approximate identity. Then for every positive linear functional $`\omega `$ one has
$$\mathrm{ker}\pi _\omega 𝒥_\omega \mathrm{ker}\omega .$$
(4.7)
Proof: This generalizes (4.5). Let $`A𝒜`$ and $`E_\alpha =E_\alpha ^{}𝒜`$ such that $`E_\alpha A=A=AE_\alpha `$. Then $`\pi _\omega (A)=0`$ implies $`0=\pi _\omega (A)\psi _{E_\alpha }=\psi _A`$, whence $`A𝒥_\omega `$. For $`A𝒥_\omega `$, we find by the Cauchy-Schwarz inequality that $`0\omega (A)\overline{\omega (A)}\omega (A^{}A)\omega (E_\alpha ^{}E_\alpha )=0`$ and hence $`A\mathrm{ker}\omega `$. $`\mathrm{}`$
Let us investigate the lattice of closed -ideals more closely. While the lattice of -ideals trivially has a minimal and a maximal element, namely $`\{0\}`$ and $`𝒜`$, the closed -ideals provide a more interesting structure. We define the minimal closed -ideal of $`𝒜`$ by
$$𝒥_{\mathrm{min}}(𝒜):=\underset{\text{ closed }\text{}\text{-ideal}}{},$$
(4.8)
which is again a closed -ideal due to Lemma 4.2. Furthermore $`𝒥_{\mathrm{min}}(𝒜)`$ is clearly the minimal element of the lattice. It is obvious that $`𝒥_{\mathrm{min}}(𝒜)=\{0\}`$ if and only if $`𝒜`$ has a faithful -representation on some pre-Hilbert space. Thus for a $`C^{}`$-algebra $`𝒜`$ one automatically has $`𝒥_{\mathrm{min}}(𝒜)=\{0\}`$. We note that the same holds for star-product algebras (see e.g. \[32, Sect. 4\]).
Since the kernel of a -representation can always be obtained by using direct sums of GNS representations, it is sufficient to consider GNS representations in order to characterize $`𝒥_{\mathrm{min}}(𝒜)`$. With (4.4), this gives
$$𝒥_{\mathrm{min}}(𝒜)=\underset{\omega }{}\mathrm{ker}\pi _\omega =\underset{\omega }{}\underset{B𝒜}{}𝒥_{\omega _B}=\underset{\omega }{}\underset{B𝒜}{}\mathrm{ker}\omega _B,$$
(4.9)
where $`\omega `$ runs over all positive linear functionals. If $`𝒜`$ has in addition an approximate identity, we obtain the following result:
###### Theorem 4.7
Let $`𝒜`$ be a -algebra over $`𝖢`$ with approximate identity. Then the minimal closed -ideal is given by
$$𝒥_{\mathrm{min}}(𝒜)=\underset{\omega }{}\mathrm{ker}\pi _\omega =\underset{\omega }{}𝒥_\omega =\underset{\omega }{}\mathrm{ker}\omega ,$$
(4.10)
where $`\omega `$ runs over all positive linear functionals. If $`A^{}A=0`$, or if $`A`$ is normal and nilpotent, or if $`zA=0`$ for $`0z𝖢`$ then $`A𝒥_{\mathrm{min}}(𝒜)`$.
Proof: Let $`A_\omega \mathrm{ker}\omega `$. Then $`\pi (A)=0`$ for all -representations $`\pi `$ since otherwise there would exist a $`\phi `$ such that $`\omega _\phi (A)=\phi ,\pi (A)\phi 0`$ which contradicts the assumption since $`\omega _\phi `$ is a positive functional. Hence $`A𝒥_{\mathrm{min}}(𝒜)`$ and by (4.7) and (4.9) the equality (4.10) follows. The other statements follow directly from \[9, Prop. 2.8\]. $`\mathrm{}`$
Note that all ‘non-$`C^{}`$-algebra-like’ elements are absorbed into $`𝒥_{\mathrm{min}}(𝒜)`$. In particular $`𝒜/𝒥_{\mathrm{min}}(𝒜)`$ has sufficiently many positive linear functionals and still an approximate identity. This can be seen as an analogue of the construction of the $`C^{}`$-enveloping algebra of a Banach -algebra, see e.g. \[12, Sect. II.7\]. The passage from $`𝒜`$ to $`𝒜/𝒥_{\mathrm{min}}(𝒜)`$ is functorial in the category of -algebras with approximate identities.
###### Proposition 4.8
Let $`𝒜`$, $``$ be -algebras over $`𝖢`$ with approximate identities and let $`\mathrm{\Phi }:𝒜`$ be a -homomorphism. Then $`\mathrm{\Phi }(𝒥_{\mathrm{min}}(𝒜))𝒥_{\mathrm{min}}()`$ and thus there exists a unique -homomorphism $`\varphi :𝒜/𝒥_{\mathrm{min}}(𝒜)/𝒥_{\mathrm{min}}()`$ such that
$$\begin{array}{ccc}𝒜& \stackrel{\mathrm{\Phi }}{}& \\ & & \\ 𝒜/𝒥_{\mathrm{min}}(𝒜)& \stackrel{\varphi }{}& /𝒥_{\mathrm{min}}()\end{array}$$
(4.11)
commutes. Thus the passage from $`𝒜`$ to $`𝒜/𝒥_{\mathrm{min}}(𝒜)`$ is functorial.
Proof: If $`\omega :𝖢`$ is a positive linear functional then $`\mathrm{\Phi }^{}\omega =\omega \mathrm{\Phi }:𝒜𝖢`$ is positive, too, whence $`\mathrm{\Phi }(𝒥_{\mathrm{min}}(𝒜))𝒥_{\mathrm{min}}()`$ follows from (4.10). Then (4.11) is clear. $`\mathrm{}`$
###### Remark 4.9
The minimal ideal as discussed here is related to the work of Handelman , as P. Ara brought to our attention. Let $`𝒜`$ be a (unital) -algebra over $`𝖢=𝖱(\mathrm{i})`$ with positive definite involution (\[15, pp. 338\]) and suppose $`𝖱`$ (so that $`𝖱`$ is archimedean-ordered). Let $`J^{}(𝒜)`$ be the set of bounded elements of $`𝒜`$ with zero seminorm (\[15, Sect.1\]). We note that, in general, $`J^{}(𝒜)𝒥_{\mathrm{min}}(𝒜_b)`$. But if $`𝖢=`$, then in fact $`J^{}(𝒜)=𝒥_{\mathrm{min}}(𝒜_b)`$, since in this case the -algebra $`𝒜_b/J^{}(𝒜)`$ can be completed to a $`C^{}`$-algebra and hence admits a faithful -representation.
## 5 Morita Equivalent -Algebras
In this section, we shall discuss the behavior of the lattice of (closed) -ideals under algebraic Rieffel induction and prove that the lattice of (closed) -ideals is a ‘formal Morita invariant’.
Let $`{}_{}{}^{}𝔛_{𝒜}^{}`$ be a $`(`$-$`𝒜)`$-bimodule with an $`𝒜`$-valued inner product (see Section 2). We define the map $`\mathrm{\Phi }_𝔛:2^𝒜2^{}`$ by
$$\mathrm{\Phi }_𝔛():=\{B|x,y{}_{}{}^{}𝔛_{𝒜}^{}:x,By_𝒜\}.$$
(5.1)
If we consider (closed) -ideals of $`𝒜`$, then we obtain the following properties of $`\mathrm{\Phi }_𝔛`$:
###### Proposition 5.1
Let $`𝒜`$, $``$ be -algebras over $`𝖢`$ and let $`{}_{}{}^{}𝔛_{𝒜}^{}`$ be a $`(`$-$`𝒜)`$-bimodule with $`𝒜`$-valued inner product. Then $`\mathrm{\Phi }_𝔛`$ maps -ideals of $`𝒜`$ to -ideals of $``$ and satisfies
$$\mathrm{\Phi }_𝔛()\mathrm{\Phi }_𝔛(𝒥)=\mathrm{\Phi }_𝔛(𝒥)\text{and}𝒥\mathrm{\Phi }_𝔛()\mathrm{\Phi }_𝔛(𝒥)$$
(5.2)
for all -ideals $``$, $`𝒥`$ of $`𝒜`$. If in addition the inner product is positive semi-definite and satisfies P, then $`\mathrm{\Phi }_𝔛`$ maps closed -ideals to closed -ideals, satisfies (5.2) with respect to the lattice of closed -ideals, and one has
$$\mathrm{ker}(_𝔛\pi _𝒜)=\mathrm{\Phi }_𝔛(\mathrm{ker}\pi _𝒜)$$
(5.3)
for any -representation $`\pi _𝒜`$ of $`𝒜`$.
Proof: The first part is a simple computation thus assume that the inner product is positive semi-definite and satisfies P. Let $`\pi _𝒜`$ be a -representation of $`𝒜`$ on $``$ and let $`\pi _{}=_𝔛\pi _𝒜`$ be the induced -representation of $``$. Then $`\pi _{}(B)=0`$ if and only if $`[By\varphi ]=0`$ for all $`y{}_{}{}^{}𝔛_{𝒜}^{}`$ and $`\varphi `$ since the equivalence classes of those $`y\varphi `$ span $`𝔎`$. But since the inner product of $`𝔎`$ is non-degenerate this is equivalent to $`0=[x\psi ],[By\varphi ]_𝔎=\psi ,\pi _𝒜(x,By_𝒜)\varphi _{}`$ for all $`x,y{}_{}{}^{}𝔛_{𝒜}^{}`$ and $`\varphi ,\psi `$. Thus it is equivalent to $`\pi _𝒜(x,By_𝒜)=0`$ and hence $`B\mathrm{\Phi }_𝔛(\mathrm{ker}\pi _𝒜)`$ which implies (5.3). Hence closed -ideals are mapped to closed -ideals and since $``$ and $``$ are given by $``$ and $``$ the proof is completed. $`\mathrm{}`$
Note that from (5.2) one obtains that $`\mathrm{\Phi }_𝔛`$ is almost a homomorphism of lattices, i.e. we have
$$\mathrm{\Phi }_𝔛()\mathrm{\Phi }_𝔛(𝒥)\mathrm{\Phi }_𝔛(𝒥)$$
(5.4)
for all -ideals or closed -ideals, respectively. Note that in general we only can guarantee ‘$``$’ since there may exist more and perhaps smaller (closed) -ideals of $``$ which are not in the image of $`\mathrm{\Phi }_𝔛`$.
In order to discuss Morita equivalence, we need the following technical lemma, which does not use any positivity requirements but only the -structures.
###### Lemma 5.2
Let $`𝒜`$, $``$ be -algebras over $`𝖢`$ and $`{}_{}{}^{}𝔛_{𝒜}^{}`$ a $`(`$-$`𝒜)`$-bimodule with compatible $`𝒜`$\- and $``$-valued inner products and let $`𝒜`$ be a -ideal. Then $`A\mathrm{\Phi }_{\overline{𝔛}}(\mathrm{\Phi }_𝔛())`$ if and only if for all $`x,y,z,w{}_{}{}^{}𝔛_{𝒜}^{}`$ one has
$$x,y_𝒜Az,w_𝒜.$$
(5.5)
Proof: By definition $`A\mathrm{\Phi }_{\overline{𝔛}}(\mathrm{\Phi }_𝔛())`$ if and only if one has $`y,\overline{x},A\overline{z}_{}w_𝒜`$ for all $`\overline{x},\overline{z}{}_{𝒜}{}^{}\overline{𝔛}_{}^{}`$ and $`y,w{}_{}{}^{}𝔛_{𝒜}^{}`$. Then it is a simple computation to obtain the condition (5.5). $`\mathrm{}`$
If in addition the $`𝒜`$-valued inner product is full, then condition (5.5) means that the ideal generated by $`A`$ must be contained in $``$. If furthermore $`𝒜`$ has an approximate identity, then we have $`\mathrm{\Phi }_{\overline{𝔛}}(\mathrm{\Phi }_𝔛())=`$ for all -ideals $``$. By symmetry in $`𝒜`$ and $``$, we end up with the following proposition.
###### Proposition 5.3
Let $`𝒜`$, $``$ be -algebras over $`𝖢`$ with approximate identities and let $`{}_{}{}^{}𝔛_{𝒜}^{}`$ be a $`(`$-$`𝒜)`$-bimodule with compatible full $``$\- and $`𝒜`$-valued inner products. Then the lattices of -ideals are isomorphic via $`\mathrm{\Phi }_𝔛`$ and $`\mathrm{\Phi }_{\overline{𝔛}}`$.
Proof: The fact that $`\mathrm{\Phi }_𝔛`$ and $`\mathrm{\Phi }_{\overline{𝔛}}`$ are inverse to each other follows from the above considerations. Hence it remains to show that $`\mathrm{\Phi }_𝔛`$ (and thus $`\mathrm{\Phi }_{\overline{𝔛}}`$) is a lattice homomorphism. But this follows from bijectivity and Proposition 5.1. $`\mathrm{}`$
In order to make statements about closed -ideals, we need the positivity requirements for the inner products to guarantee that $`_𝔛\pi _𝒜`$ is indeed a -representation.
###### Theorem 5.4
Let $`𝒜`$, $``$ be idempotent -algebras over $`𝖢`$ and let $`{}_{}{}^{}𝔛_{𝒜}^{}`$ be an equivalence bimodule. Then the lattices of closed -ideals are isomorphic via $`\mathrm{\Phi }_𝔛`$ and $`\mathrm{\Phi }_{\overline{𝔛}}`$.
Proof: Let $`=\mathrm{ker}\pi _𝒜`$. We may assume, due to Lemma 4.3, that $`\pi _𝒜`$ is strongly non-degenerate. Hence, by \[9, Lem. 5.7\], we know that $`_{\overline{𝔛}}_𝔛\pi _𝒜`$ is unitarily equivalent to $`\pi _𝒜`$. So they have the same kernel and by (5.3) it follows that $`\mathrm{\Phi }_{\overline{𝔛}}(\mathrm{\Phi }_𝔛())=`$. Analogously, one obtains $`\mathrm{\Phi }_𝔛(\mathrm{\Phi }_{\overline{𝔛}}(𝒥))=𝒥`$ for all closed -ideals of $``$ and hence $`\mathrm{\Phi }_𝔛`$ and $`\mathrm{\Phi }_{\overline{𝔛}}`$ are inverse to each other. The homomorphism properties follow as in Proposition 5.3. $`\mathrm{}`$
Note that the so-called Rieffel correspondence theorem (\[26, Thm. 3.24\]) can be recovered as a consequence of the above theorem. Indeed, if $`𝒜`$ and $``$ are strongly Morita equivalent $`C^{}`$-algebras, then by Theorem 3.7 it follows that $`𝒫_𝒜`$ and $`𝒫_{}`$ are formally Morita equivalent. So $`𝔏(𝒫_𝒜)`$ and $`𝔏(𝒫_{})`$ are isomorphic by Theorem 5.4 and therefore $`𝔏(𝒜)`$ and $`𝔏()`$ are isomorphic by Proposition 4.4. We also observe the following corollary.
###### Corollary 5.5
Let $`𝒜`$ and $``$ be non-degenerate and idempotent -algebras over $`𝖢`$ which are formally Morita equivalent. Then $`𝒥_{\mathrm{min}}(𝒜)=0`$ if and only if $`𝒥_{\mathrm{min}}()=0`$.
Proof: Let $`{}_{}{}^{}𝔛_{𝒜}^{}`$ be an equivalence bimodule. Note that if $`𝒜`$ is non-degenerate and idempotent, then $`{}_{}{}^{}x,yA=0`$ for all $`x,y𝔛`$ implies that $`A=0`$ (see the proof of Lemma 3.1). Hence the conclusion follows. $`\mathrm{}`$
Let us finally discuss the general phenomena of -algebras having equivalent (even isomorphic) categories of (strongly non-degenerate) -representations but which are not formally Morita equivalent. The following discussion generalizes the example of $`𝖢`$ and $`^{}(𝖢^n)`$ from \[9, Cor. 5.20\].
Consider the categories of -representations of $`𝒜`$ and $`𝒜/𝒥_{\mathrm{min}}(𝒜)`$. More generally, we shall consider a -ideal $`𝒥𝒥_{\mathrm{min}}(𝒜)`$ and the -algebra $`𝒜/𝒥`$. Denoting by $`\rho :𝒜𝒜/𝒥`$ the -homomorphism $`A[A]`$ we can pull-back -representations $`\pi :𝒜/𝒥𝔅()`$ by
$$\rho ^{}\pi (A):=\pi (\rho (A))=\pi ([A]),$$
(5.6)
and obtain a -representation $`\rho ^{}\pi `$ of $`𝒜`$ on the same pre-Hilbert space $``$. This is clearly functorial, i.e. compatible with intertwiners. Moreover, if $`\pi `$ is also strongly non-degenerate then $`\rho ^{}\pi `$ is strongly non-degenerate. On the other hand, let $`\pi :𝒜𝔅()`$ be a -representation of $`𝒜`$ then $`𝒥𝒥_{\mathrm{min}}(𝒜)\mathrm{ker}\pi `$. Hence we can also push-forward the -representation and obtain a -representation $`\rho _{}\pi `$ of $`𝒜/𝒥`$ on the same representation space
$$\rho _{}\pi ([A]):=\pi (A).$$
(5.7)
Again this is functorial and if $`\pi `$ is strongly non-degenerate then $`\rho _{}\pi `$ is also strongly non-degenerate. Collecting these results, we obtain
###### Proposition 5.6
Let $`𝒜`$ be a -algebra over $`𝖢`$ and $`𝒥𝒥_{\mathrm{min}}(𝒜)`$ a -ideal. Then the functors $`\rho ^{}`$ and $`\rho _{}`$ implement an isomorphism of the categories $`{}_{}{}^{}\text{-}\mathrm{𝗋𝖾𝗉}(𝒜)`$ and $`{}_{}{}^{}\text{-}\mathrm{𝗋𝖾𝗉}(𝒜/𝒥)`$ as well as of $`{}_{}{}^{}\text{-}\mathrm{𝖱𝖾𝗉}(𝒜)`$ and $`{}_{}{}^{}\text{-}\mathrm{𝖱𝖾𝗉}(𝒜/𝒥)`$.
We also note that $`𝒜`$ and $`𝒜/𝒥`$ have isomorphic lattices of closed -ideals, for any -ideal $`𝒥𝒥_{\mathrm{min}}(𝒜)`$. Let $`=\mathrm{ker}\pi 𝔏(𝒜)`$. Since $`𝒥𝒥_{\mathrm{min}}(𝒜)\mathrm{ker}\pi `$, it follows that $`\rho ()=\mathrm{ker}\rho _{}\pi `$. Thus $`\rho `$ induces a map $`𝔏(𝒜)𝔏(𝒜/𝒥)`$ which can be easily checked to be a bijection which preserves ordering, and hence is a lattice isomorphism. We can then state
###### Proposition 5.7
Let $`𝒜`$ be a -algebra over $`𝖢`$ and $`𝒥𝒥_{\mathrm{min}}(𝒜)`$ a -ideal. Then the natural projection $`\rho :𝒜𝒜/𝒥`$ induces an isomorphism of $`𝔏(𝒜)`$ and $`𝔏(𝒜/𝒥)`$. In particular, $`𝒥_{\mathrm{min}}(𝒜/𝒥)=𝒥_{\mathrm{min}}(𝒜)/𝒥`$.
We finally observe that this discussion provides a large class of examples of -algebras over $`𝖢`$ having equivalent categories of -representations, isomorphic lattice of -ideals but which are not formally Morita equivalent.
###### Proposition 5.8
Let $`𝒜`$ be a -algebra over $`𝖢`$ with approximate identity and let $``$ be a -ideal contained in $`𝒥_{\mathrm{min}}(𝒜)`$ but $`𝒥_{\mathrm{min}}(𝒜)`$. Then $`𝒜/`$ and $`𝒜/𝒥_{\mathrm{min}}(𝒜)`$ have isomorphic categories of (strongly non-degenerate) -representations and isomorphic lattices of closed -ideals, but they are not formally Morita equivalent.
Proof: The first statement follows directly from Propositions 5.6 and 5.7 by transitivity. It also follows from Proposition 4.4 that $`𝒜/𝒥_{\mathrm{min}}(𝒜)`$ has trivial minimal ideal, while $`𝒥_{\mathrm{min}}(𝒜/)=𝒥_{\mathrm{min}}(𝒜)/0`$. So the result follows by Corollary 5.5. $`\mathrm{}`$
## 6 Examples from Differential Geometry
In this section, we describe closed -ideals and minimal ideals of -algebras arising in differential geometry. We start with a description of the closed -ideals in $`C^{\mathrm{}}(M)`$, the algebra of complex-valued smooth functions on an $`m`$-dimensional smooth manifold $`M`$.
Let us first recall that given a positive Borel measure $`\mu `$ on $`M`$, with compact support, we can define a positive linear functional $`\omega _\mu `$ on $`C^{\mathrm{}}(M)`$ by integration with respect to this measure and consider the corresponding GNS representation of $`C^{\mathrm{}}(M)`$, denoted by $`\pi _\mu `$. We observe that in this case, $`\mathrm{ker}\pi _\mu =\{fC^{\mathrm{}}(M)\text{ with }f|g|^2𝑑\mu =0\text{ for all }gC^{\mathrm{}}(M)\}`$ coincides with the Gel’fand ideal $`\{fC^{\mathrm{}}(M)\text{ with }|f|^2𝑑\mu =0\}`$. To each closed set $`FM`$, there naturally corresponds a -ideal of $`C^{\mathrm{}}(M)`$ given by
$$_F=\{fC^{\mathrm{}}(M)\text{ such that }f|_F=0\},$$
(6.1)
called the *vanishing ideal* of $`F`$. The following proposition shows that these ideals are exactly the closed ideals of $`C^{\mathrm{}}(M)`$.
###### Proposition 6.1
A -ideal $`C^{\mathrm{}}(M)`$ is closed if and only if it is the vanishing ideal of some closed set $`FM`$. Moreover, $`=\mathrm{ker}\pi _\omega `$, for some GNS representation $`\pi _\omega `$, if and only if $`=_K`$ for some compact set $`KM`$.
Proof: Let $`F\mathrm{}`$ be a closed set in $`M`$. Denote by $`_F`$ the set of positive compactly supported Borel measures on $`M`$ with support in $`F`$. By the continuity of functions in $`C^{\mathrm{}}(M)`$, it follows that $`_F=_{\mu _F}\mathrm{ker}\pi _\mu `$. Thus $`_F`$ is closed. Conversely, if $``$ is closed, then it is the intersection of kernels of GNS representations of some positive linear functionals. But the positive linear functionals in $`C^{\mathrm{}}(M)`$ are given by positive Borel measures, see e.g. \[9, App. B\]. So we can write $`=\mathrm{ker}\pi `$, for $`\pi =_\alpha \pi _{\mu _\alpha }`$ and define $`F=\overline{_\alpha supp\mu _\alpha }`$. It is then easy to check that $`=_F`$. The last assertion follows from the fact that if $`K`$ is compact, then there is a positive Borel measure with support exactly $`K`$. $`\mathrm{}`$
We recall that the vanishing ideals of closed sets in $`C^{\mathrm{}}(M)`$ are exactly the closed ideals with respect to the (weak) $`C^0`$-topology induced from $`C(M)`$, that is, the locally convex topology generated by the family of semi-norms $`\rho _K(f)=sup_{xK}|f(x)|`$, for $`KM`$ compact (this follows from \[35, Thm. 3\]). Hence we have
###### Corollary 6.2
A -ideal $``$ in $`C^{\mathrm{}}(M)`$ is closed (in the algebraic sense) if and only if it is closed with respect to the (weak) $`C^0`$-topology.
We remark that similar results hold for the algebra $`C_0^{\mathrm{}}(M)`$ of compactly supported smooth complex-valued functions on $`M`$.
###### Proposition 6.3
Let $`C_0^{\mathrm{}}(M)`$ be a -ideal. Then $``$ is closed if and only if $``$ is a vanishing ideal if and only if $``$ is the kernel of a GNS representation.
It is clear that the minimal ideal of $`C^{\mathrm{}}(M)`$ is trivial. We will now consider an example with nontrivial minimal ideal, namely the sections of the Grassmann algebra bundle over $`M`$. This generalizes the example of the Grassmann algebra discussed in the introduction. We denote by $`^{}T^{}M`$ the complexified Grassmann algebra bundle over $`M`$. There are several possibilities to define fibrewise -involutions for $`^{}T^{}M`$. We do it by declaring the real one-forms to be Hermitian and extending this to an antilinear anti-automorphism of the Grassmann algebra bundle. Then $`\mathrm{\Gamma }^{\mathrm{}}(^{}T^{}M)`$ becomes a -algebra over $``$. Denote by $`^+T^{}M`$ the bundle of forms of positive degree. We have the following characterization of the minimal ideal of $`\mathrm{\Gamma }^{\mathrm{}}(^{}T^{}M)`$.
###### Proposition 6.4
The minimal ideal of $`\mathrm{\Gamma }^{\mathrm{}}(^{}T^{}M)`$ is given by $`\mathrm{\Gamma }^{\mathrm{}}(^+T^{}M)`$. Thus the representation theories of $`\mathrm{\Gamma }^{\mathrm{}}(^{}T^{}M)`$ and $`C^{\mathrm{}}(M)`$ are equivalent but the algebras are not formally Morita equivalent. The closed ideals of $`\mathrm{\Gamma }^{\mathrm{}}(^{}T^{}M)`$ are of the form $`\mathrm{\Gamma }^{\mathrm{}}(^+T^{}M)`$, where $``$ is a closed ideal of $`C^{\mathrm{}}(M)`$.
Proof: Due to the nilpotence properties of forms with positive degree, it is easy to see that $`\mathrm{\Gamma }^{\mathrm{}}(^+T^{}M)`$ has to be contained in the minimal ideal. On the other hand, $`C^{\mathrm{}}(M)=\mathrm{\Gamma }^{\mathrm{}}(^0T^{}M)`$ has sufficiently many positive linear functionals and thus the first statement follows. The other statement is clear from the above results and Section 5. $`\mathrm{}`$
Let $`EM`$ be a complex Hermitian vector bundle over $`M`$, with $`k`$-dimensional fibers and fiber metric $`h`$. We write $`\mathrm{𝖤𝗇𝖽}(E)M`$ for the corresponding endomorphism bundle and $`\mathrm{\Gamma }^{\mathrm{}}(E)`$, $`\mathrm{\Gamma }^{\mathrm{}}(\mathrm{𝖤𝗇𝖽}(E))`$ for the spaces of smooth sections of $`E`$ and $`\mathrm{𝖤𝗇𝖽}(E)`$, respectively. We denote the trivial bundle $`M\times ^kM`$ by $`t(^k)`$. The final part of this section shows that $`\mathrm{\Gamma }^{\mathrm{}}(E)`$ is a ($`\mathrm{\Gamma }^{\mathrm{}}(\mathrm{𝖤𝗇𝖽}(E))`$-$`C^{\mathrm{}}(M)`$)-equivalence bimodule and gives a description of the closed -ideals in $`\mathrm{\Gamma }^{\mathrm{}}(\mathrm{𝖤𝗇𝖽}(E))`$.
We recall that $`E`$ can always be regarded as a subbundle of some trivial bundle $`t(^N)`$ over $`M`$ (the proof in \[16, Ch. I, Thm. 6.5\] works for arbitrary manifolds since in this case vector bundles are of finite type, see \[20, Sect. 5\], \[21, Lem. 2.7\]). Denote $`C^{\mathrm{}}(M)`$ by $`𝒜`$. The natural Hermitian metric in $`t(^N)`$ induces an $`𝒜`$-valued inner product on $`\mathrm{\Gamma }^{\mathrm{}}(t(^N))𝒜^N`$ and the embedding $`i:Et(^N)`$ can be assumed to preserve metrics (\[16, Ch. I, Thm. 8.8\]). Identifying $`\mathrm{\Gamma }^{\mathrm{}}(E)`$ with a submodule of $`𝒜^N`$, we can write $`\mathrm{\Gamma }^{\mathrm{}}(E)\mathrm{\Gamma }^{\mathrm{}}(E)^{}𝒜^N`$ and hence there is a projection $`QM_N(𝒜)`$ so that $`\mathrm{\Gamma }^{\mathrm{}}(E)Q𝒜^N`$ as right $`𝒜`$-modules. Note also that $`QM_N(𝒜)Q\mathrm{\Gamma }^{\mathrm{}}(\mathrm{𝖤𝗇𝖽}(E))`$ as -algebras. It is clear that $`\mathrm{\Gamma }^{\mathrm{}}(E)`$ is a ($`\mathrm{\Gamma }^{\mathrm{}}(\mathrm{𝖤𝗇𝖽}(E))`$-$`C^{\mathrm{}}(M)`$)-bimodule equipped with $`\mathrm{\Gamma }^{\mathrm{}}(\mathrm{𝖤𝗇𝖽}(E))`$ and $`C^{\mathrm{}}(M)`$ valued inner products given by $`(u,v)\mathrm{\Theta }_{u,v}\mathrm{\Gamma }^{\mathrm{}}(\mathrm{𝖤𝗇𝖽}(E))`$, where $`\mathrm{\Theta }_{u,v}(w)=uh(v,w)`$, and $`(u,v)h(u,v)C^{\mathrm{}}(M)`$, for $`u,v,w\mathrm{\Gamma }^{\mathrm{}}(E)`$. We finally observe that this bimodule structure on $`\mathrm{\Gamma }^{\mathrm{}}(E)`$ (with the algebra valued inner products) is just the one considered for $`Q𝒜^N`$ in \[9, Sect. 6\] after the aforementioned identifications. Recall that a projection $`QM_N(𝒜)`$ is called *full* if $``$-span $`\{TQS;T,SM_N(𝒜)\}=M_N(𝒜)`$. The following observation is well-known.
###### Lemma 6.5
Let $`QM_N(C^{\mathrm{}}(M))`$ be a non-zero projection. Then $`Q`$ is full.
It then follows from \[9, Sect. 6\] that $`\mathrm{\Gamma }^{\mathrm{}}(E)`$ satisfies all the properties of a ($`\mathrm{\Gamma }^{\mathrm{}}(\mathrm{𝖤𝗇𝖽}(E))`$-$`C^{\mathrm{}}(M)`$)-equivalence bimodule, except possibly for property Q, which we now show holds in general. Let $`(\pi ,)`$ be a -representation of $`QM_N(𝒜)Q\mathrm{\Gamma }^{\mathrm{}}(\mathrm{𝖤𝗇𝖽}(E))`$ and $`\stackrel{~}{𝔎}=\overline{Q𝒜^N}_{QM_N(𝒜)Q}`$. Note that $`\overline{Q𝒜^N}`$ is again just $`\mathrm{\Gamma }^{\mathrm{}}(E)`$ (with left and right actions given by the adjoint of the previous ones). Given $`s_i\mathrm{\Gamma }^{\mathrm{}}(E)`$ and $`\psi _i`$, we must show that
$$\underset{i}{}s_i\psi _i,\underset{j}{}s_j\psi _j_{\stackrel{~}{𝔎}}0,$$
(6.2)
where $`s_1\psi _1,s_2\psi _2_{\stackrel{~}{𝔎}}=\psi _1,\pi (s_1,s_2)\psi _2_{}`$. We then have
###### Lemma 6.6
$`\mathrm{\Gamma }^{\mathrm{}}(E)`$ satisfies property Q.
Proof: Assuming $`E=t(^k)`$, the $`s_i`$’s are just $`^k`$-valued functions on $`M`$ and $`\mathrm{𝖤𝗇𝖽}(E)=t(M_k())`$. Note that if $`e_j`$ denote the canonical section basis of $`E`$, then for each $`i`$ one can find a $`T_it(M_k())`$ so that $`e_1T_i=s_i`$ (recall that $`e_1T_i`$ means $`T_{i}^{}{}_{}{}^{}e_1`$, where the last expressions is just the usual matrix multiplication). For instance, define $`T_i`$ to be the matrix valued function with $`s_i`$ as the first row and zero everywhere else. Now the proof follows the usual trick: $`_is_i\psi _i,_js_j\psi _j_{\stackrel{~}{𝔎}}=_ie_1T_i\psi _i,_je_1T_j\psi _j_{\stackrel{~}{𝔎}}=e_1_iT_i\psi _i,e_1_jT_j\psi _j_{\stackrel{~}{𝔎}}0`$.
To prove the result in general, let $`U_k`$, $`k=1,\mathrm{},r`$ be an open cover of $`M`$ so that $`E|_{U_k}`$ is trivial. Let $`\chi _k`$ be a quadratic partition of unity subordinated to this cover. Then observe that for $`s_1,s_2\mathrm{\Gamma }^{\mathrm{}}(E)`$, we have that $`s_1,s_2_{QM_N(𝒜)Q}=_k\chi _ks_1,\chi _ks_2_{QM_N(𝒜)Q}`$ (since $`_k\chi _ks_1,\chi _ks_2(x)=_k\mathrm{\Theta }_{\chi _ks_1(x),\chi _ks_2(x)}=_k\chi _k(x)\overline{\chi _k}(x)\mathrm{\Theta }_{s_1(x),s_2(x)}=s_1,s_2(x)`$). Hence,
$$\underset{i}{}s_i\psi _i,\underset{i}{}s_i\psi _i_{\stackrel{~}{𝔎}}=\underset{i,j}{}\psi _i,s_i,s_j\psi _j_{\stackrel{~}{𝔎}}=\underset{i,j,k}{}\psi _i,\chi _ks_i,\chi _ks_j\psi _j_{\stackrel{~}{𝔎}}.$$
Now note that the last expression is just $`_k_i(\chi _ks_i)\psi _i,_i(\chi _ks_i)\psi _i`$. But since $`supp\chi _ks_iU_k`$, and $`E|_{U_k}`$ is trivial, the result follows from the previous discussion. $`\mathrm{}`$
###### Proposition 6.7
If $`E`$ is a complex hermitian vector bundle over $`M`$, then $`\mathrm{\Gamma }^{\mathrm{}}(E)`$ is a $`(C^{\mathrm{}}(M)`$-$`\mathrm{\Gamma }^{\mathrm{}}(\mathrm{𝖤𝗇𝖽}(E)))`$-equivalence bimodule.
One can also consider compactly supported functions and sections. In this case, $`C_0^{\mathrm{}}(M)`$ and $`\mathrm{\Gamma }_0^{\mathrm{}}(\mathrm{𝖤𝗇𝖽}(E))`$ are no longer unital in general, but have approximate identities. So the results of \[9, Sect. 6\] can still be applied and the same proof used for Prop. 6.7 shows that
###### Proposition 6.8
$`\mathrm{\Gamma }_0^{\mathrm{}}(E)`$ is a $`(C_0^{\mathrm{}}(M)`$-$`\mathrm{\Gamma }_0^{\mathrm{}}(\mathrm{𝖤𝗇𝖽}(E)))`$-equivalence bimodule.
We finally remark that combining the results of Section 5, one easily arrives at a characterization of the closed -ideals of $`\mathrm{\Gamma }^{\mathrm{}}(\mathrm{𝖤𝗇𝖽}(E))`$ (and analogously for $`\mathrm{\Gamma }_0^{\mathrm{}}(\mathrm{𝖤𝗇𝖽}(E))`$).
###### Corollary 6.9
A -ideal $`\mathrm{\Gamma }^{\mathrm{}}(\mathrm{𝖤𝗇𝖽}(E))`$ is closed if and only it is a vanishing ideal, i.e. of the form $`\{B\mathrm{\Gamma }^{\mathrm{}}(\mathrm{𝖤𝗇𝖽}(E))|B(x)=0xF\}`$ for some closed $`FM`$.
## Acknowledgements
We would like to thank Prof. Ara for his clarifying remarks and for pointing out to us the paper and Martin Bordemann for many useful discussions, in particular concerning the minimal ideal. We would also like to thank Pierre Bieliavsky, Michel Cahen, Laura DeMarco, Diogo Gomes, Simone Gutt, Daniel Markiewicz, and Alan Weinstein for helpful remarks and discussions.
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# Multiwavelength study of the nuclei of a volume-limited sample of galaxies I: X-ray observations
## 1 Introduction
Evidence has been growing that both (i) weak nuclear activity, and (ii) the presence of quiescent black holes, are common features of normal galaxies. The best evidence to date for (i) comes from the very thorough high S/N spectroscopic survey of 486 nearby galaxies by ?. They found that 86% of galaxies have emission lines; roughly a third have Seyfert 2 or LINER spectra, indicating but not proving some kind of weak AGN; and 10% have weak broad emission lines, so are almost certainly AGN (?). These latter objects are, however, three orders of magnitude less luminous than classic Seyfert galaxies and six orders of magnitude less luminous than the most luminous quasars. They have become known as ‘dwarf AGN’. The evidence for (ii) comes from dynamical studies. The review by ? found evidence for ‘Massive Dark Objects’ (MDOs) in 20% of galaxies, but this can be considered a lower limit because of the difficulty of rigorously demonstrating the requirement for a dark mass. With a more relaxed analysis, ? have claimed that there is strong evidence for an MDO in essentially all galaxies.
Most interestingly, ? and ? also claim a correlation between the MDO mass and the stellar bulge mass, with $`M_{MDO}/M_{bulge}0.0030.006`$. Likewise, various authors have shown evidence that AGN activity correlates with galaxy size (??????). A weakness with all these studies is that only fairly luminous galaxies are considered. As well as investigating whether feeble AGN are present in normal galaxies, it is very important to know whether AGN occur at all in feeble galaxies. Is there a minimum galaxy size below which quasar-like nuclei are simply not present? There are two known examples of AGN in relatively small galaxies: a Seyfert 1 nucleus in NGC 4395 (??) and a Seyfert 2 nucleus in G1200-2038 (?). Both these galaxies have $`M_B18`$, an order of magnitude less luminous than typical $`L^{}`$ galaxies and two orders of magnitude less luminous than the giant ellipticals which dominate the ? sample. Are they rare freaks, or the tip of the iceberg? Furthermore, in the local field there are are plenty of spiral galaxies an order of magnitude less luminous still than NGC 4395 and below that we have true dwarf galaxies, dwarf irregulars, dwarf ellipticals, and dwarf spheroidals, which continue all the way down to $`M_B7`$. Where does AGN activity stop?
To address the issues of (a) how common weak AGN are, and (b) how activity is connected to host galaxy size and type, we have undertaken a multi-wavelength study of a complete volume-limited sample of 46 very nearby AGN, with (broad-band and emission line) optical, IR, X-ray and radio imaging, and long slit spectroscopy. In terms of survey size and spectroscopic quality it would be hard to improve on the ? survey. Instead, our study has several key features. (i) The multiwavelength approach gives the greatest chance of finding a weak AGN, being sensitive to a variety of signatures, and gives us diagnostic information on the nature of the objects seen. (ii) All our galaxies are closer than 7 Mpc, maximizing spatial resolution, and minimizing the detectable luminosity. (iii) The volume-limited nature (as opposed to magnitude-limited) of the sample means that we have a fairly even distribution of galaxy sizes and types, going down to extremely small galaxies. (iv) In many cases we found that the location of the true nucleus is not entirely obvious. Multi-wavelength imaging gave us a ‘shooting list’ of multiple targets for spectroscopy to ensure we did not miss the true nucleus.
This paper (Paper I) focuses on the analysis of the ROSAT HRI X-ray data for the galaxies in our sample and presents some preliminary scientific results. Section 2 describes the main characteristics of the sample. In section 3 the acquisition and analysis of the X-ray data are described. The atlas of X-ray images is found in section 4. In Section 5 we describe detailed results for individual galaxies. In section 6 we summarize and discuss some of the most interesting results and examine the correlation of X-ray luminosity and host galaxy luminosity.
In subsequent papers we will present optical-IR imaging and radio maps (Paper II), long-slit spectroscopy (Paper III), and a final analysis of results (Paper IV).
## 2 The volume-limited sample
To study the occurrence of nuclear activity in galaxies to the lowest possible level, a volume-limited sample of galaxies was defined. The sample contains the nearest examples of various morphological types of galaxies, and a representative range of intrinsic luminosities (Johnson 1997; see also Paper II).
The starting point was the Kraan-Korteweg & Tamman Catalogue (??) with all galaxies with V $`<`$ 500 km/s giving a distance-complete sample within $`d=0.35\times d_{\mathrm{Virgo}}`$ (after applying a virgocentric flow model with an infall velocity of the Local Group equal to 220 km s<sup>-1</sup>). With $`d_{\mathrm{Virgo}}=21.5`$ Mpc, the sample is complete for nearby galaxies within 7.6 Mpc. The catalogue is not complete for intrinsically faint galaxies without velocity information, but it includes all galaxies down to $`M_B13`$. Comparing the distribution of absolute magnitudes of the galaxies in the KKT with the galaxies in the Revised Shapley Ames (magnitude-limited) catalogue (?) it is found that the RSA sample of galaxies presents a strong peak around $`M_B=21`$ and has less than 10% of the objects below $`M_B=18`$, while the KKT sample shows a broad distribution with a maximum around $`M_B=16`$ and a significant number of galaxies down to magnitude $`M_B=13`$.
The observations were restricted to objects with $`\delta >35^{}`$. Because the aim was to observe objects with well-formed nuclei, galaxies classified as Sdm, Sm and Irr were removed from the sample. This gave a basic sample of 46 nearby galaxies (26 spirals, 5 E/S0, 12 dwarf ellipticals and 3 blue compact dwarfs, of which two are classified as Sm spirals and one is a peculiar object). The assumed distances for the galaxies are based on the observed recession velocities and the assumed Virgo flow model.
Table 1 shows the resulting sample on which our study is based, the properties of the galaxies concerned, and a summary of the observations available to date. Note that only 34/46 galaxies have available X-ray data. However, we can define a very useful sub-sample of all those galaxies with $`M_B<14`$ (these are all disk systems, except for the elliptical galaxy NGC 221). Of these, 29/31 have X-ray data – only UGC 7321 and NGC 7793 are missing. In addition this paper describes observations of five further dwarf galaxies: NGC 147, NGC 185, UGC 6456, Leo B, and A 0951+68. Only one of these has a detected X-ray source. The analysis we present in section 6 is based on the almost-complete subsample with $`M_B<14`$.
## 3 High resolution X-ray imaging
### 3.1 Observations
The data reported here come from pointed observations carried out with the ROSAT HRI. Table 2 shows those galaxies in the volume-limited sample with available data. Observations for 17 galaxies were awarded to this project through Announcement of Opportunity (AO) calls. For 13 objects data were retrieved from the public archives using arnie (a World Wide Web interface to the databases and catalogues supported by the Leicester Database and Archive Service, ledas). Galaxies are identified in table 2 with an ‘N’ (new data allocated to this project) or a ‘P’ (public data retrieved from the archive). Since NGC 4826 has both new and archive observations a total of 29 objects are covered. The ROSAT HRI data have been reduced and analyzed as described below.
For several galaxies more than one exposure was available. This could be because several observations were requested, or because a single request was scheduled as several observations months apart. In table 2 observations of the same galaxy associated with different requests are distinguished by their ROR (ROSAT Observation Request) numbers. If a single request was scheduled in more than one observation only one ROR is given but the different observing dates distinguish between the individual observations.
### 3.2 Data analysis
The aim of the analysis is to extract fluxes and upper limits for all the sources within a $`6\mathrm{}\times 6\mathrm{}`$ region centered on the position of the galaxy. In this way the X-ray field matches the size of the optical frames obtained using the Jacobus Kapteyn Telescope (JKT) (Johnson 1997; Paper II). For each galaxy an isointensity contour map was overlaid onto the optical image. No attempt has been made to systematically study the temporal behavior of the sources, although some interesting individual cases will be mentioned later (section 5).
The reduction and analysis of the ROSAT HRI data was done using pros, an X-ray analysis software system designed to run under the Image Reduction and Analysis Facility (iraf).
For each galaxy with more than one exposure, the different images were coadded, merging their respective livetimes. Using the sources in the field, the images were inspected and corrected for systematic shifts in the positions before co-adding.
#### 3.2.1 Source identification
For the identification of all point sources in the images we used the ROSAT Standard Analysis Software System (SASS) source list as a starting point. The following steps were followed to identify all source candidates in the X-ray images:
1. the images were binned into 2″ pixels and smoothed using gaussians with $`\sigma =2\mathrm{},4\mathrm{}`$ and $`8\mathrm{}`$; the smoothed images were inspected visually to evaluate the existence of sources not reported by SASS and to exclude obvious spurious sources;
2. whenever more than one exposure was available for a target the different files were compared so that variable sources that might appear weak in the coadded image could be recognized in the individual frames where they might have been more luminous;
3. optical images with X-ray isocontours were generated in order to look for weak sources with optical counterparts.
For each source candidate a background subtracted count number was obtained from the coadded imaged. An aperture of 10″ was adopted for all point sources which should encircle $``$ 99% of the photons at 0.2 keV and $``$ 86% of the photons at 1.7 keV for nearly on-axis sources (?). The pixel coordinates were obtained from the SASS report or by using a separate centroid algorithm. To estimate the background one or two large circular regions free of evident X-ray sources and away from the galaxy were used. These regions normally lay outside the $`6\mathrm{}\times 6\mathrm{}`$ central image.
A final list of sources for each galaxy was created with all sources that comply with at least one of the following criteria:
1. have a signal-to-noise ratio above 2.5 in the coadded image,
2. have a signal-to-noise ratio above 2.5 in at least one of the individual images,
3. or have a low signal-to-noise ratio (between 1.5 and 2.5) and an optical counterpart.
For sources with an extended component an additional, larger aperture was used to estimate its contribution. The size of the aperture was determined from the radial profile of the source. For galaxies without nuclear X-ray source detections, upper limits were established using a 10″ aperture located at the best estimate of the nuclear position from our optical images. They were computed as $`2\sigma `$ limits (ie., 0.9772 CL) assuming Bayesian statistics (?). To find the count rates the total livetime of the coadded images was used.
#### 3.2.2 Astrometry checks
In order to check the absolute astrometry of the HRI observations we used bright X-ray sources in the field of view of each image and looked for optical counterparts that would signal the presence of large shifts in the astrometric solutions. However, in many cases the only bright sources available were found at large off-axis angles. Given the dependancy of the PSF with distance from the centre of the HRI optical axis (?), it is important to establish if the position of these sources could be reliably recovered.
To establish the accuracy with which the position of X-ray sources could be recovered we examined two HRI observations for which a large number of optical couterparts for the X-ray sources has been found: one of the Lockman Hole observations (ROR 701867) and one of the ROSAT Deep Surveys centered in the direction $`\alpha =13^h34^m36^s`$ and $`\delta =37\mathrm{°}54\mathrm{}36\mathrm{}`$ (ROR 900717). The optical identifications of the X-ray sources were obtained from ? and ?.
We used a total of 12 sources from the Lockman Hole and the ROSAT Deep Survey image, with off-axis angles ranging from 4 to 14.5 arcmin, and count numbers ranging from $``$ hundred to a few thounsands. It was found that the position of the X-ray sources could be recovered with an accuracy better than 5 arcsec (average of 2.8 arcsec) regardless of the count numbers and off-axis angle of the source. Therefore, by using optical identifications (IDs) of the X-ray sources found in the field of view of our HRI observations, it should be possible to detect shifts in the astrometric solutions larger than 5 arcsec.
We checked the astrometry of our HRI observations for all those images that showed X-ray sources that could be associated with the galaxies. To look for optical counterparts large DSS images were used and inspected carefully to identify as many candidate IDs as possible. The celestial positions for the IDs were obtained usign a centroid algorithm and compared with the positions of the X-ray sources. The results can be grouped into two categories as follows: fields for which at least two secure optical IDs were found (NGC 247, NGC 2976, NGC 5457, NGC 6503, NGC 6946, UGC 6456); fields for which only one secure ID was found (IC 342, NGC 221, NGC 404, NGC 2403, NGC 4136, NGC 4150, NGC 4236, NGC 4395, NGC 4736, NGC 4826, NGC 5204, NGC 5236).
The analysis showed that most astrometric solutions were accurate to within 5 arcsec. Exceptions were NGC 2976, for which the observation obtained on the 12/04/1995 presented a shift of $`6`$ arcsec, NGC 4236 for which both observations presented a significant shift (notice however, that this result is based on only one optical ID), and NGC 6946 – ROR 600718, which was corrected by a shift $`7`$ arcsec.
#### 3.2.3 X-ray fluxes and luminosities
The conversion of the HRI count rates to fluxes was done assuming a Bremsstrahlung spectrum with $`kT=5`$ keV, an energy range 0.1–2.4 keV, and a Galactic line of sight absorption derived from the 21 cm line of atomic hydrogen (?), listed in table 3. This choice of parameters is justified by typical galactic soft X-ray spectral properties (?). Fluxes were converted to luminosities assuming the distances in table 1. Fluxes and luminosities inferred from a Bremsstrahlung spectrum with $`kT=1.0`$ or 0.1 keV, or a power law model with index $`\alpha =1.0`$ or 1.5, can be calculated using the conversion factors in table 3 for each galaxy.
### 3.3 Contour maps
Isointensity contour maps were produced for all the coadded images. All the files were binned into $`2\mathrm{}\times 2\mathrm{}`$ pixels and then smoothed using a gaussian with $`\sigma =4\mathrm{}`$. The binning and gaussian sizes were chosen so that the representation of the sources was consistent with the statistical asssesment of their strength. Contours were drawn at $`2.5^n`$ times the standard deviation in the smoothed background, where $`n=1,2,3`$, and so on. With this selection of contour intensities bright sources do not present contour crowding. For each map, the background fluctuation was calculated for the same region used earlier (when measuring source fluxes in the raw data).
## 4 The atlas
In this section we present an atlas of X-ray images, and tabulate all X-ray sources found satisfying the conditions described in the previous section, for the 29 galaxies with ROSAT HRI data found in table 2.
The atlas consists of maps of isointensity contour levels overlaid on optical images. The size of the images is $`6\mathrm{}\times 6\mathrm{}`$ centered on the position of the galaxy. Whenever an optical JKT image was not available, a $`6\mathrm{}\times 6\mathrm{}`$ Digital Sky Survey plate was used. The atlas is ordered by increasing values of right ascension.
All detected X-ray point sources fluxes are given in table 5. The sources are ordered by increasing values of right ascension. For each source in table 5 successive columns list the measured signal to noise ratio, the flux corrected for Galactic absorption in the 0.1–2.4 keV band, and the luminosity found assuming the distances shown in table 1. Comments on particular sources are given in the last column (nuclear sources are also noted in this way). When there is evidence of an extended component, the source is listed twice: the flux measured in a 10″ radius aperture is followed by the total flux observed in a larger aperture (the radius of the second aperture is given in the comments).
Upper limits ($`2\sigma `$) were obtained for galaxies without detected nuclear sources using a 10″ aperture located at the nuclear positions and can be found in table 6.
In addition to the galaxies listed in table 2 the results for four galaxies with comprehensive studies of their X-ray emission from HRI observations have been obtained directly from the literature. Luminosities for the detected point sources in these galaxies (and located within a $`6\mathrm{}\times 6\mathrm{}`$ region centered on the nuclei) are shown in table 7. These values will be used during the analysis in chapter 6.
## 5 Notes on individual objects
In this section we briefly review the results on all 34 objects with high resolution X-ray data. This includes objects for which ROSAT HRI data is analyzed in this paper (tables 5 and 6), objects for which we have used ROSAT HRI results from the literature (table 7), and Leo B, which was observed with the Einstein HRI but has not been observed by the ROSAT HRI.
NGC 147: NGC 147 was observed for the first time in X-rays using the ROSAT HRI and no emission was detected. These observations have also been reported by ?. They give a $`2\sigma `$ upper limit for the flux from a point source located in the galaxy of $`6\times 10^{14}`$ ergs s<sup>-1</sup> cm<sup>-2</sup> in the 0.1–2.5 keV band-pass, which is in good agreement with the upper limit reported in table 6.
NGC 185: As with NGC 147, this galaxy was observed for the first time in X-rays and no emission was detected. These observations have also been reported by ?. A $`2\sigma `$ upper limit for the flux of a point source was found to be $`4\times 10^{14}`$ ergs s<sup>-1</sup> cm<sup>-2</sup> in the 0.1–2.5 keV band-pass, again in good agreement with the value quoted in table 6.
NGC 205 - M 110: This small elliptical galaxy was observed with the Einstein HRI and no X-ray sources were detected with a flux $`>1.8\times 10^{12}`$ ergs s<sup>-1</sup> cm<sup>-2</sup> in the 0.5–4.0 keV band-pass (??). The ROSAT HRI observations of NGC 205 reported here give an upper limit of $`6.78\times 10^{14}`$ ergs s<sup>-1</sup> cm<sup>-2</sup> for the flux of any point source in the nuclear region. The only source seen in figure 3 lies 2.7′ away from the nucleus and is probably not associated with the galaxy.
NGC 221 - M 32: A strong off-nuclear source in NGC 221 has been found with the HRI on both Einstein and ROSAT. The flux reported here is consistent with the Einstein measurements ($`F_X=9.1\times 10^{13}`$ ergs s<sup>-1</sup> cm<sup>-2</sup> in the 0.5–4.0 keV band-pass (?)). The X-ray source lies $`7\mathrm{}`$ away from the NGC 221 nucleus and has no optical counterpart. ASCA observations of this source reported by ? reveal a hard spectrum and a flux decrease of 25 percent in two weeks, favoring the identification of the source as a single XRB. Variability by a factor 3 has also been detected during ROSAT PSPC observations (?). An upper limit for the nuclear X-ray emission is given in table 6. The position of the aperture used to compute the upper limit was shifted slightly from the exact position of the galactic nucleus to avoid contamination from the off-nuclear source.
Dynamical studies of the stellar rotation velocities in this galaxy have revealed the presence of a central dark massive object, probably a black hole, with mass $`3\times 10^6M_{}`$ (??), which corresponds to an Eddington luminosity of $`10^{44}`$ ergs s<sup>-1</sup>. The X-ray upper limit in table 6 shows that the central object is emitting at most at $`10^8L_{\mathrm{Edd}}`$.
NGC 224 - M 31: No analysis of X-ray data for this galaxy has been done in this paper. Einstein and ROSAT HRI observations of M 31 have detected over 100 individual sources with luminosities $`10^{36}L_x10^{38}`$ ergs s<sup>-1</sup> (??). The ROSAT PSPC, which is a more sensitive than the HRI (but with worse spatial resolution) gives an upper limit for a diffuse component associated with the galactic bulge of $`2.6\times 10^{38}`$ ergs s<sup>-1</sup> (?).
The point sources are associated with two components, the disk and the bulge, and are strongly concentrated towards the centre. PSPC data show that the bulge and disk components account for about one and two thirds of the total emission, respectively (?). Ginga (2–20 keV) and BeppoSAX (0.1–10 keV) observations have shown that this emission from the galaxy is consistent with a population of LMXRB (??). A list of the sources found within $`6\mathrm{}`$ from the nucleus is shown in table 7. The source coincident with the galactic nucleus has been reported to vary (?) and is probably an XRB.
NGC 247: Two strong X-ray sources are seen in the southern region of the ROSAT HRI image (see figure 5). A $`2\sigma `$ upper limit for a point source located in the nuclear region of the galaxy can be found in table 6. ? report the detection of 5 ROSAT PSPC sources associated with the galaxy, 3 of which lie outside our optical image. The remaining two correspond to our sources X1 and X2. No obvious optical couterparts are seen at these positions, although X1 is located at the edge of a bright star-forming region. Notice that the astrometry of the observations is thought to be accurate (section 3.2.2).
From a fit in the 0.1–2.0 keV range, ? found that the PSPC spectrum of X1 indicates a very soft and obscured source ($`kT=0.12`$ keV, $`N_H=6.4\times 10^{21}`$ cm<sup>-2</sup>). Using their spectral parameters we find that the HRI count rate corresponds to a flux in the 0.1–2.0 keV energy range, (corrected by Galactic absorption only) of $`5.0\times 10^{13}`$ ergs s<sup>-1</sup> cm<sup>-2</sup>, compared with the PSPC flux of $`2.1\times 10^{13}`$ ergs s<sup>-1</sup> cm<sup>-2</sup> given by ?. This corresponds to an increase in the source luminosity by more than a factor 2 between the PSPC and HRI observations. From the HRI count rate, the intrinsic luminosity of the source is found to be $`10^{42}`$ ergs s<sup>-1</sup> for an assumed distance to NGC 247 of 3.69 Mpc, although variations in the adopted parameters can change this significantly. For example, changing the hydrogen column and temperature by 1$`\sigma `$ each (using the error bars in ?, ie., $`kT=0.15`$ keV, $`N_H=4.7\times 10^{21}`$ cm<sup>-2</sup>) decreases the luminosity by an order of magnitude. The very soft spectral distribution, extremely high intrinsic luminosity, and detected variability make this source a good candidate for a stellar accreting black hole.
? also report PSPC observations and find a faint nuclear source with $`L_X=1\times 10^{36}`$ ergs s<sup>-1</sup>, which is well below the detection limit of our HRI observations. However, due to the poor spatial resolution of the PSPC, it is not clear whether this source is coincident with the galaxy nucleus.
NGC 253: No analysis of X-ray data for this galaxy has been done in this paper. This starburst galaxy has very complex X-ray emission. Thirty one point sources have been detected with the Einstein and ROSAT HRI (?) with luminosities of up to a few times $`10^{38}`$ ergs s<sup>-1</sup> (see table 7). Analysis of HRI and PSPC observations show that the source located in the nuclear region is extended, soft and possibly variable (?). Spectral fits to some disk sources using PSPC observations show that they are consistent with the spectra of absorbed XRBs (?).
Extended emission is detected well above and below the galactic plane of this galaxy (see figure 11 and 12 in ?) and NGC 253 is a prototype object for the study of X-ray emission from starburst galaxies. Combined ASCA (FWHM $`3\mathrm{}`$) and PSPC observations show that three components are required to fit the integral spectrum of the galaxy: a power law ($`\mathrm{\Gamma }1.9`$) and a two-temperatur plasma ($`kT0.3`$ keV and $`kT0.7`$) (?). The determination of more detailed physical parameters is hampered by the difficulty of measuring multiple absorbing columns, as discussed by ?. Their analysis of the halo diffuse emission detected with the PSPC shows that the spectrum is well fitted by a two-temperatur plasma ($`kT0.1`$ keV and $`kT0.7`$) if solar abundances are assumed. BeppoSAX observations of NGC 253 show the presence of a $`300`$ eV Fe K line at 6.7 KeV, probably of thermal origin (?). Other prominent lines are also resolved with ASCA (?).
NGC 404: A weak X-ray nuclear source ($`F_X=7.5\times 10^{12}`$ ergs s<sup>-1</sup>) has been detected by the ROSAT HRI in the LINER nucleus of this galaxy. An ASCA 2–10 keV upper limit of $`3\times 10^{13}`$ ergs s<sup>-1</sup> cm<sup>-2</sup> has been reported by ?, which implies that the ROSAT flux is either dominated by very soft emission, or that the source is variable. The former idea is supported by recent UV observations of the NGC 404 nucleus which show that the spectrum is dominated by stellar absorption features from massive young stars (?) and not by the blue, featureless continuum expected from an active nucleus.
NGC 598 - M 33: No analysis of X-ray data for this galaxy has been done in this paper. Several point sources have been detected with the Einstein and ROSAT HRI in this spiral galaxy (??). Those sources located within the central $`6\mathrm{}\times 6\mathrm{}`$ region are listed in table 7. The most striking source is the nucleus, with a luminosity of $`10^{39}`$ ergs s<sup>-1</sup> which makes it a good AGN candidate. However, the nucleus is not detected at radio wavelengths and shows very little line emission in the optical (Schulman & Bregman 1995, Paper III). Nuclear dynamical studies also give a strict limit for a central black hole mass of $`5\times 10^4M_{}`$ (?). ASCA observations show that the X-ray emission from this source is much softer than the typical AGN spectrum and it does not show signs of variability to within 10% on time scales of 100 minutes (?). This almost certainly excludes the possibility of an active nucleus in NGC 598. In fact, a disk blackbody fit to the ASCA data shows that its emission is similar to Galactic black hole candidates (??). Diffuse emission around the nucleus has been detected in ROSAT HRI as well as PSPC observations (???).
IC 342: We have re-analyzed the ROSAT HRI data already discussed by ?. They found evidence for a diffuse nuclear component that could be explained as a hot interstellar medium generated by a very young nuclear starburst.
Figure 7 shows that three sources lie within the limits of the optical JKT image of the galaxy. The brightest source is coincident with the nucleus and visual inspection of the image suggests that it is marginally resolved. A comparison of the azimuthally averaged profile of the source and the HRI model PSF can be seen in figure 30 and a clear deviation from the model PSF is visible at a radius of between 7 and 13 arcsec from the centre. For comparison, an off-nuclear source is also shown, which is in good agreement with the model PSF, confirming the diffuse nuclear component.
X3 is located close to a faint knot of optical emission while X1 has no obvious optical counterpart. Notice however, that only one source was used to check the astrometry of this observation (section 3.2.2). X1 has a S/N of 2.3 (see table 5), but is included here because it was reported by ? (they found a S/N of 2.8, probably due to a different background estimate).
? analyzed Einstein IPC observations of the nuclear region in IC 342 and argued that the emission is consistent with starburst activity. Unfortunately, the IPC was not able to resolve the three sources detected with the HRI.
NGC 1560: This galaxy was observed for the first time in X-rays and no emission was detected in the ROSAT HRI data. A $`2\sigma `$ upper limit for a point source located in the nuclear region of the galaxy can be found in table 6.
NGC 2366: No emission was detected in the ROSAT HRI observations of this galaxy. A $`2\sigma `$ upper limit for a point source located in the nuclear region of the galaxy can be found in table 6.
NGC 2403: This galaxy was observed with the Einstein HRI and IPC instruments (?). Three prominent sources were identified in the IPC observations with no emission from the nuclear region of the galaxy.
The ROSAT HRI observations reported here show a total of 4 point sources associated with the galaxy. Sources X1 and X3 in figure 10 correspond to two of the sources reported by ? (their third source lies outside the JKT image). Sources X2 and X4 are about one order of magnitude fainter than X1 and X3, and were probably below the sensitivity threshold of the Einstein observations. It is confirmed that no nuclear source is present in the galaxy and a $`2\sigma `$ upper limit for a nuclear point source is given in table 6.
A search for SNRs in NGC 2403 has yielded 35 detections (?). The position of remnant number 15 in table 2 of ? is coincident with the X-ray source X3 reported here (see table 5). A very faint optical counterpart is observed at this position. If the identification is correct the SNR would belong to a class of super-luminous (probably young) remnants (?).
Source X4 is coincident with a giant HII region. A photometric study of this region (N2403-A) reveals more than 1400 detected stars, among them 800 O-type stars and a lower limit of 23 WR stars (?).
The coincidence of X3 and X4 with previously known sources in NGC 2403 gives support to the astrometric checks carried out for this galaxy (section 3.2.2).
NGC 2976: This galaxy was observed for the first time in X-rays and no emission was detected in the ROSAT HRI data. A $`2\sigma `$ upper limit for a point source located in the nuclear region of the galaxy can be found in table 6.
The only detected X-ray source has no obvious counterpart. Notice, however, that the astrometry of the observations is highly uncertain given that only one optical identification was found in the field of view of the HRI and a significant shift was applied to one of the images (section 3.2.2).
A 0951+68: This galaxy has been observed in X-rays for the first time. As can be seen in figure 12, the only detected source is probably a foreground or background object.
NGC 3031: No analysis of X-ray data for this galaxy has been done in this paper. Nine and 30 point sources have been detected in this galaxy from observations with the Einstein and ROSAT HRI, respectively (??). All Einstein sources within the galaxy D<sub>25</sub> isophote were detected by ROSAT with the exception of the Einstein source X1. ROSAT sources located within the central $`6\mathrm{}\times 6\mathrm{}`$ region are listed in table 7.
The central emission is dominated by the nuclear source, coincident with a low-luminosity active nucleus, with a luminosity $`10^{40}`$ ergs s<sup>-1</sup> in the 0.2–4.0 keV band-pass. Einstein IPC data show that the spectrum of the nuclear source is soft, with a good fit given by thermal emission with $`kT1`$ keV or by a power law with index $`\alpha 2`$ (?). Broad band observations obtained with BBXRT (0.5–10 keV) and ASCA (0.2–10 keV) show, however, that the nuclear emission is consistent with a power law distribution with index $`\alpha 1`$ (???). The ASCA data also revealed the presence of a broad iron K emission line similar to those seen in more luminous Seyfert galaxies (??). It must be kept in mind that due to the coarse spatial resolution of ASCA some contamination is expected from nearby sources, although ? estimated this to be less than 10 percent. X-ray long term and fast variability by significant factors have been reported for the nuclear source (???).
Leo B: No analysis of X-ray data for this galaxy has been done in this paper. Leo B was observed with the Einstein HRI. ? report that no sources were detected.
UGC 6456: This is a blue compact galaxy. These galaxies are characterized by low metallicities, high gas content and vigorous star formation. In UGC 6456, evidence of both a recent episode of strong star formation (600–700 Myr) and an older stellar population has been found (?).
The JKT optical images show that the galaxy has numerous bright knots of emission surrounded by a low surface brightness outer envelope (Johnson 1997; Paper II). The knots of emission are displaced south from the geometrical center of the envelope and it is not clear whether they represent the true nuclear region of the galaxy (figure 31).
ROSAT PSPC observations of this galaxy show a central X-ray core and three extended structures connected to the central source (?). This morphology was interpreted as outflows from the central region of the galaxy, powered by starburst activity. The total PSPC flux within a circular aperture of radius 3′ is $`1.6\times 10^{13}`$ ergs s<sup>-1</sup> cm<sup>-2</sup>.
The high resolution data reported here show a strong X-ray source located at the north-most limit of the optical emission knots, but displaced to the west with respect to the geometrical centre of the outer envelope (figure 31). Notice also that the astrometry of the HRI observation is fairly well established (section 3.2.2). There is no evidence of extended X-ray emission in the observations, probably due to the lower sensitivity of the HRI. The HRI flux of the point source is $`9\times 10^{13}`$ ergs s<sup>-1</sup> cm<sup>-2</sup>, about 5 times more luminous than the PSPC observations. A very luminous XRB ($`L_X10^{38}`$ ergs s<sup>-1</sup>) could be responsible for this flux variation.
NGC 3738: This galaxy has an irregular optical appearance with several bright knots of emission, although the outer parts are quite regular. None of the bright knots seems to coincide with the geometrical center of the galaxy (Johnson 1997; Paper II).
NGC 3738 has been observed in X-rays for the first time and no point sources associated with the galaxy have been found. An upper limit for the X-ray emission can be found in table 6. Source X1 in figure 14 is coincident with a faint knot of optical emission and is probably a foreground or background object. Source X2 is coincident with a bright point-like object and probably corresponds to a foreground star.
NGC 4136: This spiral galaxy has been observed for the first time in X-rays. No sources associated with the nuclear region have been found. An upper limit to the flux from a nuclear point source can be found in table 6. A strong X-ray source is coincident with one of the spiral arms of the galaxy where several knots of emission can be seen in the optical image (see figure 15). The remnant of the historic Type II supernova SN 1941C seen in NGC 4136 is not located close to this X-ray source (?).
NGC 4144: This galaxy was observed for the first time in X-rays and no emission was detected in the ROSAT HRI data. A $`2\sigma `$ upper limit for a point source located in the nuclear region of the galaxy can be found in table 6.
NGC 4150: A strong point-like source coincident with this galaxy was detected in the ROSAT All-Sky Survey with $`F_X=6\times 10^{13}`$ ergs s<sup>-1</sup> cm<sup>-2</sup> and a photon index $`\mathrm{\Gamma }=1.41`$. The emission was assumed to be from the nucleus of NGC 4150 (??). The high resolution image seen in figure 17 shows, however, that the X-ray source is more than 15″ away from the galactic nucleus and has a position consistent with a knot of optical emission. Spectroscopy of the optical counterpart shows that the source is a background quasar at redshift 0.52 (figure 32).
Figure 17 shows that the X-ray contours of the source are elongated in the north west direction, suggesting that some emission might be coming from the nuclear region of the galaxy. An estimate of the nuclear emission was obtained using a 10″ radius aperture located as shown in figure 33. Although the observed counts have a S/R $`2.6`$ (and so would be considered a significant detection by the criteria defined in section 3.2.2) the measurement will be treated as an upper limit because of contamination from the nearby quasar.
NGC 4236: This galaxy has a very low surface brightness and no obvious nucleus (Johnson 1997; Paper II). From the ROSAT HRI observations reported here no X-ray sources have been found in the central region of the galaxy. The only detected source (X1) is located in the galactic plane and might have a faint optical counterpart. Notice, however, that the astrometry of the observations is highly uncertain given that only one optical identification was found in the field of view of the HRI and a significant shift was applied to both coadded images (section 3.2.2). A $`2\sigma `$ upper limit for a point source located in the nuclear region of the galaxy can be found in table 6.
NGC 4244: No emission was detected in the ROSAT HRI observations of this galaxy. A $`2\sigma `$ upper limit for a point source located in the nuclear region of the galaxy can be found in table 6.
NGC 4395: This galaxy contains the faintest and nearest Seyfert 1 nucleus known today (??). Its nuclear X-ray source is highly variable and the continuum is well fitted by a power-law distribution with photon index $`\mathrm{\Gamma }=1.7`$ (?). The bright source seen in figure 20 (X2) has no obvious optical counterpart.
NGC 4605: This galaxy was observed for the first time in X-rays and no emission was detected in the ROSAT HRI data. A $`2\sigma `$ upper limit for a point source located in the nuclear region of the galaxy can be found in table 6.
NGC 4736 - M 94: Strong X-ray emission is associated with the LINER nucleus of this galaxy. An extended nuclear source can be seen in figure 22. The azimuthally averaged profile of this source and the HRI model PSF are shown in figure 34.
The galaxy was previously imaged with the Einstein HRI. A nuclear flux of $`2.0\times 10^{12}`$ ergs s<sup>-1</sup> cm<sup>-2</sup> was measured within an aperture of 60″ radius in the 0.5–4.0 keV band-pass (?). The fluxes measured from the ROSAT HRI observation using a small (r = 10″) and a large aperture (r = 100″) are $`1.4\times 10^{12}`$ and $`4.7\times 10^{12}`$ ergs s<sup>-1</sup> cm<sup>-2</sup> respectively (see table 5). Since the emission is highly concentrated towards the nucleus, the difference between the 60″ Einstein and 100″ ROSAT fluxes cannot be explained by the different aperture sizes alone. The change in the total observed flux could be explained, however, if the central source is variable or if a substantial fraction of the emission is radiated in the very soft X-rays ($`kT0.5`$ keV).
? have recently reported on ASCA and ROSAT PSPC and HRI observations for NGC 4736. They find that the nuclear emission is consistent with an unresolved source plus an extended component. PSF modeling shows that the unresolved source accounts for more than 50% of the detected emission (?). The 0.1–10 keV ASCA spectrum of the emission is consistent with a power-law with index $`\alpha 1`$ plus a softer thermal component ($`kT0.10.6`$ keV) which dominates below 2 keV.
NGC 4826: The nucleus of NGC 4826 has been classified as a transition object (a combination of a LINER and an HIIR nucleus) by ?. The galaxy was observed with the Einstein IPC and a nuclear flux of $`7.89\times 10^{13}`$ ergs s<sup>-1</sup> cm<sup>-2</sup> was measured within a 4.5′ radius aperture (?), in good agreement with the flux found in table 5. Figure 35 shows the azimuthally averaged profile as observed by the ROSAT HRI. Significant extended emission is observed within $`20\mathrm{}`$ of the central peak.
NGC 5204: Einstein IPC observations of this galaxy show strong X-ray emission with a total flux of $`9.6\times 10^{13}`$ ergs s<sup>-1</sup> cm<sup>-2</sup> (?) in good agreement with the measurement given here (see table 5). Population I XRBs were thought to be responsible for this emission since the number of OB stars inferred from IUE observations were in agreement with the X-ray luminosity (?). However, the HRI image (see figure 24) shows that the X-ray emission is consistent with a single off-nuclear point source $`17\mathrm{}`$ away from the nucleus. Since only one optical ID was found in the field of view of the HRI, the astrometric solution of these observations is quite uncertain. However, the very good agreement in the pointing of the two coadded exposures argues in favor of a fairly accurate astrometry.
Three SNRs have been identified in NGC 5204, but none of them is coincident with the position of the X-ray source (?). Although there is no obvious optical counterpart for the source, several nearby optical knots can be seen in figure 24. The spectrum of one of the candidates, which corresponds to a star forming region, will be presented in Paper III.
The extremely high luminosity of the X-ray source, together with the lack of variability observed between the Einstein and the ROSAT observations, favor an identification as a SNR from a class of super-luminous remnants (the SASS report does not find conclusive evidence of variability during the HRI observations, either). However, the null detection by ? argues against this hypothesis, unless the remnant has an unusually low SII/H$`\alpha `$ ratio (a value $`0.45`$ was adopted by ? as selection criteria to identify SNRs) or unsually high $`L_X`$/H$`\alpha `$ ratio (the super-luminous remnant in NGC 6946 - see below - has $`L_X`$/H$`\alpha 15`$, which is the highest value seen in these type of objects; the SNR in NGC 5204 would require $`L_X`$/H$`\alpha >550`$).
NGC 5236 - M 83: The very complex X-ray emission in this galaxy can be appreciated in figure 25. Several X-ray knots are distributed on top of bright, uneven extended emission. Comparison between the two ROSAT HRI observations reveals that at least three of the eight detected point sources are variable. Table 8 shows the count rates observed in January 1993 and July 1994. Sources X1, X7 and X8 are not detected during the observations obtained in July 1994, but are among the brightest objects seen in January 1993.
? reported on the Einstein HRI observations of this galaxy. From their observations (obtained in January 1980 and February 1981) only 3 sources were detected in the nuclear region of the galaxy. Two correspond to the ROSAT sources X3 (the nucleus) and X6, which are the brightest sources observed in the ROSAT HRI data (see table 8). Their fluxes in the Einstein (0.5–3.0 keV) band-pass are in good agreement with the fluxes given in table 5. The third source had a luminosity of $`2.3\times 10^{38}`$ ergs s<sup>-1</sup> in the Einstein band-pass (for a distance to NGC 5236 of 3.75 Mpc as assumed by ?), and it is not detected in the ROSAT data. It probably corresponds to a transient XRB.
? report on ROSAT PSPC observations of NGC 5236 obtained between January 1992 and January 1993. They find a luminosity for the nuclear source of $`7\times 10^{13}`$ ergs s<sup>-1</sup> cm<sup>-2</sup>, in good agreement with the measurement in table 5. They also detect all the point sources seen in the HRI observations, with the exception of X5, which corresponds to the faintest source in the nuclear region. This suggests that the variable sources X1, X7 and X8 were detectable for at least a year (from the beginning of 1992 until the beginning of 1993), before fading away, becoming undetectable by July 1994.
The diffuse emission from the central region of NGC 5236 can be appreciated in figure 36. From the PSPC observations ? find that the soft 0.1-0.4 keV diffuse component accounts for almost half of the total X-ray emission and argue that most of it is due to hot gas in a super bubble with radius $`1015`$ kpc. Evidence of vigorous starburst activity comes from observations of the nuclear and circumnuclear regions of NGC 5236 which have intricate morphologies in the UV, optical and infrared (???).
Finally, several historic supernovae have been observed in this galaxy, but none of them is consistent with the positions of the point X-ray sources.
NGC 5238: This galaxy was observed for the first time in X-rays and no emission was detected in the ROSAT HRI data. A $`2\sigma `$ upper limit for a point source located in the nuclear region of the galaxy can be found in table 6.
NGC 5457 - M 101: Observations of this galaxy with the Einstein IPC and ROSAT PSPC have been widely reported (???). Results from an ultra-deep (229 ks) ROSAT HRI observation have become available recently (?). The data reported here were the result of combining 2 of the 4 images used to produce the ultra-deep observation reported by ?.
? find a total of 51 point sources down to fluxes $`6\times 10^{15}`$ ergs s<sup>-1</sup> cm<sup>-2</sup>, of which about half are thought to be associated with the galaxy. The X-ray emission beautifully traces the spiral arms of the galaxy and 5 of the individual sources are associated with giant HIIRs (?), but they are located outside the $`6\mathrm{}\times 6\mathrm{}`$ optical JKT image shown in figure 27.
From IPC data ? find a nuclear X-ray flux of $`3\times 10^{13}`$ ergs s<sup>-1</sup> cm<sup>-2</sup> within a circle of 90″ radius for the Einstein (0.2–4.0 keV) band-pass. From the ROSAT HRI observations a nuclear flux of $`4\times 10^{14}`$ ergs s<sup>-1</sup> cm<sup>-2</sup> is obtained, an order of magnitude fainter than the IPC flux. The difference can be explained if sources X2 and X3, which are not resolved by the IPC, were contained within the large aperture used by ?, or by the effect of luminous and highly variable XRBs. The ROSAT HRI count rate found by ? for the nuclear source is identical to our value (0.67 counts ks<sup>-1</sup>), although the inferred fluxes disagree by a factor $`1.5`$, probably because different spectral models were assumed.
The presence of a soft diffuse component is discussed by ? and ?. They find conclusive evidence in ROSAT PSPC observations for extended emission within the inner 7′ of the galaxy. The radial profile obtained from the ROSAT HRI observations shown in figure 37 suggests some patchiness in the X-ray. The better signal to noise data analyzed by ? confirm this result.
An astonishing total of 93 SNRs have been identified in the galaxy (?). Remnants 57 and 54 are consistent with the positions of sources X2 and X3 seen in figure 27. ? detected X-ray emission from two further remnants, but they fall outside our JKT image.
NGC 6503: The nucleus of this galaxy has been classified as a combination of a transition object and a Seyfert 2 nucleus by ?. The ROSAT HRI observations in figure 28 show extended and elongated emission close to the central region of the galaxy. This diffuse emission can be better appreciated in figure 38 where the X-ray image has been smoothed using a Gaussians with $`\sigma =8\mathrm{}`$.
The extended source lies $`10\mathrm{}`$ away from the galactic nucleus. The astrometric solution of the HRI observation has been confirmed using two bright X-ray sources in the field with optical counterparts and therefore the location of the diffuse emission is probably off-nuclear. Notice, however, that the astrometric checks carried out in section 3.2.2 assumed a point source distribution of counts, while the source seen in NGC 6503 is clearly extended, implying possible errors in the determination of the centroid of the emission. The shape and energetics of the diffuse emission could be explained within the super-wind model for starburst galaxies. More detailed observations are necessary to confirm this hypothesis.
NGC 6946: Einstein IPC observations showed X-ray emission associated with the whole body of this galaxy. Two peaks of emission were detected, one coincident with its starburst nucleus and the other associated with a prominent northern spiral arm (?). The spectral fit to these data was consistent with a soft and a hard component, probably associated with diffuse emission and individual accreting sources respectively. The presence of diffuse emission across the disk of the galaxy was confirmed by ROSAT PSPC observations (?) and can also be appreciated in our HRI image (figure 29). The PSPC also resolved the nucleus into three sources which correspond to X3, X4 and X7 in the HRI data. Figure 39 shows the profile of the nuclear source X3 and of an off-nuclear source (X7). Comparing both plots it is clear that the nuclear source is extended up to $`20\mathrm{}`$ away from the central peak. From the analysis of ASCA and PSPC data, ? found that a composite spectrum (a Raymond-Smith plasma with $`kT1`$ keV plus a power-law with photon index $`\mathrm{\Gamma }2.5`$) was a good fit to the observations. It should be remembered, however, that due to the poor spatial resolution of ASCA most of the sources seen in figure 29 were probably contained in the extraction apertures of the X-ray spectra.
Six historic SNs have been seen in this spiral galaxy (SNs 1917A, 1939C, 1948B, 1968D, 1969P, and 1980K) (?). A total of 27 remnants (not including the historic SNs) have been detected by ?. SN 1980K has been observed in X-rays (?), but the source lies outside the $`6\mathrm{}\times 6\mathrm{}`$ JKT optical image of the galaxy. The extreme northern source seen in this galaxy (X8 in figure 29) seems to belong to an extreme group of SNRs with $`10^{39}`$ ergs s<sup>-1</sup> (?). A faint, red ($`R=18.81;BV=1.43`$) counterpart is seen at the position of this source in our optical images, which also corresponds to the object No 16 in the ?’s list of SNRs. Recent PC2 HST images show that the optical source has an intricate morphology, with a small bright shell and what seem to be two outer loops or arcs, as can be seen in the narrow filter image (F673N) shown in figure 40. Based on this morphology, ? suggested that the source could be explained as two interacting remnants of different age, with the smaller shell being a young remnant colliding with the outer, older shells. However, ? have found no high velocity components in their optical high-dispersion spectra of the source, which is not consistent with the presence of a very young and compact SNR. From the H$`\alpha `$ flux they derived a kinetic energy of $`7\times 10^{50}`$ ergs for the shocked gas. Although this value is somewhat larger than what is normally assumed for supernova explosions, it is still within the range of normal events.
## 6 Discussion
First we summarize the most striking results, before discussing some of the key issues more carefully.
(i) Of the 34 galaxies with X-ray data, 12 have X-ray sources associated with their nuclear regions: a 35% success rate. The first result, then, is that nuclear X-ray sources are very common. (It should be noticed that due to the astrometric uncertainties discussed in section 3.2.2, the aligment of the sources with the galactic nuclei is a tentative result for NGC 404, IC 342, and NGC 4395. However, the presence of an active nucleus in NGC 4395 – see section 5 – gives further support to this identification).
(ii) However, it is clear that the detection rate changes markedly with host galaxy luminosity. All 12 detections are in the sub-sample of 29 galaxies with $`M_B<14`$. This effect is discussed in more detail below.
(iii) Of the galaxies spectroscopically classified as Seyfert, LINER, or transition object, 5/7 have detected nuclear X-ray sources, whereas only 7/22 objects with HIIR or absorption-line spectra are detected (spectroscopic classifications will be discussed in Paper III, but the same result holds using the classifications of ?). However, this apparent preference for AGN only reflects the fact that nearly all the smaller galaxies are classified as HIIR. Amongst detected nuclear X-ray sources, 5/12 are classified as AGN of some kind, and 7/12 as HIIR or absorption line.
(iv) Some of the nuclear sources are only just at the limit of detectability, but many are considerably more luminous, up to $`10^{40}`$ erg s<sup>-1</sup>. Of these more luminous sources, a large fraction are clearly extended on a scale of 10 arcsec or more, corresponding to $`>150`$ pc or so at the typical distance of our sample galaxies.
(v) Many sources outside the nucleus are detected – most interestingly nine off-nuclear sources are found with luminosities exceeding $`10^{39}`$ erg s<sup>-1</sup>, not easily explained as individual X-ray binaries or SNRs.
### 6.1 Dwarf galaxies
Of the twelve galaxies classified as dwarfs in the sample (see table 1), only five have Einstein or ROSAT HRI observations (NGC 147, NGC 185, A 0951+68, Leo B and UGC 6456). Of these only one has a positive detection of an X-ray source (UGC 6456) and even this one is not nuclear. For the rest of section 6, we do not include these dwarfs in the analysis, concentrating on the sub-sample of 29 galaxies with $`M_B<14`$.
### 6.2 Super-luminous X-ray off-nuclear sources
Outside the nucleus, remarkably luminous sources ($`L>10^{38}`$ erg s<sup>-1</sup>) are seen very frequently in galaxies spanning 2 orders of magnitude in luminosity. These are plotted against host galaxy luminosity in figure 41. The most luminous sources ($`L>10^{39}`$ erg s<sup>-1</sup>) have been labeled with numbers and brief comments about them can be found in table 9. Figure 41 shows that there is no obvious correlation with galaxy size.
The population of X-ray sources in the Milky Way and M 31 does not include luminosities above $`10^{38}`$ ergs s<sup>-1</sup>. The presence of luminous ($`L_X10^{38}`$ ergs s<sup>-1</sup>) sources in the Magellanic Clouds was assumed to be a metallicity effect (??). However, as more galaxies were surveyed using the capabilities of the Einstein and ROSAT satellites it became clear that extremely luminous objects ($`L_X10^{39}`$ ergs s<sup>-1</sup>) were not rare (?).
Next we assess the possibility that the ultra-luminous sources are actually background objects. Using the results from the ROSAT Deep Survey in the Lockman Field (?) we estimate that no more than three background sources with fluxes of $`10^{12}10^{13}`$ ergs s<sup>-1</sup> cm<sup>-2</sup> would be found in a 1 deg<sup>2</sup> field. If these targets are distant, and so obscured by the intervening galaxy, their intrinsic luminosities must be higher and the corresponding number densities even lower. Given the average projected size of the galaxies on the sky (see table 3 for the area contained within the D<sub>25</sub> ellipses) the probability that all the sources correspond to background objects is extremely small.
At least one super-luminous source has been identified as a multiple object formed by interacting SNRs in the disk of NGC 6946 (?) (see figure 40). ? has recently suggested that some of the X-ray bright remnants observed in M 101 could correspond to ‘Hypernova Remnants’, a term introduced by Paczyński (1998) for super-energetic $`\gamma `$-ray bursts and the associated afterglow event - but see also ?. Other multiple object systems could be formed from XRBs, SNRs, and diffuse emission from hot bubbles of interstellar gas. Alternatively, the sources could be super-Eddington XRBs with luminosities several times greater than the Eddington limit for a $`1.4M_{}`$ neutron star. The high luminosities in this case are explained as anisotropic emission from neutron stars in binary systems with very strong magnetic fields (?), or as black hole XRBs.
? have recently analyzed ASCA data from three nearby spiral galaxies. One of the galaxies is M 33 and these observations have already been discussed in section 5. The other two galaxies (NGC 1313 and NGC 5408) each harbor a luminous source ($`L_X10^{39}`$ ergs s<sup>-1</sup>) displaced $`50\mathrm{}`$ from the nucleus. The fit to the X-ray spectra of these sources shows that at least two components, a steep power law and a ‘Disk Black Body’, are required to explain the observations, suggesting that these are accretion driven systems. The parameters from the model imply BH masses of $`10010,000M_{\mathrm{}}`$, intermediate between those seen in XRB and AGN. Given the observed X-ray luminosities and the estimated masses, these objects are not super Eddington sources.
It is clear, then, that super-luminous point sources are common. However, unlike nuclear sources discussed in the next section, there seems to be no correlation between the probability of having a super-luminous off-nuclear source and the brightness of the parent galaxy.
### 6.3 Correlation between nuclear X-ray luminosity and host galaxy luminosity
Figure 42 shows the nuclear X-ray luminosities and upper limits for the galaxies in the sample as a function of the blue absolute magnitude of the host. Different symbols have been used to show the optical classification of the nuclear emission as starlight, HII regions, LINERs, transition objects (i.e., a HIIR and LINER composite), and Seyferts (Paper III). As was suggested before, it is found that the probability of detection of a nuclear X-ray source correlates strongly with host galaxy luminosity: the rate of detection is 0% (0/7) below $`M_B17`$, 25% (3/12) for $`17<M_B<19`$, and 90% (9/10) for galaxies with $`M_B<19`$.
? have collected data from ROSAT HRI observations for 39 nearby spiral and elliptical galaxies, of which 9 objects are common to our sample. Their results confirm what is seen in figure 42: for those sources detected within $`6`$ arcsecs from the optical center of the galaxy, and therefore likely to be nuclear sources, the detection rate is nearly 100 percent for hosts with $`M_B20`$.
The distribution of points does not suggest a straightforward correlation, but rather a large spread with an upper envelope. In other words, large galaxies can have luminous or feeble nuclei, but small galaxies can only have feeble nuclei. A similar upper-envelope effect has been claimed for the host galaxy luminosities of quasars and Seyferts (?). Given this correlation, the lack of detection of nuclear sources in small galaxies does not necessarily mean they cannot form AGN at all – it may be that they are exceedingly weak nuclear sources.
It is not clear how to model the envelope-like relationship between nuclear and host galaxy luminosities. However, a simple visual comparison with other related samples reveals a rather striking effect. The upper envelope for our sample seems to go roughly as $`L_xL_{\mathrm{host}}^{1.5}`$. Figures 43, 44 and 45 show nuclear emission line data taken from the study of ?. The nuclear H$`\alpha `$ luminosity from galaxies classified as HII regions shows a very similar correlation with host galaxy luminosity (figure 43) and once again an upper envelope with slope 1.5 is a good fit. Likewise for dwarf Seyferts, the broad H$`\alpha `$ component shows a correlation with a slope of 1.5 (figure 44). However, for the narrow-line component of H$`\alpha `$ for Seyfert nuclei, the correlation is much steeper, with a slope of 3.0 (figure 45). Figure 46 shows data from the radio survey of elliptical galaxies by ?. Here we find that nuclear radio luminosity shows a very strong correlation with an upper envelope slope of 3.0. (? do a more careful analysis showing that the 30th percentile goes as $`L_{\mathrm{host}}^{2.2}`$). So it seems that narrow emission lines in AGN are intimately connected with radio emission. On the other hand, the relation seen for our X-ray sources is equally consistent with either AGN or star formation.
### 6.4 Nuclear X-ray sources as disk sources
Some of the nuclear X-ray sources we have seen are consistent with the most luminous known X-ray binaries. Furthermore, we have seen that super-luminous off-nuclear sources do occur quite often. Is it possible, then, that the nuclear sources are simply examples of the general X-ray source population that happen to be located in the nucleus? Where we see a clear extended source, this cannot be the case. This argument assumes that the observed extended emission is truly diffuse in nature. Of course, we cannot exclude a contribution of unresolved sources to the extended component, but this cannot account for the bulk of the emission for the more conspicuous cases (see ?, ?, ?, and ?). It could also be that more luminous galaxies are simply more likely to have at least one example of a particularly luminous source. The strongest argument against this is that amongst the off-nuclear sources we do not see the upper-envelope correlation with host galaxy luminosity (figure 41). We can also crudely estimate the likely size of such an effect if we know the luminosity distribution of X-ray sources.
The combined luminosity distribution of disk (ie, off-nuclear) X-ray sources from a sample of 83 spiral galaxies has been established recently by ?. The sample was defined as those objects surveyed by ? that had archive ROSAT HRI observations. The distribution determined by ? reaches luminosities for off-nuclear sources just below $`10^{41}`$ ergs s<sup>-1</sup>, ie, it extends well into the realm of the super-luminous sources discussed in the previous section. The low luminosity end of the distribution was found using the ROSAT PSPC deep observations of NGC 224 (M 31) published by ?. ? find that the (differential) luminosity distribution is well fitted by a power law slope of $`1.8`$, with a flattening for luminosities $`10^{36}`$ ergs s<sup>-1</sup>. It is not clear whether this change in the slope is due to incompleteness or to an intrinsic variation in the faint source population.
Since the luminosity distribution determined by ? corresponds to an average distribution a direct comparison between the properties of the populations in individual galaxies and their hosts is not possible. Moreover, as their sample is drawn from the flux-limited sample surveyed by ?, with an under-representation of low luminosity galaxies, it is not possible to confirm the lack of correlation between the probability of finding a super-luminous off-nuclear source and the luminosity parent galaxy seen in figure 41.
The luminosity distribution of off-nuclear sources predicts that about 8 disk sources with $`L_X10^{37}`$ ergs s<sup>-1</sup> will be found in a $`10^{10}\times L_{\mathrm{}}`$ ($`M20`$) galaxy. Of these sources only one will have a luminosity $`10^{38}`$ ergs s<sup>-1</sup>. The probability of that source being located within the nucleus of the parent galaxy can be expressed as the ratio between the total luminosity and the luminosity of a ‘nuclear region’ of 100 pc (equal to the HRI spatial resolution at a distance of 4 Mpc). Using the well established exponential disk profile in spirals, and assuming a scale length for the disk equal to 4 kpc (?) this ratio is only $`10^4`$.
For more luminous sources the situation is even worse, since the ? luminosity distribution predicts that about $`8\times 10^{10}\times L_{\mathrm{}}`$ galaxies have to be surveyed in order to find a $`L_X10^{39}`$ ergs s<sup>-1</sup> X-ray source. This prediction does not agree with our results, however. In our sample of 29 galaxies we have detected 9 sources with luminosities above $`10^{39}`$ ergs s<sup>-1</sup> and all but one of the host galaxies (NGC 2403) have $`M_B>20`$ (see figure 41). A closer look at the results reveals that two effects could be responsible for this difference: (1) statistic fluctuations introduced in the results drawn from our smaller sample due to the distances adopted for the parent galaxies; (2) the fairly large threshold radius (25 arcsecs) adopted by ? to discriminate between nuclear and off-nuclear sources (as an example, the ultra-luminous off-nuclear source in NGC 5204 was labeled as nuclear by ?). As pointed by ?, individual galaxies can also show important deviations from the determined luminosity distribution. The X-ray population in NGC 5457, for example, shows the same power-law index determined by ? for the combined sample, but the normalization is 4 times larger.
In summary, even with a large error (or fluctuation) in the normalization of the luminosity distribution of off-nuclear sources, it is clear that nuclear sources cannot be explained as disk sources located by chance in the nuclei of galaxies. There is another possibility, however. Nuclear sources could be explained if massive X-ray emitting systems, such as black hole binaries, suffer considerable dynamical friction against field stars and so migrate to the centre of their parent galaxies. In this scenario the upper envelope seen in figure 42 could be the result of smaller galaxies having shallower potential wells and therefore being less efficient in dragging the heavy objects to their centres. However, it is well known that the scale time for a significant change in the orbit of any stellar system is comparable to the local relaxation time which, in the case of disks, is much longer than the Hubble time (?). Therefore, galactic disks are too young for a source to have migrated significantly from its original location.
### 6.5 Nuclear X-ray sources as bulge sources
An alternative to the scenario described above is that the massive XRBs could have originally been located in the bulge of the galaxy, where larger background stellar densities make mass segregation a much more efficient mechanism. However, the anticorrelation between the time it takes for a massive object to migrate to the galactic nucleus and the mass of the object implies that an XRB would need a mass of over $`100M_{\mathrm{}}`$ to fall to the galactic centre from a distance of only 10 pc within the lifetime of the parent galaxy (1 Gyr) (?). Therefore, dynamical friction will only be efficient for very heavy or very nearby stars. Indeed, ? show that only extremely massive objects such as $`10^6M_{\mathrm{}}`$ globular clusters would have spiralled into the centre of their parent galaxies in less than a Hubble time from a significantly extra-nuclear (2 kpc) distance.
It could also be argued that the X-ray emission from the central regions corresponds to bulge sources which appear coincident with the galactic nuclei. This is unlikely, however, given the spatial resolution of our observations (ranging from $`16`$ pc for NGC 224 to $`190`$ pc for NGC 5457, with most galaxies with detected nuclear sources being observed at a resolution $`100`$ pc), implying that a very large number of highly centrally concentrated bulge sources would be necessary to explain the observations. NGC 224 indeed shows a high concentration of bulge sources when compared with our own Galaxy, with 4 times more sources seen in NGC 224 within the central 5 arcmins (?). However, only two sources are found within the innermost 100 pc, giving a total luminosity of $`4\times 10^{37}`$ ergs s<sup>-1</sup>. The probability of those sources having $`L_X10^{38}`$ ergs s<sup>-1</sup> is extremely small, as is shown in the bulge luminosity distribution determined of ?. Also, as we have already shown, the hypothesis that particularly heavy (and luminous) binary systems are located within this region because of mass segregation is not viable. Therefore, most of the galaxies with $`M_B<18`$ seen in figure 42 would need to harbor a much larger population of bulge sources that the one seen in NGC 224 in order to explain their observed nuclear luminosities, which is a very unlikely scenario.
### 6.6 The nature of the nuclear X-ray sources
The very different correlations between X-ray luminosity and host galaxy absolute magnitude seen in figures 41 and 42 for nuclear and off-nuclear sources imply two different populations. Nuclear sources are not disk or bulge sources located in the nuclear region by chance. Instead, nuclear sources have a particular nuclear nature.
The next step is to try to unveil this nature. What are these sources? Are they connected with stellar processes in the nuclear region of the galaxies? Could they be an expression of nuclear activity that somehow escapes detection at optical wavelengths?
An interesting pattern can be seen in figure 42 between the optical (spectral) classification of the galaxies and their X-ray luminosities. Both LINER objects (NGC 404 and NGC 4258) lie close to the upper envelope of the distribution of detected sources, as does the Seyfert galaxy NGC 3031. The second Seyfert nucleus, NGC 4395, is heavily absorbed below 3 keV and its intrinsic soft X-ray luminosity is estimated to be an order of magnitude larger than the values found from ROSAT observations (?). Introducing this last correction all objects classifies as AGN would be found near the top of the upper envelope in figure 42.
The distribution of X-ray sources with galaxy absolute magnitude, shown in figure 42, can be bounded by a nearly linear upper envelope. A similar linear upper envelope has been found in the correlation between nuclear (AGN) luminosity and host galaxy luminosity for samples of Seyfert galaxies and low redshift quasars at IR and optical wavelengths (??). This suggests that there is a maximum allowed AGN luminosity which is an increasing function of the luminosity of the parent galaxy. The limit could be the result of a correlation between the total mass of the host galaxy and the mass of the central engine. The observed range of nuclear luminosities would be given by different accretion rates, with sources located at the bound of the region, showing the maximum possible luminosity. Indeed, from their study of nearby quasars ? and ? show that if the correlation between black hole mass and galaxy size claimed by ? is correct, then the upper envelope is consistent with the most luminous AGN radiating at roughly 20% of the Eddington limit. If the smaller ratio of black hole mass to galaxy mass suggested by ? is used then the closeness of the upper envelope to the Eddington prediction is even more remarkable.
However, in the case of the comparatively weak nuclear X-ray sources we have found here, it is not clear what process determines the maximum possible luminosity. If the observed X-ray emission is driven by accretion of matter onto a BH, then the Eddington limit for the central mass seems the most straight-forward mechanism to explain the existence of the linear upper envelope. However, this cannot be the case. The only two Seyfert galaxies in the sample, NGC 3031 (M 81) and NGC 4395, are thought to be extremely sub-Eddington (??), with their bolometric luminosities being only $`10^410^3\times L_{\mathrm{Edd}}`$.
Figure 47 makes this point graphically. We use the relation given by ?, $`\mathrm{log}M_{BH}=1.96+\mathrm{log}L`$, where $`L`$ is the $`V`$-band luminosity of the bulge and all quantities are expressed in solar units. The bulge magnitudes for the galaxies in our sample can be found from the total magnitudes shown in table 1 using the relantionship obtained by ? as a function of Hubble Type, while a colour $`BV=1.0`$ was assumed for all objects to convert the blue magnitudes to visual magnitudes. Using the relationship above we can infer the X-ray output for a central BH emitting at the Eddington limit as a function of its host galaxy bulge luminosity, assuming that the soft X-ray emission corresponds to $`10\%`$ of the bolometric luminosity. It is found that the predicted X-ray luminosities (seen as a dotted line in figure 47) are about 4 orders of magnitude brighter than what is observed.
The nuclear-host correlation we have seen cannot, therefore, be simply interpreted as the upper envelope relation seen amongst luminous AGN, because the X-ray sources are so feeble. However, it should be borne in mind that the objects discussed in ? have been pre-selected as AGN, whereas we are looking at galaxies as a whole. For example, it is possible that the probability distribution $`P(L_x/L_{\mathrm{Edd}})`$ is a universal function for all galaxies, declining continuously from large values at small $`L_x/L_{\mathrm{Edd}}`$ to small values at large $`L_x/L_{\mathrm{Edd}}`$, but with a cutoff at $`L_x/L_{\mathrm{Edd}}=1`$. The envelope seen in AGN samples then traces the cutoff, but the envelope seen in complete galaxy samples is a statistical effect, tracing in each host-galaxy-bin the nuclear luminosity beyond which the expected number of objects is just less than one. With a larger sample, the density of points would be greater, so the envelope would move up, but the slope might stay the same.
Most of the preceding discussion has centred round the possibility that the weak nuclear X-ray sources are AGN. However, this is not obvious. They seem to show no particular preference for AGN-like optical spectra, and the correlation with host galaxy luminosity is exactly like that shown by the emission line luminosity for HII nuclei (see figure 43). Furthermore many of the most luminous sources are extended on scales of hundreds of parsecs or more. It seems equally likely that the X-ray emission is connected with star formation, but with the current data we cannot go much further. In later papers in this series we examine the correlation between X-ray, optical, and emission-line luminosities, and test specific AGN and star formation models.
Although it is unclear which physical process causes the observed distributions, the upper envelopes give strong evidence for a correlation between the host galaxy luminosity and the level of nuclear activity. The fact that an upper envelope is also observed in the distribution of the nuclear HII regions might suggest that the correlations are governed by the amount of gas available in the nuclear regions. If bigger galaxies are more efficient in dragging gas to their nuclear regions they could show more vigorous star formation and, potentially, feed their massive BHs more efficiently.
## 7 Summary
X-ray sources in the nuclei of galaxies are very common and the luminosity of these sources is strongly connected to the luminosity of the host galaxy. The highest luminosities reach $`10^{40}`$ ergs s<sup>-1</sup> which is $`10^4`$ times less than the Eddington limit for the massive black holes that may be present in such galaxies. The most luminous nuclear X-ray sources are frequently extended on scales of hundreds of parsecs. It is, so far, unclear whether these nuclear X-ray sources are miniature AGN of some kind or a phenomenon connected with normal star formation. Outside the nuclei of galaxies, extremely luminous sources ($`L>10^{39}`$ erg s<sup>-1</sup>, compared with ‘normal’ $`10^{37}`$ erg s<sup>-1</sup> SNRs or $`10^{38}`$ erg s<sup>-1</sup> XRBs) are also quite common, occurring in a third of all galaxies. The nature of these sources is currently also unclear.
## Acknowledgements
We thank Gordon Stewart and John Arabadjis for insightfull discussion on this paper. This research has made use of data obtained from the Leicester Database and Archive Service at the Department of Physics and Astronomy, Leicester University, UK. The Jacobus Kapteyn Telescope is operated on the island of La Palma on behalf of the UK Particle Physics and Astronomy Research Council (PPARC) and the Nederlanse Organisatie voor Wetenschappelijk Onderzoek (NWO) and is the located at the Spanish Observatorio del Roque de los Muchachos. The observatory is hosted by the Instituto de Astrofísica de Canarias. The DSS plates used in this paper were based on photographic data of the National Geographic Society – Palomar Observatory Sky Survey (NGS-POSS) obtained using the Oschin Telescope on Palomar Mountain. The NGS-POSS was funded by a grant from the National Geographic Society to the California Institute of Technology. The plates were processed into the present compressed digital form with their permission. The Digitized Sky Survey was produced at the Space Telescope Science Institute under US Government grant NAG W-2166. The data presented in this paper were reduced using Starlink facilities.
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# An Interpretation of the Evidence for TeV Emission from Gamma-Ray Burst 970417a
## 1. Introduction
Gamma-ray bursts (GRBs) have been the most mysterious astronomical phenomenon in the universe for about 30 years after the discovery (Klebesadel, Strong, & Olson 1973). One of the reasons why the GRB phenomenon is difficult to understand is that GRBs had been observed only in the soft gamma-ray band for a long time. Detection of GRBs in other wavelengths is quite valuable for the progress of the GRB study, as proved by the dramatic progress in recent years following the discovery of afterglows in longer wavelengths of X, optical and radio bands (see, e.g., Piran 1999 for a recent review). Detection of gamma-rays harder than the ordinary sub-MeV band is also important information for better understanding of GRBs. It has been confirmed that emission from GRBs extends up to $``$ 10 GeV, as seen in the famous long-duration GeV emission from GRB 940217 (Hurley et al. 1994). There have been some suggestive results for the emission beyond TeV range (Amenomori et al. 1996; Padilla et al. 1998), although these results were not claimed as firm detections of GRBs.
Recently the Milagro group reported evidence for TeV emission from one (GRB 970417a) of the 54 BATSE GRBs in the field of view of their detector, Milagrito (proto-type of Milagro) (Atkins et al. 2000). An excess of gamma-rays above background is clearly seen during the duration of this burst in the BATSE error circle, and the chance probability of such an event after examining 54 GRBs is estimated as $`1.5\times 10^3`$, giving stronger evidence for TeV emission compared with the earlier observations in this energy band. If this signal is truly from the GRB, the TeV fluence must be at least 10 times greater than the sub-MeV fluence of this GRB without taking into account the intergalactic absorption of TeV gamma-rays. The impact on the GRB energetics would be quite strong.
In this letter I discuss theoretical implications of this interesting event, assuming that the signal observed by the Milagrito is truly from the GRB 970417a. I first try to estimate a likely value of the redshift and energetics of this GRB from the observed energy fluence and detection rate (1 per 54 GRBs), under the assumption that this GRB is not a peculiar GRB. I then discuss whether this extreme phenomenon can be explained in a reasonable theoretical framework. I show that the proton-synchrotron model of GRBs (Totani 1998b, 1999) gives a possible explanation for the Milagrito result.
Throughout this letter I use the isotropic energy for the total energy of GRBs. My analysis does not depend on the unknown collimation factor of GRBs, and the actual energy emitted by the central engine can be much smaller than the isotropic energy if GRBs are strongly collimated, jet-like explosions.
## 2. Redshift and Energetics of GRB 970417a
The basic assumption is that GRB 970417a is the nearest GRB to us among the 54 GRBs observed by the Milagrito. The unknown luminosity function of GRBs in the TeV range makes this assumption less reliable, but this is reasonable because detectability of TeV gamma-rays from cosmological GRBs very rapidly decreases with increasing redshift, due to the well-known effect of intergalactic absorption by the cosmic infrared background radiation (see, e.g., Salamon & Stecker 1998; Primack et al. 1999). Another argument supporting this assumption will be given in §4.
We can estimate the fraction of GRBs within a given redshift $`z`$ of all GRBs detectable by the BATSE satellite, $`N\left(<z\right)`$, if we know the sub-MeV luminosity function of GRBs, GRB rate history as a function of $`z`$, and the threshold flux of the BATSE. Let $`L_p`$ be the peak photon-number luminosity (i.e., not energy) in the restframe 50–300 keV range, and $`P`$ be the observed peak photon flux in the BATSE range of 50–300 keV. These two are related as:
$$P(L_p,z)=\frac{\left(1+z\right)^{2\alpha }}{4\pi d_L^2}L_p,$$
(1)
where $`d_L`$ is the standard luminosity distance and I have assumed a power-law spectrum of GRBs in the sub-MeV range with a spectral index $`\alpha `$, as $`dL_p/d\epsilon \epsilon ^\alpha `$. The observed rate of GRBs with a given set of $`(L_p,z)`$ is given as
$$\frac{d^2N}{dL_pdz}=\varphi \left(L_P\right)\frac{R_{\mathrm{GRB}}\left(z\right)}{\left(1+z\right)}\frac{dV}{dz},$$
(2)
where $`\varphi \left(L_p\right)`$ is the luminosity function of GRBs, $`R_{\mathrm{GRB}}`$ is the comoving rate density of GRBs as a function of redshift, and $`dV/dz`$ is the comoving volume element of the universe. The factor $`\left(1+z\right)^1`$ comes from the cosmological time dilation effect. I assume that the GRB rate evolution traces the star formation history in the universe (Totani 1997). The form of the star formation rate evolution is modeled based on the data compiled by Madau, Pozzetti, & Dickinson (1998). Then $`N\left(<z\right)`$ is given as
$$N\left(<z\right)=\frac{_0^z𝑑z𝑑L_p\left(d^2N/dL_pdz\right)\mathrm{\Theta }\left[P(L_p,z)P_{\mathrm{th}}\right]}{_0^{\mathrm{}}𝑑z𝑑L_p\left(d^2N/dL_pdz\right)\mathrm{\Theta }\left[P(L_p,z)P_{\mathrm{th}}\right]},$$
(3)
where $`\mathrm{\Theta }`$ is the step function \[i.e., $`\mathrm{\Theta }\left(x\right)=1`$ and 0 for $`x>0`$ and $`x<0`$, respectively\], and $`P_{\mathrm{th}}`$ is the detection threshold of the BATSE ($`P_{\mathrm{th}}0.3\mathrm{photons}\mathrm{cm}^2\mathrm{sec}^1`$, Meegan et al. 1996). In this letter I assume a form of the luminosity function of GRBs as logarithmically flat \[$`\varphi \left(L_p\right)L_p^1`$\] within a range of $`(L_{p,\mathrm{min}},L_{p,\mathrm{max}})=(10^{57},10^{59})`$ \[photons s<sup>-1</sup>\], based on the observed luminosity distribution of GRBs with secure redshifts (see, e.g., Table 1 of Lamb & Reichart 2000). The spectral index $`\alpha `$ is set to 1, which is a rough average of GRB spectra (Mallozzi, Pendleton, & Paciesas 1996). I use a reasonable set of cosmological parameters of $`H_0`$ = 70 km/s/Mpc and $`(\mathrm{\Omega }_0,\mathrm{\Omega }_\mathrm{\Lambda })=(0.3,0.7)`$.
In Fig. 1, I show $`N\left(<z\right)`$ as a function of redshift by the dot-dashed line. \[See the right-hand-axis for the scale of $`N\left(<z\right)`$\]. It becomes consistent with the detection rate of GRBs by the Milagrito detector, 1/54, at redshift $`0.7`$. (The vertical solid line shows the redshift corresponding to the detection rate of 1/54, and the shaded region shows a redshift range in which the detection rate is consistent within a factor of two.) Also shown by the dashed line is the total isotropic energy emitted in the BATSE range, $`E_{\mathrm{MeV}}`$, which is estimated by the BATSE fluence of the GRB 970417a, $`3.9\times 10^7\mathrm{erg}\mathrm{cm}^2`$, in all the four energy channels of the BATSE ($`>20`$ keV). \[The latest BATSE catalog at MSFC is available at http://cossc.gsfc.nasa.gov/cossc/BATSE.html\]. The isotropic energy in the sub-MeV band then becomes $`\stackrel{<}{}10^{51}`$ erg in the likely redshift range. The observed total sub-MeV energy of the GRBs with secure redshifts is widely distributed in a range $`10^{51}`$$`10^{54}`$ erg, and GRB 970417a belongs to a class of the least energetic GRBs in the sub-MeV band.
In order to estimate the total energy emitted in the TeV range, the absorption optical depth of TeV gamma-rays due to the cosmic infrared background is necessary. I have calculated this optical depth as a function of the source redshift and observed photon energy, by using a model of luminosity density evolution of stellar lights in the universe (Totani, Yoshii, & Sato 1997). The standard formulation for the calculation of optical depth from the luminosity density evolution in the universe is given in e.g., Salamon & Stecker (1998). The dust-emission component is calculated assuming that the dust emission spectrum is the same as that in the solar neighborhood and the fraction of stellar light absorbed by dust is determined to reproduce the far infrared background radiation measured by the COBE satellite (Hauser et al. 1998). I have checked that this model of optical depth is quantitatively consistent with earlier publications within the model uncertainties (e.g., Salamon & Stecker 1998; Primack et al. 1999). Fig. 2 shows the observed photon energy ($`\epsilon `$) corresponding to several values of optical depth \[$`\tau (z,\epsilon )`$\] as a function of the source redshift.
The events observed by the Milagrito are considered to be gamma-rays above 50 GeV (Atkins et al. 2000), and an estimate of the TeV fluence of the GRB 970417a is given as a function of the upper cut-off energy and spectral index (see Fig. 4 of Atkins et al. 2000). Here I assume the spectral index of 1.5, which is the standard photon index of synchrotron radiation with particle index of 2. If the TeV spectrum is harder or softer than this, the estimate of the total energy in TeV range becomes smaller or larger, respectively (see Fig. 4 of Atkins et al. 2000). I use the photon energy corresponding to the intergalactic optical depth $`\tau `$=1, denoted as $`\epsilon _{\tau =1}\left(z\right)`$, as the upper cut-off energy, unless this energy is smaller than 50 GeV. Then we can estimate the total isotropic energy emitted in the TeV range as
$$E_{\mathrm{TeV}}=\frac{4\pi d_L^2}{\left(1+z\right)}F_{\mathrm{TeV}}\left(\epsilon _{\mathrm{cut}}\right)\mathrm{exp}\left[\tau (z,\epsilon _{\mathrm{cut}})\right],$$
(4)
where the upper cut-off energy is $`\epsilon _{\mathrm{cut}}=\mathrm{max}(50\mathrm{GeV},\epsilon _{\tau =1})`$, and $`F_{\mathrm{TeV}}`$ is the observed TeV fluence which is a function of the upper cut-off energy. The result is shown by the solid line in Fig. 1 as a function of redshift.
From these results, a possible interpretation is that GRB 970417a was located at $`z`$ 0.7 emitting isotropic energies of $`\stackrel{>}{}10^{54}`$ and $`\stackrel{<}{}10^{51}`$ erg in the TeV and sub-MeV range, respectively. It is interesting to note that the energy emitted in the TeV range is similar to that in the sub-MeV range of the most energetic GRB observed to date: GRB 990123, whose total isotropic energy was estimated as $`3\times 10^{54}`$ erg (Kulkarni et al. 1999). Therefore the extremely large isotropic energy in the TeV range for GRB 970417a is not too large for the total energy budget of GRBs. In fact, we do not know why most of the energy of GRBs is emitted in the sub-MeV range, and there is no robust theoretical argument which excludes a possibility that most of the total GRB energy is emitted in other photon energy bands. In fact, such an extreme phenomenon has been predicted by the proton-synchrotron model of GRBs (Totani 1998b; 1999). In the rest of this letter I discuss whether this model can explain the Milagrito result.
## 3. Interpretation by the Proton-Synchrotron Model
Full description of this model has already been given in Totani (1998b; 1999), and here I summarize the qualitative feature of the model.
Currently the most popular explanation for the GRB phenomenon is dissipation of the kinetic energy of ultra-relativistic bulk motion with a Lorentz factor of $`\mathrm{\Gamma }\stackrel{>}{}10^{23}`$, in internal shocks which are generated by relative velocity differences of relativistic shells ejected from the central engine. All the total energy ejected as relativistic bulk motion cannot be dissipated in internal shocks, and hence the total energy truly emitted as kinetic motion ($`E_{\mathrm{iso}}`$) should be larger than the observed total energy of gamma-rays, $`E_{\gamma ,\mathrm{iso}}`$, at least by a factor of several. Therefore, the most energetic class of GRBs, such as GRB 990123, must emit quite a large amount of energy, $`E_{\mathrm{iso}}\stackrel{>}{}10^{55}`$ erg. If the efficiency of the internal shock is not so high, we may have to consider an isotropic energy reaching $`10^{56}`$ erg. Therefore, if GRBs are produced by stellar death events, GRBs must be strongly collimated at least by a factor of ($`4\pi /\mathrm{\Delta }\mathrm{\Omega })`$ 100 to reduce the actual energy emitted from the central engine.
Since the origin of the GRB energy is relativistic bulk motion, protons should carry a much larger amount of energy than electrons by a factor of $`m_p/m_e2,000`$ in the initial stage of the internal shock generation. It is very uncertain what fraction of the proton energy is converted into electrons, but the simplest Coulomb interaction cannot transfer the proton energy into electrons within the typical time scale of GRBs (Totani 1998a, 1999). The soft gamma-rays are generally considered to be generated by electrons, because of the short time variability of GRBs. Therefore it is not unreasonable that, in some GRBs, only $`\left(m_e/m_p\right)10^3`$ of the total kinetic energy is carried by electrons and then emitted as soft gamma-rays. If the hidden energy carried by protons is directly emitted in the TeV range, then much more energy can be radiated in the TeV range than in the sub-MeV range by a factor of almost one thousand.
GRBs are known as a possible site for the acceleration of protons up to $`10^{20}`$ eV, which are observed on the Earth as ultra-high energy cosmic rays (Waxman 1995; Vietri 1995; Milgrom & Usov 1995). I have already shown (Totani 1998b) that, when $`E_{\mathrm{iso}}\stackrel{>}{}10^{55}`$ erg and the magnetic field is as strong as the energy density of the shocked region, synchrotron radiation of protons accelerated to $`10^{20}`$ eV can be an efficient emission process because the cooling time of such protons is comparable with the typical GRB duration ($``$ 10 sec) in the observer’s frame. The energy of these synchrotron photons for an observer is about 1–10 TeV, and strong TeV emission from GRBs is possible. I suggest that the TeV gamma-rays possibly detected by the Milagrito were produced by this mechanism.
On the other hand, it may also be possible that a physical process works as an energy conveyor from the hidden energy reservoir (i.e., protons) into electrons (or positrons). If the energy transfer is almost complete in a GRB, a significant fraction of $`E_{\mathrm{iso}}`$ can be radiated as gamma-rays in the sub-MeV range. I have pointed out (Totani 1999) that $`e^\pm `$-pair creation by TeV photons of proton-synchrotron might work as the new energy channel for the energy transfer from protons into electrons and positrons, giving an explanation for the energetic sub-MeV GRB phenomenon such as GRB 990123. Proton-synchrotron photons interact with low energy electron-synchrotron photons and create $`e^\pm `$ pairs. It can also be shown that the photon energy range of the synchrotron radiation of the created pairs becomes about MeV, i.e., consistent with the BATSE range.
Then what is the crucial parameter which determines whether a GRB is bright in TeV or MeV? The GRB luminosity in the sub-MeV range is determined by the efficiency of energy transfer from protons into $`e^\pm `$ pairs, i.e., the opacity of pair-production reaction for the proton-synchrotron TeV photons. Based on the typical fireball parameters of GRBs, this pair-production opacity is typically of order unity, and strongly depends on the bulk Lorentz factor of GRBs by the special relativistic effect as $`\tau \mathrm{\Gamma }^5`$ in a simple internal shock model (Totani 1999; see also Baring & Harding 1997; Böttecher & Dermer 1998). Because of this strong dependence of the pair-creation optical depth on $`\mathrm{\Gamma }`$, a modest dispersion in $`\mathrm{\Gamma }`$ by a factor of 3–4 from one GRB to another results in drastic change in the sub-MeV energetics of GRBs by a factor of up to $`10^3`$ (Totani 1999). A large value of $`\mathrm{\Gamma }`$ results in negligible optical depth to pair-creation and hence a GRB strong in TeV such as GRB 970417a, while a small $`\mathrm{\Gamma }`$ in the inverse case of a GRB strong in MeV such as GRB 990123.
This mechanism gives a natural explanation for the wide dispersion in the observed total GRB energy in the sub-MeV range, with almost no correlation with the afterglow luminosity (Totani 1999). This model assumes a relatively uniform distribution of the total kinetic energy emitted from the central engine, and sub-MeV luminosity of GRBs is not correlated with the kinetic energy injected into interstellar/circumstellar medium. It may be similar to a see-saw between sub-MeV and TeV energies, in which the total kinetic energy of GRBs is roughly the same for all GRBs and difference of GRB energetics is whether dominant emission is in TeV or MeV bands.
To summarize, GRB 970417a observed by the Milagrito can be understood as a GRB with the isotropic kinetic energy $`E_{\mathrm{iso}}\stackrel{>}{}10^{55}`$ erg ejected from the central engine, a significant fraction of which is radiated in the TeV range by the synchrotron radiation of ultra-high-energy protons, with almost no energy transfer from protons into electrons due to a relatively large value of $`\mathrm{\Gamma }`$.
## 4. Discussion
Here I discuss some observational signatures expected in future experiments from the interpretation presented in this letter.
The proton-synchrotron model predicts that the ratio of TeV/MeV luminosities drastically changes from burst to burst by the difference of energy transfer efficiency from protons into electrons and positrons. When the optical depth to the pair-creation reaction is much larger than unity, it is possible that the TeV energy fluence is much weaker than the sub-MeV fluence in contrast to the GRB 970417a, even after the intergalactic absorption of TeV gamma-rays is corrected. However, a generic prediction of this model is that the total of isotropic TeV and sub-MeV energies is about $`10^{5455}`$ erg for all GRBs with much smaller scatter than that of the sub-MeV or TeV isotropic energies.
The assumption that the GRB 970417a is the closest burst to us may be wrong if TeV luminosity of GRBs drastically changes from burst to burst, as expected in the proton-synchrotron model. However, I emphasize that this assumption is, in this case, conservative from a theoretical point of view. The argument that the redshift of the closest GRB in the 54 BATSE GRBs is about $`z0.7`$ is based only on the sub-MeV luminosity function of GRBs and GRB rate evolution. Then, if the GRB 970417a is not the closest but intrinsically even brighter GRB in TeV, its distance must be larger than $`z0.7`$. Therefore this estimate can be considered as a lower limit of the redshift. If the redshift is significantly larger than 0.7, the TeV isotropic energy of this burst would be much larger than $`10^{54}`$ erg, which would be quite difficult to explain by the stellar death models even if we invoke a strongly collimated jet-like explosion. Hence I consider the estimate of redshift and energetics presented here is reasonable.
The Milagro detector, which has significantly increased the sensitivity to GRBs between 0.1 and 10 TeV, is now operating to search for GRBs (Atkins et al. 2000). If the signal observed by the Milagrito is truly from GRB 970417a, the Milagro detector would detect an event similar to GRB 970417a with better signal-to-noise ratio. It is also important to detect GRBs more distant than GRB 970417a to increase the detection rate, because $`N\left(<z\right)`$ rapidly increases with redshift. It is unfortunate that the intergalactic optical depth of TeV gamma-rays also rapidly increases, and hence improvement of detector sensitivity does not significantly extend the distance of marginally detectable GRBs. The cut-off photon energy by the intergalactic absorption falls below $``$ 50 GeV at $`z1.2`$, where $`N\left(>z\right)0.1`$. Therefore it seems unlikely that the detection rate in the TeV range is increased to more than 1 per 10 BATSE-detected GRBs, even if the detector sensitivity is significantly improved.
Another improvement of the Milagro over Milagrito is the improved of spectral information. If we can measure the spectral cut-off by the intergalactic absorption for GRBs with known redshifts, it would provide information on the flux of the cosmic infrared background radiation as well as its evolution to $`z\stackrel{>}{}0.5`$, which would be valuable for the study of galaxy formation and evolution.
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# The Effect of a Small Mixing Angle in the Atmospheric Neutrinos 11footnote 1Talk presented at the 1st Workshop on Neutrino Oscillations and their Origin, February 11-13, 2000, Fujiyoshida, Japan
## Abstract
The effect of matter enhanced neutrino oscillations on atmosheric neutrinos is investigated systematically in the framework of one mass dominant model of three neutrinos. The resonance conditions of neutrino crossing the earth are determined by the three parameters, namely, the zenith angle, $`\mathrm{\Delta }m^2/E`$, and the mixing angle $`\theta _3`$ of the electron neutrinos with tau neutrinos. The values of the triplet under the resonance is found numerically.
## 1 Introduction
It was almost ten years ago that the atmospheric neutrino anomaly could be explained by the oscillations of muon neutrinos into tau neutrinos . This interpretation has been confirmed by the observation of the zenith angle dependence of the neutrino fluxes by the SuperKamiokande. While the zenith angle dependence of the electron-like events is consistent with the theoretical predictions, that of the muon-like events disagree with the theory, especially for the up-going events.
The neutrino state with flavor basis and that of mass basis is related to each other by the mixing matrix as follows,
$$\left(\begin{array}{c}\nu _e\\ \nu _\mu \\ \nu _\tau \end{array}\right)=\left(\begin{array}{ccc}c_1c_3& s_1c_3& s_3\\ s_1c_2c_1s_2s_3& c_1c_2s_1s_2s_3& s_2c_3\\ s_1s_2c_1c_2s_3& c_1s_2s_1c_2s_3& c_2c_3\end{array}\right)\left(\begin{array}{c}\nu _1\\ \nu _2\\ \nu _3\end{array}\right),$$
(1)
where $`s_i=\mathrm{sin}\theta _i,c_i=\mathrm{cos}\theta _i\mathrm{for}i=1,2,\mathrm{and}3.`$ The Super-Kamiokande shows that
$$P(\nu _\mu \nu _\mu )+P(\nu _\mu \nu _\tau )1.$$
(2)
The equation (2) means that $`s_30,`$ which is consistent with the CHOOS results.
Once the the matter enhanced neutrino oscillations was proposed , it has soon been applied to the atmospheric neutrinos crossing the earth. Recently much attention has been devoted to this subject. However, it seems that no systmatic study has not been made yet, especially for the three flavor model of neutrinos.
The purpose of this paper is to reveal the conditions that amplify the effect of $`\theta _3`$ through matter oscillations so that we can determine this small mixing angle.
## 2 Matter Enhanced Atmospheric Neutrinos
I assume that the masses of neutrinos are hierarchical, namely, $`m_1<`$ $`m_2<m_3`$, Furthermore I assume that the solar neutrino oscillations are attributed to $`\delta m_{21}^2m_2^2m_1^2`$. The atmospheric neutrino oscillations are derived by $`\delta m_{32}^2`$ $`m_3^2m_2^2`$, the value of which is $`𝒪(10^3)eV^2`$. The solar neutrino problem may be solved by the matter enhanced, or just so oscillations. In either case, $`\delta m_{21}^2`$ is much smaller than $`\delta m_{32}^2`$, and irrelevant to our concerns. In the following, we assign $`\delta m_{21}^2=0`$, and $`\mathrm{\Delta }m^2\delta m_{31}^2=\delta m_{32}^2`$. The parameters contained in the model are $`\theta _2`$, $`\theta _3`$, and $`\mathrm{\Delta }m^2`$.
In calculating the matter effect of the earth, I use the density profile given by the “Preliminary reference Earth model”(PREM).
The prominent feature of our model is that the survival probability of electron neutrinos $`P(\nu _e\nu _e)`$ is completely determined by the three parameters, i.e., $`\theta _3`$, $`\mathrm{\Delta }m^2/E`$ where $`E`$ is the neutrino energy, and the zenith angle $`z`$. It is natural to define the resonance condition as
$$P(\nu _e\nu _e)=0$$
(3)
The resonance states form curves in the the three parameter space $`(\theta _3,\mathrm{\Delta }m^2/E,`$ $`\mathrm{cos}z)`$ corresponding to discretization of the resonance wavelength. I solve eq.(3) numerically, and show in Fig.1 and Fig.2 the values of the three parameters at the resonance. Fig.1 and Fig.2 correspond to the longest and the second longest wavelength, respectively.
The other probabilities $`P(\nu _e\nu _\mu )`$, $`P(\nu _e\nu _\tau )`$, $`P(\nu _\mu \nu _\mu )`$, and $`P(\nu _\mu `$ $`\nu _\tau )`$ also depend on the $`\nu _\mu \nu _\tau `$ mixing angle $`\theta _2`$.
## 3 Summary and Discussions
I have investigated the resonance conditions of neutrinos passing through the earth in a wide range of parameters, although only a small mixing angle $`\theta _3<13^{}`$ has a practical meaning as has been suggested by the CHOOS experiments.
The results show that it is necessary to observe the high energy electron neutrinos crossing the depth of the earth to investgate the small mixing angle $`\theta _3`$. For example, the resonance occurs at $`\theta _3=6.4^{}`$, $`\mathrm{\Delta }m^2/E=6.1\times 10^4eV^2/GeV`$, and $`\mathrm{cos}z=0.9`$. The energy of electron neutrino is $`5GeV`$ if $`\mathrm{\Delta }m^2=3\times 10^3eV^2`$. This indicates the need to assemble the data of the zenith angle distribution of the electron neutrinos whose energy is of order $`5GeV`$ in order to determine the lepton mixing angle $`\theta _3`$.
This research was supported in part by the Grant-in-Aid for Scientific Research of the Ministry of Education, Scinece and Culture of Japan No. 10640281.
Figure Captions
Fig.1 The resonace conditions(I) as functions of $`\mathrm{cos}z`$. The solid line represents $`\theta _3`$, and the dot-dashed line represents $`\mathrm{\Delta }m^2/E`$.
Fig.2 The resonace conditions(II) as functions of $`\mathrm{cos}z`$. The solid linerepresents $`\theta _3`$, and the dot-dashed line represents $`\mathrm{\Delta }m^2/E`$.
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# 1 Introduction
## 1 Introduction
Linked cluster expansions (LCEs) have a long tradition in statistical physics. Originally applied to classical fluids, later to magnetic systems (,, and references therein), they were generalized to applications in particle physics in the eighties . There they have been used to study the continuum limit of a lattice $`\mathrm{\Phi }^4`$ field theory in 4 dimensions at zero temperature. In they were further generalized to field theories at finite temperature, simultaneously the highest order in the expansion parameter was increased to 18. Usually the analytic expansions are obtained as graphical expansions. Because of the progress in computer facilities and the development of efficient algorithms for generating the graphs, it is nowadays possible to handle of the order of billions of graphs. The whole range from high temperatures down to the critical region becomes available, and thermodynamic quantities like critical indices and critical temperatures are determined with high precision (the precision is comparable or even better than in corresponding high quality Monte Carlo results) -. An extension of LCEs to a finite volume in combination with a high order in the expansion parameter turned out to be a particularly powerful tool for investigating the phase structure of systems with first and second order transitions by means of a finite size scaling analysis .
Linked cluster expansions are series expansions of the free energy and connected correlation functions about an ultralocal, decoupled theory in terms of a hopping parameter $`K`$. The corresponding graphical representation is a sum in terms of connected graphs. The value of $`K`$ parametrizes the strength of interactions between fields at different lattice sites. Usually they are chosen as nearest neighbours. In contrast to the ultralocal terms of a generic interaction we will sometimes refer to hopping terms as non-ultralocal.
In this paper we develop dynamical linked cluster expansions (DLCEs). These are linked cluster expansions with hopping parameter terms that are endowed with their own dynamics. Such systems are realized in spin glasses with (fast) spins and (slow) interactions -. They also occur in variational estimates for SU(N)-gauge-Higgs systems, cf. . Like LCEs they are expected to converge for a large class of interactions.
Formally DLCEs amount to a generalization of an expansion scheme from 2-point to point-link-point-interactions. These are interactions between fields associated with two points and with one pair of points called link. The points and links are not necessarily embedded on a lattice, and the links need not be restricted to nearest neighbours. We have developed a new multiple-line graph theory in which a generalized notion of connectivity plays a central role. Standard notions of equivalence classes of graphs like 1-line irreducible and 1-vertex irreducible graphs have been generalized, and new notions like 1-multiple-line irreducible graphs were defined in order to give a systematic classification.
The paper is organized as follows. In Sec. 2 we specify the models that admit a DLCE. We introduce multiple-line graphs and explain the idea behind the abstract notions of multiple-line graph theory. Detailed definitions of multiple-line graphs , related notions and the computation of weights are given in Sect. 3. Sect. 4 treats the issue of renormalization in the sense of suitable resummations of graphs. Applications to spin glasses are presented in Sect. 5. There it is of particular interest that DLCEs allow for the possibility of avoiding the replica trick. In the quenched limit we derive that DLCEs must be restricted to a subclass of the corresponding graphical expansion, so-called quenched DLCEs (QDLCEs). We also list some examples for models whose phase structure is accessible to QDLCEs. To these belongs in particular the bond diluted Ising model. Sect. 6 contains the summary and conclusions.
## 2 A Short Primer to DLCEs
In this section we first specify the class of models for which we develop dynamical linked cluster expansions. Next we illustrate some basic notions of multiple-line graph theory, in particular the need for a new notion of connectivity.
By $`\mathrm{\Lambda }_0`$ we denote a finite or infinite set of points. One of its realizations is a hypercubic lattice in $`D`$ dimensions, infinite or finite in some directions with the topology of a torus. $`\mathrm{\Lambda }_1`$ denotes the set of unordered pairs $`(x,y)`$ of sites $`x,y\mathrm{\Lambda }_0`$, $`xy`$, also called unoriented links, and $`\overline{\mathrm{\Lambda }}_1`$ a subset of $`\mathrm{\Lambda }_1`$.
We consider physical systems with a partition function of the generic form
$`Z(H,I,v)\mathrm{exp}W(H,I,v)`$
$`=𝒩{\displaystyle 𝒟\varphi 𝒟U\mathrm{exp}(S(\varphi ,U,v))\mathrm{exp}(\underset{x\mathrm{\Lambda }_{}}{}H(x)\varphi (x)+\underset{l\overline{\mathrm{\Lambda }}_1}{}I(l)U(l))},`$ (1)
with measures
$$𝒟\varphi =\underset{x\mathrm{\Lambda }_0}{}d\varphi (x),𝒟U=\underset{l\overline{\mathrm{\Lambda }}_1}{}dU(l)$$
(2)
and action
$$S(\varphi ,U,v)=\underset{x\mathrm{\Lambda }_{}}{}S^{}(\varphi (x))+\underset{l\overline{\mathrm{\Lambda }}_1}{}S^1(U(l))\frac{1}{2}\underset{x,y\mathrm{\Lambda }_{}}{}v(x,y)\varphi (x)U(x,y)\varphi (y),$$
(3)
with non-ultralocal couplings
$`v(x,y)=v(y,x)\mathrm{\hspace{0.33em}0}\text{only for }(x,y)\overline{\mathrm{\Lambda }}_1,`$
$`\text{in particular}v(x,x)=\mathrm{\hspace{0.33em}0}.`$ (4)
For later convenience the normalization via $`𝒩`$ is chosen such that $`W[0,0,0]=0`$.
The field $`\varphi (x)`$ is associated with the sites $`x\mathrm{\Lambda }_0`$ and the field $`U(l)`$ lives on the links $`l\overline{\mathrm{\Lambda }}_1`$, and we write $`U(x,y)=U(l)`$ for $`l=(x,y)`$. For definiteness and for simplicity of the notation here we assume $`\varphi (x)𝐑`$ and $`U(l)𝐑`$. In our actual applications to spin glasses the $`\varphi `$s are the (fast) Ising spins and the $`U`$s $`𝐑`$ are the (slow) interactions. The action is split into two ultralocal parts, $`S^{}`$ depending on fields on single sites, and $`S^1`$ depending on fields on single links $`l\overline{\mathrm{\Lambda }}_1`$. For simplicity we choose $`S^1`$ as the same function for all links $`l\overline{\mathrm{\Lambda }}_1`$. We may identify $`\overline{\mathrm{\Lambda }}_1`$ with the support of $`v`$,
$$\overline{\mathrm{\Lambda }}_1=\{l=(x,y)|v(x,y)0\}.$$
(5)
The support of $`v(x,y)`$ need not be restricted to nearest neighbours, also the precise form of $`S^{}`$ and $`S^1`$ does not matter for the generic description of DLCEs, $`S^{}`$ and $`S^1`$ can be any polynomials in $`\varphi `$ and $`U`$, respectively. The only restriction is the existence of the partition function.
Note that the interaction term $`v(x,y)\varphi (x)U(x,y)\varphi (y)`$ contains a point-link-point-interaction and generalizes the 2-point-interactions $`v(x,y)\varphi (x)\varphi (y)`$ of usual hopping parameter expansions. The effective coupling of the $`\varphi `$ fields has its own dynamics governed by $`S^1(U)`$, the reason why we have called our new expansion scheme dynamical LCE.
Dynamical linked cluster expansions are induced from a Taylor expansion of $`W(H,I,v)=\mathrm{ln}Z(H,I,v)`$ about $`v=0`$, the limit of a completely decoupled system. We want to express the series for $`W`$ in terms of connected graphs. Let us consider the generating equation
$`W/v(xy)`$ $`=`$ $`1/2<\varphi (x)U(x,y)\varphi (y)>`$ (6)
$`=`$ $`1/2(W_{H(x)I(x,y)H(y)}+W_{H(x)H(y)}W_{I(x,y)}`$
$`+`$ $`W_{H(x)I(x,y)}W_{H(y)}+W_{I(x,y)H(y)}W_{H(x)}`$
$`+`$ $`W_{H(x)}W_{H(y)}W_{I(x,y)}).`$
Here $`<>`$ denotes the normalized expectation value w.r.t. the partition function of Eq. (2). Subscripts $`H(x)`$ and $`I(x,y)=I(y,x)=I(l)`$ denote the derivatives of $`W`$ w.r.t. $`H(x)`$ and $`I(x,y)`$, respectively.
Next we would like to represent the right hand side of Eq. (6) in terms of connected graphs. Once we have such a representation for the first derivative of $`W`$ w.r.t. $`v`$, grapical expansions for the higher derivatives can be traced back to the first one.
For each $`W`$ in Eq. (6) we draw a shaded bubble, for each derivative w.r.t. $`H`$ a solid line, called a $`\varphi `$-line, with endpoint vertex $`x`$, and for each derivative w.r.t. $`I`$ a dashed line, called a $`U`$-line, with link label $`l=(x,y)`$. The main graphical constituents are shown in Fig. 1. Two $`\varphi `$-lines with endpoints $`x`$ and $`y`$ are then joined by means of a dashed $`U`$-line with label $`l`$, if the link $`l`$ has $`x`$ and $`y`$ as its endpoints, i.e. $`l=(x,y)`$. According to these rules Eq. (6), multiplied by $`v(x,y)`$ and summed over $`x`$ and $`y`$, is represented by Fig. 2. Note that, because of the Taylor operation, each solid line from $`x`$ to $`y`$ carries a factor $`v(x,y)`$.
Since the actual need for a new type of connectivity is not quite obvious from Fig. 2, because Eq. (6) does not contain higher than first order derivatives w.r.t. $`I`$, let us consider a term
$$W_{H(x)}W_{H(y)}W_{H(r)}W_{H(s)}W_{I(x,y)I(r,s)}$$
(7)
occurring in the second derivative of $`W`$ w.r.t. $`v(x,y)`$, $`v(r,s)`$. According to the above rules this term would be represented as shown in Fig. 3a. While the 2 vertices in the last term of Fig. 2 are connected in the usual sense via a common (solid) line (the dashed line with an attached bubble could be omitted in this case), the graph in Fig. 3a would be disconnected in the old sense, since neither $`x`$ nor $`y`$ are line-connected with $`r`$ and $`s`$, but -as a new feature of DLCE graphs- the lines from $`x`$ to $`y`$ and from $`r`$ to $`s`$ are connected via the dashed lines emerging from a common bubble shown in the middle of the graph. As we see from Fig. 3a, we need an additional notion of connectivity referring to the possibility of multiple-line connectivity. While the analytic expression is fixed, it is a matter of convenience to further simplify the graphical notation of Fig. 3a at $`v=0`$. Two possibilities are shown in Fig. 3b and Fig. 3c. To Fig. 3b we later refer in the formal definition of the new type of multiple-line connectivity. In the familiar standard notion of connectivity two vertices of a graph are connected via lines. The vertices are line-connected. Already there, in a dual language, one could call two lines connected via vertices. The second formulation is just appropriate for our need to define when two lines are connected. The corresponding vertices mediating the connectivity of lines are visualized by tubes, in Fig. 3b we have just one of them. The tubes should be distinguished from the former type of vertices represented as full dots which are connected via bare $`\varphi `$-lines. In Fig. 3c we show a simplified representation of Fig. 3b that we actually use in graphical expansions.
The derivative terms have to be evaluated at $`v=0`$. For $`v=0`$ we have a decomposition of $`W`$ according to
$$W(H,I,v=0)=\underset{x\mathrm{\Lambda }_0}{}W^{}(H(x))+\underset{l\overline{\mathrm{\Lambda }}_1}{}W^1(I(l))$$
(8)
with
$$\mathrm{exp}W^{}(H)Z^{}(H)=\frac{_{\mathrm{}}^{\mathrm{}}𝑑\varphi \mathrm{exp}(S^{}(\varphi )+H\varphi )}{_{\mathrm{}}^{\mathrm{}}𝑑\varphi \mathrm{exp}(S^{}(\varphi ))}$$
(9)
and
$$\mathrm{exp}W^1(l)Z^1(I)=\frac{_{\mathrm{}}^{\mathrm{}}𝑑U\mathrm{exp}(S^1(U)+IU)}{_{\mathrm{}}^{\mathrm{}}𝑑U\mathrm{exp}(S^1(U))}.$$
(10)
In Eq.s (9,10) we have omitted any single site or single link dependence, because we assume that $`S^{}`$ and $`S^1`$ are the same for all $`x\mathrm{\Lambda }_0`$ and $`l\overline{\mathrm{\Lambda }_1}`$, respectively. Therefore, at $`v=0`$, the only non-vanishing derivatives of $`W`$ are
$$W_{H(x_1)H(x_2)\mathrm{}H(x_n)}|_{v=0}=\frac{^nW^{}(H(x_1))}{H(x_1)^n}\delta _{x_1,x_2,\mathrm{}x_n}$$
(11)
and
$$W_{I(l_1)I(l_2)\mathrm{}I(l_m)}|_{v=0}=\frac{^mW^1(I(l_1))}{I(l_1)^m}\delta _{l_1,l_2,\mathrm{}l_m},$$
(12)
but mixed derivatives w.r.t. H and I vanish. As anticipated in Fig.s 3b and 3c, for $`v=0`$ we replace the dashed bubbles and graphically distinguish between bubbles with $`\varphi `$-lines and bubbles with $`U`$-lines. We define
$$\text{}\}n=v^c_n=(\frac{^nW^{}(H)}{H^n})_{H=0}$$
(13)
for a connected n-point vertex with $`n1`$ bare $`\varphi `$-lines emerging from it and
$$\text{}\}\nu =m^{1c}_\nu =(\frac{^\nu W^1(I)}{I^\nu })_{I=0}$$
(14)
for a connected $`\nu `$-line consisting of $`\nu `$ bare lines. If $`\nu =1`$, we often omit the dashed line. If the bare lines of a $`\nu `$-line are internal $`\varphi `$-lines, they get vertices attached to their endpoints, if they are external $`U`$-lines, no vertices will be attached.
Let $`V`$ denote the lattice volume in $`D`$ dimensions. The Taylor expansion of $`W`$ about $`v=0`$ to second order in $`v`$ then reads
$`W(H,I,v)`$ $`=`$ $`W(H,I,v=0)`$
$`+`$ $`{\displaystyle \underset{x,y\mathrm{\Lambda }_0}{}}v(x,y){\displaystyle \frac{1}{2}}W_{H(x)}W_{H(y)}W_{I(x,y)}`$
$`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{x,y,r,s\mathrm{\Lambda }_0}{}}{\displaystyle \frac{1}{4}}v(x,y)v(r,s)`$
$``$ $`(`$ $`4W_{H(y)}W_{H(s)}W_{H(r)H(x)}W_{I(x,y)}W_{I(r,s)}`$ (15)
$`+`$ $`2W_{H(x)H(r)}W_{H(y)H(s)}W_{I(x,y)}W_{I(r,s)}`$
$`+`$ $`4W_{H(y)}W_{H(s)}W_{H(r)H(x)}W_{I(x,y)I(r,s)}`$
$`+`$ $`2W_{H(r)H(x)}W_{H(y)H(s)}W_{I(r,y)I(x,s)}`$
$`+`$ $`W_{H(x)}W_{H(y)}W_{H(r)}W_{H(s)}W_{I(x,y)I(r,s)})_{v=0}`$
$`+`$ $`O(v^3),`$
where we have used that $`v(x,x)=0`$. For each $`W`$ in the products of $`W`$s we now insert Eq.s (11),(12).
If we choose $`v`$ in a standard way as next-neighbour couplings
$$v(x,y)=2K\underset{\mu =0}{\overset{D1}{}}(\delta _{x+\widehat{\mu },y}+\delta _{x\widehat{\mu },y})$$
(16)
with $`\widehat{\mu }`$ denoting the unit vector in $`\mu `$-direction, Eq. (2) becomes in a graphical representation at $`H=I=0`$
$`{\displaystyle \frac{W(0,0,v)}{V}}`$ $`=`$ $`(2K){\displaystyle \frac{1}{2}}(2D)\text{}`$
$`+`$ $`(2K)^2\{{\displaystyle \frac{1}{2}}(2D)^2\text{}+{\displaystyle \frac{1}{4}}(2D)\text{}`$
$`+{\displaystyle \frac{1}{2}}(2D)^2\text{}+{\displaystyle \frac{1}{4}}(2D)\text{}`$
$`+{\displaystyle \frac{1}{8}}\mathrm{\hspace{0.33em}2}(2D)\text{}\}`$
$`+`$ $`O(K^3).`$
For clarity, here we have written explicitly the topological symmetry factors and the lattice embedding numbers. (Usually graphs represent their full weights including these factors.) Note that the first two graphs of the second order contribution also occur in a usual LCE with frozen $`U`$-dynamics, the next two differ by an additional dashed 2-line and the last one becomes even disconnected without the dashed line.
As usual, graphical expansions for correlation functions, in particular susceptibilities, are generated from $`W(H,I,v)`$ by taking derivatives w.r.t. the external fields $`H`$ and $`I`$. Graphically this amounts to attaching external $`\varphi `$-lines and $`U`$-lines with
(18)
In passing we remark that the conventional LCE is included as a special case of the DLCE, if the $`U`$-dynamics is ”frozen” to some value $`U_00`$, so that
$`W^1(I)`$ $`=`$ $`S_1(U_0)+IU_0,`$
$`{\displaystyle \frac{W^1(I)}{I}}`$ $`=`$ $`U_0,`$ (19)
$`{\displaystyle \frac{^nW^1(I)}{I^n}}`$ $`=`$ $`0\text{for all }n>1,`$
i.e., no n-lines do occur with $`n>1`$. In this case it becomes redundant to attach dashed lines to bare lines. As mentioned above, in an LCE only the first three contributions would be left in Eq. (2).
## 3 Graphical expansion
### 3.1 Multiple-line graph theory
The definition of a multiple-line graph as it will be given here is adapted to the computation of susceptibilities and the free energy, where the sum is taken over all possible locations of the fields. The definition easily generalizes to correlation functions.
For details of the standard definiton of graphs in the framework of linked cluster expansions and related notions we refer e.g. to . Here, for convenience, we briefly recall the very definition of a graph to point out the new properties of multiple-line graphs as defined below in this section.
A (standard LCE) graph or diagram is a structure
$$\stackrel{~}{\mathrm{\Gamma }}=(\stackrel{~}{}_\mathrm{\Gamma },\stackrel{~}{}_\mathrm{\Gamma },\stackrel{~}{E}_\mathrm{\Gamma },\stackrel{~}{\mathrm{\Phi }}_\mathrm{\Gamma }),$$
(20)
where $`\stackrel{~}{}_\mathrm{\Gamma }`$ and $`\stackrel{~}{}_\mathrm{\Gamma }\mathrm{}`$ are disjoint sets of internal lines and vertices of $`\stackrel{~}{\mathrm{\Gamma }}`$, respectively. $`\stackrel{~}{E}_\mathrm{\Gamma }`$ is a map
$`\stackrel{~}{E}_\mathrm{\Gamma }:\stackrel{~}{}_\mathrm{\Gamma }`$ $``$ $`\{0,1,2,\mathrm{}\},`$
$`v`$ $``$ $`\stackrel{~}{E}_\mathrm{\Gamma }(v)`$ (21)
that assigns to every vertex $`v`$ the number of external lines $`\stackrel{~}{E}_\mathrm{\Gamma }(v)`$ attached to it. Finally, $`\stackrel{~}{\mathrm{\Phi }}_\mathrm{\Gamma }`$ is the incidence relation that assigns internal lines to their two endpoints.
A multiple-line graph or multiple-line diagram is a structure
$$\mathrm{\Gamma }=(_\mathrm{\Gamma },_\mathrm{\Gamma },_\mathrm{\Gamma },E_\mathrm{\Gamma }^{(\varphi )},E_\mathrm{\Gamma }^{(U)},\mathrm{\Phi }_\mathrm{\Gamma },\mathrm{\Psi }_\mathrm{\Gamma }).$$
(22)
$`_\mathrm{\Gamma }`$, $`_\mathrm{\Gamma }`$ and $`_\mathrm{\Gamma }`$ are three mutually disjoint sets,
$`_\mathrm{\Gamma }`$ $`=`$ $`\text{set of bare internal lines of }\mathrm{\Gamma },`$ (23)
$`_\mathrm{\Gamma }`$ $`=`$ $`\text{set of multiple lines of }\mathrm{\Gamma },`$ (24)
$`_\mathrm{\Gamma }`$ $`=`$ $`\text{set of vertices of }\mathrm{\Gamma }.`$ (25)
$`E_\mathrm{\Gamma }^{(\varphi )}`$ is a map
$`E_\mathrm{\Gamma }^{(\varphi )}:_\mathrm{\Gamma }`$ $``$ $`\{0,1,2,\mathrm{}\},`$
$`v`$ $``$ $`E_\mathrm{\Gamma }^{(\varphi )}(v)`$ (26)
that assigns to every vertex $`v`$ the number of bare external $`\varphi `$-lines $`E_\mathrm{\Gamma }^{(\varphi )}(v)`$ attached to $`v`$. Every such $`\varphi `$-line represents a field $`\varphi `$. The number of external $`\varphi `$-lines of $`\mathrm{\Gamma }`$ is denoted by $`E_\mathrm{\Gamma }^{(\varphi )}=_{v_\mathrm{\Gamma }}E_\mathrm{\Gamma }^{(\varphi )}(v)`$. Similarly, $`E_\mathrm{\Gamma }^{(U)}`$ is a map
$`E_\mathrm{\Gamma }^{(U)}:_\mathrm{\Gamma }`$ $``$ $`\{0,1,2,\mathrm{}\},`$
$`m`$ $``$ $`E_\mathrm{\Gamma }^{(U)}(m)`$ (27)
that assigns to every multiple line $`m`$ the number of external $`U`$-lines $`E_\mathrm{\Gamma }^{(U)}(m)`$ attached to $`m`$. Every such $`U`$-line represents a field $`U`$ associated with a lattice link. The number of external $`U`$-lines of $`\mathrm{\Gamma }`$ is given by $`E_\mathrm{\Gamma }^{(U)}=_{m_\mathrm{\Gamma }}E_\mathrm{\Gamma }^{(U)}(m)`$.
Furthermore, $`\mathrm{\Phi }_\mathrm{\Gamma }`$ and $`\mathrm{\Psi }_\mathrm{\Gamma }`$ are incidence relations that assign bare internal lines to their endpoint vertices and to their multiple lines, respectively. We treat lines as unoriented. The generalization to oriented lines is easily done. More precisely, let $`\overline{(_\mathrm{\Gamma }\times _\mathrm{\Gamma })}^{}`$ be the set of unordered pairs of vertices $`(v,w)`$ with $`v,w_\mathrm{\Gamma }`$, $`vw`$. (The bar implies unordered pairs, the prime the exclusions of $`(v,v)`$, $`v_\mathrm{\Gamma }`$.) As in the standard linked cluster expansion, self-lines are excluded. Every bare internal line is then mapped onto its pair of endpoints via
$$\mathrm{\Phi }_\mathrm{\Gamma }:_\mathrm{\Gamma }\overline{(_\mathrm{\Gamma }\times _\mathrm{\Gamma })}^{}.$$
(28)
We say that $`v`$ and $`w`$ are the endpoint vertices of $`l_\mathrm{\Gamma }`$ if $`\mathrm{\Phi }_\mathrm{\Gamma }(l)=(v,w)`$. If there is such an $`l_\mathrm{\Gamma }`$, $`v`$ and $`w`$ are called neighbours. Similarly, $`\mathrm{\Psi }_\mathrm{\Gamma }`$ is a map
$`\mathrm{\Psi }_\mathrm{\Gamma }:_\mathrm{\Gamma }`$ $``$ $`_\mathrm{\Gamma },`$
$`l`$ $``$ $`\mathrm{\Psi }_\mathrm{\Gamma }(l)`$ (29)
that maps every bare internal line to a multiple line. A multiple line $`m_\mathrm{\Gamma }`$ is composed of bare internal lines $`l_\mathrm{\Gamma }`$ which belong to $`m`$ in the sense that $`\mathrm{\Psi }_\mathrm{\Gamma }(l)=m`$. $`l__\mathrm{\Gamma }(m)`$ is the total number of bare internal lines belonging to $`m`$. With $`\nu =l__\mathrm{\Gamma }(m)+E_\mathrm{\Gamma }^{(U)}(m)`$, $`m`$ is called a $`\nu `$-line. We always require that $`\nu 1`$. On the other hand, every bare internal line belongs to one and only one multiple line. For simplicity we often identify a $`1`$-line with the only one bare line that belongs to it. Next we introduce some further notions that will be used later. External vertices are vertices having external $`\varphi `$-lines attached,
$$_{\mathrm{\Gamma },ext}=\{v_\mathrm{\Gamma }|E_\mathrm{\Gamma }^{(\varphi )}(v)0\},$$
(30)
whereas internal vertices do not, $`_{\mathrm{\Gamma },int}=_\mathrm{\Gamma }_{\mathrm{\Gamma },ext}`$. Similarly, external multiple lines have external $`U`$-lines attached,
$$_{\mathrm{\Gamma },ext}=\{m_\mathrm{\Gamma }|E_\mathrm{\Gamma }^{(U)}(m)0\},$$
(31)
and the complement in $`_\mathrm{\Gamma }`$ are the internal multiple lines, $`_{\mathrm{\Gamma },int}=_\mathrm{\Gamma }_{\mathrm{\Gamma },ext}`$.
For every pair of vertices $`v,w_\mathrm{\Gamma }`$, $`vw`$, let $`\overline{\mathrm{\Phi }}^1(v,w)`$ be the set of lines with endpoint vertices $`v`$ and $`w`$, and $`|\overline{\mathrm{\Phi }}^1(v,w)|`$ the number of these lines. Thus $`\overline{\mathrm{\Phi }}^1(v,w)`$ is the set of lines $`v`$ and $`w`$ have in common. With $`E_\mathrm{\Gamma }^{(\varphi )}(v)`$ denoting the number of external $`\varphi `$-lines attached to $`v_\mathrm{\Gamma }`$,
$$t__\mathrm{\Gamma }(v)=\underset{w_\mathrm{\Gamma }}{}|\overline{\mathrm{\Phi }}^1(v,w)|+E_\mathrm{\Gamma }^{(\varphi )}(v)$$
(32)
is the total number of bare lines attached to $`v`$.
Some topological notions and global properties of graphs will be of major interest in the following. A central notion is the connectivity of a multiple-line graph . Recall that we want to consider the DLCE expansion of the free energy and of truncated correlation functions as an expansion in connected graphs. As indicated in section 2, the main generalization compared to the common notion of connectivity of a graph which is required here is that an additional type of connectivity is provided by multiple-lines. To define the connectivity of a multiple-line graph $`\mathrm{\Gamma }`$, $`\mathrm{\Gamma }`$ first is mapped to a (standard) LCE graph $`\overline{\mathrm{\Gamma }}`$ to which the standard notion of connectivity applies. There are various equivalent ways to define such a map. We choose the following one.
* For every multiple-line $`m_\mathrm{\Gamma }`$ define a new vertex $`w(m)`$. Let $`\stackrel{~}{}_\mathrm{\Gamma }=\{w(m)|m_\mathrm{\Gamma }\}`$ and define $`\overline{}=_\mathrm{\Gamma }\stackrel{~}{}_\mathrm{\Gamma }`$ as the union of the vertices of $`\mathrm{\Gamma }`$ and the new set of vertices originating from the multiple-lines.
* For every bare internal line $`l_\mathrm{\Gamma }`$ define two new internal lines $`l_1,l_2`$ and incidence relations
$`\overline{\mathrm{\Phi }}(l_1)`$ $`=`$ $`(v_1,w(\mathrm{\Psi }_\mathrm{\Gamma }(l))),`$
$`\overline{\mathrm{\Phi }}(l_2)`$ $`=`$ $`(v_2,w(\mathrm{\Psi }_\mathrm{\Gamma }(l))),`$ (33)
where $`v_1`$ and $`v_2`$ are the two endpoint vertices of $`l`$. The set of all lines $`l_1,l_2`$, for all $`l_\mathrm{\Gamma }`$, is denoted by $`\overline{}`$.
* Define the external incidence relations
$`\overline{E}:\overline{}`$ $``$ $`\{0,1,2,\mathrm{}\},`$
$`\overline{E}(v)`$ $`=`$ $`E_\mathrm{\Gamma }^{(\varphi )}(v),\text{for }v_\mathrm{\Gamma },`$ (34)
$`\overline{E}(v)`$ $`=`$ $`E_\mathrm{\Gamma }^{(U)}(m),\text{for }v=w(m)\stackrel{~}{}_\mathrm{\Gamma }.`$
Now, $`\overline{\mathrm{\Gamma }}`$ is defined by
$$\overline{\mathrm{\Gamma }}=(\overline{},\overline{},\overline{E},\overline{\mathrm{\Phi }}).$$
(35)
Having defined the standard LCE graph $`\overline{\mathrm{\Gamma }}`$ for any multiple-line graph $`\mathrm{\Gamma }`$, we call $`\mathrm{\Gamma }`$ multiple-line connected or just connected if $`\overline{\mathrm{\Gamma }}`$ is connected (in the usual sense). In Fig. 4 we have given two examples for a connected (upper graph) and a disconnected (lower graph) multiple-line graph .
The next important notion is the topological equivalence of two multiple-line graphs . Two multiple-line graphs
$$\mathrm{\Gamma }_i=(_i,_i,_i,E_i^{(\varphi )},E_i^{(U)},\mathrm{\Phi }_i,\mathrm{\Psi }_i),i=1,2$$
(36)
are called (topologically) equivalent if there are three invertible maps
$`\varphi _1:_1`$ $``$ $`_2,`$
$`\varphi _2:_1`$ $``$ $`_2,`$ (37)
$`\varphi _3:_1`$ $``$ $`_2,`$
such that
$`\mathrm{\Phi }_2\varphi _2`$ $`=`$ $`\overline{\varphi }_1\mathrm{\Phi }_1,`$
$`\mathrm{\Psi }_2\varphi _2`$ $`=`$ $`\varphi _3\mathrm{\Psi }_1,`$ (38)
and
$`E_2^{(\varphi )}\varphi _1`$ $`=`$ $`E_1^{(\varphi )},`$
$`E_2^{(U)}\varphi _3`$ $`=`$ $`E_1^{(U)}.`$ (39)
Here $``$ means decomposition of maps, and
$`\overline{\varphi }_1:\overline{_1\times _1}^{}`$ $``$ $`\overline{_2\times _2}^{}`$
$`\overline{\varphi }_1(v,w)`$ $`=`$ $`(\varphi _1(v),\varphi _1(w)).`$ (40)
A symmetry of a multiple-line graph $`\mathrm{\Gamma }=(,,,E^{(\varphi )},E^{(U)},\mathrm{\Phi },\mathrm{\Psi })`$ is a triple of maps $`\varphi _1:`$, $`\varphi _2:`$ and $`\varphi _3:`$ such that
$`\mathrm{\Phi }\varphi _2`$ $`=`$ $`\overline{\varphi }_1\mathrm{\Phi },`$
$`\mathrm{\Psi }\varphi _2`$ $`=`$ $`\varphi _3\mathrm{\Psi },`$ (41)
and
$`E^{(\varphi )}\varphi _1`$ $`=`$ $`E^{(\varphi )}`$
$`E^{(U)}\varphi _3`$ $`=`$ $`E^{(U)}.`$ (42)
The number of these maps is called the symmetry number of $`\mathrm{\Gamma }`$.
We denote by $`𝒢_{E_1,E_2}(L)`$ the set of equivalence classes of connected multiple-line graphs with $`L`$ bare internal lines, $`E_1`$ external $`\varphi `$-lines and $`E_2`$ external $`U`$-lines. Furthermore we set
$$𝒢_{E_1,E_2}:=\underset{L0}{}𝒢_{E_1,E_2}(L).$$
(43)
A multiple line graph $`\mathrm{\Gamma }`$ does not need to have a vertex. If $`_\mathrm{\Gamma }=0`$, we have $`_\mathrm{\Gamma }=0`$ as well. If in addition $`\mathrm{\Gamma }`$ is connected, $`_\mathrm{\Gamma }`$ consists of only one element, with all external $`U`$-lines attached to it. (We anticipate that $`\mathrm{\Gamma }`$ is 1-multiple-line irreducible (1MLI) by definition. For the definition of 1MLI cf. section 4 below.) The only graph of $`𝒢_{0,E}(L=0)`$ is given by
$$\mathrm{\Gamma }=\text{}\}E.$$
(44)
It represents the leading term of the susceptibility
$$\chi _{0,E}=\frac{1}{VD}\underset{l_1,\mathrm{},l_E\overline{\mathrm{\Lambda }}_1}{}<U(l_1)\mathrm{}U(l_E)>^c$$
(45)
and is given by $`^EW^1(I)/I^E|_{I=0}`$. The index $`c`$ in (45) stands for truncated (connected) correlation.
By removal of a $`\nu `$-line $`m_\mathrm{\Gamma }`$ we mean that $`m`$ is dropped together with all bare internal lines and all external $`U`$-lines that belong to $`m`$. This notion is explained in Fig. 5a. (It is used in section 4 for $`1`$-lines to define 1-particle irreducible (1PI) and 1-line irreducible (1LI) multiple-line graphs .)
On the other hand, by decomposition of a $`\nu `$-line $`m_\mathrm{\Gamma }`$ we mean that $`m`$ is dropped together with the external $`U`$-lines of $`m`$, but all bare internal lines that belong to $`m`$ are kept in the graph, being identified now with $`1`$-lines. This notion will be used below to define 1MLI and renormalized multiple-line moments. It is illustrated in Fig. 5b.
Similarly, decomposition of a vertex $`v_\mathrm{\Gamma }`$ means to remove the vertex $`v`$ and to attach the free end of every line that entered $`v`$ before to a new vertex, a separate one for each line. This notion is used to define 1-vertex-irreducible (1VI) and renormalized vertex moments for multiple-line graphs . For an example see Fig. 5c.
### 3.2 Susceptibilities and weights
In the last section we have defined multiple-line graphs and the notions of connectivity and equivalence of such graphs. The definition is chosen in such a way that the series expansions of the free energy and of truncated correlation functions are obtained as a sum over equivalence classes of connected multiple-line graphs . The number $`L`$ of bare internal lines of a multiple-line graph $`\mathrm{\Gamma }`$ counts the order in the expansion parameter $`v(x,y)`$ to which $`\mathrm{\Gamma }`$ contributes. If $`v(x,y)`$ is of the form
$$v(x,y)=\mathrm{\hspace{0.33em}2}K\underset{z𝒩(x)}{}\delta _{y,z},$$
(46)
with $`𝒩(x)`$ any finite $`x`$-dependent set of lattice sites, the contribution of $`\mathrm{\Gamma }`$ is a multiple of $`(2K)^L`$. Often used special cases are the nearest neighbour interactions
$$v(x,y)=\mathrm{\hspace{0.33em}2}K\underset{\mu =0}{\overset{D1}{}}\left(\delta _{x,y+\widehat{\mu }}+\delta _{x,y\widehat{\mu }}\right)$$
(47)
and the uniform interaction
$$v(x,y)=\mathrm{\hspace{0.33em}2}K\left(1\delta _{x,y}\right),$$
(48)
which is used in models of spin glasses and partially annealed neural networks. Susceptibilities of the $`\varphi `$ and $`U`$ fields will be represented as
$`\chi _{E_1,E_2}`$ $`=`$ $`{\displaystyle \frac{1}{VD}}{\displaystyle \underset{x_1,\mathrm{},x_{E_1}\mathrm{\Lambda }_0}{}}{\displaystyle \underset{l_1,\mathrm{},l_{E_2}\overline{\mathrm{\Lambda }}_1}{}}<\varphi (x_1)\mathrm{}\varphi (x_{E_1})U(l_1)\mathrm{}U(l_{E_2})>^c`$ (49)
$``$ $`{\displaystyle \frac{1}{VD}}{\displaystyle \underset{x_1,\mathrm{},x_{E_1}\mathrm{\Lambda }_0}{}}{\displaystyle \underset{l_1,\mathrm{},l_{E_2}\overline{\mathrm{\Lambda }}_1}{}}{\displaystyle \frac{^{E_1+E_2}W(H,I,v)}{H(x_1)\mathrm{}H(x_{E_1})I(l_1)\mathrm{}I(l_{E_2})}}|_{H=I=0}`$
$`=`$ $`{\displaystyle \underset{L0}{}}(2K)^L{\displaystyle \underset{\mathrm{\Gamma }𝒢_{E_1,E_2}(L)}{}}w(\mathrm{\Gamma })`$
with lattice volume $`V`$ and dimension $`D`$. Similar representations hold for higher moments $`\mu `$.
The weight $`w(\mathrm{\Gamma })`$ of a multiple-line graph $`\mathrm{\Gamma }𝒢_{E_1,E_2}(L)`$ is given as the product of the following factors
* for every vertex $`v_\mathrm{\Gamma }`$ a factor
$$v_n^c=\left(\frac{^nW(H)^{}}{H^n}\right)_{H=0},$$
(50)
where $`n=t__\mathrm{\Gamma }(v)`$ is the total number of bare lines attached to $`v`$.
* for every multiple line $`m_\mathrm{\Gamma }`$ a factor
$$m_\nu ^{1c}=\left(\frac{^\nu W^1(I)}{I^\nu }\right)_{I=0},$$
(51)
where $`\nu =l__\mathrm{\Gamma }(m)+E_\mathrm{\Gamma }^{(U)}(m)`$, that is $`m`$ is a $`\nu `$-line,
* a factor $`1/S_\mathrm{\Gamma }`$, where $`S_\mathrm{\Gamma }`$ is the topological symmetry number of $`\mathrm{\Gamma }`$,
* a factor counting the permutation symmetry of external $`\varphi `$-lines,
$$\frac{E_\mathrm{\Gamma }^{(\varphi )}!}{_{v_\mathrm{\Gamma }}E_\mathrm{\Gamma }^{(\varphi )}(v)!},$$
(52)
* a factor counting the permutation symmetry of external $`U`$-lines,
$$\frac{E_\mathrm{\Gamma }^{(U)}!}{_{m_\mathrm{\Gamma }}E_\mathrm{\Gamma }^{(U)}(m)!},$$
(53)
* the lattice embedding number of $`\mathrm{\Gamma }`$, which is the number of ways $`\mathrm{\Gamma }`$ can be embedded on a lattice of given geometry, e. g. on a hypercubic lattice. To this end, the vertices of $`\mathrm{\Gamma }`$ (if any) are placed onto lattice sites. One arbitrary vertex is placed at a fixed lattice site, in order to account for the volume factor $`1/V`$ in (49). A priori there is no exclusion principle. This means that any number of vertices can be placed at the same lattice site. (This is sometimes called free embedding.) Two restrictions apply to the embeddings. The first constraint results from the fact that a bare internal line represents a hopping propagator $`v(x,y)`$, with lattice sites $`x`$ and $`y`$ at which the two endpoint vertices of the line are placed at. A reasonable computation of the embedding number takes into account the particular form of $`v(x,y)`$ from the very beginning. The second constraint is that bare lines of the same multiple-line have to be mapped on the same pair of sites.
For example, if $`v(x,y)`$ is the nearest neighbour interaction (47), two vertices which have at least one line in common are to be placed at nearest neighbour lattice sites. On the other hand, a propagator $`v(x,y)`$ of the form (48) implies a rather weak constraint in that $`x`$ and $`y`$ must be different, but otherwise can be freely placed over the lattice.
We remark that in case of a non-trivial internal symmetry (such as considered in section 7) the expressions of Eq.s (49)-(51) must be modified appropriately. In particular, the weight (51) of a multiple-line does no longer take such a simple form.
## 4 Renormalization
Truncated correlation functions, susceptibilities and other moments are obtained as sums over multiple-line graphs that are connected. Their number rapidly grows with increasing order, that is with increasing number of bare internal lines. The procedure of ”renormalization” means that the connected moments are represented in terms of reduced ones. The reduced moments are obtained by summation over multiple-line graph classes which are more restricted than just by their property of being connected. Of course the number of graphs of such classes is smaller. Only the most restricted multiple-line graph classes must be constructed. The subsequent steps towards the moment computation are most conveniently done by operating analytically with the reduced moments. In particular, it is no longer necessary to generate all connected and the corresponding intermediate multiple-line graph classes.
A connected multiple-line graph $`\mathrm{\Gamma }`$ is called 1-particle irreducible (1PI) if it satisfies the following condition. Remove an arbitrary $`1`$-line of $`\mathrm{\Gamma }`$. There is at most one connected component left that has external lines attached. (This notion is the same as the one used in the context of Feynman graphs.) On the other hand, if in addition the remaining graph is still connected, then $`\mathrm{\Gamma }`$ is called 1-line irreducible (1LI). In many cases it is sufficient to use only the second notion. It is for instance sufficient that all vertices are constrained to have only an even number of lines attached, or more generally, if graphs and subgraphs with one external line are forbidden. For notational simplicity we assume in the following that this is the case and henceforth refer only to the notion 1LI <sup>4</sup><sup>4</sup>4In ref. the term 1PI was used instead.. The generalization to the case in which 1LI and 1PI graphs must be distinguished goes along the same lines as for LCEs, which was discussed in .
By $`𝒢_{E_1,E_2}^{1LI}(L)`$ we denote the subset of multiple-line graphs $`\mathrm{\Gamma }𝒢_{E_1,E_2}(L)`$ that are 1LI. 1LI-susceptibilities are defined as series in the hopping parameter similarly as in (49) by restricting the summation to 1LI graphs,
$$\chi _{E_1,E_2}^{1LI}=\underset{L0}{}(2K)^L\underset{\mathrm{\Gamma }𝒢_{E_1,E_2}^{1LI}(L)}{}w(\mathrm{\Gamma }).$$
(54)
Susceptibilities are easily obtained in a closed form in terms of 1LI-susceptibilities $`\chi ^{1LI}`$. It can be shown that the $`\chi ^{1LI}`$s can be obtained by an appropriate Legendre transform. For instance
$`\chi _{2,0}`$ $`=`$ $`{\displaystyle \frac{\chi _{2,0}^{1LI}}{1\stackrel{~}{v}(0)\chi _{2,0}^{1LI}}},`$
$`\chi _{2,1}`$ $`=`$ $`{\displaystyle \frac{\chi _{2,1}^{1LI}}{(1\stackrel{~}{v}(0)\chi _{2,0}^{1LI})^2}},`$ (55)
where $`\stackrel{~}{v}(k)`$ is the Fourier transform of the hopping propagator $`v(x,y)`$,
$$v(x,y)=_\pi ^\pi \frac{d^Dk}{(2\pi )^D}e^{ik(xy)}\stackrel{~}{v}(k).$$
(56)
In LCEs the second important resummation comes from so called vertex renormalizations. This means partial resummation of graphs with specific properties such as having one external vertex only. These sums then are considered as ”renormalized vertices” replacing the vertices of graphs with complementary properties. The procedure naturally leads to the notion of 1-vertex irreducibility (1VI) and renormalized moments.
In DLCE we follow this procedure. The very definition of 1VI has to be modified slightly for multiple-line graphs because of the enhanced connectivity properties due to multiple-lines. In addition, as a natural generalization, we supplement vertex renormalization by multiple-line renormalization.
A multiple-line graph $`\mathrm{\Gamma }`$ is called 1-vertex irreducible (1VI) if it satisfies the following condition. Decompose an arbitrary vertex $`v_\mathrm{\Gamma }`$. Every connected component of the remaining graph has then at least one external line attached. It can be a $`\varphi `$-line or a $`U`$-line. We write
$$𝒢_{E_1,E_2}^{1VI}(L)=\{\mathrm{\Gamma }𝒢_{E_1,E_2}^{1LI}(L)|\mathrm{\Gamma }\text{ is 1VI}\}$$
(57)
for the set of equivalence classes of graphs that are both 1LI and 1VI, with $`E_1`$ external $`\varphi `$-lines, $`E_2`$ external $`U`$-lines and $`L`$ bare internal lines.
The renormalized vertex moment graphs are 1LI graphs that have precisely one external vertex and no external multiple line,
$$Q_k(L)=\{\mathrm{\Gamma }𝒢_{k,0}^{1LI}(L)|\text{there is }v_\mathrm{\Gamma }\text{ with }E_\mathrm{\Gamma }^{(\varphi )}(v)=k\}.$$
(58)
A multiple-line graph $`\mathrm{\Gamma }`$ is called 1-multiple-line irreducible (1MLI) if it satisfies the following criterion. Decompose an arbitrary multiple-line $`m_\mathrm{\Gamma }`$. Every remaining connected component has then at least one external line attached. It can be a $`\varphi `$-line or a $`U`$-line. We write
$$𝒢_{E_1,E_2}^{1MLI}(L)=\{\mathrm{\Gamma }𝒢_{E_1,E_2}^{1LI}(L)|\mathrm{\Gamma }\text{ is 1MLI}\}.$$
(59)
The renormalized multiple-line moment graphs are graphs that are 1LI and have precisely one external multiple-line, but no external vertex,
$$R_k(L)=\{\mathrm{\Gamma }𝒢_{0,k}^{1LI}(L)|\text{there is }m_\mathrm{\Gamma }\text{ with }E_\mathrm{\Gamma }^{(U)}(m)=k\}.$$
(60)
The equivalence classes of graphs that are both 1VI and 1MLI are denoted by
$$S_{E_1,E_2}(L)=𝒢_{E_1,E_2}^{1VI}(L)𝒢_{E_1,E_2}^{1MLI}(L).$$
(61)
With the renormalized moment graphs as defined above, the 1LI-susceptibilities are now obtained in the form
$$\chi _{E_1,E_2}^{1LI}=\underset{L0}{}(2K)^L\underset{\mathrm{\Gamma }𝒮_{E_1,E_2}(L)}{}\stackrel{~}{w}(\mathrm{\Gamma }).$$
(62)
The weights $`\stackrel{~}{w}(\mathrm{\Gamma })`$ are given as a product of factors as described in the last subsection, with the following two exceptions.
* The vertex coupling constants $`v_n^c`$ are replaced by the renormalized vertex moments
$$v_n^cv_n^c=\underset{L0}{}(2K)^L\underset{\mathrm{\Gamma }Q_n(L)}{}w(\mathrm{\Gamma }).$$
(63)
* The multiple line coupling constants $`m_\nu ^{1c}`$ are replaced by the renormalized multiple line moments
$$m_\nu ^{1c}m_\nu ^c=\underset{L0}{}(2K)^L\underset{\mathrm{\Gamma }R_\nu (L)}{}w(\mathrm{\Gamma }).$$
(64)
In the series representations above, the $`w(\mathrm{\Gamma })`$ are computed according to the rules of subsection 3.2.
In we have described an algorithmic construction of graphs that is the first step for an automatic computer aided generation. In exceptional cases DLCEs can be summed up in a closed form. Otherwise, a computer aided generation of graphs is unavoidable, because the proliferation of DLCE graphs is pronounced even compared to LCE graphs. For LCE graphs we know that billions of graphs must be included to obtain the critical exponents with the accuracy of some per mil.
## 5 Applications to spin glasses
In this section we consider applications of DLCEs to disordered systems, in particular to spin glasses with ”‘slow”’ interactions coupled to ”‘fast”’ spins. The interactions $`J`$ are assumed to be in equilibrium with a thermal heat bath of inverse temperature $`\beta ^{}`$, while the spins $`\sigma `$ are equilibrated according to a second inverse temperature $`\beta `$. Both systems need not be mutually in equilibrium. Let $`Z_\beta (J)`$ be the partition function that describes the equilibrium distribution of the spins for given $`J`$s,
$$Z_\beta (J)=\underset{\{\sigma _i=\pm 1\}}{}\mathrm{exp}(\beta \underset{i<j}{}J_{(i,j)}\sigma _i\sigma _j).$$
(65)
The sum runs over pairs $`(i,j)`$ that need not be restricted to nearest neighbours only.
We further assume that the dynamics of the time evolution of the slow interactions $`J`$ is governed by a Langevin equation
$$N\frac{d}{dt}J_{(i,j)}=\frac{}{J_{(i,j)}}(J)+\sqrt{N}\eta _{ij}(t)$$
(66)
with
$$(J)=\frac{1}{\beta }\mathrm{ln}Z_\beta (J)+\frac{1}{2}\mu N\underset{i<j=1}{\overset{N}{}}J_{(i,j)}^2.$$
(67)
Here $`Z_\beta (J)`$ is given by (65), $`N`$ is the total number of spins, $`\mu `$ is a positive constant and $`\eta _{ij}`$ is a stochastic gaussian white noise of zero mean and correlation
$$<\eta _{ij}(t)\eta _{kl}(t^{})>=\frac{2}{\beta ^{}}\delta _{(ij),(kl)}\delta (tt^{}).$$
(68)
Such a Langevin equation for the $`J`$s can be derived from an ansatz which is motivated by neural networks -. Moreover, since the time evolution of the $`J`$s is determined by a dissipative Langevin equation, the equilibrium distribution of the slow variables is again a Boltzmann distribution, governed now by the second temperature $`\beta _{}^{}{}_{}{}^{1}`$,
$$Z_\beta ^{}^{}=𝒩_{\mathrm{}}^{\mathrm{}}\underset{i<j=1}{\overset{N}{}}dJ_{(i,j)}\mathrm{exp}(\beta ^{}(J)),$$
(69)
with $`𝒩`$ some normalization that will be specified below. The effective Hamiltonian $``$ of $`J`$ is given by (67).
It is these equilibrium aspects of coupled systems of fast spins and slow interactions that we can treat analytically with DLCEs, as we will show below.
Let us first rewrite $`Z_\beta ^{}^{}`$ in the form
$`Z_{x\beta }^{}`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \underset{i<j=1}{\overset{N}{}}}\left(\sqrt{{\displaystyle \frac{QN}{2\pi }}}dJ_{(i,j)}\right)\mathrm{exp}({\displaystyle \frac{1}{2}}QN{\displaystyle \underset{i<j}{}}J_{(i,j)}^2)Z_\beta (J)^x`$ (70)
$``$ $`[[Z_\beta (J)^x]],`$
where we have introduced $`Q=\beta \mu `$ and real $`x=\beta ^{}/\beta `$ as the ratio of two temperatures. The normalization has been chosen such that $`[[1]]=1`$. In the limit of $`x0`$ for fixed $`\beta `$, i.e. $`\beta ^{}0`$, we have a quenched system. The $`J`$ only feel the infinitely high temperature $`\beta _{}^{}{}_{}{}^{1}`$, but are decoupled from the spins. The time scale of fluctuations of the spins is assumed to be so short that the $`J`$ are only sensitive to averages of the $`\sigma `$. Vice versa, the spin dynamics does depend on the $`J`$s. Therefore the quantity of physical interest is not
$$\mathrm{ln}\left[_{\mathrm{}}^{\mathrm{}}\underset{i<j=1}{\overset{N}{}}\left(\sqrt{\frac{QN}{2\pi }}dJ_{(i,j)}\right)\left(\underset{\{\sigma _i\}}{}\mathrm{exp}(\beta \underset{i<j=1}{\overset{N}{}}J_{(i,j)}\sigma _i\sigma _j)\right)^x\mathrm{exp}(\frac{1}{2}QN\underset{i<j}{}J_{(i,j)}^2)\right]$$
(71)
where fluctuations of the $`\sigma `$s do influence the $`J`$s, but
$$_{\mathrm{}}^{\mathrm{}}\underset{i<j=1}{\overset{N}{}}\left(\sqrt{\frac{QN}{2\pi }}dJ_{(i,j)}\right)\mathrm{ln}\left[\underset{\{\sigma _i\}}{}\mathrm{exp}(\beta \underset{i<j=1}{\overset{N}{}}J_{(i,j)}\sigma _i\sigma _j)\right]\mathrm{exp}(\frac{1}{2}QN\underset{i<j}{}J_{(i,j)}^2),$$
(72)
or, in a shorthand notation,
$$𝒟J\mathrm{ln}Z_\beta (J)[[\mathrm{ln}Z_\beta (J)]].$$
(73)
Usually one rewrites
$`{\displaystyle 𝒟J\mathrm{ln}Z_\beta (J)}={\displaystyle 𝒟J\underset{x0}{lim}\frac{Z_\beta (J)^x1}{x}}`$
$`=\underset{x0}{lim}{\displaystyle 𝒟J\frac{Z_\beta (J)^x1}{x}}=\underset{x0}{lim}{\displaystyle \frac{\mathrm{ln}\{1+([[Z_\beta (J)^x]]1)\}}{x}}`$
$`=\underset{x0}{lim}{\displaystyle \frac{\mathrm{ln}Z_{x\beta }^{}}{x}}.`$ (74)
For the second equality sign one has assumed that $`𝒟J`$ commutes with $`lim_{x0}`$, in the third one that $`lim_{x0}[[Z_\beta (J)^x]]=1`$. So far, $`x`$ as the ratio of two temperatures is real. Rewriting the left hand side of (5) according to the right hand side is called the replica trick . The uncontrolled approximation that usually enters the replica trick is that now the right hand side is evaluated for positive integer $`xn`$ and extrapolated to $`n=0`$. Clearly a function that is known only for positive integer $`n`$ does not have a unique extrapolation to $`n=0`$ without further assumptions. Nevertheless, this approximation is made, because it is rather convenient. For integer $`n`$, $`Z_\beta (J)^n`$ is the partition function of an $`n`$ times replicated system of which the logarithm is taken after the integration over the $`J`$s. It is seen as follows. We rewrite
$$Z_\beta (J)^n=\underset{\{\sigma _i^{(a)}\}}{}\mathrm{exp}(\beta \underset{a=1}{\overset{n}{}}\underset{i<j=1}{\overset{N}{}}J_{(i,j)}\sigma _i^{(a)}\sigma _j^{(a)}),$$
(75)
with $`a=1,\mathrm{},n`$ labelling the replicated spin variables, so that
$`Z_{n\beta }^{}`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \underset{i<j=1}{\overset{N}{}}}dJ_{(i,j)}{\displaystyle \underset{\{\sigma _i^{(a)}=\pm 1\}}{}}\mathrm{exp}(S(J,\sigma ^{(a)})),`$
$`S(J,\sigma ^{(a)})`$ $`=`$ $`\beta {\displaystyle \underset{a=1}{\overset{n}{}}}{\displaystyle \underset{i<j=1}{\overset{N}{}}}J_{(i,j)}\sigma _i^{(a)}\sigma _j^{(a)}+{\displaystyle \frac{1}{2}}QN{\displaystyle \underset{i<j}{}}J_{(i,j)}^2.`$ (76)
Linear terms in $`\sigma `$ and $`J`$ may be included according to
$$S_{lin}=h\underset{a=1}{\overset{n}{}}\underset{i=1}{\overset{N}{}}\sigma _i^{(a)}+c\underset{i<j=1}{\overset{N}{}}J_{(i,j)}$$
(77)
with constant external fields $`h`$ and $`c`$.
Apparently, because of integer $`n`$, $`Z_{n\beta }^{}`$ has the form of models to which DLCE applies, with a hopping term
$$S_{hop}(J,\sigma ^{(a)})=\beta \underset{a=1}{\overset{n}{}}\underset{i<j=1}{\overset{N}{}}J_{(i,j)}\sigma _i^{(a)}\sigma _j^{(a)},$$
(78)
a single link action
$$S^1(J_{(i,j)})=cJ_{(i,j)}+\frac{1}{2}QNJ_{(i,j)}^2,$$
(79)
and a single site action
$$S^{}(\sigma _i^{(a)})=h\underset{a=1}{\overset{n}{}}\sigma _i^{(a)}.$$
(80)
Depending on $`n`$ we distinguish the following cases.
* $`n=1`$. First we note that for $`n=1`$ we can directly apply DLCE to $`\mathrm{ln}Z_{\beta ^{}=\beta }^{}`$ and to derived quantities to obtain their series expansions in $`\beta `$. But from a physical point of view, in a disordered system one is not interested in $`n=1`$, because $`n=1`$ corresponds to the completely annealed situation, in which the fast spins and the slow interactions are in mutual equilibrium. (In contrast, in particle physics one is interested in the $`n=1`$ case, cf. our applications of DLCEs in the framework of variational cumulant expansions of the $`SU(2)`$ Higgs model .)
* $`n>1`$, integer. Again we apply DLCE to $`\mathrm{ln}Z_{\beta ^{}=n\beta }^{}`$, but have to account for the permutation symmetry between the replicas. Formally, the replica symmetry plays a role similar to an internal symmetry, e.g. an $`O(N)`$ symmetry in a scalar Higgs model. DLCEs with nontrivial internal symmetries have been discussed in connection with the $`SU(2)`$ Higgs model . Thus we can study ”‘unquenched”’ equilibrium aspects of systems with two temperatures and compare the results from DLCEs adapted to ”internal” replica symmetry with Monte Carlo simulations for the same $`n`$ .
* $`n=0`$, the quenched limit. As we will show in the next section, in order to discuss the $`x0`$ limit, we need not refer to $`n`$ times replicated systems $`Z_\beta (J)^n`$ characterized by (78)-(80), but just to $`Z_\beta (J)`$ given by (78)-(80) with $`n=1`$. By means of special DLCEs , so-called quenched DLCEs (QDLCEs), we directly calculate the left hand side of (5). Therefore, setting $`n=1`$ in (78)-(80) in QDLCEs does not imply the completely annealed case, because we first take the logarithm of $`Z_\beta (J)`$ and then average over the $`J`$s.
### 5.1 Avoiding the replica trick
First we adapt the notation to section 2 to include more general cases. $`\mathrm{\Lambda }_0`$ denotes the support of the spins, that is the set of lattice sites, with $`V=|\mathrm{\Lambda }_0|`$ denoting their total number. $`\overline{\mathrm{\Lambda }}_1`$ $`\mathrm{\Lambda }_1`$ are the pairs of sites whose spins interact. In accordance with (70), we write for the normalized link-average of a function $`f(J)`$
$$[[f(J)]]=𝒟Jf(J)$$
(81)
with
$`𝒟J`$ $`=`$ $`{\displaystyle \underset{l\overline{\mathrm{\Lambda }}_1}{}}d\mu (J(l)),`$
$`d\mu (J)`$ $`=`$ $`𝒩_1dJ\mathrm{exp}(S^1(J)),{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\mu (J)=1.`$ (82)
It is convenient to introduce the single link expectation values
$$<g(J)>_1𝑑\mu (J)g(J)$$
(83)
and the generating function $`W^1(I)`$ by
$$\mathrm{exp}W^1(I)<\mathrm{exp}(IJ)>_1.$$
(84)
The way in which the replica trick can be avoided is examplified for the free energy density $`W_{sp}/V`$ of the spin system averaged over the link couplings. The partition function of the spin system for a given distribution of the link interactions $`J(x,y)`$ is given by
$$\mathrm{exp}W_{sp}(J)=𝒩_{sp}𝒟\sigma \mathrm{exp}(S_{sp}(\sigma ,J)),$$
(85)
where $`W_{sp}(0)=0`$ and
$`S_{sp}(\sigma ,J)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{x,y\mathrm{\Lambda }_0}{}}v(x,y)\sigma (x)\sigma (y)J(x,y),`$
$`𝒟\sigma `$ $`=`$ $`{\displaystyle \underset{x\mathrm{\Lambda }_0}{}}d\sigma (x)\mathrm{exp}(S^{}(\sigma (x))).`$ (86)
Without loss of generality we identify the support of the interaction $`v(x,y)=v(y,x)`$ with the set $`\overline{\mathrm{\Lambda }}_1`$ of lattice sites where $`𝒟J`$ is supported,
$$\overline{\mathrm{\Lambda }}_1=\{l=(x,y)\overline{\mathrm{\Lambda }_0\times \mathrm{\Lambda }_0}|v(x,y)0\}.$$
(87)
For simplicity we assume $`v(x,y)`$ to be of the form (46), so that $`K`$ is a measure of the strength of the interactions $`v(x,y)`$.
The free energy density of the spin system allows for a series expansion in the standard LCE sense, with the link field $`J(l)`$ playing the role of a ”background field”,
$$\frac{1}{V}W_{sp}(J)=\underset{L0}{}(2K)^L\underset{\mathrm{\Gamma }𝒢_0^{sp}(L)}{}w^{sp}(\mathrm{\Gamma },J).$$
(88)
Here $`𝒢_E^{sp}(L)`$ (with $`E=0`$) denotes the set of equivalence classes of connected LCE graphs with $`E`$ external lines and $`L`$ internal lines. The spin-weights $`w^{sp}(\mathrm{\Gamma },J)`$ are of the form
$$w^{sp}(\mathrm{\Gamma },J)=R^{sp}(\mathrm{\Gamma })\underset{_\mathrm{\Gamma }\overline{\mathrm{\Lambda }}_1}{\overset{}{}}\underset{l\overline{\mathrm{\Lambda }}_1}{}J(l)^{m(l)}.$$
(89)
The sum is taken over all non-vanishing lattice embeddings of the graph $`\mathrm{\Gamma }`$. It runs over all maps of internal lines of the graph $`\mathrm{\Gamma }`$ to pairs of lattice sites of $`\overline{\mathrm{\Lambda }}_1`$ that are consistent with the graph topology in the sense discussed in section 3. For every $`l\overline{\mathrm{\Lambda }}_1`$, $`m(l)`$ denotes the number of lines of $`\mathrm{\Gamma }`$ that are mapped onto the link $`l`$ by the embedding. All other factors that contribute to the weight are collected in the prefactor $`R^{sp}(\mathrm{\Gamma })`$, including the inverse topological symmetry number of $`\mathrm{\Gamma }`$.
Next we want to express $`[[W_{sp}(J)]]`$ as a series in $`K`$ by means of DLCE. Toward this end we set $`f(J)=W_{sp}(J)`$ and insert the series (88) with (89) into (81). At this stage we are not concerned with question of (uniform or dominated) convergence and obtain
$`[[{\displaystyle \frac{1}{V}}W_{sp}(J)]]`$ $`=`$ $`{\displaystyle \underset{L0}{}}(2K)^L{\displaystyle \underset{\mathrm{\Gamma }𝒢_0^{sp}(L)}{}}{\displaystyle 𝒟Jw^{sp}(\mathrm{\Gamma },J)}`$ (90)
$`=`$ $`{\displaystyle \underset{L0}{}}(2K)^L{\displaystyle \underset{\mathrm{\Gamma }𝒢_0^{sp}(L)}{}}R^{sp}(\mathrm{\Gamma }){\displaystyle \underset{_\mathrm{\Gamma }\overline{\mathrm{\Lambda }}_1}{\overset{}{}}}{\displaystyle \underset{l\overline{\mathrm{\Lambda }}_1}{}}<J(l)^{m(l)}>_1.`$
The next step is to express the full single link expectation values in terms of the connected ones. They are related by
$$<J^m>_1=\underset{\mathrm{\Pi }𝒫(\underset{¯}{m})}{}\underset{P\mathrm{\Pi }}{}<J^{|P|}>_1^c,$$
(91)
where $`𝒫(\underset{¯}{m})`$ denotes the set of all partitions of $`\underset{¯}{m}=\{1,\mathrm{},m\}`$ into non-empty, mutually disjoint subsets of $`\underset{¯}{m}`$. $`|P|`$ is the number of elements of the set $`P`$. The relation (91) is equivalent to the partition of all lines of $`\mathrm{\Gamma }`$ that are mapped to the same lattice link into multiple lines, with every multiple line contributing a factor
$$<J^{|P|}>_1^c=\frac{^{|P|}W^1(I)}{I^{|P|}}|_{I=0}=m_{|P|}^{1c}.$$
(92)
Using (91), (92) we rewrite (90) as
$$[[\frac{1}{V}W_{sp}(J)]]=\underset{L0}{}(2K)^L\underset{\mathrm{\Gamma }𝒢_0^{sp}(L)}{}R^{sp}(\mathrm{\Gamma })\underset{\mathrm{\Pi }𝒫(_\mathrm{\Gamma })}{}\left(\underset{P\mathrm{\Pi }}{}m_{|P|}^{1c}\right)\left(\underset{\mathrm{\Pi }\overline{\mathrm{\Lambda }}_1}{\overset{}{}}\underset{l\overline{\mathrm{\Lambda }}_1}{}1\right).$$
(93)
The last summation in (93) is over all maps $`_\mathrm{\Gamma }\overline{\mathrm{\Lambda }}_1`$ of the lines of $`\mathrm{\Gamma }`$ to the lattice links of $`\overline{\mathrm{\Lambda }}_1`$ subject to the constraint that all lines that belong to the same multiple-line corresponding to some $`P\mathrm{\Pi }`$ are mapped onto the same lattice link.
Finally we rewrite (93) as a sum over multiple-line graphs . To this end, we first observe that for every $`\mathrm{\Gamma }𝒢_0^{sp}(L)`$, every partition $`\mathrm{\Pi }𝒫(_\mathrm{\Gamma })`$ of the lines of $`\mathrm{\Gamma }`$ into multiple-lines generates a multiple-line graph $`\mathrm{\Delta }=(\mathrm{\Gamma },\mathrm{\Pi })`$ in the obvious way. Let us denote by $`\overline{𝒢}_{0,0}(L)`$ the subset of multiple-line graphs of $`𝒢_{0,0}(L)`$ that stay connected after decomposition of all multiple lines. (These are the multiple-line graphs which stay connected in the usual graph theoretical sense, when the dashed lines are omitted.) For every $`\mathrm{\Delta }\overline{𝒢}_{0,0}(L)`$ there is a unique $`\mathrm{\Gamma }(\mathrm{\Delta })𝒢_0^{sp}(L)`$ and at least one $`\mathrm{\Pi }𝒫(_{\mathrm{\Gamma }(\mathrm{\Delta })})`$ such that $`(\mathrm{\Gamma }(\mathrm{\Delta }),\mathrm{\Pi })=\mathrm{\Delta }`$. Let $`n_\mathrm{\Delta }`$ be the (uniquely determined) number of partitions $`\mathrm{\Pi }𝒫(_{\mathrm{\Gamma }(\mathrm{\Delta })})`$ with $`(\mathrm{\Gamma }(\mathrm{\Delta }),\mathrm{\Pi })=\mathrm{\Delta }`$, and $`\mathrm{\Pi }(\mathrm{\Delta })`$ such an arbitrary partition. Eq. (93) then becomes
$$[[\frac{1}{V}W_{sp}(J)]]=\underset{L0}{}(2K)^L\underset{\mathrm{\Delta }\overline{𝒢}_{0,0}(L)}{}n_\mathrm{\Delta }R^{sp}(\mathrm{\Gamma }(\mathrm{\Delta }))\left(\underset{P\mathrm{\Pi }(\mathrm{\Delta })}{}m_{|P|}^{1c}\right)\left(\underset{\mathrm{\Pi }(\mathrm{\Delta })\overline{\mathrm{\Lambda }}_1}{\overset{}{}}1\right).$$
(94)
The last bracket of (94) is the lattice embedding factor of the multiple-line graph $`\mathrm{\Delta }`$. The second bracket from the right does not depend on the choice of $`\mathrm{\Pi }(\mathrm{\Delta })`$ and is the product of the multiple-line coupling constants as defined in section 3. Finally, $`n_\mathrm{\Delta }R^{sp}(\mathrm{\Gamma }(\mathrm{\Delta }))`$ is precisely the remaining part of the weight of $`\mathrm{\Delta }`$ that was described in detail in section 3, endowed with the correct inverse topological symmetry number of the multiple-line graph $`\mathrm{\Delta }`$ (because of the factor $`n_\mathrm{\Delta }`$).
In summary, we obtain the series expansion of the link-averaged free energy density in terms of DLCE graphs,
$$[[\frac{1}{V}W_{sp}(J)]]=\underset{L0}{}(2K)^L\underset{\mathrm{\Delta }\overline{𝒢}_{0,0}(L)}{}w(\mathrm{\Delta }).$$
(95)
The weight $`w(\mathrm{\Delta })`$ of a multiple-line graph $`\mathrm{\Delta }`$ is defined and computed according to the rules given in section 3.
Eq. (95) is the series representation of the link-averaged free energy density of the spin system, i.e. the free energy density of the $`n=0`$ replica system, in terms of DLCE graphs. It looks much like the series representation of the $`1`$-replica system, which is given by
$$\frac{1}{V}W_{1repl}\frac{1}{V}\mathrm{ln}[[\mathrm{exp}W_{sp}(J)]]=\underset{L0}{}(2K)^L\underset{\mathrm{\Delta }𝒢_{0,0}(L)}{}w(\mathrm{\Delta })$$
(96)
according to the discussion of section 2. We recall that $`𝒢_{0,0}(L)`$ is the set of DLCE vacuum graphs with $`L`$ bare lines that are connected in the generalized DLCE sense. Comparing (95) and (96), the transition from $`n=1`$ to $`n=0`$ replicas is achieved by keeping only the subset $`\overline{𝒢}_{0,0}(L)𝒢_{0,0}(L)`$ of multiple-line graphs that are connected in the original (LCE) sense.
We emphasize that the restriction of DLCEs to QDLCEs is not an ad hoc (or intuitively motivated) assumption but a derived consequence of the fact that the logarithm is taken before the integration $`𝒟J`$. This procedure accounts for all graphs that contribute to a given order in $`K`$. Thus we do have to truncate the series unless the series can be completely summed up, as it happens in exceptional cases.
We expect that the series (95) are convergent for a large class of interactions $`S^1(J)`$ and $`v(x,y)`$ if the coupling constant $`K`$ is sufficiently small. For special interactions most of the multiple-line graphs yield vanishing contributions so that we can further restrict the sum to a subset of $`\overline{𝒢}_{0,0}(L)`$. An example is given by the mean field type of interaction of the Sherrington-Kirkpatrick model, cf. Sect. 5.2.
### 5.2 Applications of QDLCEs
In the following we list some examples for systems of which we can study the phase structure by means of QDLCEs. Their actions are special cases of (78)-(80) with $`n=1`$ (as explained above) and the following choice of variables.
* Infinite range models. Choose $`J_{(i,j)}`$ as before, $`\sigma _i\pm 1`$, $`i,j\{1,\mathrm{},N\}`$,
$`v(x,y)`$ $`=`$ $`K(1\delta _{x,y}),`$
$`\overline{\mathrm{\Lambda }}_1`$ $``$ $`\mathrm{\Lambda }_1\text{is the set of all pairs of sites,}`$ (97)
$`S^1(J)`$ $`=`$ $`N{\displaystyle \frac{1}{2}}J^2.`$
Now the sum over the sites in (78) runs over arbitrary pairs $`(i,j)`$, $`i<j`$, and we obtain the infinite range Sherrington-Kirkpatrick model. For infinte range and in the thermodynamic limit ($`N\mathrm{}`$), the phase structure can be solved by replica mean field theory, cf. e.g. . For QDLCEs the infinite range and $`N\mathrm{}`$ limits imply that only tree graphs of 2-lines contribute to the series of the free energy density, such as
(98)
The reason for that is that each 2-line gets a factor of $`1/N`$ from the $`S_1`$-part of the action, but each vertex gets a factor $`(N1)`$ from the embedding onto a lattice $`\mathrm{\Lambda }_0`$. The contribution of every tree graph to the free energy is proportional to $`N`$. If a chain of 2-lines connecting the vertices gets closed, forming a loop, there is one $`(N1)`$ less in the total embedding factor. Thus the contribution is suppressed by $`1/(N1)`$ for every loop and vanishes in the thermodynamic limit. Because of the simple tree structure there is a chance for summing up the series. This is currently under investigation.
* Finite range connectivity. The sum $`_{i<j}`$ of the spins is now restricted to next-neighbours or, more generally, to a finite number of pairs. Rather than specifying $`S_1(J_{(i,j)})`$ of (79), it is sufficient for DLCEs and QDLCEs to choose $`\mathrm{exp}(S_1(J_{(i,j)}))`$. Let
$$exp(S_1(J))=(1p)\delta (J)+p\delta (J1)$$
(99)
with $`p[0,1]`$. The variables $`J_{(i,j)}\{0,1\}`$ can then be interpreted as occupation numbers of the bonds. Furthermore, if we choose $`\sigma _i\{\pm 1\}`$ we obtain a
+ bond-diluted Ising model.
Choosing $`\sigma _i_q`$ we obtain a
+ bond-diluted $`q`$-state Potts model.
If $`\sigma _iS_q`$, we obtain a
+ bond-diluted Heisenberg model.
QDLCEs provide a systematic analytic expansion for disordered systems with bond dilution in a quenched limit. Coming from the high temperature (small $`\beta `$) region one can study the phase structure as a function of the degree of dilution. Work in this direction is in progress.
## 6 Summary and Conclusions
In this paper we have introduced a new expansion scheme for 3-point interactions or, more precisely, for point-link-point interactions. This scheme generalizes linked cluster expansions for 2-point interactions by including hopping parameter terms endowed with their own dynamics. In chapters 3-4 we have developed a multiple-line graph theory with an additional new type of multiple-line connectivity. We have introduced appropriate equivalence classes of graphs and discussed the issue of renormalization. These notions are required for an algorithmic generation of graphs. Because of the fast proliferation of graphs already at low orders in the expansion, a computer aided implementation becomes unavoidable, if one is interested in higher orders of the expansion than we have computed so far.
In Sect. 5 we have shown how to avoid the replica trick for calculating the free energy of disordered systems in the quenched limit. DLCEs are a systematic expansion method to study the phase structure of disordered systems. It is systematic in the sense that we do not restrict the expansion to certain subclasses of graphs that can be summed up, but we identify and keep all graphs that contribute to a given order in the expansion parameter. DLCEs provide an analytic tool for studying systems in situations in which it has been impossible so far.
## Acknowledgment
We would like to thank Reimar Kühn (Heidelberg) for discussions.
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# \quad1.3.1. Theorem
MIRROR SYMMETRY AND QUANTIZATION
OF ABELIAN VARIETIES
Yu. I. Manin
Max–Planck–Institut für Mathematik, Bonn
0. Introduction
0. Plan of the paper. This paper consists of two sections discussing various aspects of commutative and non–commutative geometry of tori and abelian varieties.
In the first section, we present a new definition of mirror symmetry for abelian varieties and, more generally, complex and $`p`$–adic tori, that is, spaces of the form $`T/B`$ where $`T`$ is an algebraic group isomorphic to a product of multiplicative groups, $`K`$ is a complete normed field, and $`BT(K)`$ is a discrete subgroup of maximal rank in it. We also check its compatibility with other definitions discussed in the literature.
In the second section, we develop an approach to the quantization of abelian varieties first introduced in \[Ma1\], namely, via theta functions on non–commutative, or quantum, tori endowed with a discrete period lattice. These theta functions satisfy a functional equation which is a generalization of the classical one, in particular, involve a multiplier. Since multipliers cease to be central in the quantum case, one must decide where to put them. In \[Ma1\] only one–sided multipliers were considered. As a result, the product of two theta functions in general was not a theta function. Here we suggest a partial remedy to this problem by introducing two–sided multipliers. The resulting space of theta functions possesses partial multiplication and has sufficiently rich functorial properties so that rudiments of Mumford’s theory (\[Mu\]) can be developed. Main results of \[Ma1\] are reproduced here in a generalized form, so that this paper can be read independently.
We will now briefly describe a broader picture into which this work fits.
1. On mirror symmetry. In this paper mirror symmetry is understood as a binary relation between (weak) Calabi–Yau manifolds endowed with some additional data. A projective (or compact complex) manifold $`V`$ is a weak Calabi–Yau, if it admits nowhere vanishing global volume form. Additional data which are commonly considered are of two types.
A) A symplectic or complexified symplectic structure $`\omega _V`$ on $`V`$, which is “sufficiently large” in the case of complex base field.
B) A cusp $`c_V`$ in the moduli space (or rather stack) of deformations of $`V`$, that is, a neighborhood of a point of “maximal degeneration” to which $`V`$ belongs.
The mirror partnership relation between $`(V,c_V,\omega _V)`$ and $`(W,c_W,\omega _W)`$ consists in a host of identifications (partly conjectural) of various structures that can be produced starting with such triples. In the case of strict Calabi–Yau’s, this includes an identification of two Frobenius manifolds: quantum cohomology of $`(V,\omega _V)`$ and a germ of the extended moduli space of $`W`$ with its flat structure determined by $`c_W`$, and similarly with roles of $`V`$ and $`W`$ reversed. Generally, one expects also a representation of $`V`$ and $`W`$ as dual real Lagrangian torus fibrations over a common base, with a rich structure of Fourier–Mukai transform connecting Lagrangian/complex analytic objects on both sides. For more details on this, see original works \[MirS1\], \[MirS2\], \[StYZ\], \[Ko\], \[Giv\], \[Bar\], \[LYZ\], and a report \[Ma3\].
Most important testing ground for all levels of mirror correspondence is furnished by toric mirrors introduced and studied by V. Batyrev. In the simplest version, this construction looks as follows (see \[Bat\]).
Let $`T`$ be an $`n`$–dimensional algebraic torus, that is an algebraic group which is isomorphic to a product of $`n`$ multiplicative groups $`𝔾_m`$. It determines (and is functorially determined by any of) the two free abelian groups of rank $`n`$: its character group $`M_T:=\mathrm{Hom}(T,𝔾_m)`$ and its group of one–parametric subgroups $`N_T:=\mathrm{Hom}(𝔾_m,T)`$. These groups are naturally dual to each other. Denote by $`T^t`$ the dual torus, whose character group is $`M_{T^t}=N_T`$ and respectively $`N_{T^t}=M_T`$. In the main text of paper, one of these groups is denoted $`H`$ and another $`H^t.`$
Among various toric compactifications of $`T`$ we are interested in those for which their anticanonical system is ample. Anticanonical divisors on them are $`n`$–dimensional Calabi–Yau manifolds.
Such toric compactifications $`\overline{T}_\mathrm{\Delta }`$ are naturally indexed by reflexive polyhedra $`\mathrm{\Delta }`$ in $`N_T`$. Via standard convex duality, each such polyhedron determines the dual polyhedron $`\mathrm{\Delta }^t`$ in $`N_{T^t}`$ which is also reflexive. According to Batyrev, families of anticanonical hypersurfaces in the respective toric compactifications are expected to be mirror partners:
$$\overline{T}_\mathrm{\Delta }\left|K_{\overline{T}_\mathrm{\Delta }}\right|\left|K_{\overline{T}_{\mathrm{\Delta }^t}^t}\right|\overline{T}_{\mathrm{\Delta }^t}^t.$$
$`(0.1)`$
The relevant maximally degenerate CY’s are simply divisors at infinity in these toric compactifications.
Supplementing $`(T,\mathrm{\Delta })`$ with additional combinatorial structure, one can generalize this picture to some Calabi–Yau complete intersections.
The main goal of the first section of this paper is to provide a new definition of mirror symmetry for abelian varieties (and more generally, complex and $`p`$–adic tori) similar to (0.1). We use a “multiplicative uniformization” which goes back to Jacobi and which represents $`𝒜`$ as a quotient of $`n`$–dimensional algebraic torus $`T`$ by a multiplicative discrete lattice $`BT(K).`$ Consider now two dual algebraic tori endowed with period lattices which are explicitly identified, that is a diagram of the form:
$$(i,i^t):TBT^t.$$
$`(0.2)`$
We will say that pairs $`(𝒜:=T(K)/i(B),i^t)`$ and $`(:=T^t(K)/i^t(B),i)`$ are mirror dual to each other. Over complex field, we will relate $`i^t`$ (resp. $`i`$) to a structure similar to $`\omega _V`$ (resp. $`\omega _W`$) above and compare our mirror relation with that of \[Gr1\], \[Gr2\], \[AP\] and \[GolLO\]. In particular, we will see how this diagram gives rise to two dual fibrations of the relevant abelian varieties (or complex tori) by mutually dual Lagrangian real tori over the same base.
The choice of multiplicative uniformization is not unique, and it provides the environment for a partial compactification of the relevant moduli space and choice of a maximally degenerate point at the boundary of moduli space to which $`𝒜`$ is close, that is, of the relevant cusp of the moduli space. Roughly speaking, at the boundary some generators of the period lattice are forced to vanish so that the rank of the image of $`B`$ drops. To make it more precise, we can choose a fan $`\mathrm{\Phi }`$ in $`N_T`$, construct the dual fan $`\mathrm{\Phi }^t`$ and study the moduli space of the refined diagrams
$$(i,i^t):\overline{T}_\mathrm{\Phi }(K)B\overline{T}_{\mathrm{\Phi }^t}^t(K).$$
$`(0.3)`$
In this context, it turns out that the fibers of the dual mirror fibrations lie in the monodromy invariant homology class $`\tau `$ of middle dimension, as in the strict Calabi–Yau case.
0.3. Commutative and non–commutative tori and theta functions. In the second section of this paper we address the problem of constructing “quantum abelian varieties”. As we already mentioned, what we actually construct is a linear space of quantized theta functions (on a noncommutative torus with a period lattice) endowed with partial multiplication. We regard this semiring as a quantum deformation of the universal multigraded ring $`\mathrm{\Gamma }(L),L\mathrm{Pic}𝒜,[1]^{}(L)L`$, which makes sense for any abelian variety $`𝒜`$ (or indeed for any projective variety).
Since the whole construction is algebraic, it can be performed over any complete normed field, for example, $`p`$–adic field, and applied to the classical abelian varieties as well. Only those $`p`$–adic abelian varieties admit a multiplicative uniformization which have maximally degenerate stable reduction modulo $`p`$. This is the definition of the $`p`$–adic cusp. Diagrams of the type (0.2) make full sense in this context as well, but $`i^t`$ admits no straightforward interpretation as anything like symplectic form. It would be interesting to investigate the meaning of $`p`$–adic $`B`$–fields for strict Calabi–Yau manifolds.
We will now briefly discuss issues of non–commutative geometry involved in our construction of quantized theta functions. The standard approach is via deformation of classical function rings, and this intuition guided our initial construction in \[Ma1\] and to a certain degree its extension presented in this paper.
A complementary paradigm, made explicit on many occasions in Connes’ papers and the book \[Co\] is the natural appearance of non–commutative rings as objects encoding commutative spaces “with bad geometric properties”, typically quotients of commutative spaces by non–separated equivalence relations.
For example, in the context of multiplicative uniformization of abelian varieties $`𝒜()=T()/B`$, the multiplicative period lattice $`B`$ can degenerate also by ceasing to be discrete, although keeping its rank constant, and in this case it becomes natural to interpret $`T()/i(B)`$ in the realm of non–commutative geometry.
To be concrete, consider the case of one–dimensional commutative torus $`T.`$ Its maximal toric compactification is $`^1=𝔾_m\{0\}\{\mathrm{}\}`$. Choose $`B=`$ so that $`𝒜=^{}/(q^{})`$. On the level of diagram (0.3) we can choose here any $`q^1()`$. When $`|q|0,1,\mathrm{}`$, we get an elliptic curve $`𝒜=E_q`$ fibered by images of real tori $`|z|=\mathrm{const}`$ (where $`z`$ is the coordinate on $`^{}`$) over the circle $`/\mathrm{log}|q|`$. Dualizing this fibration, we will get the mirror dual elliptic curve. The points $`0,\mathrm{}`$ become divisors on a modular curve and provide the familiar degeneration picture in algebraic geometry. Points where $`q`$ is a root of unity also become visible in algebraic geometry as cusp points of modular curves of higher levels. However, $`q`$ of infinite order with $`|q|=1`$ are not considered in algebraic geometry at all. One way to interprete the space $`𝔾_m/(q^{})`$ in this case is to identify it with (one of the versions of) the two–dimensional non–commutative torus $`T_q`$, whose function ring is generated by $`x^{\pm 1},y^{\pm 1}`$ satisfying the commutation relation $`xy=qyx.`$ See Appendix for an informal explanation of this in the context of noncommutative geometry à la Alain Connes.
One remarkable trace of this origin of $`T_q`$ for $`q=e^{2\pi i\tau },\tau ,`$ as a limiting elliptic curve is the re–appearance of the modular group $`SL(2,Z)`$ as a symmetry group in the non–commutative situation: acting on $`\tau `$, it produces Morita equivalent function rings of $`T_q`$’s. See \[RiS\] for a thorough discussion of this in arbitrary dimension.
What I want to stress here, is a somewhat neglected complementary aspect of this picture: namely, that even if the group $`q^{}`$ is discrete, $`T_q`$ still can be viewed as a legitimate incarnation of the elliptic curve $`E_q`$ in the non–commutative world. A systematic treatment of the correspondence between, say, coherent sheaves on $`E_q`$ and modules over $`T_q`$ remains a problem for future, but see the recent preprint \[BEG\] for some precise facts about this correspondence, and many interesting suggestions are contained in \[So\].
This remark throws some light on the problems left open with our approach to quantized theta–functions, for example, that of functional equations corresponding to a change of cusp. Some of our non–commutative abelian varieties $`T(H,\alpha )/B`$ where $`T(H,\alpha )`$ is a non–commutative torus with quantization parameter $`\alpha `$ (see 2.1) can be understood as a result of taking a quotient of a commutative torus $`T(H^{},1)`$ of halved rank, by a non–discrete period subgroup with too many generators. This might provide a bridge between our construction and that of Weinstein (\[We\]) remaining entirely in the realm of commutative geometry.
Acknowledgement. I am grateful to D. Orlov for illuminating correspondence on Abelian mirrors and comments on the paper \[GolLO\]. I learned from Y. Soibelman the philosophy of treating boundaries of various moduli spaces as bridges to the non–commutative realm. A. Polishchuk has drawn my attention to the paper \[Fu\] and sent me a copy of it.
1. Mirror symmetry for complex tori and abelian varieties
1.1. Toric formalism. Since we will have to work with several different types of tori which must be carefully distinguished, we start with some terminological conventions.
Let $`K`$ be a ground field, $`H`$ a free abelian group of rank $`n`$ which we will always write additively. Algebraic torus $`T`$ with character group $`H`$ is the affine spectrum of the group ring of $`H`$, and $`H=\mathrm{Hom}(T,𝔾_m)`$. Put $`H^t:=\mathrm{Hom}(𝔾_m,T)`$. Groups $`H,H^t`$ are connected by the canonical duality map $`H\times H^t\mathrm{Aut}𝔾_m=.`$
An element $`hH`$ considered as a character of $`T`$ will be denoted $`e(h)`$ ($`e`$ for exponential). We have $`e(h+h^{})=e(h)e(h^{}).`$
In the next section, we will consider noncommutative algebraic tori as well, for which the multiplication rule for $`e(h)`$ is twisted: see (2.3) below.
1.1.1. Duality of algebraic tori. If $`T`$ is an algebraic torus with character group $`H_T`$, the dual algebraic torus $`T^t`$ has the character group $`H_{T^t}:=H_T^t=\mathrm{Hom}(H_T,)`$.
1.1.2. Periods and abstract tori. Let $`H`$ be the character group of $`T`$. Denote by $`B`$ another free abelian group of the same rank $`n`$, and by $`i:BT(K)`$ a homomorphism which we will call “period map”. Since $`T(K)=\mathrm{Hom}(H,K^{})`$, to give $`i`$ is the same as to give a pairing $`H\times BK^{}`$ such that $`(h,b)`$ is the value of $`e(h)`$ at the point $`i(b)`$.
We will refer to the quotient space $`T/i(B)`$ as an abstract torus, and to $`T`$ as its covering algebraic torus. One may imagine $`T/i(B)`$ simply as a functor of points on $`K`$–algebras $`RT(R)/i(B).`$ In the case when $`K=`$ and $`i`$ is an injection with discrete image, $`T()/i(B)`$ is a complex torus, which may admit a structure of abelian variety (it is then unique). However, a large part of our elementary formalism will not depend on additional assumptions. Thus, as explained in the Introduction, we will be able to include into our mirror picture $`p`$–adic tori and abelian varieties, and eventually non–commutative tori viewed as models of $`T()/i(B)`$ in non–commutative geometry.
In order not preclude the eventual interpretation of the space $`𝒜=T/i(B)`$ we will consistently identify $`𝒜`$ with the triple
$$(H_𝒜,B_𝒜,(,)_𝒜:H_𝒜\times B_𝒜K^{})$$
$`(1.1)`$
as above.
1.1.3. Poincaré dual abstract tori. By definition, Poincaré duality interchanges characters and periods. More precisely, $`𝒜`$ and $`\widehat{𝒜}`$ are Poincaré dual if
$$H_{\widehat{𝒜}}=B_𝒜,B_{\widehat{𝒜}}=H_𝒜,(b,h)_{\widehat{𝒜}}=(h,b)_𝒜^1\mathrm{for}hH_𝒜,bB_𝒜.$$
$`(1.2)`$
For abelian varieties this agrees with the classical definition. Notice however, that a choice of the covering algebraic torus is an additional structure, and we explicitly extend Poincaré duality to this context.
1.1.4. Framed tori. A framing of the abstract torus (1.1) is a map $`i^t:B_𝒜T_𝒜^t(K)`$. A framed abstract torus is a pair $`(𝒜,i^t)`$.
A framing of the complex torus or abelian variety consists of its representation as an abstract torus and framing of that abstract torus.
1.1.5. Mirror dual framed abstract tori. Two framed abstract tori $`(𝒜,i^t)`$ and $`(,i)`$ are called mirror dual, or mirror partners, if their covering algebraic tori are dual, and their periods are explicitly identified. More precisely, the relation of mirror partnership is provided by diagrams of the form (0.2). If one thinks about $`i,i^t`$ in terms of the respective character/period pairings (1.1), the mirror duality is provided by pairings
$$H_𝒜\times B_𝒜K^{},H_𝒜^t\times B_𝒜K^{}.$$
$`(1.3)`$
A framing is called non–degenerate if both kernels of the respective pairing are trivial.
This notion of mirror duality is the main definition of this section. We start with studying it for $`K=.`$
1.2. Complex tori. Assume that $`i(B)`$ is discrete in $`T()`$. Put
$$\mathrm{\Gamma }=\mathrm{\Gamma }_{}=\pi _1(T()/i(B),0)=H_1(𝒜,),$$
$`\mathrm{\Gamma }_{}=\mathrm{\Gamma }_{}`$ and similarly for $`\mathrm{\Gamma }_{}`$. If $`T()/i(B)`$ is denoted $`𝒜`$, we may write $`\mathrm{\Gamma }_𝒜,`$ $`\mathrm{\Gamma }_{𝒜,}`$, etc.
The real space $`\mathrm{\Gamma }_{}`$ can be identified with the Lie algebra of $`𝒜`$, and the exponential map $`\mathrm{exp}:\mathrm{\Gamma }_{}𝒜`$ with kernel $`\mathrm{\Gamma }_{}`$ is the universal covering of $`𝒜`$.
Let $`H=H_𝒜`$ be the character group of $`T`$. The map $`h{\displaystyle \frac{1}{2\pi i}}{\displaystyle \frac{de(h)}{e(h)}}`$ induces canonical identification $`H=H^1(T(),)`$. Hence we have an exact sequence
$$0H_𝒜^tH_1(𝒜,)B_𝒜0$$
$`(1.4)`$
where the third arrow is induced by $`i`$. Similarly, we have for a mirror dual framed torus $`(=T^t()/i^t(B),i)`$
$$0H_𝒜H_1(,)B_𝒜0$$
$`(1.5)`$
where the third arrow is now induced by $`i^t`$.
1.2.1. Mirror partners as dual real torus fibrations. In this subsection we will show that our definition of mirror partners $`(𝒜,)`$ over $``$ naturally fits into the general context of Lagrangian/complex duality: see \[AP\], \[StYZ\], \[Gr2\].
We start with a brief description of this context.
Let $`(X,\omega )`$ be a $`C^{\mathrm{}}`$ symplectic manifold, endowed with a submersion $`p_X:XU`$ whose fibers are Lagrangian tori. We will fix also a Lagrangian section $`0_X:UX.`$
Using $`\omega `$, we can identify the bundle of Lie algebras of the tori $`p_X^1(u),uU,`$ with the cotangent bundle $`T_U^{}.`$ Hence we have a canonical isomorphism $`X=T_U^{}/H`$ where $`H`$ is a Lagrangian sublattice in $`T_U^{}`$ with respect to the lift of $`\omega `$ which is the standard symplectic form on the cotangent bundle. There exists also a canonical flat symmetric connection on $`T_U^{}`$ for which $`H`$ is horizontal.
The local system $`H^t=om(H,)`$ is embedded as a sublattice into $`T_U`$, and we can define the mirror partner of $`(p_X:XU,\omega ,0_X)`$ as the toric fibration $`Y:=T_U/H^t`$ endowed with the projection to the same base $`p_Y:YU`$ and the zero section $`0_Y`$.
Passing from $`X`$ to $`Y`$ we have lost the symplectic form. To compensate for this loss, we have acquired a complex structure $`J:T_YT_Y`$ which can be produced from $`(p:XU,\omega ,0_X)`$ in the following way. The flat connection on $`T_U`$ obtained by the dualization from $`T_U^{}`$ produces a natural splitting $`T_Y=p_Y^{}(T_U)p_Y^{}(T_U).`$ With respect to this splitting, $`J`$ acts as $`(t_1,t_2)(t_2,t_1).`$
Conversely, suppose that we have a complex manifold $`Y`$ endowed with a fibration by real tori $`YU`$ with zero section, such that the operator of complex structure along the zero section identifies $`T_U`$ with the bundle of Lie algebras of fibers. Then we can consecutively construct the lattice $`H^tT_U`$, the dual fibration $`X:=T_U^{}/H`$ and the symplectic form on $`X`$ coming from the cotangent bundle.
Now we can return to complex tori.
Put $`S^1=\{|z|=1|z\}.`$ We have the Lie group isomorphism $`^{}S^1\times :`$ $`z(z/|z|,\mathrm{log}|z|)`$. This induces an isomorphism
$$(\alpha ,\lambda ):T()\mathrm{Hom}(H,S^1)\times \mathrm{Hom}(H,).$$
$`(1.6)`$
If $`i(B)`$ is discrete of maximal rank which I will assume, then $`\lambda i(B)`$ is an additive lattice in the real space $`\mathrm{Hom}(H,)`$. Thus (1.6) produces a real torus fibration of $`T()`$ over the base which is as well a real torus of the same dimension:
$$0\mathrm{Hom}(H,S^1)T()/i(B)\mathrm{Hom}(H,)/\lambda i(B)0.$$
$`(1.7)`$
Similarly, we have
$$0\mathrm{Hom}(H^t,S^1)T^t()/i^t(B)\mathrm{Hom}(H^t,)/\lambda ^ti^t(B)0$$
$`(1.8)`$
where $`\lambda ^t`$ is defined for $`T^t`$ in the same way as $`\lambda `$ for $`T`$. Let us identify linear real spaces $`H_{}`$ with $`H_{}^t`$ in such a way that lattice points $`\lambda i(b)`$ and $`\lambda ^ti^t(b)`$ are identified for all $`bB`$. Then (1.7) and (1.8) become dual real torus fibrations over the common base.
The relevant complex structures in our context come from covering tori. We have to introduce symplectic forms. Let us construct, say, $`\omega _𝒜`$. From (1.8) one sees that $`𝒜=T()/i(B)`$ can be obtained as quotient space of the tangent bundle of the base by a lattice. The tangent bundle (and the lattice) is canonically trivialized, and its fiber is $`H_{}.`$ Using the two framings, we have identified $`H_{}`$ with $`H_{}^t,`$ that is, tangent bundle with cotangent bundle. The canonical symplectic form on the cotangent bundle becomes our $`\omega _𝒜.`$ Clearly, fibers of (1.7) are Lagrangian tori. It remains to check that $`\omega _𝒜`$ determines $`I_{}`$ as above, but this is quite straightforward.
1.2.2. Maximal degeneration point and monodromy. Let us consider now the situation, described by (one half of) the diagram (0.3). More precisely, consider the space of maps of $`B`$ to a neighbourhood of a point of maximal degeneration in some toric compactification $`\overline{T}_\mathrm{\Phi }.`$ Such a point is a zero–dimensional orbit of $`T`$, thus it corresponds to a maximal cone in $`\mathrm{\Phi }`$. Assume for simplicity that it is the simplicial cone generated by a basis of $`H^t`$. This means that we identify $`H^t`$ with $`^n`$, $`T()`$ with $`(^{})^n`$, and choose as partial compactification the imbedding $`(^{})^n^n`$. Let $`D_r`$ be the $`r`$–th coordinate hyperplane in $`^n`$.
Choosing a basis of $`B`$ as well, that is, identifying it with $`^n`$, we see that the region of the partially compactified moduli space of multiplicatively uniformized complex tori that we are interested in can be identified with an open subspace of the matrix space $`^{n\times n}`$ whose columns generate a multiplicative sublattice. The discriminant locus in this region consists of its intersection with $`D_{rs}`$ where $`D_{rs}`$ is the $`r`$–th coordinate hyperplane in the $`s`$–th copy of $`^n`$ times other copies. The origin (intersection of all $`D_{rs}`$) is the maximum degeneration point.
Let $`q=(q_{rs})`$ be a point of the moduli space, $`𝒜_q`$ the respective torus. Denote by $`\gamma _rH_1(𝒜_q,)`$ the image of the $`r`$–th $`S^1`$ in (1.7) with counterclockwise orientation. Let $`M_{rs}`$ be the monodromy action of a small counterclockwise loop around $`D_{rs}`$ in the moduli space. All cycles $`\gamma _r`$ are monodromy invariant. Let $`\beta _s`$ be any lift to $`H_1(𝒜_q,)`$ of the $`s`$–th 1–cycle in the base torus in (1.7). Then $`M_{rs}`$ transforms $`\beta _s`$ into $`\beta _s+\gamma _r`$ and leaves other 1–cycles invariant.
Thus, the homology class $`\tau `$ of any fiber of (1.7) generates the cyclic group of invariant cycles of middle dimension.
This statement holds independently of the choice of the simplicial fan and a maximal degenerating cone in it. In this sense, the choice of a covering torus alone encodes essential information about large complex structure.
1.3. Framings and well–becoming pairs. In this subsection we will compare our construction of mirror partners with that of \[GolLO\]. In that paper, the additional structure on $`𝒜`$ is a complex–valued 2–form $`\omega `$ rather than symplectic form as above. I will show that this as well can be related to an appropriate framing.
Consider first a pair of framed complex tori $`(𝒜,i^t)`$ and $`(,i)`$ which are mirror partners as above. In \[GolLO\], sec. 10, the authors use decompositions
$$\mathrm{\Gamma }_𝒜=\mathrm{\Gamma }_{1,𝒜}\mathrm{\Gamma }_{2,𝒜},\mathrm{\Gamma }_{}=\mathrm{\Gamma }_{1,}\mathrm{\Gamma }_{2,}.$$
We will call such decompositions compatible with our choice of covering tori, if $`\mathrm{\Gamma }_{1,𝒜}=H_𝒜^t,`$ $`\mathrm{\Gamma }_{1,}=H_𝒜`$ as in (1.4), (1.5). Thus compatible decompositions are simply splittings of (1.4) and (1.5). The spaces $`\mathrm{\Gamma }_{𝒜,},`$ $`\mathrm{\Gamma }_,`$ are endowed respectively with complex structures $`I_𝒜,I_{}`$ coming from covering tori. Clearly,
$$\mathrm{\Gamma }_{𝒜,}=H_{𝒜,}^tI_𝒜H_{𝒜,}^t$$
so that compatible splittings satisfy the condition 10.3.1 (2) of \[GolLO\].
Consider now only $`𝒜`$, but equipped with a class $`\omega H^2(𝒜,)`$ interpreted as an antisymmetric complex–valued form on $`\mathrm{\Gamma }_𝒜`$. Assume that the extension of $`\omega `$ to $`\mathrm{\Gamma }_{𝒜,}`$ is $`I_𝒜`$–invariant (see \[GolLO\], last lines of 1.4 for an explanation of this condition). Assume moreover that there exists a compatible decomposition of $`\mathrm{\Gamma }_𝒜`$ such that $`\mathrm{\Gamma }_{1,𝒜}`$ and $`\mathrm{\Gamma }_{2,𝒜}`$ are $`\omega `$–isotropic. This means that $`(𝒜,\omega )`$ is a well–becoming pair in the sense of \[GolLO\], 10.3.1. Looking at (1.4), we see that $`\omega `$ can be uniquely reconstructed from its restriction which we also denote $`\omega `$
$$\omega :H_𝒜^t\times B_𝒜.$$
$`(1.9)`$
Exponentiating (1.9) we produce a framing $`i^t:BT^t()`$ in the form (1.3), that is
$$e(h)(i^t(b))=e^{2\pi i\omega (h,b+H_𝒜^t)},hH_𝒜^t,bB.$$
$`(1.10)`$
Clearly, $`i^t`$ remains the same, if we choose another compatible isotropic splitting. Moreover, it does not change if we add to $`\omega `$ another pairing taking values in $`2\pi i`$.
In this way we get a map from the set of all well–becoming pairs $`(𝒜,\omega )`$ admitting compatible isotropic splittings in the sense of \[GolLO\] to the set of framed abstract tori $`(𝒜,i^t)`$ with non–degenerate framings in our sense.
We can now complete the comparison of our definition of mirror duality with that of \[GolLO\].
###### \quad1.3.1. Theorem
Let $`(𝒜,i^t)`$, $`(,i)`$ be a mirror dual pair of framed complex abstract tori, admitting lifts $`(𝒜,\omega _𝒜)`$, $`(,\omega _{})`$ to well–becoming pairs. Then $`(𝒜,\omega _𝒜)`$, $`(,\omega _{})`$ are mirror dual in the sense of \[GolLO\].
Proof. We will compare our setting with that of \[GolLO\], 10.4 and 10.4.1. Groups $`\mathrm{\Gamma }_1,\mathrm{\Gamma }_2`$ in \[GolLO\] are our $`H_𝒜^t,B`$, fundamental group of the mirror dual torus is $`\mathrm{\Gamma }_1^t\mathrm{\Gamma }_2`$. This means that the real covering torus of their $``$ is the same as ours, that is $`T^t`$. It remains to compare the complex structures. In our case it is simply induced from $`T^t().`$ In \[GolLO\] it is described with the help of the complex structure operator $`I`$ acting upon $`(\mathrm{\Gamma }_1^t\mathrm{\Gamma }_2)`$ produced from $`\omega `$ in two steps: via formula (14) and subsequent projection described in 10.4.
In order to check that they coincide, we will reproduce a part of the argument in \[GolLO\] in the form which hopefully clarifies the meaning of their crucial formula (14).
Since the following construction must be considered in two different situations, described in \[GolLO\] 9.2 and 10.4 respectively, we slightly change the scope of our notation. From now on, $`\mathrm{\Gamma }_1,\mathrm{\Gamma }_2`$ will denote two abstract free abelian groups of the same finite rank, $`\mathrm{\Gamma }=\mathrm{\Gamma }_1\mathrm{\Gamma }_2`$, $`\mathrm{\Gamma }^{}=\mathrm{\Gamma }_1^t\mathrm{\Gamma }_2`$. Introduce two real tori
$$𝒞=(\mathrm{\Gamma }_1\mathrm{\Gamma }_2)_{}/\mathrm{\Gamma }_1\mathrm{\Gamma }_2,𝒞^{}=(\mathrm{\Gamma }_1^t\mathrm{\Gamma }_2)_{}/\mathrm{\Gamma }_1^t\mathrm{\Gamma }_2.$$
$`(1.11)`$
Consider the data of two types.
(i) Complex structures on $`𝒞^{}`$ described by the operators $`I`$ on $`(\mathrm{\Gamma }_1^t\mathrm{\Gamma }_2)_{}`$ such that their crossover components
$$I_{12}:\mathrm{\Gamma }_{2,}\mathrm{\Gamma }_{1,}^t,I_{21}:\mathrm{\Gamma }_{1,}^t\mathrm{\Gamma }_{2,}$$
are bijective.
(ii) Forms $`\omega ^2\mathrm{\Gamma }_{}^t`$ for which $`\mathrm{\Gamma }_{1,}`$ and $`\mathrm{\Gamma }_{2,}`$ are maximal isotropic.
We will establish a bijection between them in the following way.
Let us start with a complex structure $`I`$ in $`𝒞^{}`$. It determines (and is determined by) the space of invariant holomorphic 1–forms on $`𝒞^{}`$. Integrating them over $`\mathrm{\Gamma }_1^tH_1(𝒞^{},)`$, we will get all additive maps $`\mathrm{\Gamma }_1^t`$, in particular, all elements of $`\mathrm{\Gamma }_1`$. So we have an embedding $`\mathrm{\Gamma }_1H^0(𝒞^{},\mathrm{\Omega }^1):\gamma \nu _\gamma `$ such that for all $`\beta \mathrm{\Gamma }_1^t,\gamma \mathrm{\Gamma }_1,`$
$$(\beta ,\gamma )=_\beta \nu _\gamma .$$
$`(1.12)`$
This allows one to define a non–degenerate scalar product $`,:\mathrm{\Gamma }_2\mathrm{\Gamma }_1`$:
$$\gamma _2,\gamma _1=_{\gamma _2}\nu _{\gamma _1}.$$
$`(1.13)`$
Finally, we can extend it to a complex skew–symmetric form $`\omega `$ on $`\mathrm{\Gamma }`$ declaring $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$ to be isotropic:
$$\omega ((\gamma _1,\gamma _2),(\gamma _1^{},\gamma _2^{}))=\gamma _2,\gamma _1^{}\gamma _2^{},\gamma _1.$$
$`(1.14)`$
If we choose a basis of $`\mathrm{\Gamma }_1^t,\mathrm{\Gamma }_2`$ and a basis of holomorphic 1–forms whose period matrix over $`\mathrm{\Gamma }_1^t`$ is the identity $`E`$, then the Gram matrix of the pairing $`,`$ will be just the second half of the total period matrix. Let us denote it $`\tau `$. Then $`e^{2\pi i\tau }`$ is the matrix generating the multiplicative period lattice in the covering complex torus $`\mathrm{\Gamma }_{}^{}/\mathrm{\Gamma }_1^t`$ which in view of (1.14) agrees with (1.10).
Arguing now in reverse direction, we will show that knowing $`\tau `$ we can reconstruct the operator $`I`$ in the same basis and get essentially the \[GolLO\] formula. In fact, $`I`$ is uniquely determined by the requirement that for all $`\gamma \mathrm{\Gamma }^{}`$ and all holomorphic $`\nu `$ we have $`_{I\gamma }\nu =i_\gamma \nu .`$ Hence to find $`I`$ we must solve the matrix system
$$(\mathrm{Re}\tau +i\mathrm{Im}\tau ,E)\left(\begin{array}{cc}X& Y\\ U& V\end{array}\right)=(\mathrm{Im}\tau +i\mathrm{Re}\tau ,iE)$$
which gives
$$I=\left(\begin{array}{cc}X& Y\\ U& V\end{array}\right)=\left(\begin{array}{cc}(\mathrm{Im}\tau )^1\mathrm{Re}\tau & (\mathrm{Im}\tau )^1\\ \mathrm{Im}\tau \mathrm{Re}\tau (\mathrm{Im}\tau )^1\mathrm{Re}\tau & \mathrm{Re}\tau (\mathrm{Im}\tau )^1\end{array}\right)$$
$`(1.15)`$
For the first application of (1.15), let us choose a real torus $`𝒜`$ with period lattice $`\mathrm{\Gamma }_𝒜`$ and put $`𝒞=𝒜\widehat{𝒜}`$ where $`\widehat{𝒜}`$ is the Poincaré dual real torus, that is $`\mathrm{\Gamma }_{\widehat{𝒜}}=\mathrm{\Gamma }_𝒜^t`$ (this does not contradict our multiplicative description (1.2) although it may not be immediately obvious). Put $`\mathrm{\Gamma }_1=\mathrm{\Gamma }_𝒜,`$ $`\mathrm{\Gamma }_2=\mathrm{\Gamma }_𝒜^t.`$ Comparing our formula (1.15) with \[GolLO\] (14) in this situation, we see that our $`I`$ coincides with their $`I_\omega `$. In \[GolLO\], 8.4, $`𝒜`$ additionally possesses a complex structure, which produces the canonical complex structures on $`\widehat{𝒜}`$, $`𝒞`$ and $`𝒞^{}`$, say, $`J.`$ Moreover, $`\omega `$ is restricted to lie in the complexified Néron–Severi group, and as a result $`I`$ commutes with the inherited complex structure $`J`$ on $`𝒞^{}.`$
The setup which we are discussing in the Theorem 1.3.1 is that of \[GolLO\] 10.4 and 10.5. The relevant torus $`𝒞`$ is now $`𝒜`$, and its homology lattice is now split by the choice of a compatible decomposition like in 1.3. Extension of this splitting to $`𝒜\widehat{𝒜}`$ produces the period matrix $`\tau `$ which is block diagonal and consists of two blocks. Formula (1.15) is still valid for the mirror complex structure, when one replaces $`\tau `$ in it by the respective block. Putting everything together we see that our complex structure determined essentially by (1.10) indeed agrees with that of \[GolLO\].
1.3.2. Remark. \[GolLO\] contains several tentative descriptions of mirror dual pairs differing mostly by the exact choice of the additional structure that should be added to $`𝒜`$. Our Theorem 1.3.1 together with the Theorem 10.5 in \[GolLO\] indicates that the notion of a well–becoming pair endowed with a choice of one half of an isotropic decomposition captures just right amount of information. Replacing this structure by that of framing, we make explicit the important aspect of “large complex structure” in the case $`K=`$ and simultaneously extend the definition to abstract tori over arbitrary fields.
2. Quantized theta–functions and abelian varieties.
2.1. Category of non–commutative tori. Let $`H`$ be a free abelian group of finite rank and $`\alpha :H\times HK^{}`$ an alternating pairing: for all $`h,gH`$
$$\alpha (h,g)=\alpha (g,h)^1,\alpha (h_1+h_2,g)=\alpha (h_1,g)\alpha (h_2,g).$$
$`(2.1)`$
A morphism $`f:(H_1,\alpha _1)(H_2,\alpha _2)`$ is a group homomorphism $`f:H_1H_2`$ such that for all $`h,gH_1`$ we have
$$\alpha _2^2(f(h),f(g))=\alpha _1^2(h,g).$$
$`(2.2)`$
The bilinear form
$$\epsilon _f(h,g):=\alpha _1(h,g)\alpha _2^1(f(h),f(g))$$
with values in $`\{\pm 1\}`$ is called the characteristic of $`f`$ (and of $`F`$).
Any such pair $`(H,\alpha )`$ will be called the character group of the non–commutative torus $`T(H,\alpha )`$ whose ring of algebraic functions $`Al(H,\alpha )`$ is the linear space spanned over $`K`$ by the symbols $`e(h),hH,`$ with multiplication rule
$$e(h)e(h^{})=\alpha (h,h^{})e(h+h^{}).$$
$`(2.3)`$
We may write $`e_{H,\alpha }(h)`$ for $`e(h)`$ if need be.
Notice that $`\epsilon (h):=\alpha (h,h)`$ is a character of $`H`$ taking values $`\pm 1`$, and that from (2.3) we get the following formulas:
$$e(h_1)e(h_2)e(h_3)=\alpha (h_1,h_2)\alpha (h_1,h_3)\alpha (h_2,h_3)e(h_1+h_2+h_3),$$
$`(2.4)`$
$$e(h)^1=\epsilon (h)e(h).$$
$`(2.5)`$
We can also consider the two-sided $`Al(H,\alpha )`$–module of formal functions $`Af(H,\alpha )`$ consisting of infinite linear combinations $`_ha_he(h),a_hK`$, and, in the case of a complete normed field $`K`$ and an unitary quantization parameter $`\alpha `$ (that is, $`|\alpha |=1`$) the ring of analytic functions $`An(H,\alpha )`$ consisting of those formal functions for which $`|a_h|h^N0`$ for any $`N`$ as $`h\mathrm{}`$, $`h`$ being any Euclidean norm on $`H`$.
The form $`\alpha `$ can be called the quantization parameter. When $`\alpha 1`$, we get the usual notions of commutative geometry, so that $`T(H,1)`$ is the algebraic torus with character group $`H`$.
A morphism $`F:T(H,\alpha )T(H^{},\alpha ^{})`$, by definition, is given by the contravariant $`K`$–algebra homomorphism $`F^{}:Al(H^{},\alpha ^{})Al(H,\alpha )`$.
###### \quad2.1.1. Proposition
a) The set of invertible elements of $`Al(H,\alpha )`$ is $`\{ae(h)|aK^{},hH\}.`$ If $`F:T(H_2,\alpha _2)T(H_1,\alpha _1)`$ is a morphism of non–commutative tori, then the induced map $`f=[F]:H_1H_2`$ determined by $`F^{}(e(h))=a_he(f(h)),a_hK^{},`$ is additive and satisfies (2.2) and thus is a morphism of character groups.
b) The set of all morphisms $`F`$ with fixed $`f=[F]`$ is either empty, or has a natural structure of the principal homogeneous space over the group $`T(H_1,1)(K)=\mathrm{Hom}(H_1,K^{}).`$ In particular, if the characteristic of $`f`$ is 1, then $`F^{}:e(h)e(f(h))`$ extends to a uniquely defined morphism of rings of algebraic functions.
c) Any morphism $`F^{}`$ extends to $`Af`$ by $`F^{}(a_he(h))=a_hF^{}(e(h)).`$ If $`K`$ is normed and $`\alpha `$ unitary, then this extension maps analytic functions to analytic.
Proof. The first statement follows from the fact that $`H`$ can be endowed with the structure of a well–ordered group. For any such structure, the highest (resp. lowest) terms of a product are products of the highest (resp. lowest) terms, so that an invertible element coincides with its highest and lowest term.
To prove the second statement, rewrite the equality
$$F^{}(e(h)e(g))=F^{}(e(h))F^{}(e(g))$$
using (2.3). Comparing the $`e`$–terms, we see that $`f(h+g)=f(h)+f(g).`$ Comparing the coefficients, we get
$$a_ha_ga_{h+g}^1=\alpha _1(h,g)\alpha _2^1(f(h),f(g)).$$
$`(2.6)`$
The left hand side is a symmetric form in $`h,g`$, whereas the right hand side is alternate. Therefore this form takes values $`\{\pm 1\}`$. Hence (2.2) holds, and (2.6) is the characteristic $`\epsilon _f(h,g)`$ of $`f`$.
Finally, let $`f`$ be a morphism of character groups with characteristic $`\epsilon `$. Then ring morphisms $`F^{}`$ with $`[F]=f`$ bijectively correspond to the solutions $`\{a_h|hH_1\}`$ of the equations $`a_ha_ga_{h+g}^1=\epsilon _f(h,g)`$. If one such solution exists, then all others are of the form $`a_hc(h)`$ where $`c:H_1K^{}`$ is an arbitrary homomorphism.
The remaining statements are straightforward.
2.1.2. Direct product. By definition, the ring of algebraic functions of $`T(H_1,\alpha _1)\times T(H_2,\alpha _2)`$ is the tensor product of the respective rings. We can write
$$T(H_1,\alpha _1)\times T(H_2,\alpha _2)=T(H_1H_2,\alpha _1\alpha _2)$$
by identifying
$$e_{H_1,\alpha _1}(h_1)e_{H_2,\alpha _2}(h_2)=e_{H_1H_2,\alpha _1\alpha _2}((h_1,h_2)).$$
2.1.3. Some standard morphisms. (i) Shifts. Any point $`bT(H,1)(K)=\mathrm{Hom}(H,K^{})`$ determines an automorphism $`b^{}`$ of $`T(H,\alpha )`$:
$$b^{}(e(h)):=h(b)e(h),$$
$`(2.7)`$
where from now on we denote by $`h(b)`$ the value of $`e_{H,1}(h)`$ at the point $`b`$.
(ii) Multiplication by $`n`$. This is the morphism
$$[n]:T(H,\alpha )T(H,\alpha ^{n^2})$$
defined by
$$[n]^{}(e_{H,\alpha ^{n^2}}(h))=e_{H,\alpha }(nh).$$
$`(2.8)`$
For $`n=1`$ it is an endomorphism of $`T(H,\alpha )`$. It is also an endomorphism, if $`\alpha `$ takes values in roots of unity of degree $`d`$ and $`n^21\mathrm{mod}d.`$
The commutation rule with shifts is
$$b^{}[n]^{}=[n]^{}(nb)^{}.$$
$`(2.9)`$
(iii) External multiplication. It is the morphism
$$m_{\alpha ,\beta }:T(H,\alpha )\times T(H,\beta )T(H,\alpha \beta )$$
defined by
$$m_{\alpha ,\beta }^{}(e_{H,\alpha \beta }(h))=e_{H,\alpha }(h)e_{H,\beta }(h).$$
$`(2.10)`$
(iv) Mumford’s morphism. This is the morphism
$$M:T(HH,\alpha \alpha )T(HH,\alpha ^2\alpha ^2),$$
$$M^{}(e(h,g))=e(h+g,hg).$$
$`(2.11)`$
It is well defined, because
$$(\alpha \alpha )[(h+g,hg),(h^{}+g^{},h^{}g^{})]$$
$$=\alpha (h+g,h^{}+g^{})\alpha (hg,h^{}g^{})=\alpha ^2(h,h^{})\alpha ^2(g,g^{}).$$
2.2. Periods. We choose and fix an abelian group of periods $`BT(H,1)(K).`$ The period group is written additively; it acts upon $`T(H,\alpha )`$ by shifts.
Trying to make sense of the quotient $`T(H,\alpha )/B`$ we will study formal or analytic functions on $`T(H,\alpha )`$ with automorphic properties with respect to the the group $`\{b^{}|bB\}.`$
###### \quad2.3. Definition
A (two–sided) theta multiplier $``$ for the non–commutative torus $`T(H,\alpha )`$ and period group $`B`$ consists of the data $`=(h_l,h_r,\psi ,(,))`$ where
(i) $`h_l,h_r:BH`$ are two group homomorphisms.
We also put $`h^\pm :=h_l\pm h_r`$ and denote the image of $`bB`$ with respect to $`h_l`$ (resp. $`h_r,h^\pm )`$ as $`h_{b,l}`$ (resp. $`h_{b,r},h_b^\pm )`$.
(ii) $`\psi :BK^{}`$ is also a group homomorphism.
(iii) $`(,):B\times BK^{}`$ is a symmetric pairing.
These data must satisfy the following condition: for all $`b_1,b_2B`$
$$h_{b_2}^{}(b_1)=(b_1,b_2)^2\alpha (h_{b_1,l},h_{b_2,l})\alpha (h_{b_1,r},h_{b_2,r})^1.$$
$`(2.12)`$
2.3.1. Remark. Moduli space of quotients $`T(H,\alpha )/B`$ locally splits into a product of the classical moduli space of commutative tori $`T(H,1)/B`$ and the space of quantization parameters $`\alpha `$ (which in a sense also are “hidden periods”: cf. our discussion of Connes’ treatment of bad equivalence relations).
When $`K=`$, existence of sufficiently many theta functions is equivalent to the algebraicity of $`T(H,1)/B`$ which becomes an abelian manifold. Multipliers of such theta functions satisfy Riemann symmetry and positivity conditions.
Relations (2.12) represent an extension of the symmetry conditions to our enlarged moduli space. For a quantum version of positivity conditions, see Theorem 2.6.1 b) below.
2.4. Automorphy factors. For any theta multiplier $``$ and period $`bB`$, the automorphy factor $`j_{}(b)`$ is, by definition, the following linear endomorphism of any of the function spaces $`Al,Af,An`$ of $`T(H,\alpha )`$:
$$j_{}(b):\mathrm{\Phi }\psi (b)(b,b)e(h_{b,l})\mathrm{\Phi }e(h_{b,r})^1.$$
$`(2.13)`$
Clearly, it is invertible.
###### \quad2.5. Proposition
For $``$ fixed, the map $`bj_{}(b)^1b^{}`$ is a group homomorphism.
It is injective if $`\mathrm{Ker}h^{}=0.`$
Proof. We must check that
$$j_{}(b_1+b_2)^1(b_1+b_2)^{}=j_{}(b_1)^1b_1^{}j_{}(b_2)^1b_2^{}.$$
$`(2.14)`$
We have in view of (2.13), (2.4), (2.5):
$$j_{}(b)^1(e(h))=\psi (b)^1(b,b)^1\epsilon (h_{b,l})\alpha (h,h_b^+)\alpha (h_{b,r},h_{b,l})e(hh_b^{}).$$
$`(2.15)`$
Now apply both sides of (2.14) to the arbitrary $`e(h)`$ and then compare them using (2.12). A somewhat lengthy but straightforward calculation gives (2.14).
When $`h^{}`$ is injective, $`h_b^{}0`$ for non–zero $`b`$, so that $`j_{}(b)^1b^{}(e(h))e(h).`$
2.6. Theta functions and theta types. A (quantized) theta function with multiplier $``$ is a formal series $`\theta Af(T(H,\alpha ))`$ invariant with respect to the transformation group
$$\{j_{}(b)^1b^{}|bB\}.$$
In other words, $`\theta `$ must satisfy the functional equations
$$b^{}(\theta )=\psi (b)(b,b)e(h_{b,l})\theta e(h_{b,r})^1$$
$`(2.16)`$
for all $`bB.`$ Clearly, theta functions with multiplier $``$ form a linear space which we denote $`\mathrm{\Gamma }()`$. This notation is supposed to remind the case of usual abelian varieties where we deal with invertible sheaves and their sections.
Actually, different multipliers may produce the same space of theta functions or even homomorphisms $`bj_{}(b)^1b^{}`$. Consider, for example, the case of a commutative torus, where $`\alpha 1.`$ Then $`\mathrm{\Gamma }()`$ depends on $`h_l,h_r`$ only via their difference $`h^{}=h_lh_r.`$ Moreover, if $`(,)^{}`$ is another symmetric pairing such that $`\phi (b_1,b_2):=(b_1,b_2)^{}(b_1,b_2)^1`$ takes values in $`\{\pm 1\}`$, then $`\phi (b,b)`$ is multiplicative in $`b`$, and we may replace $`(\psi ,(,))`$ by $`(\psi ^{},(,)^{})`$ where $`\psi ^{}(b)=\varphi (b,b)\psi (b).`$
Generally, from (2.15) one sees that if $`j_{^{}}(b)^1b^{}=j_{}(b)^1b^{}`$ for all $`b`$, then $`h^{}=h^{}`$ and moreover,
$$(b,b)^1\alpha (h_{b,r},h_{b,l})=\pm (b,b)^1\alpha (h_{b,r}^{},h_{b,l}^{}),$$
$$\psi (b)^1\alpha (h,h_b^+)=\pm \psi ^{}(b)^1\alpha (h,h_b^+).$$
We will call two multipliers equivalent, if they have the same space of theta–functions. (This definition is reasonable only for ample multipliers, see below). An equivalence class of multipliers $`L`$ will be called a theta type. The space $`\mathrm{\Gamma }()`$ depends only on this class and can be denoted also $`\mathrm{\Gamma }(L).`$
###### \quad2.6.1. Theorem
a) We have
$$\mathrm{dim}\mathrm{\Gamma }()=[H:h^{}(B)].$$
$`(2.17)`$
b) Assume that $`K`$ is a normed field. Then all theta functions of type $``$ are analytic if $`[H:h^{}(B)]<\mathrm{}`$ and
$$\mathrm{log}|(b,b)\alpha (h_{b,l},h_{b,r})|$$
$`(2.18)`$
is a positively defined quadratic form on $`B`$.
In particular, assume that $`B`$ is free and the quantization parameter $`\alpha `$ is unitary, $`|\alpha (h,h^{})|1.`$ Then this condition means that $`\mathrm{rk}B=\mathrm{rk}H`$, $`B`$ is discrete in $`T(H,1)(K)`$ and $`\mathrm{log}|(b,b)|`$ is positively defined.
Proof. Let $`\theta =_{hH}a_he(h),a_hK,`$ $`bB.`$ We have
$$b^{}(\theta )=\underset{hH}{}a_hh(b)e(h)=\underset{hH}{}a_{h+h_b^{}}(h+h_b^{})(b)e(h+h_b^{}),$$
$`(2.19)`$
whereas the right hand side of (2.16) is:
$$\psi (b)(b,b)\underset{hH}{}a_h\epsilon (h_{b,r})\alpha (h_{b,r},h_{b,i})\alpha (h_b^+,h)e(h+h_b^{}).$$
$`(2.20)`$
In (2.19) we replace $`h_b^{}(b)`$ by $`(b,b)^2\epsilon (h_{b,l})\epsilon (h_{b,r})`$ (see (2.12)). Comparing coefficients of (2.19) and (2.20), we see that (2.16) is equivalent to
$$a_{h+h_b^{}}=a_h\psi (b)h(b)^1(b,b)^1\epsilon (h_{b,l})\alpha (h_{b,r},h_{b,l})\alpha (h_b^+,h)$$
$`(2.21)`$
for all $`hH`$ and $`bB.`$ Thus one can arbitrarily choose values $`a_h`$ for all $`h`$ in a system of representatives of $`H/h^{}(B)`$ and then uniquely reconstruct $`\theta `$. This proves the first statement of the theorem. It also shows that if $`\mathrm{\Gamma }()`$ is not finite dimensional, it necessarily contains non–analytic functions.
Assume now that $`\mathrm{\Gamma }()`$ is finite dimensional. Then on each coset $`h+h^{}(B)`$ we have
$$\mathrm{log}|a_{h+h_b^{}}|=\mathrm{log}|a_h|\mathrm{log}|(b,b)\alpha (h_{b,r},h_{b,l})|+\mathrm{log}|\psi (b)h(b)^1\alpha (h_b^+,h)|.$$
The second summand in the right hand side is quadratic in $`b`$ whereas the third is linear. Hence analyticity follows from the positive definiteness of (2.18).
2.6.2. Multiplication of theta functions. Generally, we can multiply analytic functions, but not formal ones. We will call a theta multiplier $``$ analytic, if $`\mathrm{\Gamma }()`$ consists of analytic functions, and ample, if $``$ satisfies conditions of Theorem 2.6.1 b).
We will call two analytic theta multipliers $`_i=(h_l^{(i)},h_r^{(i)},\psi _i,(,)_i),i=1,2,`$ composable (in this order) if $`h_l^{(2)}=h_r^{(1)}.`$ Define their product as
$$_1_2=:=(h_l^{(1)},h_r^{(2)},\psi _1\psi _2,(,)_1(,)_2).$$
A straightforward calculation shows that if $`_i`$ are ample, $``$ is ample as well, and the product of theta functions produces a well defined map
$$\mathrm{\Gamma }(_1)\mathrm{\Gamma }(_2)\mathrm{\Gamma }():\theta _1\theta _2\theta _1\theta _2.$$
$`(2.22)`$
One can call two ample theta types $`L_1`$ composable, if they contain pairs of composable multipliers. The multiplication in (2.22) does not depend on the choice of such a pair, but I did not check that the product type $`L`$ cannot change.
2.6.3. Quantized abelian varieties. Assume that $`(T(H,\alpha ),B)`$ admits ample theta multipliers. Then we consider $`_L\mathrm{\Gamma }(L)`$ where $`L`$ runs over all theta types with $`\psi (b)\{\pm 1\}`$, together with partial multiplication defined above, as a quantized version of the graded coordinate ring of an abelian variety. In the classical case, it is graded by the symmetric elements of $`\mathrm{Pic}`$ lying in the effective cone (see 2.7.2 below).
2.7. Functorial properties of theta functions. Consider a morphism of non–commutative tori $`F:T(H_2,\alpha _2)T(H_1,\alpha _1)`$. As in Proposition 2.1.1 , let $`F^{}(e(h))=a_he(f(h)),a_hK^{}`$, $`f:H_1H_2.`$ The map $`f`$ induces a morphism of commutative tori which we also denote $`F`$: $`T(H_2,1)T(H_1,1)`$. Let $`B_iT(H_i,1)(K)`$ be two period lattices such that $`F(B_2)B_1.`$
Choose a theta multiplier $`_1=(h_l,h_r,\psi ,(,))`$ for $`T(H_1,\alpha _1)`$ and $`B_1.`$ We will show, how to produce from it a new theta multiplier
$$_2=F^{}(_1)=(h_l^{},h_r^{},\psi ^{},(,)^{})$$
for $`T(H_2,\alpha _2),B_2,`$ if the following condition holds:
the characteristic of $`F`$ is 1, that is, $`a_{h+g}=a_ha_g`$ (cf. end of the proof of Proposition 2.2.1), and $`f`$ is compatible with $`\alpha _1`$ and $`\alpha _2`$, and not just their squares.
The map $`f:H_1H_2`$ induces a map $`F:\mathrm{Hom}(H_2,K^{})\mathrm{Hom}(H_1,K^{})`$ whose restriction on $`B_2`$ sends it to $`B_1`$. Put for $`b,b_1,b_2B_2`$:
$$h_{l,r}^{}=fh_{l,r}F:B_2H_2,$$
$`(2.23)`$
so that $`h_{b,l}^{}=f(h_{F(b),l})`$ and similarly for $`h_r^{},h^\pm ,`$
$$\psi ^{}(b)=\psi (F(b))a_{h_{F(b)}^{}},$$
$`(2.24)`$
$$(b_1,b_2)^{}=(F(b_1),F(b_2)).$$
$`(2.25)`$
###### \quad2.7.1. Theorem
The data (2.23)–(2.25) constitute a theta multiplier for $`T(H_2,\alpha _2)`$, $`B_2`$ such that $`F^{}(\mathrm{\Gamma }())\mathrm{\Gamma }(F^{}()).`$
Proof. We have first to check that (2.12) holds for $`F^{}()`$. For $`b_1,b_2B_2`$, the left hand side becomes
$$h_{b_2}^{}(b_1)=f(h_{F(b_2)}^{})(b_1)=h_{F(b_2)}^{}(F(b_1)).$$
Furthermore, if the characteristic of $`F`$ is 1, the right hand side can be rewritten as
$$(b_1,b_2)^2\alpha _2(f(h_{F(b_1),l}),f(h_{F(b_2),l}))\alpha _2(f(h_{F(b_1),r}),f(h_{F(b_2),r}))^1$$
$$=(F(b_1),F(b_2))^2\alpha _1(h_{F(b_1),l},h_{F(b_2),l})\alpha _1(h_{F(b_1),r},h_{F(b_2),r})^1.$$
Applying (2.12) to $`F(b_1),F(b_2)`$ in lieu of $`b_1,b_2`$, we see that both expressions coincide.
Let us now check that $`F^{}(\mathrm{\Gamma }())\mathrm{\Gamma }(F^{}()).`$
Choose a formal function
$$\theta =\underset{hH_1}{}c_he(h),e=e_{H_1,\alpha _1}.$$
According to (2.21), it belongs to $`\mathrm{\Gamma }()`$ iff the following conditions are satisfied for all $`hH_1,bB_1`$:
$$c_{h+h_b^{}}=c_h\psi (b)h(b)^1(b,b)^1\epsilon _1(h_{b,l})\alpha _1(h_{b,r},h_{b,l})\alpha _1(h_b^+,h).$$
$`(2.26)`$
Notice that our notation slightly differs from (2.21): $`a_h`$ is now reserved for $`F^{}(e(h))=a_he(f(h))`$ as in the first lines of 2.7.
Put now
$$F^{}(\theta )=\underset{gH_2}{}C_ge(g),e=e_{H_2,\alpha _2}.$$
We have $`C_g=0`$, if $`gf(H_1)`$; otherwise
$$C_g=\underset{hf^1(g)}{}c_ha_h.$$
$`(2.27)`$
To prove that $`F^{}(\theta )\mathrm{\Gamma }(F^{}())`$, we will check that $`C_g`$ satisfy analogs of relations (2.26) written for all $`gH_2`$ and all $`bB_2.`$
Consider first the case $`gf(H_1).`$ Then the relevant analog of (2.26) says that $`C_{g+h_b^{}}`$ must be proportional to $`C_g`$ that is, zero. This is indeed true, because in view of (2.23), $`h_b^{}=f(h_{F(b)}^{})`$ and thus $`g+h_b^{}f(H_1).`$
Now assume that $`gf(H_1)`$, fix also $`bB_2`$, and write down separately both sides of the relevant case of (2.26). Since characteristic of $`F`$ is 1, the left hand side can be written as
$$C_{g+h_b^{}}=\underset{hf^1(g)}{}c_{h+h_{F(b)}^{}}a_ha_{h_{F(b)}^{}}$$
$$=\underset{hf^1(g)}{}c_ha_h\psi (F(b))a_{h_{F(b)}^{}}h(F(b))^1(F(b),F(b))^1$$
$$\times \epsilon _1(h_{F(b),l})\alpha _1(h_{F(b),r},h_{F(b),l})\alpha _1(h_{F(b)}^+,h).$$
$`(2.28)`$
Here we have rewritten $`c_{h+h_{F(b)}^{}}`$ using (2.26).
The right hand side is
$$\left(\underset{hf^1(g)}{}c_ha_h\right)\psi ^{}(b)g(b)^1(F(b),F(b))^1\epsilon _2(h_{b,l}^{})\alpha _2(h_{b,r}^{},h_{b,l}^{})\alpha _2(h_b^+,g).$$
$`(2.29)`$
We can now compare (2.28) and (2.29) term by term using (2.23)–(2.25), and convince ourselves that they coincide. In particular, we use the identities $`g(b)=h(F(b))`$ and $`g=f(h)`$ for any $`hf^1(g).`$
2.7.2. Examples. All morphisms of tori described in 2.1.3 have characteristic 1. Therefore, complementing these tori by compatible period lattices, we obtain quantized versions of many standard morphisms of abelian varieties.
In particular, $`[1]^{}`$ acts on theta multipliers for $`T(H,\alpha ),B`$ by simply inverting $`\psi `$. Hence symmetric theta multipliers correspond to $`\psi :B\{\pm 1\}.`$
Mumford’s morphism for abelian varieties, induced by the toric morphism (2.11), is the starting point for his study of homogeneous coordinate rings of abelian varieties. Our preparatory work allows us to hope that at least part of this study can be extended to the quantum case.
Appendix. Commutative Geometry as Noncommutative Geometry
The main goal of this Appendix is to illustrate Connes approach to noncommutative geometry in an algebraic geometric context and to make convincing our claim that an elliptic curve “is” a two–dimensional noncommutative torus.
We take as our starting point Connes explanations about how to treat as a non–commutative space quotient of a “commutative space” by an equivalence relation or equivalence groupoid: see \[Co\], II.2–II.5. This viewpoint is complementary to the more popular in algebraic geometry deformation paradigm.
First, a reminder about groupoids.
Let $`U`$ be a set. Classically, an equivalence relation $``$ on $`U`$ is given by its graph $`RU\times U,R:=\{(a,b)|ab\}`$ which satisfies the three conditions:
Reflexivity:
$$aa\mathrm{\Delta }_UR;$$
Symmetry:
$$(abba)s_{12}(R)=R;$$
Transitivity:
$$((ab)\&(bc)ac)\mathrm{pr}_{13}[(R\times U)(U\times R)]R.$$
All of this can be rephrased as follows: there exists a category with the set of objects $`U,`$ set of morphisms $`R,`$ such that $`RU\times U`$ is the map $`f`$ (source of $`f`$, target of $`f`$), and in addition, every morphism is an isomorphism, and all automorphism groups are trivial.
Consider now a diagram $`RU\times U`$ satisfying this description with the last condition deleted so that the automorphism groups can now be arbitrary. We will call such a diagram an equivalence groupoid (on the set $`U`$). Of course, an equivalence groupoid $`RU\times U`$ comes together with the identity map $`UR:a\mathrm{id}_a`$ and the associative multiplication map $`R\times _URR`$ satisfying the usual categorical axioms which reduce to the reflexivity, symmetry, and transitivity for the usual equivalence relations. Notice that the image of $`R`$ is in fact an equivalence relation, and the respective quotient is the set of isomorphism classes of objects.
Thus the basic difference between equivalence groupoids and equivalence relations on sets can be demonstrated on one–point sets $`U=\{\}`$: in this case $`R`$ is simply a group. In the framework of homotopy theory, the respective quotient object $`\{\}/R`$ is represented by the classifying space $`BR`$. Stacks provide a categorical context for constructing such quotients in algebraic geometry.
In fact, the notion of equivalence groupoid was formulated in such a way that it readily generalizes to the case when $`RU\times U`$ is a diagram in an arbitrary category with products, e.g. schemes. One should imagine $`U`$ as an atlas and $`R`$ as gluing rules, so that the geometric object we are interested in is symbolically $`U/R`$.
Connes prescription, roughly speaking, consists in reducing geometric study of $`X=U/R`$ to the representation theory of a non–commutative ring $`_X`$. The heuristic rule for constructing $`_X`$ can be stated as follows. Let us write elements (points) of $`R`$ as morphisms $`j:uu^{}`$ so that $`RU\times U`$ maps such a point to $`(u,u^{}).`$ Then $`_X`$ consists of certain functions $`f`$ on $`R`$ endowed with convolution multiplication:
$$(fg)(k:uu^{\prime \prime })=\underset{(i,j):ij=k}{}f(j:uu^{})g(i:u^{}u^{\prime \prime }).$$
$`(A.1)`$
Of course, correct choice of the class of functions and making sense of the convolution multiplication may present a problem, but in algebraic geometry it is clear what to start with at least when $`U`$ is affine.
Let us illustrate such a setup by several examples.
Example 1. If $`X`$ is the quotient $`\{\}/G`$ where $`G`$ is a finite group, then $`_X`$ is the group ring of $`G`$. Representation theory of $`_X`$ is essentially $`K`$–theory of the classifying stack of $`G`$, in accordance with the common wisdom.
Example 2. If $`X`$ is the quotient of an affine scheme $`U=\mathrm{Spec}A`$ by the action of a group $`G`$, then $`_X`$ should contain at least the twisted product of $`A`$ with a version of the group ring of $`G`$.
When $`U=𝔾_m`$ with coordinate $`y`$ and $`G=`$ with generator $`x`$ acting on $`y`$ as multiplication by $`q`$, then in this twisted product we have $`xyx^1=qx`$ that is, the basic relation of the noncommutative torus. On the other hand, $`U/G`$ makes sense as an elliptic curve in analytic geometry, if $`|q|1.`$
Example 3. A projective scheme $`\mathrm{Proj}A`$ is the quotient of its cone $`\mathrm{Spec}A`$ (with vertex deleted) by the action of $`𝔾_m`$ determining the grading. Studying representations of the twisted product of $`A`$ with the group ring of $`G_m`$ is equivalent to studying graded $`A`$–modules. Deleting the vertex boils down to taking the quotient of the category of graded modules by the subcategory of modules with finite number of nonvanishing components. This is the familiar Serre’s picture which is archetypal in the following sense: after finding an appropriate ring $`_X`$, one proceeds to establish an equivalence of categories $`_X`$$`\mathrm{Mod}\mathrm{Coh}_X`$, or its localized and/or derived version.
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# Depletion potential in hard-sphere mixtures: theory and applications
## I Introduction
Two big colloidal particles immersed in a fluid of smaller colloidal particles or non-adsorbing polymers or micelles experience an attractive depletion force when the separation $`h`$ of the surfaces of the big particles is less than the diameter of the small ones. The expulsion or depletion of the small particles gives rise to anisotropy of the local pressure which results in the effective attractive force between the big ones. Asakura and Oosawa and, independently, Vrij used excluded volume arguments to determine the effective potential between two big hard spheres (modeling the colloids) assuming that the small particles or polymers form a mutually non-interacting fluid whose centers are excluded from the surfaces of the colloids by a distance $`R_s`$ . The resulting depletion potential is attractive for $`h<2R_s`$ and is zero for $`h2R_s`$; it increases monotonically with $`h`$ from its value at contact, $`h=0`$, and is proportional to $`\eta _s`$, the packing fraction of the small particles \[see, c.f., Eq. (14)\]. Much attention has been paid to depletion induced attraction within colloid science since it provides an important driving force for phase separation and flocculation phenomena in mixtures of colloids and in colloid-polymer mixtures. From a statistical mechanics viewpoint depletion forces are of considerable interest since they arise from purely entropic effects because the bare interactions between the particles are hard-sphere-like. Formally, it is the integrating out of the degrees of freedom of the small particles which gives rise to the effective interaction between two big ones.
Although colloid-polymer mixtures can, under favorable circumstances, be modeled by a binary mixture of hard spheres and ideal, non-interacting polymers (we term this the Asakura-Oosawa model), for mixtures of colloids or colloids and micelles a more appropriate zeroth-order model is a binary hard-sphere mixture, i.e., the small particles are not interpenetrating but experience mutual hard-sphere repulsion. In this case it becomes a key question as to how the depletion potential between two big hard spheres is influenced by interactions between the small spheres. For high packing fractions $`\eta _s`$, one might suppose that the small spheres exhibit pronounced short-ranged correlations (layering) leading to significant changes in the depletion potential. This would, in turn, have repercussions for the phase behavior of the bulk mixture, making this significantly different from that of the Asakura-Oosawa model. Such considerations have prompted several recent theoretical investigations of phase behavior based on an effective one-component depletion potential description of model colloidal mixtures . The crucial ingredient in such investigations is an accurate depletion potential.
Having a proper understanding of depletion potentials is not only relevant for bulk phase behavior; it is of intrinsic interest. In recent years a variety of experimental techniques have been developed which measure, directly or indirectly, the depletion potential between a colloidal particle, immersed in a sea of small colloids or polymers, and a fixed object such as a planar wall . Video microscopy has also been used to determine depletion forces for a single big colloid in a solution of small colloids inside a vesicle – a system which resembles hard spheres inside a hard cavity . Very recently Crocker et al. measured the depletion potential between two big PMMA spheres immersed in a sea of small polystyrene spheres for a range of packing fractions of the latter (see, c.f. Sec V B). At low values of $`\eta _s`$ the measured depletion potential is well-described by the Asakura-Oosawa result but at higher packing fractions the potential exhibits a repulsive barrier and for $`\eta _s0.26`$ the depletion potential is damped oscillatory with a wavelength that is of the order of the small particle diameter. As experiments grow in sophistication and in resolution it is likely that further details of depletion potentials will be revealed whose interpretation will require a reliable and versatile theoretical approach. Such an approach should be able to tackle experimental situations where $`\eta _s`$ is rather high and to treat general ‘confining’ geometries. The latter include a big particle near a planar wall or in a wedge or cavity, as well as the case of a big particle near another, fixed big particle. In this paper we describe such a theory for the depletion potential based on a density functional treatment (DFT) of a fluid mixture. Our treatment avoids the limitations of the virial expansion (in powers of $`\eta _s`$) and the uncontrolled nature of the Derjaguin approximation which are inherent in recent approaches to depletion forces in hard-sphere mixtures. It is less cumbersome than the alternative integral equation treatments and more easily adopted to different geometries. A key feature of our treatment is that it does not require the calculation of the total free energy of the inhomogeneous fluid or of the local density of the small particles in contact with the big particle . The method is much easier to implement than a direct minimization of the free energy functional which is numerically very demanding when any symmetry of the density profile of the small spheres is broken by the presence of the big particle.
The paper is arranged as follows: Subsection II A defines the depletion potential in an arbitrary mixture, showing how this is related to the one-body inhomogeneous direct correlation function of the big particles. In Subsec. II B we use this result to derive an explicit formula for the the depletion potential in the low-density limit, where the densities of all species approach zero. For the particular case of a binary hard-sphere mixture in this limit we recover the Asakura-Oosawa result. Subsection II C describes the general asymptotic behavior of the depletion potential for $`h\mathrm{}`$ while Subsec. II D describes the implementation of the theory for a binary hard-sphere mixture using the DFT of Rosenfeld . In Section III we present several comparisons of our hard-sphere DFT results, for both sphere-sphere and (planar) wall-sphere depletion potentials, with those of computer simulation . Our theory performs well for all size ratios $`sR_s/R_b`$ and packing fractions $`\eta _s`$ for which simulation results are available. We show that the leading-order asymptotic result for the depletion potential provides an excellent account of the oscillations in the calculated potential not only at longest range but also at intermediate separations of the big spheres. Section IV is concerned with assessing the regime of validity of the well-known Derjaguin approximation which relates the force between two big objects to the integral of the solvation force, or excess pressure, of the small particles confined between two planar walls \[see, c.f., Eq. (40)\]. We argue that this approximation is not reliable for the hard-sphere mixture even when the size ratio $`s`$ is very small. In Subsec. V A we describe a simple but accurate parametrization of the depletion potential suitable for a big hard sphere near a planar hard wall and for the potential between two big hard spheres. Such a parametrized form should prove useful for effective Hamiltonian studies of phase behavior and colloid structure . Subsection V B presents results for the depletion potential in a binary hard-sphere mixture where the size ratio is chosen to mimic the system considered in the experiments of Ref. . We conclude in Sec.VI with a discussion and summary of our results.
## II The depletion potential
### A General theory
We consider a general mixture of $`\nu `$ components in which each species $`i`$ ($`i=1,\mathrm{},\nu `$), characterized by its radius $`R_i`$, is coupled to a reservoir with chemical potential $`\mu _i`$ and is subject to an external potential $`V_i(𝐫)`$. The mixture at thermodynamic equilibrium can be described by the set of number density profiles $`\{\rho _i(𝐫)\}`$. For such a mixture we wish to calculate the depletion potential, or the depletion force, between an object fixed at position $`𝐫_1`$ and a second one fixed at $`𝐫_2`$. Without loss of generality the position $`𝐫_1`$ of the first object is chosen as the origin of the coordinate system. This fixed object then exerts an external potential on the particles constituting the mixture. The external potential can represent a planar hard wall or a fixed particle of the mixture, or more generally, a curved surface or soft planar walls . If the depletion potential between two particles of the mixture is to be calculated either particle can be chosen to act as the external potential and this point will be addressed in more detail in a later section.
In the following the second object is a test particle of a species denoted as $`b`$. The grand potential of the mixture when the test particle is fixed at the position $`𝐫_b`$ in the presence of the fixed object exerting the external potential $`V_b(𝐫)`$ is denoted by $`\mathrm{\Omega }_{tb}(𝐫_b;\{\mu _i\};T)`$. $`W_t(𝐫_b)`$, the quantity of interest here, is defined as the difference of grand potential between a configuration in which the test particle is in the vicinity of the fixed object and one in which the test particle is deep in the bulk, i.e., $`𝐫_b\mathrm{}`$:
$$W_t(𝐫_b)=\mathrm{\Omega }_{tb}(𝐫_b;\{\mu _i\};T)\mathrm{\Omega }_{tb}(𝐫_b\mathrm{};\{\mu _i\};T).$$
(1)
In order to calculate this difference the test particle can be moved along any path from one configuration to the other. A particular path which simplifies the calculation is via the reservoir. This path can be divided into two steps. In the first step the test particle is removed from the bulk at $`𝐫_b\mathrm{}`$ and put into the reservoir. In the second step the test particle is taken from the reservoir and is inserted back into the mixture but now at $`𝐫_b`$. The formal means to describe particle insertion in a general mixture is the potential distribution theorem and we employ this in the grand ensemble .
The potential distribution theorem provides an expression for the partition function $`\mathrm{\Xi }_{tb}(𝐫_b;\{\mu _i\};T)`$ of the mixture after a test particle of species $`b`$ is inserted at position $`𝐫_b`$ in terms of the partition function of the mixture $`\stackrel{~}{\mathrm{\Xi }}(\{\mu _i\};T)`$ and the number density profile $`\rho _b(𝐫)`$ of species $`b`$ before the particle insertion:
$$\mathrm{\Xi }_{tb}(𝐫_b;\{\mu _i\};T)=\mathrm{exp}\left(\beta (V_b(𝐫_b)\mu _b)\right)\mathrm{\Lambda }_b^3\rho _b(𝐫_b)\stackrel{~}{\mathrm{\Xi }}(\{\mu _i\};T),$$
(2)
with $`\beta ^1=k_BT`$ and $`\mathrm{\Lambda }_b`$ the thermal wavelength of species $`b`$. Together with a well-known result from density functional theory (DFT) ,
$$\mathrm{\Lambda }_b^3\rho _b(𝐫)=\mathrm{exp}\left(\beta (\mu _bV_b(𝐫))+c_b^{(1)}(𝐫;\{\mu _i\})\right),$$
(3)
it follows that the one-body direct correlation function $`c_b^{(1)}`$ of species $`b`$ can be written as
$`c_b^{(1)}(𝐫_b;\{\mu _i\})`$ $`=`$ $`\mathrm{ln}\left(\mathrm{\Xi }_{tb}(𝐫_b;\{\mu _i\};T)/\stackrel{~}{\mathrm{\Xi }}(\{\mu _i\};T)\right)`$ (4)
$`=`$ $`\beta \stackrel{~}{\mathrm{\Omega }}(\{\mu _i\};T)\beta \mathrm{\Omega }_{tb}(𝐫_b;\{\mu _i\};T),`$ (5)
i.e., $`\beta c_b^{(1)}(𝐫_b;\{\mu _i\})`$ describes the change in the grand potential of the whole system due to insertion of a test particle. The grand potential difference defined by Eq. (1) can now be expressed in terms of the difference of one-body direct correlation functions:
$$\beta W_t(𝐫_b)=c_b^{(1)}(𝐫_b\mathrm{};\{\mu _i\})c_b^{(1)}(𝐫_b;\{\mu _i\}).$$
(6)
As the potential distribution theorem, Eq. (2), is a general result, valid for any number of components, for arbitrary densities of all components and, in fact, for any inter-particle potential function the same generality holds for Eq. (6). No approximations have been made so far. However, in order to use Eq. (6) to calculate $`\beta W_t(𝐫)`$ an explicit procedure that can treat a mixture must be applied. Simulations provide such a procedure as does density functional theory. We shall consider both here.
We emphasize that the direct correlation function entering Eq. (6) depends on the equilibrium density profiles before the test particle of species $`b`$ is inserted at position $`𝐫_b`$. This observation simplifies the calculation of $`\beta W_t(𝐫)`$ dramatically because the symmetry of the relevant density profiles $`\{\rho _i(𝐫)\}`$ is determined solely by the symmetry of the external potentials and therefore depends only on the nature of the object that is fixed at the origin. If this object is a structureless planar wall and in the absence of spontaneous symmetry breaking such as prefreezing or crystalline layer formation the density profiles of all species reduce to one-dimensional profiles $`\{\rho _i(z)\}`$ with $`z`$ the distance perpendicular to the wall. For a fixed spherical or cylindrical wall or particle the density profiles $`\{\rho _i(r)\}`$ depend only on the radial distance. Even if the fixed object is a wall of more general shape, so that there is no simple symmetry involved in the problem, calculating the density profiles before particle insertion is much easier than after insertion, when the broken symmetry due to the presence of the test particle leads to a more complex dependence of the profiles on the coordinates.
While Eq. (6) can be evaluated for arbitrary densities of species $`b`$ within the present DFT approach, a particular limit in which the density of species $`b`$ goes to zero is considered now. This dilute limit is especially important since it arises in the context of measuring depletion forces and in formal procedures for deriving effective Hamiltonians for big particles by integrating out the degrees of freedom of the small particles. For example, if in a binary mixture the degrees of freedom of the small particles are integrated out the resulting effective one-component fluid can be described by an effective Hamiltonian containing a volume term, to which only the small particles contribute, a one-body term, in which a single big particle in a ’sea’ of small particles contributes, a two-body term, a three-body term and so on . For highly asymmetric mixtures the most important contributions come from the volume and the one- and two-body terms. This assumption is substantiated by the results of calculations of three-body contributions reported in Ref. for a size ratio $`s=0.1`$. Three-body contributions also seem to be small for $`s=0.2`$ . Note that the two-body term describes an effective pairwise interaction potential between two big particles which turns out to be precisely the depletion potential, i.e., $`\beta W_t(𝐫)`$ evaluated in the dilute limit .
In the grand ensemble the dilute limit can be obtained by taking the limit in which the chemical potential of species $`b`$, $`\mu _b\mathrm{}`$, with the chemical potentials of all other species $`\{\mu _{ib}\}`$ kept fixed. The depletion potential is then given by
$`\beta W(𝐫)`$ $``$ $`\underset{\mu _b\mathrm{}}{lim}\beta W_t(𝐫)`$ (7)
$`=`$ $`c_b^{(1)}(𝐫\mathrm{};\{\mu _{ib}\},\mu _b\mathrm{})c_b^{(1)}(𝐫;\{\mu _{ib}\},\mu _b\mathrm{}),`$ (8)
which contains no explicit dependence on the external potentials that are present, i.e., the depletion potential depends only on the intrinsic change of the grand potential.
Although in the dilute limit both the density profile $`\rho _b(𝐫)`$ and the bulk density $`\rho _b^{bulk}=\rho _b(\mathrm{})`$ of species $`b`$ vanish, the ratio stays finite and the depletion potential can also be obtained from the result
$$\beta W(𝐫)=\underset{\mu _b\mathrm{}}{lim}\mathrm{ln}\left(\frac{\rho _b(𝐫)}{\rho _b(\mathrm{})}\right)V_b(𝐫)+V_b(\mathrm{}).$$
(9)
which takes a more familiar form if we re-write Eq. (9) as $`p(𝐫)/p(\mathrm{})=\mathrm{exp}[\beta (W(𝐫)+V_b(𝐫))]`$, where $`p(𝐫)`$ is the probability density of finding the particle of species $`b`$ at a position $`𝐫`$ and we assume $`V_b(\mathrm{})=0`$. This route to the depletion potential was employed successfully in a grand canonical Monte Carlo simulation of a big sphere in a sea of small hard spheres near a hard wall . It is also the route used to obtain $`W(𝐫)`$ from experiment . Note that in the same limit $`\mu _b\mathrm{}`$ the density profiles of all other species $`\{\rho _{ib}(𝐫)\}`$ reduce to those of a $`\nu 1`$ component mixture.
### B The low density limit
In order to implement the formal result in Eq. (7) a way of determining the direct correlation function $`c_b^{(1)}`$ is required. It is convenient to adopt a DFT perspective. In density functional theory the intrinsic Helmholtz free energy functional can be divided into an ideal gas contribution plus an excess over the ideal gas contribution. While the former is known exactly, in general, only approximations are available for the excess part . One important exception is the excess free energy functional for a general mixture in the low density limit, i.e., in the limit of all densities going to zero. By means of a diagrammatic expansion it can be shown that the exact excess free energy functional in this limit is given by
$$\underset{\{\mu _i\mathrm{}\}}{lim}\beta _{ex}[\{\rho _i\}]=\frac{1}{2}\underset{i,j}{}\mathrm{d}^3r\mathrm{d}^3r^{}\rho _i(𝐫)\rho _j(𝐫^{})f_{ij}(𝐫𝐫^{}),$$
(10)
where $`f_{ij}`$ is the Mayer bond between a particle of species $`i`$ and one of species $`j`$.
For a binary mixture in the low density limit the depletion potential acting on a big particle $`b`$ can be calculated from Eq. (10) using the definition of the one-body direct correlation function given within density functional theory,
$$c_b^{(1)}(𝐫;\{\mu _i\})=\beta \frac{\delta _{ex}[\{\rho _i\}]}{\delta \rho _b(𝐫)},$$
(11)
and we obtain
$$\beta W(𝐫)=\mathrm{d}^3r^{}(\rho _s(𝐫^{})\rho _s(\mathrm{}))f_{bs}(𝐫𝐫^{}),$$
(12)
where $`s`$ refers to the small particles. In the same limit the density profile of the small particles reduces to the density profile of an ideal gas in the external potential $`V_s(𝐫)`$, i.e., $`\rho _s(𝐫)=\rho _s(\mathrm{})\mathrm{exp}(\beta V_s(𝐫))`$ and the depletion potential can be written as
$$\beta W(𝐫)=\rho _s(\mathrm{})\mathrm{d}^3r^{}(\mathrm{exp}[\beta V_s(𝐫^{})]1)f_{bs}(𝐫𝐫^{}).$$
(13)
This result is more familiar for the case of a binary hard-sphere mixture with sphere radii $`R_b`$ and $`R_s`$ where $`f_{bs}(𝐫𝐫^{})=\mathrm{\Theta }((R_b+R_s)|𝐫𝐫^{}|)`$, where $`\mathrm{\Theta }`$ is the Heaviside function. Then Eq. (13) reduces to the well-known Asakura-Oosawa depletion potential . As an example we consider the depletion potential between two big spheres; in this case $`\mathrm{exp}[\beta V_s(r)]1=f_{bs}(r)`$ and the sphere-sphere depletion potential can be expressed as
$`\beta W_{bb}^{AO}(h)`$ $`=`$ $`\rho _s(\mathrm{})\pi (2R_sh)\left\{R_s[R_b+{\displaystyle \frac{2}{3}}R_s]{\displaystyle \frac{h}{2}}[R_b+{\displaystyle \frac{R_s}{3}}]{\displaystyle \frac{h^2}{12}}\right\}\text{for }h<2R_s`$ (14)
$`=`$ $`0\text{for }h>2R_s,`$ (15)
where $`h`$ is the separation between the surfaces of the big hard spheres. As a second example we consider a big sphere near a planar, structureless hard wall. The depletion potential is then
$`\beta W_{wb}^{AO}(h)`$ $`=`$ $`2\rho _s(\mathrm{})\pi (2R_sh)\left\{R_s[R_b+{\displaystyle \frac{R_s}{3}}]{\displaystyle \frac{h}{2}}[R_b{\displaystyle \frac{R_s}{3}}]{\displaystyle \frac{h^2}{6}}\right\}\text{for }h<2R_s`$ (16)
$`=`$ $`0\text{for }h>2R_s,`$ (17)
where $`h`$ is the separation between the surface of the big hard sphere and the hard wall.
It is important to recognize that Eq. (13) provides the exact low density expression for the depletion potential even if the interactions between the species or between the wall and species $`s`$ are soft and possibly contain an attractive part so that the Mayer $`f`$ functions cannot be expressed in terms of the Heaviside function $`\mathrm{\Theta }`$. In general there is also a direct interaction potential between two big particles, or between a single big particle and a wall, so that the total effective potential, after integrating out the degrees of freedom of the small particles, is the sum of the intrinsic contribution – the depletion potential – and the direct interaction potential $`V_b(𝐫)`$, i.e.,
$$\mathrm{\Phi }_{tot}(𝐫)=W(𝐫)+V_b(𝐫).$$
(18)
For example the total effective potential between two big hard spheres in the sea of small hard spheres is $`\mathrm{\Phi }_{tot}(r)=W(r)+V_b(r)`$, with $`V_b(r)`$ the hard-sphere potential between the two big ones, and it is $`\mathrm{\Phi }_{tot}(r)`$ which constitutes the effective pair potential in the effective one-component Hamiltonian for the big spheres .
It is instructive to note that the functional given by the r.h.s. of Eq. (10) generates the appropriate depletion potential for the original Asakura-Oosawa model of a mixture of colloids and ideal, non-interacting polymers. This model binary mixture is specified by $`f_{cc}`$, $`f_{cp}`$, and $`f_{pp}`$, the Mayer $`f`$ functions describing the pairwise interactions between two colloids, between a colloid and a polymer, and between two polymers, respectively. $`f_{pp}`$ is set to zero in order to describe the ideal, non-interacting polymer coils. The resulting depletion potential is still given by Eq. (13) but this result now holds for all polymer densities $`\rho _s(\mathrm{})`$ not just in the dilute limit $`\rho _s(\mathrm{})0`$, because the polymer is taken to be ideal. The total effective potential between two colloids is then given by Eq. (18), with the Asakura-Oosawa result (14) for $`W(r)`$, which may be employed in an effective Hamiltonian for the colloids .
### C Asymptotic behavior
In the previous subsection we showed that for the case of hard spheres the depletion potential reduces in the low density limit to the Asakura-Oosawa result. Examination of Eqs. (14) and (16) shows that this potential is identically zero for separations $`h`$ between the spheres or between the sphere and the wall that are greater than 2$`R_s\sigma _s`$. Outside the low density limit of the small particles this is no longer valid. From the general theory of the asymptotic decay of correlations it is known that for systems in which the interatomic forces are short-ranged, i.e., excluding power-law decay, the density profiles of both components of a binary mixture exhibit a common damped oscillatory form in the asymptotic regime, far from the wall or fixed particle, which is determined fully by the pole structure of the total pair correlation functions $`h_{ij}(r)`$ of the bulk mixture. The depletion potential is related to the density profile of species $`b`$ via Eq. (9) and therefore its asymptotic behavior should be related directly to that of this density profile. In order to understand this connection in more detail we first recall some arguments from Ref..
A bulk binary mixture consisting of small particles of density $`\rho _s^{bulk}`$ and big particles of density $`\rho _b^{bulk}`$ is considered. The total correlation functions in the bulk, $`h_{ij}(r)`$, with $`i,j=s,b`$ are related to the radial distribution functions $`g_{ij}(r)`$ via $`h_{ij}(r)=g_{ij}(r)1`$ and to the two-body direct correlation functions $`c_{ij}^{(2)}(r)`$ via the Ornstein-Zernike relation for mixtures. In Fourier space the latter can be expressed as
$$\widehat{h}_{ij}(q)=\frac{\widehat{N}_{ij}(q)}{\widehat{D}(q)}$$
(19)
where $`\widehat{h}_{ij}(q)`$ is the 3 dimensional Fourier transform of $`h_{ij}(r)`$, the numerator is given by
$$\begin{array}{ccc}\widehat{N}_{aa}(q)\hfill & =& \widehat{c}_{aa}^{(2)}(q)+\rho _b^{bulk}(\widehat{c}_{ab}^{(2)}(q)^2\widehat{c}_{aa}^{(2)}(q)\widehat{c}_{bb}^{(2)}(q)),\hfill \\ \widehat{N}_{bb}(q)\hfill & =& \widehat{c}_{bb}^{(2)}(q)+\rho _a^{bulk}(\widehat{c}_{ba}^{(2)}(q)^2\widehat{c}_{aa}^{(2)}(q)\widehat{c}_{bb}^{(2)}(q)),\hfill \\ \widehat{N}_{ab}(q)\hfill & =& \widehat{c}_{ab}^{(2)}(q),\hfill \end{array}$$
(20)
and
$$\widehat{D}(q)=(1\rho _s^{bulk}\widehat{c}_{ss}^{(2)}(q))(1\rho _b^{bulk}\widehat{c}_{bb}^{(2)}(q))\rho _s^{bulk}\rho _b^{bulk}\widehat{c}_{sb}^{(2)}(q)^2$$
(21)
is a common denominator. The total correlation function in real space can be obtained by taking the inverse Fourier transform:
$$rh_{ij}(r)=\frac{1}{2\pi ^2}\underset{0}{\overset{\mathrm{}}{}}dqq\mathrm{sin}(qr)\widehat{h}_{ij}(q)$$
(22)
which can then be evaluated by means of the residue theorem. If $`q_n`$ denotes the $`n`$-th pole in the upper complex q half-plane and $`R_n`$ the corresponding residue of $`q\widehat{h}_{ij}(q)`$, the total correlation function can be written as :
$$rh_{ij}(r)=\frac{1}{2\pi }\underset{n}{}e^{iq_nr}R_n.$$
(23)
From this equation it becomes clear that the asymptotic behavior of $`h_{ij}(r)`$ is dominated by the pole or poles $`q_n`$ with the smallest imaginary part, since this gives rise to the slowest exponential decay.
For all pairs $`i,j=b,s`$ the poles are determined by the condition $`\widehat{D}(q)=0`$ . For a binary mixture in the dilute limit of the big particles, i.e., $`\rho _b^{bulk}0`$, the general theory of the asymptotic decay simplifies considerably and from Eq. (21) we see that the pole structure of all three total correlation functions can be obtained from the solutions of the equation
$$1\rho _s^{bulk}\widehat{c}_{ss}^{(2)}(q)=0,$$
(24)
with $`\widehat{c}_{ss}^{(2)}(q)`$ referring to the fluid of pure $`s`$ at density $`\rho _s^{bulk}`$. In general there will be an infinite number of solutions of Eq. (24), but only the solution $`q_nq=a_1+ia_0`$ with the smallest imaginary part $`a_0`$ is important for the following. The asymptotic behavior of the radial distribution functions $`\rho _i(r)/\rho _i^{bulk}`$ of species $`i`$ around a fixed particle, which can be a small or a big one, can be ascertained and it follows that the density profiles exhibit asymptotic decay of the form
$$\rho _i(r)\rho _i^{bulk}\frac{A_{pi}}{r}\mathrm{exp}(a_0r)\mathrm{cos}(a_1r\mathrm{\Theta }_{pi}),r\mathrm{},$$
(25)
with a common characteristic inverse decay length $`a_0^1`$ and wavelength of oscillations $`2\pi /a_1`$ for both species $`i=s,b`$. Remarkably, exactly the same inverse decay length and wavelength also characterize the asymptotic decay of the density profiles close to a planar wall. This is given by
$$\rho _i(z)\rho _i^{bulk}A_{wi}\mathrm{exp}(a_0z)\mathrm{cos}(a_1z\mathrm{\Theta }_{wi}),z\mathrm{},$$
(26)
for $`i=s,b`$. The amplitudes $`A_{pi}`$ and $`A_{wi}`$ and the phases $`\mathrm{\Theta }_{pi}`$ and $`\mathrm{\Theta }_{wi}`$ do depend on species $`i`$ and whether a particle or wall is the source of the external potential. Note that from Eq. (24) it follows that in the dilute limit for the big particles $`a_0`$ and $`a_1`$ are functions of the packing fraction of the small particles only; thus they do not depend on the size ratio.
The asymptotic behavior of the depletion potential can now be obtained from Eq. (9). Assuming that the external potential acting on the big spheres is of finite range the depletion potential between two big spheres has an asymptotic behavior of the form
$`\beta W(r)`$ $``$ $`\mathrm{ln}\left(1+{\displaystyle \frac{A_{pb}}{r}}\mathrm{exp}(a_0r)\mathrm{cos}(a_1r\mathrm{\Theta }_{pb})/\rho _b^{bulk}\right)`$ (27)
$``$ $`{\displaystyle \frac{A_{pb}}{r}}\mathrm{exp}(a_0r)\mathrm{cos}(a_1r\mathrm{\Theta }_{pb})/\rho _b^{bulk},r\mathrm{},`$ (28)
and that between a single big sphere and a planar wall takes the form
$$\beta W(z)A_{wb}\mathrm{exp}(a_0z)\mathrm{cos}(a_1z\mathrm{\Theta }_{wb})/\rho _b^{bulk},z\mathrm{}.$$
(29)
We shall see later that our DFT results for the density profiles and the depletion potential conform with these asymptotic results at very large separations and, strikingly, at intermediate separations.
### D A density functional approach for hard spheres
For the system of primary interest, namely the mixture of hard spheres, a very reliable DFT exists, namely the Rosenfeld fundamental measures functional . While, in principle, this functional also can treat generally shaped convex hard particles , its application has been restricted to the particular cases of hard spheres and parallel hard cubes .
In the low density limit the Rosenfeld functional reduces to the exact excess free energy functional of Eq. (10). For arbitrary densities it has the following structure:
$$_{ex}[\{\rho _i\}]=\mathrm{d}^3r\mathrm{\Phi }(\{n_\alpha (𝐫)\}),$$
(30)
where $`\mathrm{\Phi }`$ is a function of a set of weighted densities $`\{n_\alpha \}`$ which are defined by
$$n_\alpha (𝐫)=\underset{i=1}{\overset{\nu }{}}\mathrm{d}^3r^{}\rho _i(𝐫^{})\omega _i^\alpha (𝐫𝐫^{}).$$
(31)
The weight functions $`\omega _i^\alpha `$ in Eq. (31) depend only on the geometrical features, the so-called fundamental measures, of species $`i`$. Explicit expressions for the weight functions of hard-sphere mixtures and for $`\mathrm{\Phi }`$ can be found in Ref. and in Ref. . The Rosenfeld functional has the following properties: i) the free-energy of the homogeneous mixture is identical to that from Percus-Yevick or scaled-particle theory and ii) the pair direct correlation functions of the homogeneous hard-sphere mixture, generated by functional differentiation of $`_{ex}`$, are identical to those of Percus-Yevick theory. The index $`\alpha `$ labels 4 scalar plus 2 vector weights . While the original functional given in Ref. did not account for the freezing transition of pure hard spheres, more sophisticated extensions do account for freezing; the weight functions remain the same but $`\mathrm{\Phi }`$ is changed slightly. For the depletion potential problems under consideration the different versions give almost identical results for bulk packing fractions $`\eta _s0.3`$ . At higher packing fractions the density profiles of the small spheres $`\rho _s(z)`$ close to a hard planar wall, or $`\rho _s(r)`$ close to a fixed particle, do display small deviations between the different versions of the Rosenfeld functional. Moreover, when calculating the depletion potential for size ratios of $`s=0.1`$ or smaller, these deviations are amplified and one observes slightly smaller amplitudes of oscillation for the more sophisticated versions of the theory. An example is given in Subsection V B (see, c.f., Fig. 11).
The one-body direct correlation function, defined within density functional theory by Eq. (11), can be written as
$$c_b^{(1)}(𝐫;\{\mu _i\})=\underset{\alpha }{}\mathrm{d}^3r^{}\left(\frac{\beta \mathrm{\Phi }(\{n_\alpha \})}{n_\alpha }\right)_𝐫^{}\omega _b^\alpha (𝐫^{}𝐫)$$
(32)
for the Rosenfeld functional. In the limit where all species have the same radius it is easy to check that the weighted densities $`n_\alpha `$ in Eq. (31), and hence $`\mathrm{\Phi }`$, reduce to the corresponding quantities for the pure fluid and, since the weight functions $`\omega _b^\alpha `$ in Eq. (32) reduce to the weight functions of the pure system, $`c_b^{(1)}`$ reduces to the one-body direct correlation function of the pure ($`s`$) fluid. The depletion potential is then given by $`\beta W(r)=\mathrm{ln}(\rho _s(r)/\rho _s(\mathrm{}))`$, which is the correct result .
Defining functions $`\mathrm{\Psi }^\alpha `$ as
$$\mathrm{\Psi }^\alpha (𝐫^{})\left(\frac{\beta \mathrm{\Phi }(\{n_\alpha \})}{n_\alpha }\right)_𝐫^{}\left(\frac{\beta \mathrm{\Phi }(\{n_\alpha \})}{n_\alpha }\right)_{\mathrm{}},$$
(33)
the grand potential difference in Eq. (6) can be written as a sum of convolutions of these functions with the weight functions of species $`b`$:
$$\beta W_t(𝐫)=\underset{\alpha }{}\mathrm{d}^3r^{}\mathrm{\Psi }^\alpha (𝐫^{})\omega _b^\alpha (𝐫^{}𝐫).$$
(34)
This expression is valid for arbitrary densities. The dilute limit of species $`b`$ can now be taken, within the Rosenfeld functional, by considering the weighted densities Eq. (31) which in this case reduce to
$$n_\alpha ^{dilute}(𝐫)=\underset{ib}{}\mathrm{d}^3r^{}\rho _i(𝐫^{})\omega _i^\alpha (𝐫𝐫^{}),$$
(35)
where the set of density profiles $`\{\rho _i(𝐫)\}`$ that enter Eq. (35) is that of the $`\nu 1`$ component fluid, i.e., the one obtained after taking the limit. It follows that the Helmholtz free energy in Eq. (30) and, consequently, the functions $`\mathrm{\Psi }^\alpha `$ in Eq. (33) are those of a $`\nu 1`$ component mixture. Species $`b`$ enters into the calculation of the depletion potential, i.e., the dilute limit of Eq. (34), only through its geometry, i.e., via the weight functions $`\omega _b^\alpha `$.
This feature of the theory becomes especially important if the number of components $`\nu `$ is small. In the particular case of a binary mixture, $`\nu =2`$, the minimization of the functional in the dilute limit reduces to the minimization of the functional of a pure fluid and the weighted densities depend only on the density profile $`\rho _s(𝐫)`$ of the small spheres:
$$n_\alpha ^{dilute}(𝐫)=\mathrm{d}^3r^{}\rho _s(𝐫^{})\omega _s^\alpha (𝐫𝐫^{}).$$
(36)
Although the direct approach of calculating the depletion potential via evaluating grand potential differences by brute force requires only a functional describing the pure fluid, the above considerations demonstrate that our present approach based on the one-body direct correlation function of the big spheres in a sea of small ones requires a functional that describes the binary mixture.
## III Results from the DFT approach and assessment of their accuracy
In this section we examine the accuracy of some of the approximations inherent in the present DFT approach by comparing our DFT results for the depletion potential with those of simulations and with the predictions of the general asymptotic theory given in Subsection II C.
### A Consistency check
We consider first the results of two separate routes to obtaining the dilute limit for the case of a binary hard-sphere mixture. In the first route both components of the mixture are treated on equal footing so that one calculates both $`\rho _b(𝐫)`$ and $`\rho _s(𝐫)`$ and obtains $`W_t(𝐫)`$ using Eq. (34). By requiring the chemical potential of the big spheres $`\mu _b`$ to become more and more negative, the bulk density $`\rho _b(\mathrm{})`$ of this component approaches zero and the dilute limit is taken numerically. For all the mixtures we investigated, a bulk packing fraction of the big spheres of $`\eta _b=10^4`$ was sufficiently small to ensure that the density profile of the small spheres is indistinguishable from that of a pure fluid at the same $`\eta _s`$. Moreover the convergence of $`\rho _s(𝐫)`$ and $`W_t(𝐫)`$ to their limiting values is rather fast; an explicit example is given in Fig. 1 of an earlier Letter . Using the second route, employing the weighted densities of Eq. (36), the dilute limit is taken directly in the functional. In Fig. 1 depletion potentials corresponding to both routes are shown for a big hard sphere near a planar wall and a size ratio $`s=R_s/R_b=0.1`$. The bulk packing fraction of the small spheres is $`\eta _s=0.3`$. We find excellent agreement between the two sets of results. The same level of agreement is found for a wide range of size ratios $`s`$ and packing fractions $`\eta _s`$. From this we conclude that the limit can be taken directly in the functional, which makes the calculations significantly easier to perform, and all the results for the depletion potential we present subsequently will be based on this route.
### B Comparison with simulation data
The results presented in Fig. 1 test the self-consistency of the two routes to the dilute limit within the given DFT approach. In order to test the accuracy of approximations introduced by employing the Rosenfeld functional the results of the present approach are compared with those of simulations. Fortunately some independent sets of simulation results for depletion potentials are available for both the sphere-sphere and the wall-sphere case. In Fig. 2 the depletion potentials between two big spheres in a sea of small spheres, at a size ratio of $`s=0.1`$ and various packing fractions up to $`\eta _s=0.6\pi /60.314`$, obtained from the molecular dynamics simulations of Ref. are compared with results of the present DFT approach. The agreement between the latter and the simulations is generally very good. At the higher packing fractions small deviations can be seen near contact and near the first minimum but the agreement is within the error bars of the simulations which are not indicated here. We note that in Ref. the depletion force was the quantity measured in the simulations and the depletion potential was calculated by integrating a smoothed force. For a higher packing fraction, $`\eta _s=0.7\pi /60.367`$, the agreement between the depletion potential obtained in the simulations of Ref. and our present result is poorer (not shown in Fig. 2), but for this large value of $`\eta _s`$ the error bars of the simulations are probably bigger than for small values of $`\eta _s`$ .
In Ref. the depletion potential for a single big hard sphere near a planar hard wall calculated within DFT was compared with the results of two independent sets of simulations for a size ratio $`s=0.2`$ and a packing fraction $`\eta _s=0.3`$. Very good agreement was found. In Fig. 3 we present a comparison of our results with simulation results from Ref. for size ratios $`s=0.2`$ (a) and $`s=0.1`$ (b), for various packing fractions of the small spheres up to $`\eta _s=0.3`$. The original simulation results did not oscillate around $`W=0`$ which led us to follow the procedure described in Ref. and to shift the data by a small constant amount in order to match the contact values with those of our DFT result. We note that in the simulations of Ref. the depletion force was measured and the depletion potential was obtained by integrating the force. Since the data for the force are available only for $`hh_{max}`$, the integral depends on the cut-off $`h_{max}`$.We surmise that this cut-off dependence is responsible for $`W(h)`$ not oscillating around zero. The agreement between our DFT results and those of the shifted simulation data is very good. The differences probably lie within the error bars of the simulations, for all packing fractions when $`s=0.2`$, and for $`\eta _s=0.1`$ and $`\eta _s=0.2`$ when $`s=0.1`$. However, for $`\eta _s=0.3`$ and $`s=0.1`$ clear deviations remain between our results and those of the simulations. In this case the shifted simulation data for the depletion potential are close to the DFT results for $`h<\sigma _s`$ – the height and position of the first maximum are the same – but, in contrast to the DFT results, the simulation data do not oscillate around zero. Clearly some alternative procedure for interpreting the simulation data is required.
### C Density profiles
In Fig. 4 the number density profiles of a binary hard-sphere mixture near a planar hard wall as obtained from DFT are shown for three size ratios. The packing fraction of the small spheres is $`\eta _s=0.3`$ and that of the big spheres is $`\eta _b=10^4`$. The latter is sufficiently low that the density profiles of the small spheres, $`\rho _s(z)`$, shown in Fig. 4(a), are practically equal to that corresponding to pure small spheres. Therefore, they are indistinguishable for all size ratios. Because of the hard-body interaction between the small spheres and the wall the density profile $`\rho _s(z)`$ exhibits a discontinuous fall to zero at $`z=\sigma _s/2`$. The density profiles of the big spheres for size ratios $`s=0.1`$ (full line), 0.1333 (dotted line) and 0.2 (dashed line) are shown in Fig. 4(b). These density profiles do differ significantly for different size ratios. The hard-body interaction between the wall and the big spheres does not allow their centers to encroach closer than $`z=\sigma _b/2`$ and we find that the contact value is very different in all three cases (see caption to Fig. 4). We note that the wavelength of the oscillations in both $`\rho _s(z)`$ and $`\rho _b(z)`$ is approximately $`\sigma _s`$. In order to display the asymptotic behavior of these density profiles the logarithm of the difference between each density and its bulk value is shown in Fig. 5. For $`z/\sigma _s2`$ these plots conform very closely to the asymptotic form given by Eq. (26). Straight lines joining the maxima have a common slope and the distance between adjacent maxima is the same in all cases. Only the amplitudes of the oscillations in $`\rho _b(z)`$ differ for different values of $`s`$. It follows that the decay length, $`a_0^1`$, and the wavelength of the oscillation, $`2\pi /a_1`$, are the same for both density profiles, i.e., for the big and the small spheres, and are independent of the size ratio. We have confirmed that the same values for $`a_0`$ and $`a_1`$ are obtained from plots of the density profiles of the same binary mixture in the presence of a fixed big hard sphere, i.e., our results are consistent with Eq. (25). At high packing fractions of the small spheres, e.g., $`\eta _s=0.42`$, we can easily resolve up to 25 damped oscillations. At long range we find the calculated density profiles to be in excellent agreement with the predictions of the theory of asymptotic decay. As the amplitude of the 24th oscillation is smaller than the amplitude of the first one by a factor of approximately $`5\times 10^6`$ this attests further to the high numerical accuracy of our results. In addition we confirmed numerically that the modifications of the Rosenfeld functional which we employed lead to the same asymptotic behavior of $`W(h)`$ as the original functional .
### D Asymptotic behavior
The asymptotic behavior of the depletion potential calculated within DFT is shown in Fig. 6. For $`z/\sigma _s2`$ our results conform very closely to Eq. (29): although the amplitude of the oscillations depends on $`s`$, $`W(z)`$ is characterized by the same, common decay length $`a_0^1`$ and wavelength $`2\pi /a_1`$ which describe the density profiles of the mixture. The results displayed in Figs. 5 and 6 indicate that the asymptotic behavior of the density profiles and of the depletion potential set in at rather small distances from the wall. For wall-sphere surface separations of typically $`z2\sigma _s`$, or even smaller, the asymptotic formulae are already remarkably accurate. This is in keeping with the results of earlier studies of the bulk pairwise correlation functions of hard-sphere mixtures , where leading-order asymptotics were shown to be accurate down to second-nearest neighbor separations. We shall make use of this observation in a later section in which we develop an explicit parameterized form for the depletion potential.
As a final examination of the validity of the asymptotic analysis we calculated values of $`a_0`$ and $`a_1`$ from plots of (the logarithm of) the density profiles of the small spheres $`\rho _s(z)`$ near the planar hard wall (see Fig. 5) for a range of values for $`\eta _s`$ and various size ratios from $`s=0.5`$ to $`s=0.1`$. In accordance with the above statement the results for $`a_0`$ and $`a_1`$ do not depend on $`s`$ and are shown in Fig. 7 together with the values obtained using the Percus-Yevick result for $`c_{ss}^{(2)}(r)`$ in the pure fluid to solve Eq. (24) for the poles $`q=a_1+ia_0`$; we recall that the Rosenfeld hard-sphere functional generates the Percus-Yevick two-body direct correlation functions for a bulk mixture . For small packing fractions $`\eta _s`$ the oscillations are damped very rapidly, i.e., the decay length $`a_0^1`$ is small, so that the numerical determination of the wavelength from a density profile is quite difficult. Nevertheless, the level of agreement between the two sets of results is very good, for all values of $`\eta _s`$ that were considered, confirming that the DFT results are consistent with the general predictions for the asymptotic behavior.
Our approach predicts depletion potentials for both the wall-sphere and the sphere-sphere case which are in very good agreement with simulations for distances close to contact and which are consistent with predictions of the general theory of the asymptotic decay of correlations in hard-sphere mixtures for distances away from contact. From our comparisons we conclude that our approach yields accurate results in the whole range of distances, for packing fractions up to (at least) $`\eta _s=0.3`$ and for size ratios down to (at least) $`s=0.1`$. We emphasize that the full structure of the depletion potential, which is correctly described by the present approach, is not captured by the Asakura-Oosawa approximation or by a truncated virial expansion .
### E Large asymmetries
So far it is not apparent how well our present approach will fare for extreme asymmetries, i.e., for $`s1`$. The Rosenfeld functional, which is the density functional we apply for all of the calculations of the depletion potentials, is designed to treat a multi-component hard-sphere mixture with arbitrary inhomogeneities. Whilst its accuracy in describing the density profiles for a pure fluid and for binary mixtures at moderate packing fractions and moderate size ratios has been confirmed by comparison with simulation results, a highly asymmetric binary mixture has not yet been studied systematically using this functional. Thus it is not known for which size ratios the results calculated with this functional are accurate. We recall that the Percus-Yevick approximation becomes increasingly less accurate for bulk properties as $`s0`$, but here we are interested, in particular, with the reliability of our approach for determining depletion potentials. The latter are obtained from density profiles, having taken the dilute limit of one of the species \[see Eq. (9)\].
In this context it is instructive to consider the depletion potential between a hard sphere of radius $`R_1`$ and one of radius $`R_2`$ in a sea of small hard spheres of radius $`R_s`$ at a packing fraction $`\eta _s`$. This system is formally a mixture of three components in which two are dilute. The radius ratio $`R_s/R_1`$ is chosen such that on the basis of our previous results we know that the Rosenfeld functional can treat a mixture of species 1 and $`s`$ accurately. On the other hand, the radius $`R_2`$ is chosen to be much bigger than $`R_s`$ and $`R_1`$. The depletion potential can be calculated in two different ways. In the first route, sphere $`2`$ (with large radius $`R_2`$) is fixed and enters into the calculation as an external potential for sphere 1 and species $`s`$. The density profile of the small spheres in the presence of this external potential can be calculated, and from it the depletion potential, using the theoretical approach described in Sec. II. Thus, in this calculation the Rosenfeld functional treats a mixture with a moderate size ratio $`R_s/R_1`$ exposed to an external potential. Therefore we expect these results to be very accurate. In the second route, the roles of spheres 1 and 2 are exchanged. The sphere of medium radius $`R_1`$ is fixed and acts as an external potential for the very large sphere 2 and the small species $`s`$. Now the Rosenfeld functional must treat a very asymmetric mixture. Of course, in an exact treatment of this problem it does not matter which sphere is fixed first as the depletion potential is simply the difference in the grand potential between a configuration in which spheres 1 and 2 are fixed and positioned close to each other and one in which both spheres are at infinite separation. The result of an exact treatment cannot depend upon which sphere is regarded as an external potential. However, it is not immediately obvious that the underlying symmetry is respected in an approximate DFT treatment. At first sight one way of calculation might appear to be less demanding on the theory than the other.
In Fig. 8 we show the depletion potential between the big spheres when sphere $`2`$ is fixed (solid line) and then with sphere $`1`$ fixed (symbols), for a packing fraction of the small spheres $`\eta _s=0.3`$, and $`R_1=5R_s`$. Two different values of $`R_2`$ are considered, namely $`R_s/R_2=0.02`$ (a) and $`R_s/R_1=0.01`$ (b). We find excellent agreement between the results of the two routes for both (a) and (b). Only very small differences between the curves can be ascertained and these occur for separations close to contact where numerics are most difficult. Note that the results in Figs. 8(a) and 8(b) lie very close to each other. This can be understood easily when $`R_2`$ is the fixed sphere. For cases (a) and (b), $`R_2R_1=5R_s`$ and one is effectively in the planar-wall limit so that both sets of results lie close to those in Fig. 3(a), with $`\eta _s=0.3`$. From the results shown in Fig. 8 and further comparisons for other values of $`\eta _s`$ it is evident that the Rosenfeld functional does maintain the required symmetry between 1 and 2. It is important to understand this. In the low density limit, i.e., if the ternary mixture is considered with the density of the small spheres also approaching zero, the depletion potential can be expressed in terms of Mayer $`f`$ functions and the equivalence of the two routes can be verified directly. Starting from the functional given in Eq. (10) we follow the derivation of Eq. (13) and obtain
$$\beta W_{21}(r)=\rho _s(\mathrm{})\mathrm{d}^3r^{}f_{2s}(𝐫^{})f_{s1}(𝐫𝐫^{}),$$
(37)
for the depletion potential with sphere 2 fixed and
$$\beta W_{12}(r)=\rho _s(\mathrm{})\mathrm{d}^3r^{}f_{1s}(𝐫^{})f_{s2}(𝐫𝐫^{})$$
(38)
for that with sphere 1 fixed. Here $`f_{1s}`$ and $`f_{2s}`$ are the Mayer $`f`$ functions between a small sphere and sphere 1 and 2, respectively, and it is evident that $`W_{21}(r)W_{12}(r)`$. The Rosenfeld functional will reproduce this result for packing fractions $`\eta _s0`$, since it reduces to Eq. (10) in this limit. For arbitrary values of $`\eta _s`$ it is necessary to reconsider the genesis of the functional and recognize that although the hard-sphere pairwise potentials $`\mathrm{\Phi }_{ij}(r)`$ between species $`i`$ and $`j`$ do not enter explicitly, the functional does respect the equivalence of $`\mathrm{\Phi }_{ij}(r)`$ and $`\mathrm{\Phi }_{ji}(r)`$; the Mayer functions and the weight functions which were used in constructing the functional are symmetric w.r.t. $`i`$ and $`j`$. It is straightforward to show that the equivalence of $`W_{21}(r)`$ and $`W_{12}(r)`$ is guaranteed provided the functional respects this symmetry. Thus, the two sets of results shown in Fig. 8(a) should agree with each other, as should those shown in Fig. 8(b). That there are small discrepancies reflects only numerical inaccuracies rather than any fundamental shortcoming of the DFT approach. It is pleasing that what appear to be two distinct ways of calculating the depletion potential yield the same results, even for high degrees of asymmetry. Whether other functionals, not based on fundamental measure theory, will respect the symmetry requirements remains to be ascertained. Although the present calculations should be regarded as a further test of the internal consistency of our approach rather than a formal demonstration that it is accurate for extreme asymmetries, the results, when coupled with the excellent agreement between theory and simulations for $`s=0.1`$, do suggest that the approach should remain accurate for smaller size-ratios.
## IV Derjaguin approximation
In the well-known Derjaguin approximation the force between two large convex bodies is expressed in terms of the interaction energy of two parallel plates. This approximate mapping is valid in the limit where the minimal separation of surfaces $`h`$ is much smaller than the radii of curvature and was developed assuming the force between the surfaces can be calculated by integration over all interactions between pairs of points of the two bodies. Recently the Derjaguin approximation was implemented for the depletion force between two big hard spheres in a sea of small hard spheres, employing a truncated virial expansion to calculate the excess pressure of the small spheres between planar hard walls, and results were compared with simulation data for a size ratio $`s=0.1`$ . There is, however, an important conceptual difference from earlier applications of the Derjaguin approximation as depletion effects are global effects arising from packing of the small spheres and it is not obvious that the original derivation remains applicable or what the regime of validity of the approximation should be. Some of its limitations were discussed in Ref. where it was argued that the Derjaguin approximation should not be reliable for $`s=0.1`$ if the packing fraction $`\eta _s0.3`$. Here we examine some of the key predictions of the Derjaguin approximation by making comparison with results of our DFT approach. From the arguments of Subsec. III E it is safe to assume that the present DFT approach remains reliable for rather large size ratios where the Derjaguin approximation might be expected to be valid.
There is an elegant scaling relation connecting the depletion force, $`F(h)W(h)/h`$, between two big spheres, $`F_{bb}(h)`$, in a sea of small spheres with that between a single big sphere and a planar hard wall, $`F_{wb}(h)`$. In the limit of infinite asymmetry, $`s0`$, the forces are equal except for a factor of $`2`$, i.e.,
$$2F_{bb}(h)=F_{wb}(h),$$
(39)
with $`h`$ the minimal separation of the surfaces of the two big objects. This scaling relation follows directly from the Derjaguin approximation and if it is found to be obeyed it is sometimes inferred that the Derjaguin approximation itself is valid. However, it was shown that this scaling relation follows from geometrical considerations without introducing the explicit Derjaguin approximation. This can be illustrated by comparing the explicit Asakura-Oosawa depletion potentials \[see Eqs. (14) and (16)\]. In the limit $`R_bR_s`$ both formulae reduce to $`\frac{\epsilon }{2}\rho _s(\mathrm{})\pi R_b(2R_sh)^2`$, for $`h<2R_s`$, where $`\epsilon =1`$ corresponds to the sphere-sphere and $`\epsilon =2`$ to the wall-sphere case. Thus, achieving the correct scaling property in Eq. (39) does not prove that the Derjaguin approximation,
$$F^{Derj}(h)=\epsilon \pi (R_b+R_s)\underset{h}{\overset{\mathrm{}}{}}dLf_s(L)$$
(40)
is accurate. Here $`f_s(L)`$ is the solvation force, or the excess pressure, for the small-sphere fluid confined between two planar parallel hard walls separated by a distance $`L`$ .
From our DFT calculations we find that the scaling relation is already well-obeyed at moderate size ratios. In Fig. 9 the scaled depletion force $`\beta f_{bb}^{}(h)=2\beta F_{bb}(h)R_s^2/(R_b+R_s)`$ between two big hard spheres (solid line) and that between a single big hard sphere and a planar hard wall ($`\mathrm{}`$), $`\beta f_{wb}^{}(h)=\beta F_{wb}(h)R_s^2/(R_b+R_s)`$, in a sea of small hard spheres at a packing fraction of $`\eta _s=0.3`$ is shown for size ratios $`s=0.1`$ and 0.02. While small deviations from the scaling relation in Eq. (39) are visible close to contact for $`s=0.1`$, these deviations have almost disappeared for $`s=0.05`$ (not shown in the figure) and near perfect agreement is found for $`s=0.02`$. Note also that the scaled depletion forces corresponding to the different values of $`s`$ lie close to each other.
An explicit result of the Derjaguin approximation \[Eq. (40)\] is that the depletion force between two big spheres or between a big sphere and a planar wall can be written as
$$F^{Derj}(h)=\epsilon \pi (R_b+R_s)(p(\eta _s)(h2R_s)\gamma (\eta _s)),h<2R_s,$$
(41)
where $`p(\eta _s)`$ is the bulk pressure of the small spheres and $`\gamma (\eta _s)`$ is twice the surface tension of the small-sphere fluid at a planar hard wall. The geometrical factor $`\epsilon `$ is the same as in Eq. (40). Thus, for a given size ratio the slope of the depletion force predicted by the Derjaguin approximation is constant for $`h<2R_s`$ and depends only on the equation of state of the small spheres $`p(\eta _s)`$. For the particular case of hard spheres we obtain:
$$\frac{\mathrm{d}\beta F^{Derj}(h)}{\mathrm{d}h}=\epsilon (R_b+R_s)\frac{3\eta _s}{4R_s^3}\frac{1+\eta _s+\eta _s^2\eta _s^3}{(1\eta _s)^3},h<2R_s,$$
(42)
where the quasi-exact Carnahan-Starling equation of state was used.
However, the depletion forces calculated within the present approach show a qualitatively different behavior from that predicted by Eq. (42). It was found that even for small size ratios ($`s0.05`$), only in the limit $`\eta _s0`$, in which the Asakura-Oosawa approximation becomes exact, there is agreement between the Derjaguin approximation and the results of our approach. The depletion force calculated at a packing fraction of $`\eta _s=0.3`$ does not have constant slope for $`h<2R_s`$ (see Fig. 9). This is in clear contradiction to Eq. (42). Simulation results for the depletion force also exhibit non-constant slopes for $`h<2R_s`$.
Another prediction of the Derjaguin approximation in Eq. (41) is that the contact value $`W(0)`$ of the depletion potential can be expressed simply in terms of the equation of state $`p(\eta _s)`$ and the surface tension $`\gamma (\eta _s)`$. Using the Carnahan-Starling result for $`p(\eta _s)`$ and the scaled particle result for $`\gamma (\eta _s)`$ we obtain :
$$\beta W^{Derj}(0)=\frac{\epsilon (R_b+R_s)3\eta _s}{2R_s}\frac{12\eta _s2\eta _s^2\eta _s^3}{(1\eta _s)^3}$$
(43)
which becomes positive at high packing fractions of the small spheres . Provided $`\eta _s<0.2`$ the contact values from Eq. (43) are in reasonable agreement with the results of our DFT approach for small values of $`s`$. However, the contact values obtained from Eq. (43) change sign at $`\eta _s0.3532`$, which is in complete contradiction to the results of the present approach where we find negative contact values for all packing fractions and all size ratios $`s`$ under consideration .
In Ref. it was shown that a third-order virial expansion (in powers of $`\eta _s`$) for the depletion potential calculated within the Derjaguin approximation does not yield positive contact values. However, expansion to fourth or fifth-order shows a qualitatively different behavior from third order and already indicates the onset of positive $`W(0)`$. Thus, in keeping with Ref. we conclude that the Derjaguin approximation is not very useful for the calculation of depletion forces. The good level of agreement, observed for $`h<2R_s`$, between the results of the third-order virial expansion and those of simulation for $`s=0.1`$ should be regarded fortuitous.
## V Applications
### A A parametrized form for the depletion potential of hard spheres
As mentioned in Subsec. II A, recent studies of correlation functions and phase equilibria of highly asymmetric binary mixtures have shown that it is very advantageous to map such mixtures onto effective one-component fluids . The effective pairwise potential between the big particles is then the bare pair potential between two big particles plus the depletion potential \[see Eq. (18)\]. Thus, in calculating the phase behavior of binary hard-sphere mixtures it is necessary to adopt a specific form for the depletion potential between two big hard spheres. Previous simulation studies of binary mixtures have employed the simplified third-order virial expansion formula given by Götzelmann et al. and the same potential has been used in a perturbation theory treatment of the phase behavior . Although this formula is convenient for global investigations of phase behavior, as the depletion potential is given explicitly as a function of $`\eta _s`$, clearly it would be valuable to have a simple, parametrized form for the depletion potential that (i) is better founded than the formula provided by Götzelmann et al. and (ii) captures the correct intermediate and long-range oscillatory structure as well as the important short-range features. Note that in Refs. and the effective pair potential was set equal to zero for separations $`h>2R_s`$.
We have used depletion potentials calculated within the present DFT approach, for a single big sphere near a planar hard wall and for two big spheres, to develop a suitable parameterization scheme. Although this parameterization is fairly simple it yields rather accurate fits. The depletion potential close to contact is fitted by a polynomial and is continued by the known asymptotic behavior. In the following the variable $`x`$ measures the minimal distance from contact in units of the small sphere diameter $`\sigma _s`$, i.e., $`xh/\sigma _s`$. These parametrized depletion potentials $`\overline{W}`$ are also scaled: the actual potentials $`W`$ are recovered by multiplying by a factor of $`\epsilon (R_b+R_s)/(2R_s)`$ with $`\epsilon =2`$ for the wall-sphere and $`\epsilon =1`$ for the sphere-sphere potential:
$$W=\frac{\epsilon (R_b+R_s)}{2R_s}\overline{W}.$$
(44)
Between contact at $`x=0`$ and the location $`x_0`$ of the first maximum the scaled depletion potential is fitted by a cubic polynomial:
$$\beta \overline{W}(x,\eta _s)=a(\eta _s)+b(\eta _s)x+c(\eta _s)x^2+d(\eta _s)x^3,x<x_0,$$
(45)
where the coefficients $`a`$, $`b`$, $`c`$ and $`d`$ are functions of the packing fraction of the small spheres $`\eta _s`$. More details of this polynomial and the determination of the coefficients are presented in the appendix.
In order to obtain the depletion potential for $`x>x_0`$ we assume that the asymptotic decay already sets in at the point $`x_0`$. This assumption is supported by the results presented in Fig. 6. Thus, for $`x>x_0`$ we adopt the form \[c.f. Eq. (26)\]
$$\beta \overline{W}_{asympt}^w(x,\eta _s)=A_w(\eta _s)\mathrm{exp}(a_0(\eta _s)\sigma _sx)\mathrm{cos}(a_1(\eta _s)\sigma _sx\mathrm{\Theta }_w(\eta _s)),x>x_0,$$
(46)
for the scaled depletion potential between a wall and a sphere and \[c.f. Eq. (25)\]
$$\beta \overline{W}_{asympt}^p(x,\eta _s)=\frac{A_p(\eta _s)}{s^1+x}\mathrm{exp}(a_0(\eta _s)\sigma _sx)\mathrm{cos}(a_1(\eta _s)\sigma _sx\mathrm{\Theta }_p(\eta _s)),x>x_0,$$
(47)
for the potential between two spheres. The denominator in Eq. (47) measures the separation $`\sigma _b+h`$ between the centers of the spheres in units of $`\sigma _s`$. Both forms contain the functions $`a_0(\eta _s)`$ and $`a_1(\eta _s)`$, which can be calculated from the Percus-Yevick bulk pair direct correlation function $`c_{ss}^{(2)}(r)`$ (see Subsec. III D and Fig. 7). The amplitudes $`A_j(\eta _s)`$ and phases $`\mathrm{\Theta }_j(\eta _s)`$, $`j=p,w`$, are chosen so that the depletion potential and its first derivative are continuous at $`x_0`$. $`A_p(\eta _s)`$ and $`\mathrm{\Theta }_p(\eta _s)`$ are weakly dependent on the size ratio $`s`$.
With this prescription the scaled depletion potential is completely determined. For a given packing fraction $`\eta _s`$ the coefficients $`a`$, $`b`$, $`c`$ and $`d`$ are given by Eq. (A7) and the position of the first maximum can be calculated from Eq. (A1). Using those values as input, the amplitude $`A`$ and the phase $`\mathrm{\Theta }`$ of the asymptotic decay are readily obtained from either Eqs. (A4) and (A3) or Eqs. (A6) and (A5). Thus in this parametrization the scaled depletion potential has the form
$$\beta \overline{W}(x,\eta _s)=\{\begin{array}{cc}\hfill a+bx+cx^2+dx^3,& xx_0\hfill \\ \hfill \beta \overline{W}_{asympt}^{p,w}(x,\eta _s),& x>x_0.\hfill \end{array}$$
(48)
In Fig. 10(a) fits (lines) of the form given by Eq. (48) are compared with the scaled depletion potentials between a big hard sphere and a hard wall calculated within DFT (symbols) for a size ratio $`s=0.1`$. Although the fit is relatively simple, its accuracy is high. The position of the first maximum, which depends sensitively on the packing fraction $`\eta _s`$, is reproduced very accurately. The value $`\beta \overline{W}_0=\beta W(x_0)`$ of the potential at the first maximum is also given quite accurately, and only for $`\eta _s=0.3`$ are small deviations of the fit from the full DFT results visible. Clearly the full structure of the depletion potential is reproduced well by this parametrization. In order to demonstrate the wide range of applicability of this parametrization in Fig. 10(b) we show a comparison of the parametrized scaled depletion potential ($`\mathrm{}`$) for a packing fraction $`\eta _s=0.3`$ with scaled DFT results (lines) for the depletion potential between two spheres and size ratios $`s=0.2`$ and $`s=0.05`$. Although for these size ratios the scaling relation Eq. (39) is not satisfied particularly accurately, the agreement between our parametrization and the DFT results is rather good. This gives us confidence that we have developed a satisfactory parametrized form for the depletion potential which properly incorporates all essential features.
### B Oscillatory depletion potential at high packing fractions of the small spheres
In a recent experiment by Crocker et al. the equilibrium probability distribution $`p(r)`$ for two (big) PMMA (polymethylmethacrylate) spheres of diameter $`\sigma _b=1.1\mu `$m immersed in a sea of (small) polystyrene spheres of diameter $`\sigma _s=83`$ nm was measured using line-scanned optical tweezers and digital videomicroscopy at various packing fractions in the range between $`\eta _s=0.04`$ and $`\eta _s=0.42`$. The solvent contains added salt and surfactant to prevent colloidal aggregation and the ‘bare’ interactions between the colloidal particles are expected to be screened Coulombic repulsion with a screening length of about 3 nm . Since the latter is small compared with the colloid diameters the bare interactions can be regarded, to good approximation, as hard-sphere-like. The depletion potentials $`\beta W(r)=\mathrm{ln}(p(r)/p(\mathrm{}))`$ (see Subsec. II A) obtained from these experiments are shown in Fig. 1 of Ref. . At low packing fractions, $`\eta _s=0.04`$ and 0.07, rather good agreement with the results of the Asakura-Oosawa approximation was found, after taking into account the effects of limited spatial resolution of the optical instruments. For $`\eta _s=0.15`$ and 0.21 the measured depletion potential displayed a pronounced repulsive barrier. For higher packing fractions, i.e., $`\eta _s=0.26`$, 0.34 and 0.42 damped oscillations were observed, these being particularly pronounced for the two highest packing fractions for which three maxima are clearly visible. Reference appears to be the first report of an experimental observation of an oscillatory depletion potential and, indeed, of a repulsive contribution arising from purely entropic or packing effects .
Motivated by these experiments we consider a binary hard-sphere mixture in the dilute limit with a size ratio $`s=0.0755`$, as in the experiment. (We do not attempt to include the increase of the effective radius of the spheres arising from screened Coulomb repulsion and, in keeping with the authors of Ref. , we do not include any dispersion forces.) As previously, the depletion potential between two big spheres is calculated using Eq. (34) in the dilute limit. The functions $`\mathrm{\Psi }^\alpha `$ are functionals of $`\rho _s(r)`$, the density profiles of the small spheres close to a big sphere fixed at the origin, which depend only on the radial distance $`r`$. The results are shown in Fig. 11 for the same values of $`\eta _s`$ as in the experiments. It is encouraging to find that the theoretical and experimental results have many common features. As expected, the calculated oscillations become much more pronounced as $`\eta _s`$ increases. The wavelength decreases slowly and the decay length of the envelope increases rapidly with $`\eta _s`$ – as predicted by the theory of asymptotic decay (see Fig. 7). The experimental data are consistent with both observations. Moreover the wavelength of the oscillations for $`\eta _s=0.34`$ is close to $`\sigma _s=`$ 83 nm in theory and experiment. For $`\eta _s=0.42`$ both theory and experiment yield a slightly smaller wavelength. The amplitude of the calculated oscillations is larger than in the experiment. However, we emphasize that we made no attempt to take into account effects of instrumental resolution or the polydispersity of the small polystyrene particles. Nor have we attempted to include the effects of the softness of the inter-particle potentials and any non-additivity of the effective diameters; both are likely to lead to a reduction in the amplitude of the oscillations. The qualitative agreement between the experimental results and those of our calculations persuades us that the hard-sphere model is an appropriate starting point for describing the colloidal system and that the observed oscillations do reflect the packing of the small spheres – as inferred in Ref. .
Significant deviations between our results and the experimental ones do occur, at large $`\eta _s`$, for separations near contact or near the first maximum in the depletion potential. Our results imply that the height of the first maximum and the magnitude of the contact value $`|W(0)|`$ are larger than the experimental ones by about a factor of two for $`\eta _s=0.34`$. Although the source of these differences may well reside in the experimental situation it is important to check that the particular DFT which we employ is performing reliably at these high values of $`\eta _s`$. It is precisely this regime of high density and very strong confinement of the small spheres where differences between the various DFT theories, i.e., the improvements on the original Rosenfeld version, might reveal themselves. These circumstances are reminiscent of those investigated by González et al. in their DFT studies of hard spheres in small spherical cavities. Those authors were able to ascertain that the improved theories fared better than the original version under conditions of extreme confinement.
To this end we repeated our calculations of the depletion potential with the improved versions of the Rosenfeld functional that can account for the freezing transition . At packing fractions $`\eta _s0.3`$ we obtained, as stated earlier, results almost identical to those of the original functional. At higher packing fractions, however, we find that the amplitude of the oscillations is slightly smaller than those obtained from the original functional. This is illustrated in Fig. 11 for $`\eta _s=0.42`$, using the interpolation form of the functional (dotted line). The antisymmetrized version of the functional, with $`q=3`$ , yields a depletion potential very close to that of the interpolation form. In view of the smallness of these deviations the discrepancies between the experimental findings and the theoretical results at high $`\eta _s`$ cannot be blamed on the performance of the DFT but most probably reside in differences between the actual experimental sample and the model of hard spheres.
## VI Summary and Discussion
In this paper we have developed a versatile theory for determining the depletion potential in general fluid mixtures. Our approach requires only the knowledge of the equilibrium density profile $`\rho _s(𝐫)`$ of the small particles before the big (test) particle is inserted, i.e., $`\rho _s(𝐫)`$ has the symmetry of the external potential. If the latter is exerted by a fixed particle or by a planar wall then in these cases $`\rho _s(𝐫)`$ simplifies to functions $`\rho _s(r)`$ or $`\rho _s(z)`$ of one variable. Since a one-dimensional profile can be calculated very accurately, the resulting depletion potentials can be obtained without the numerical complications and limitations that are inherent in brute-force DFT . The latter requires the calculation of the local density of the small particles around the big particles in the presence of the external potential or the calculation of the total free energy as a function of the separation of the big particles ; both calculations require considerable numerical effort due to the reduced symmetry of the density distributions. We have employed our approach in a comprehensive study of the depletion potential for hard-sphere systems, using Rosenfeld’s fundamental measure functional. The main conclusions which emerge from our study are as follows:
1. The depletion potential can be obtained by considering a liquid mixture in the limit of vanishing concentration of one of the species. Two different ways to implement this limit lead to the same result (Fig. 1).
2. Detailed comparison of our results with those of simulations, for both sphere-sphere and (planar) wall-sphere depletion potentials (see Figs.2 and 3), demonstrate that the theory is very accurate for size ratios $`s=R_s/R_b`$ as small as 0.1 and for packing fractions $`\eta _s`$ as large as 0.3. These are the most extreme cases for which reliable simulation data are presently available. The theory describes accurately the short-ranged depletion attraction, the first repulsive barrier and the subsequent oscillations in the depletion potential.
3. By performing consistency checks we argue that at least up to moderate packing fractions the predictions of the Rosenfeld DFT for depletion should be quantitatively reliable even for large asymmetries between the sizes of the solvent and the solute particles (Fig. 8). Subsection III E provides a theoretical understanding of this feature of our DFT approach.
4. Extensions of the Rosenfeld functional yield very similar results (see Fig. 11) for the cases we have studied. It would be of considerable interest to test the performance of the proposed functionals against simulation data for smaller size ratios and for higher values of $`\eta _s`$, for which more extreme packing constraints might discriminate between the various functionals.
5. Our DFT approach incorporates the correct, exponentially damped, oscillatory asymptotic ($`h\mathrm{}`$) decay of the depletion potential $`W(h)`$. This is inherent in the construction of the theory, is preserved by the approximate Rosenfeld functional and is exhibited explicitly by the numerical results (Fig. 6). The decay length $`a_0^1`$ of the oscillations increases and the wavelength $`2\pi /a_1`$ decreases with increasing $`\eta _s`$ (Fig. 7) but these quantities are independent of the size ratio $`s`$. The same values for $`a_0`$ and $`a_1`$ characterize the oscillatory decay towards the bulk values of the number density profiles of hard-sphere mixtures near a hard wall when the packing fraction of the big spheres is very small (Figs. 4 and 5).
6. We have developed simple parametrization schemes for the depletion potential between big hard sphere and a planar wall and that between two big hard spheres which provide accurate fits to our DFT results (see Fig. 10). The fitting procedure makes use of the fact that leading asymptotic behavior of $`W(h\mathrm{})`$ provides an accurate account of the oscillatory structure of the depletion potential at intermediate separations as well as at longest range. Such parametrizations are designed to provide a more accurate alternative to the third-order virial expansion formula given by Götzelmann et al. . Since these parametrization can be easily implemented we recommend that they should be employed in subsequent studies of the phase behavior of highly-asymmetric binary hard-sphere mixtures of the type reported in Refs. and .
7. In Sec. IV we investigated the regime of validity of the Derjaguin approximation \[Eq. (40)\] for the depletion potential and showed that this fails, for all but the smallest packing fractions $`\eta _s`$, for which the depletion potential reduces to the Asakura-Oosawa result. However, the scaling relation Eq. (39) connecting the depletion force between two big spheres to that between a big sphere and a planar wall – which is predicted by the Derjaguin approximation but which also follows from geometrical considerations – does remain valid even at moderate size ratios (Fig. 9).
We conclude with several remarks concerning the accuracy and usefulness of our approach. One might be surprised that a DFT which corresponds to the Percus-Yevick theory for the bulk mixture (the Rosenfeld functional yields the same bulk free-energy density and bulk pair direct correlation function) performs so well for small size ratios, for which it is known that Percus-Yevick theory becomes inaccurate. For example, it fails to predict the fluid-fluid spinodals for hard-sphere mixtures. However, our present approach involves only the calculation of a one-body direct correlation function $`c_b^{(1)}(𝐫;\{\mu _i\})`$ and, therefore, the determination of one-body density profiles. The minimization of approximate functionals can yield rather accurate one-body profiles in spite of the limitations of the underlying approximations; e.g., this is the reason why the test particle route to the bulk radial distribution function $`g(r)`$ is very successful within DFT . Furthermore, in determining the depletion potential $`W(𝐫)`$ we require only solutions of the Euler-Lagrange equation for $`\rho _s(𝐫)`$ in the limit where $`\rho _b0`$, i.e., in the absence of the big particles. The DFT is likely to be more accurate in this limiting regime than for a mixture concentrated in all species. We emphasize that taking the dilute limit of the big particles numerically, i.e., working at non-zero but very small values of $`\eta _b`$, involves more computation than taking the limit directly in the functional. Moreover, caution should be exercised in hard-sphere mixtures with extreme size ratios, $`s0.1`$, at high packing fractions $`\eta _s`$ of the small spheres since the fluid-solid phase boundary already occurs at very low packing fractions $`\eta _b`$ of the big spheres . The fluid-solid coexistence region is avoided if the dilute limit is taken directly .
Our procedure for calculating $`W_t(𝐫)`$ at arbitrary concentrations of the big particles might prove useful for interpreting (future) measurements of the effective interaction potential when the mixture is not in the dilute limit. Figure 1 of Ref. illustrates how the wall-sphere potential $`W_t(z)`$ varies with the big sphere packing fraction $`\eta _b`$ for a mixture with size ratio $`s=0.2`$ and $`\eta _s=0.2`$. For $`\eta _b=0.025`$, $`W_t(z)`$ already differs by a few percent from its dilute limit $`W(z)`$.
It is possible to calculate the depletion potential by using as input density profiles obtained by other means. In particular one might take simulation data for $`\rho _s(𝐫)`$, computed in the absence of the big test particle, and insert these into Eq. (36) to determine the weighted densities. Although such a procedure does not offer the appeal of a self-consistent approach in which both the equilibrium density profiles and the depletion potential are calculated within the same framework, in practice this could be a profitable route for complex geometries where a direct simulation of the depletion potential or force is very difficult.
Finally we mention that the techniques we have developed here are not restricted to additive, binary hard-sphere mixtures. Our general approach to the calculation of depletion potentials can be applied to hard-sphere mixtures with non-additive diameters, to ternary mixtures and to systems where the interparticle potentials are soft.
###### Acknowledgements.
We thank T. Biben and R. Dickman for providing us with their simulation data. We benefited from conversations with J.M. Brader, M. Dijkstra, H.N.K Lekkerkerker, R. van Roij, and correspondence with J.-P. Hansen, P.B. Warren, and A.G. Yodh. This research was supported in part by the EPSRC under GR/L89013.
## A Parameterizing the depletion potential
In the range between contact, $`x=h/\sigma _s=0`$, and the position $`x_0`$ of the first maximum the scaled depletion potential $`\overline{W}`$ is parametrized by a cubic polynomial \[Eq. (45)\], which is the simplest polynomial fit that remains accurate close to contact. Since $`\beta \overline{W}(x=0,\eta _s)=a(\eta _s)`$, the first coefficient is the contact value of the depletion potential. The position $`x_0`$ and the height $`W_0`$ of the first maximum can be obtained easily by differentiating Eq. (45). The cubic polynomial has two extrema, with the maximum located at
$$x_0(\eta _s)=\frac{c+\sqrt{c^23bd}}{3d}$$
(A1)
and a maximal value of
$$\beta \overline{W}_0(\eta _s)=\beta \overline{W}(x=x_0(\eta _s),\eta _s)=\frac{2c^39bcd+27ad^2+2(c^23bd)^{3/2}}{27d^2}.$$
(A2)
Beyond the position of the first maximum of the depletion potential the parametrized form is continued by imposing the known asymptotic behavior for large $`h`$. The asymptotic behaviors of the wall-sphere and sphere-sphere depletion potential are slightly different and must be considered separately. For the wall-sphere depletion potential the asymptotic behavior is given by Eq. (46) and the amplitude $`A_w`$ and the phase $`\mathrm{\Theta }_w`$ are chosen such that the function and its first derivative are continuous at $`x_0`$. From the requirement of a continuous derivative at the first maximum, i.e., $`\frac{\mathrm{d}\beta \overline{W}_{asympt}^w(x,\eta _s)}{\mathrm{d}x}|_{x=x_0}=0`$, the phase can be determined to be
$$\mathrm{\Theta }_w(\eta _s)=a_1\sigma _sx_0+\mathrm{arccos}\left(\frac{a_1}{\sqrt{a_0^2+a_1^2}}\right).$$
(A3)
From the requirement that the function is continuous at $`x_0`$, i.e., $`\overline{W}_{asympt}^w(x_0,\eta _s)=\overline{W}_0(\eta _s)`$, together with the phase from Eq. (A3), the amplitude of the asymptotic decay follows as
$$A_w(\eta _s)=\beta \overline{W}_0\mathrm{exp}(a_0\sigma _sx_0)\sqrt{\frac{a_0^2+a_1^2}{a_1^2}}.$$
(A4)
A similar calculation for the sphere-sphere case using Eq. (47) leads to slightly different expressions for the phase,
$$\mathrm{\Theta }_p(\eta _s)=a_1\sigma _sx_0+\mathrm{arccos}\left(\frac{a_1(x_0\sigma _s+\sigma _b)}{\sqrt{1+2a_0(x_0\sigma _s+\sigma _b)+(a_0^2+a_1^2)(x_0\sigma _s+\sigma _b)^2}}\right),$$
(A5)
and the amplitude
$$A_p(\eta _s)=\beta \overline{W}_0\frac{\mathrm{exp}(a_0\sigma _sx_0)}{a_1\sigma _s}\sqrt{1+2a_0(x_0\sigma _s+\sigma _b)+(a_0^2+a_1^2)(x_0\sigma _s+\sigma _b)^2}.$$
(A6)
Unlike $`\mathrm{\Theta }_w`$ and $`A_w`$, $`\mathrm{\Theta }_p`$ and $`A_p`$ depend (weakly) on the size ratio $`s=\sigma _s/\sigma _b`$.
The coefficients $`a`$, $`b`$, $`c`$ and $`d`$ are fitted to depletion potentials calculated within DFT. Scaled depletion potentials $`\overline{W}`$ obtained for a big hard sphere near a planar hard wall, for a size ratio $`s=0.1`$, are used in the range $`0\eta _s0.3`$. The resulting coefficients are given by
$$\begin{array}{ccc}\hfill a(\eta _s)& =& 2.909\eta _s,\hfill \\ \hfill b(\eta _s)& =& 6.916\eta _s4.616\eta _s^2+78.856\eta _s^3,\hfill \\ \hfill c(\eta _s)& =& 4.512\eta _s+15.860\eta _s^293.224\eta _s^3,\hfill \\ \hfill d(\eta _s)& =& \eta _s\mathrm{exp}(1.734+8.957\eta _s+1.595\eta _s^2).\hfill \end{array}$$
(A7)
There is no particular significance in the chosen form of parametrization but we note that the contact values of the scaled depletion potential $`\beta \widehat{W}(0,\eta _s)=a(\eta _s)`$ are linear in $`\eta _s`$ for this choice of parametrization. It is interesting that the coefficient $`2.909`$ is rather close to the value $`3`$ obtained from the Asakura-Oosawa result (valid as $`\eta _s0`$) in the limit of small size-ratios, see Eqs. (14) and (16) and also Eq. (43).
The quantities $`a_0`$ and $`a_1`$ are obtained by solving Eq. (24), using the Percus-Yevick pair direct correlation function $`c_{ss}^{(2)}(r)`$. The results are shown in Fig. 7. In the range $`0.05\eta _s0.4`$ they can be fitted accurately by
$$a_0(\eta _s)\sigma _s=4.674\mathrm{exp}(3.935\eta _s)+3.536\mathrm{exp}(56.270\eta _s)$$
(A8)
and
$$a_1(\eta _s)\sigma _s=0.682\mathrm{exp}(24.697\eta _s)+4.720+4.450\eta _s.$$
(A9)
These formulae specify all the ingredients for determining the parametrized form of the depletion potential.
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# Brane-World Cosmology of Modulus Stabilization with a Bulk Scalar Field
## I Introduction
One of the most striking proposals in current elementary particle theory is the existence of extra dimensions which are hidden from us, not necessarily by their smallness, but by our confinement to a four-dimensional slice (a 3-brane) of the full spacetime . Randall and Sundrum have produced a verion of this scenario which is particularly attractive because of its apparent links to deep theoretical developments: the conformal field theory/5D-anti-de Sitter correspondence and holography . On a more practical level, their idea provides an explanation for why the Planck scale is so much greater than the weak scale, independently of supersymmetry. The cosmological implications of this model have also been the subject of vigorous study .
The Randall-Sundrum proposal involves a “Planck brane” located at a position $`y=0`$ in a single additional dimension, and a second “TeV brane” located at $`y=1`$, in our conventions. The extra dimension is permeated by a negative bulk vacuum energy density, so that the space between the branes is a slice of anti-de Sitter space. Solving the 5-D Einstein equations results in the line element
$`ds^2`$ $`=`$ $`e^{2\sigma (y)}(dt^2d\stackrel{}{x}^2)b^2dy^2;`$ (1)
$`\sigma (y)`$ $`=`$ $`kby;y[0,1]`$ (2)
The constant $`k`$ is related to the 4-D and 5-D Planck masses, $`M`$ and $`M_p`$ respectively, by
$$k=\frac{M^3}{M_p^2}$$
(3)
where $`M_p^2=8\pi G_N`$ and the 5-D gravitational action is $`S=\frac{1}{2}M^3d^{\mathrm{\hspace{0.17em}5}}x\sqrt{g}R`$. The warp factor $`e^{\sigma (y)}`$ determines the physical masses of particles on the TeV brane: even if a bare mass parameter $`m_0`$ in the TeV brane Lagrangian is of order $`M_p`$, the physical mass gets rescaled by
$`m_{\mathrm{phys}}=e^{\overline{\sigma }}m_0;\overline{\sigma }\sigma (1)=kb,`$ (4)
as can be easily derived starting from a scalar field action written covariantly in terms of the metric (1), $`S=\frac{1}{2}d^{\mathrm{\hspace{0.17em}4}}xe^{4\overline{\sigma }}(e^{2\overline{\sigma }}(\varphi )^2m_0^2\varphi ^2)`$, and rescaling $`\varphi `$ so that the kinetic term becomes canonically normalized.
In order for $`m_0`$ to be of order 100 GeV, the combination $`kb`$ must be approximately 36, so that $`e^{\overline{\sigma }}10^{16}`$. Yet in the original proposal, the value of $`b`$, which is the size of the extra dimension, was completely undetermined. It is a modulus with no potential, which is phenomenologically unacceptable. For one thing, the particle associated with 4-D fluctuations of $`b`$, the radion, would couple to matter on the TeV brane similarly to gravitons, but more strongly by a factor of $`e^{\overline{\sigma }}`$ . This would lead to a fifth force which could easily have been detected. Furthermore, a massless radion leads to problems with cosmology: our brane-universe would have to have a negative energy density to expand at the expected rate, assuming that energy densities on the branes are tuned to give a static extra dimension. It was shown in refs. (see also ) that this problem disappears when the size of the extra dimension is stabilized.
Radion stabilization is therefore a crucial ingredient of the Randall-Sundrum idea. Goldberger and Wise have presented an elegant mechanism for accomplishing this , using a bulk scalar field. Self-interactions of the field on the branes forces it to take nonvanishing VEV’s, $`v_0`$ and $`v_1`$ respectively, which are generally different from each other. The field thus has a gradient in the extra dimension, and the competition between the gradient and potential energies gives a preferred value for the size of the extra dimension. In other words, a potential for the radion is generated, which has a nontrivial minimum. It is easy to get the correct brane separation using natural values of the parameters in the model.
Roughly speaking, the radion potential has the form
$$V(\varphi )\lambda \varphi ^4\left(\left(\frac{\varphi }{f}\right)^ϵ\frac{v_1}{v_0}\right)^2$$
(5)
with a nontrivial minimum at $`\varphi =f(v_1/v_0)^{1/ϵ}`$. However, there is another minimum at $`\varphi =0`$, which represents an infinitely large extra dimension. This could not describe our universe, because then $`e^{\overline{\sigma }}`$ would be zero, corresponding to vanishing particle masses on the TeV brane. In the more exact expression for the potential, we will show that $`\varphi =0`$ is actually a false vacuum—it has higher energy than the nontrivial minimum. Nevertheless, is is quite likely that the metastable state could be reached through classical evolution in the early universe. The question then naturally arises whether tunneling or thermal transitions to the desired state occurs. This is the subject of our study.
Such detailed questions about the viability of the Goldberger-Wise mechanism are important because there are few attractive alternatives at the moment. Ref. recently studied Casimir energies of fields between the branes as a possible origin of a stabilizing potential. They found that stabilization in this way is indeed possible, but that the resulting radion mass is too small to be phenomenologically consistent if the size of the extra dimension is that taken to be that dictated by the hierarchy problem.
In section 2 we derive the Goldberger-Wise potential for the radion in a slightly simplified model which allows for the exact solution of the potential. The classical evolution of the radion field is considered in section 3, where we show that for generic initial conditions, the universe reaches a state in which the radion is not stabilized, but instead the extra dimension is expanding. This is a metastable state however, and in section 4 the rate of transitions to the minimum energy state, with finite brane separation, is computed. Conclusions are given in section 5.
## II Radion Potential
Let $`\psi (y)`$ be the bulk scalar field which will be responsible for stabilizing the radion. Its action consists of a bulk term plus interactions confined to the two branes, located at coordinate positions $`y=0`$ and $`y=1`$, respectively. Using the variable $`\sigma `$ of eq. (2) in place of $`y`$, the 4-D effective Lagrangian for a static $`\psi `$ configuration can be written as
$``$ $`=`$ $`{\displaystyle \frac{k}{2}}{\displaystyle _{\overline{\sigma }}^{\overline{\sigma }}}e^{4\sigma }\left((_\sigma \psi )^2+\widehat{m}^2\psi ^2\right)𝑑\sigma `$ (6)
$``$ $`m_0(\psi _0v_0)^2e^{4\overline{\sigma }}m_1(\psi _1v_1)^2,`$ (7)
where $`\widehat{m}`$ is the dimensionless mass $`m/k`$, $`\psi _i`$ are the values of $`\psi `$ at the respective branes, and the orbifold symmetry $`\psi (\sigma )=\psi (\sigma )`$ is to be understood. In 5-D, the field $`\psi `$, as well as the VEV’s $`v_i`$ on the two branes, have dimensions of (mass)<sup>3/2</sup>, while the parameters $`m_i`$ have dimensions of mass. Denoting $`_\sigma \psi =\psi ^{}`$, the the Euler-Lagrange equation for $`\psi `$ is
$`ke^{4\sigma }(\psi ^{\prime \prime }`$ $``$ $`4\psi ^{}\widehat{m}^2\psi )=2m_0(\psi _0v_0)\delta (\sigma )`$ (8)
$`+`$ $`2e^{4\overline{\sigma }}m_1(\psi _1v_1)\delta (\sigma \overline{\sigma }).`$ (9)
The general solution has the form
$`\psi `$ $`=`$ $`e^{2\sigma }\left(Ae^{\nu \sigma }+Be^{\nu \sigma }\right);`$ (10)
$`\nu `$ $`=`$ $`\sqrt{4+\widehat{m}^2}2+ϵ.`$ (11)
To get the correct hierarchy between the Planck and weak scales, it is necessary to take $`\widehat{m}^2`$ to be small, hence the notation $`ϵ`$.
The brane terms induce boundary conditions specifying the discontinuity in $`\psi ^{}`$ at the two branes. Imposing the orbifold symmetry $`\psi (\sigma )=\psi (\sigma )`$, this implies that
$`\psi ^{}(0)`$ $`=`$ $`\widehat{m}_0(\psi _0v_0)`$ (12)
$`\psi ^{}(\overline{\sigma })`$ $`=`$ $`\widehat{m}_1(\psi _1v_1),`$ (13)
where we defined $`\widehat{m}_i=m_i/k`$. Substituting the solution (10) into (12) allows us to solve for the unknown coefficients $`A`$ and $`B`$ exactly. In this respect, the present model is simpler than that originally given in ref. . There the brane potentials were taken to be quartic functions, so that $`A`$ and $`B`$ could only be found in the approximation where the field values $`\psi _i`$ were very strongly pinned to their minimum energy values, $`v_i`$. In our model this would occur in the limit $`m_i\mathrm{}`$. However, we can easily explore the effect of leaving these parameters finite since $`A`$ and $`B`$ can be determined exactly. Let us denote
$$\widehat{\varphi }=e^{kb}=e^{\overline{\sigma }},$$
(14)
which will be convenient in the following, because $`\widehat{\varphi }`$ is proportional to the canonically normalized radion field. Then $`A`$ and $`B`$ are given by
$`A`$ $`=`$ $`\left(C_1\widehat{\varphi }^\nu +C_2\widehat{\varphi }^2\right)\widehat{\varphi }^\nu /D(\widehat{\varphi })`$ (15)
$`B`$ $`=`$ $`\left(C_3\widehat{\varphi }^\nu C_4\widehat{\varphi }^2\right)\widehat{\varphi }^\nu /D(\widehat{\varphi })`$ (16)
where
$`C_1`$ $`=`$ $`\widehat{m}_0v_0(\widehat{m}_1ϵ)`$ (17)
$`C_2`$ $`=`$ $`\widehat{m}_1v_1(\widehat{m}_0+ϵ)`$ (18)
$`C_3`$ $`=`$ $`\widehat{m}_0v_0(\widehat{m}_1+4+ϵ)`$ (19)
$`C_4`$ $`=`$ $`\widehat{m}_1v_1(\widehat{m}_04ϵ)`$ (20)
$`D(\widehat{\varphi })`$ $`=`$ $`{\displaystyle \frac{\left(C_2C_3C_1C_4\widehat{\varphi }^{2\nu }\right)}{\left(\widehat{m}_0\widehat{m}_1v_0v_1\right)}}.`$ (21)
It can be checked that in the limit $`\widehat{m}_i\mathrm{}`$, the field values on the branes approach $`\psi _iv_i`$. For finite $`\widehat{m}_i`$, the competing effect of minimizing the bulk energy causes departures from these values, however.
The solution for $`\psi `$ can be substituted back into the Lagrangian (6) to obtain the effective 4-D potential for the radion, $`V(\widehat{\varphi })`$. However, rather than substituting directly, one can take advantage of the fact that $`\psi `$ is a solution to the Euler-Lagrange equation. Doing a partial integration and using the boundary terms in (8) gives a simpler expression for $`V(\widehat{\varphi })`$. In the general case where the brane potentials are denoted by $`V_i(\psi )`$, one can easily show that
$$=\underset{i}{}e^{4\sigma _i}\left(V_i(\psi _i)\frac{1}{2}\psi _iV_i^{}(\psi _i)\right)$$
(22)
In the present case, we obtain
$`V(\widehat{\varphi })`$ $`=`$ $``$ (23)
$`=`$ $`m_0v_0(v_0\psi _0)+\widehat{\varphi }^4m_1v_1(v_1\psi _1)`$ (24)
$`=`$ $`m_0v_0(v_0(A+B))`$ (25)
$`+`$ $`\widehat{\varphi }^4m_1v_1(v_1\widehat{\varphi }^2(A\widehat{\varphi }^\nu +\widehat{\varphi }^\nu B))`$ (26)
In the following, we will be interested in the situation where $`V(\widehat{\varphi })`$ has a nontrivial minimum for very small values of $`\widehat{\varphi }10^{16}`$, as needed to address the weak scale hierarchy problem. It is therefore a good approximation to expand $`V(\widehat{\varphi })`$ near $`\widehat{\varphi }=0`$, keeping only the terms which are larger than $`\widehat{\varphi }^{2\nu }`$. This is accomplished by expanding the denominator $`D(\widehat{\varphi })`$ in eqs. (15), after which the potential can be written in the simple form
$$V(\widehat{\varphi })=\mathrm{\Lambda }\widehat{\varphi }^4\left(\left(1+ϵ_4ϵ_1\right)\widehat{\varphi }^{2ϵ}2\eta \left(1+ϵ_4\right)\widehat{\varphi }^ϵ+\eta ^2\right)$$
(27)
where we introduce the notation
$$ϵ_4=\frac{ϵ}{4};ϵ_0=\frac{ϵ}{\widehat{m}_0};ϵ_1=\frac{ϵ}{\widehat{m}_1},$$
(28)
$$\eta =\left(1+ϵ_0\right)\frac{v_1}{v_0}$$
(29)
and
$$\mathrm{\Lambda }=4kv_0^2\frac{(1+ϵ_4)\left(1+ϵ_0\right)^2}{\left(1+\frac{4}{\widehat{m}_1}+ϵ_0\right)}$$
(30)
In the following it will be convenient for us to rewrite $`V(\widehat{\varphi })`$ in the form
$$V(\widehat{\varphi })=\mathrm{\Lambda }^{}\widehat{\varphi }^4(\widehat{\varphi }^ϵc_+)(\widehat{\varphi }^ϵc_{})$$
(31)
where $`c_\pm `$ are given by
$$c_\pm =\eta \left(\frac{\left(1+ϵ_4\right)\pm \sqrt{\left(1+ϵ_4\right)^2\left(1+ϵ_4ϵ_1\right)}}{\left(1+ϵ_4ϵ_1\right)}\right)$$
(32)
and $`\mathrm{\Lambda }^{}=\mathrm{\Lambda }(1+ϵ_4ϵ_1)`$.
## III Phenomenology and Early Cosmology of the Model
The warp factor which determines the hierarchy between the weak and Planck scales can be found by minimizing the potential (31). Expanding in $`ϵ`$, it has a global minimum and a local maximum at the respective values $`\widehat{\varphi }_+`$ and $`\widehat{\varphi }_{}`$:
$`\widehat{\varphi }_\pm `$ $`=`$ $`\left({\displaystyle \frac{v_1}{v_0}}\right)^{1/ϵ}({\displaystyle \frac{1+ϵ_0}{1+ϵ_4ϵ_1}})^{1/ϵ}\left(1\pm \sqrt{ϵ_1+ϵ_4\frac{1}{2}ϵϵ_1}\right)^{1/ϵ}`$ (33)
$``$ $`\left({\displaystyle \frac{v_1}{v_0}}\right)^{1/ϵ}\mathrm{exp}\left(\pm \sqrt{\frac{1}{ϵ}(\frac{1}{\widehat{m}_1}+\frac{1}{4})}+\frac{1}{\widehat{m}_0}+\frac{1}{\widehat{m}_1}\frac{1}{4}\right)`$ (34)
The last approximation holds in the limit of small $`ϵ`$, $`ϵ_0`$ and $`ϵ_1`$; it is not always an accurate approximation for the parameter values of interest, so we will use the exact expression in any computations which might be sensitive to the actual value. The positions of the zeroes of $`V`$, $`\widehat{\varphi }=c_\pm ^{1/ϵ}`$, are slightly greater than $`\widehat{\varphi }_\pm `$, by the factor $`(1+ϵ_4)^{1/ϵ}e^{1/4}`$, as can be seen by comparing (33) with (32). The large hierarchy of $`\widehat{\varphi }_+10^{16}`$ is achieved by taking a moderately small ratio $`v_1/v_0<1`$ and raising it to a large power,<sup>*</sup><sup>*</sup>*An alternative possibility, taking $`v_1/v_0>1`$ and $`ϵ<0`$, corresponding to a negative squared mass in the bulk Lagrangian (6), does not work. Although the negative $`m^2`$ does not necessarily lead to any instability, since the field is prevented from going to infinity by the potentials on the branes, the radion potential has no nontrivial minimum in this case. $`1/ϵ`$. This leads to the mass scale which functions like the cutoff on the TeV brane,
$$\widehat{\varphi }_+M_pM_{\mathrm{TeV}}$$
(35)
where $`M_{\mathrm{TeV}}/(1\mathrm{TeV})`$ is a number of order unity, which we will take to be one of the phenomenological free parameters of the model. The choice of $`M_{\mathrm{TeV}}`$ specifies precisely where we want our cutoff scale to be. In ref. the exponential corrections to $`(v_1/v_0)^{1/ϵ}`$ in (33) were not considered; these change somewhat the values one might choose for $`v_1/v_0`$ and $`ϵ`$ to get the correct hierarchy. The factor $`e^{\pm \sqrt{(1/\widehat{m}_1+1/4)/ϵ}}`$ in particular can be significant.
Refs. and showed that the canonically normalized radion field is $`\varphi =f\widehat{\varphi }`$, where $`f=\sqrt{6M^3/k}`$ is another scale of order $`M_p`$.The choice $`f=\sqrt{24M^3/k}`$ in ref. seems to correspond to an unconventional normalization for $`M_p`$. The 4-D effective action for the radion and gravity is
$`S`$ $`=`$ $`{\displaystyle \frac{M^3}{2k}}{\displaystyle d^{\mathrm{\hspace{0.17em}4}}x\sqrt{g}\left(1\widehat{\varphi }^2\right)R}`$ (36)
$`+`$ $`{\displaystyle d^{\mathrm{\hspace{0.17em}4}}x\sqrt{g}\left(\frac{1}{2}f^2_\mu \widehat{\varphi }^\mu \widehat{\varphi }V(\widehat{\varphi })\right)},`$ (37)
where $`V(\widehat{\varphi })`$ is the potential (31) and $`R`$ is the Ricci scalar. To get the correct strength of gravity, we must therefore have
$$\frac{M^3}{k}=M_p^2=\frac{1}{8\pi G_N};f=\sqrt{6}M_p.$$
(38)
The radion mass is $`f^2`$ times the second derivative of $`V`$ evaluated at its minimum. We find the value
$$m_\varphi ^2=4\frac{\mathrm{\Lambda }}{f^2}\left(1+\frac{ϵ}{4}\frac{ϵ}{\widehat{m}_1}\right)^1\left(\sqrt{}1+\frac{4}{\widehat{m}_1}+2\frac{\sqrt{ϵ}}{\widehat{m}_1}\right)\eta ^2\widehat{\varphi }_+^2ϵ^{3/2},$$
(39)
which implies that $`m_\varphi `$ is of typically of order $`ϵ^{3/4}`$ times the TeV scale. The factor of $`ϵ^{3/4}`$ leads to the prediction that the radion will be lighter than the Kaluza-Klein excitations of the graviton, which would also be a signal of the new brane physics . However, we see that the corrections due to finite $`\widehat{m}_1`$, which were not explicitly considered in , can possibly compensate this and make the radion heavier, if $`\widehat{m}_1ϵ`$.
Now let us turn to the evolution of $`\widehat{\varphi }`$ in the early universe. For this purpose it is important to understand that the depth of the potential at its global minimum, as well as the height of the bump separating the minimum from $`\widehat{\varphi }=0`$, is set by the TeV scale. The values of the potential at these extrema are approximately (to leading order in $`ϵ_4`$, but exact in $`ϵ_1`$)
$$V(\widehat{\varphi }_\pm )2\eta ^2\mathrm{\Lambda }\widehat{\varphi }_\pm ^4\frac{ϵ_4\sqrt{ϵ_1+ϵ_4}}{1+ϵ_4ϵ_1}(1\pm \sqrt{ϵ_1}).$$
(40)
Since $`\mathrm{\Lambda }M_p^4`$, the depth at the minimum is $`V(\widehat{\varphi }_+)ϵ^{3/2}O(\mathrm{TeV})^4`$—suppressed slightly by the factor of $`ϵ^{3/2}`$. The height of the bump at $`\widehat{\varphi }_{}`$ can be considerably smaller because of the exponential factors in (33). In fact
$$\left|\frac{V(\widehat{\varphi }_{})}{V(\widehat{\varphi }_+)}\right|=\left(\frac{\widehat{\varphi }_{}}{\widehat{\varphi }_+}\right)^4\mathrm{\Omega }^4$$
(41)
where
$$\mathrm{\Omega }\left(\frac{1\sqrt{ϵ_1+ϵ_4}}{1+\sqrt{ϵ_1+ϵ_4}}\right)^{1/ϵ}\mathrm{exp}\left(\sqrt{\frac{1}{ϵ}(1+\frac{4}{\widehat{m}_1})}\right)$$
(42)
For example, if $`ϵ=0.01`$ as suggested by , $`\mathrm{\Omega }^4`$ is less than $`10^{17}`$. If the brane potential parameter $`m_1`$ is not large, so that $`\widehat{m}_1\text{ }\stackrel{<}{}\text{ }1`$, the suppression will be much greater. Figure 1 illustrates the flatness of the potential for the case $`ϵ=0.2`$, where the barrier is not so suppressed. The new mass scale $`\mathrm{\Omega }`$ TeV $``$ 1 TeV is due to the small curvature of the radion potential at the top of the barrier, and its smallness will play an important role in the following.
Thus the barrier separating the true minimum at $`\widehat{\varphi }_+`$ from the false one at $`\widehat{\varphi }=0`$ is extremely shallow. Moreover a generic initial condition for the radion is a value like $`\widehat{\varphi }1`$, quite different from the one we want to end up with, $`\widehat{\varphi }10^{16}`$. Clearly, the shape of the potential is such that, if we started with a generic initial value for $`\widehat{\varphi }`$, it would easily roll past the local minimum and the barrier, hardly noticing their presence. The point $`\widehat{\varphi }=0`$ toward which it rolls is the limit of infinite brane separation, phenomenologically disastrous since gravity no longer couples at all to the visible brane in this limit.
One might wonder whether inflation could prevent this unwanted outcome, since then there would be a damping term in the $`\varphi `$ equation of motion, possibly causing it to roll slowly:
$$\ddot{\varphi }+3H\dot{\varphi }+V_{\mathrm{eff}}^{}(\varphi )=0.$$
(43)
Indeed, with sufficiently large Hubble rate $`H`$, the motion could be damped so that $`\varphi `$ would roll to its global minimum. The condition for slow-roll is that
$$V_{\mathrm{eff}}^{\prime \prime }9H^2.$$
(44)
However, inflation is a two-edged sword in this instance, because the effective potential $`V_{\mathrm{eff}}`$ for the radion gets additional contributions from the curvature of the universe during inflation. From eq. (36) one can see that
$`V_{\mathrm{eff}}(\widehat{\varphi })`$ $`=`$ $`V(\widehat{\varphi })+{\displaystyle \frac{M^3}{2k}}R\widehat{\varphi }^2`$ (45)
$`=`$ $`V(\widehat{\varphi })+6{\displaystyle \frac{M^3}{k}}H^2\widehat{\varphi }^2`$ (46)
$`=`$ $`V(\widehat{\varphi })+H^2\varphi ^2`$ (47)
using the relation $`R=12H^2`$ which applies for de Sitter space, and eq. (38). The new term tends to destroy the nontrivial minimum of the radion potential. One can estimate the relative shift in the position of the minimum as
$$\frac{\delta \widehat{\varphi }_+}{\widehat{\varphi }_+}=\frac{\delta V^{}(\widehat{\varphi }_+)}{\widehat{\varphi }_+V^{\prime \prime }(\widehat{\varphi }_+)}=\frac{H^2}{m_\varphi ^2}$$
(48)
This should be less than unity to avoid the disappearance of the minimum altogether.
Combining the requirement that the global minimum survives with the slow-roll condition (44), evaluated near the minimum of the potential, we find the following constraint on the Hubble rate:
$$\frac{1}{9}m_\varphi ^2H^2<m_\varphi ^2.$$
(49)
This is a narrow range, if it exists at all. In fact, one never expects such a large Hubble rate in the Randall-Sundrum scenario since the TeV scale is the cutoff: $`H`$ should never exceed $`T^2/M_p10^{16}`$ TeV if the classical equations are to be valid. The problem of radion stabilization might also be exacerbated because contributions to the energy density of the universe which cause inflation can give additional terms of the type $`\widehat{\varphi }^2`$ to $`V_{\mathrm{eff}}`$ which are not considered in the above argument. For example, \[
Figure 1: $`V(\varphi )`$ versus $`\varphi `$ for $`ϵ=0.2`$, showing the smallness of the barrier (right) relative to the minimum (left). Notice the difference in vertical scales. \] a field in the bulk which does not have its endpoints fixed on the branes has a 5-D energy density which is peaked on the visible brane , $`\rho _5(y)\rho _0e^{2kby}`$, and gives a contribution to $`V_{\mathrm{eff}}`$ of
$`\delta V`$ $``$ $`\rho _0{\displaystyle _0^1}𝑑ybe^{4kby}\rho _0e^{2kby}`$ (50)
$`=`$ $`b\rho _0\widehat{\varphi }^2,`$ (51)
remembering that $`e^{kb}=\widehat{\varphi }`$. Such a contribution could destroy the nontrivial minimum even if (49) is satisfied. In any case, it does not appear to be natural to tune the Hubble rate during inflation to try to solve the stabilization problem.
## IV Phase Transition to the True Vacuum
Since the barrier of the radion potential is too small to prevent the radion from rolling into the false minimum, perhaps we can take advantage of this smallness to get tunneling or thermal transitions back into the true vacuum. The situation is quite similar to that of “old inflation” , except that in the latter, the transition was never able to complete because the universe expanded too rapidly compared to the rate of nucleation of bubbles of the true vacuum. In the present case this problem can be avoided because we are not trying to use the radion for inflation. Indeed, a small amount of inflation may take place before the tunneling occurs, since the radion potential is greater than zero at $`\varphi =0`$, but we will not insist that this be sufficient to solve the cosmological problems inflation is intended to solve—otherwise we would be stuck with the problems of old inflation. Instead we will assume that inflation is driven by some other field, and consider the transitions of the radion starting from the time of reheating. The criterion that the phase transition completes is that the rate of bubble nucleation per unit volume, $`\mathrm{\Gamma }/V`$, exceeds the rate of expansion of the universe per Hubble volume:
$$\frac{\mathrm{\Gamma }}{V}\text{ }\stackrel{>}{}\text{ }H^4$$
(52)
The reason is that the bubbles expand at nearly the speed of light, so the relevant volume is determined by the distance which light will have traveled by a given time, which is of order $`1/H`$.
### A The Euclidean Bounce
To compute the nucleation rate $`\mathrm{\Gamma }/V`$, one must construct the bounce solution which is a saddle point of the Euclideanized action , in other words, with the sign of the potential reversed. This is a critical bubble solution with the boundary conditions
$`\varphi (r)|_{r=0}`$ $`=`$ $`\varphi _0;\varphi ^{}(r)|_{r=0}=0;`$ (53)
$`\varphi (r)|_r\mathrm{}`$ $`=`$ $`0;\varphi ^{}(r)|_r\mathrm{}=0.`$ (54)
The value of $`\varphi _0`$ which ensures the desired behavior as $`r\mathrm{}`$ cannot be computed analytically because the motion of the field is damped by the term $`\varphi ^{}/r`$ in the equation of motion. We will consider bubble nucleation at finite temperature in the high $`T`$ limit, where the bounce solutions are three dimensional. The equation of motion is
$$\frac{1}{r^2}\left(r^2\varphi ^{}\right)^{}=+V_{\mathrm{eff}}^{}(\varphi ),$$
(55)
where now $`V_{\mathrm{eff}}`$ includes thermal corrections, which are much larger than the $`H^2\varphi ^2`$ term considered in eq. (45):
\[
Figure 2: The bounce solution, for the parameters $`ϵ=0.01`$, $`\widehat{m}_1=0.1`$, $`m_\varphi =T=100`$ GeV and $`M_{\mathrm{TeV}}=1`$ TeV. The rescaling $`\varphi \stackrel{~}{\varphi }`$ and $`r\stackrel{~}{r}`$ is given in eqs. (62-63). \]
$`V_{\mathrm{eff}}(\varphi )V(\varphi )`$ $`+`$ $`{\displaystyle \frac{T^2}{24}}m_{\mathrm{th}}^2(\varphi ){\displaystyle \frac{T}{12\pi }}m_{\mathrm{th}}^3(\varphi )`$ (56)
$`+`$ $`{\displaystyle \frac{c_b}{64\pi ^2}}m_{\mathrm{th}}^4(\varphi )V_0`$ (57)
$`m_{\mathrm{th}}^2(\varphi )=V^{\prime \prime }(\varphi )`$ $`+`$ $`\frac{1}{24}T^2V^{(4)}(0)`$ (58)
where $`c_b3.9`$ if the renormalization scale is taken to be equal to $`T`$. The leading thermal correction is of order $`T^2\varphi ^2`$, whereas $`H^2\varphi ^2`$ is suppressed by $`T^2/M_p^2`$ relative to this. The cubic term $`m_{\mathrm{th}}^3`$ becomes imaginary if $`V^{\prime \prime }(\varphi )+\frac{T^2}{24}V^{(4)}(0)`$ becomes negative; we take the real part only. The term $`\frac{1}{24}T^2V^{(4)}(0)`$ represents the thermal mass, which appears in the cubic and quartic terms when ring diagrams are resummed . We subtract a constant term $`V_0`$ from $`V_{\mathrm{eff}}`$ so that $`V_{\mathrm{eff}}(0)=0`$, as is needed to properly compute the action associated with the bounce solution.
The thermal corrections to the effective radion potential cause the bounce solution to fall exponentially at large $`r`$:
$$\varphi \frac{c}{r}e^{r/r_0},$$
(59)
where $`1/r_0\eta \sqrt{\mathrm{\Lambda }/f^4}T`$ if $`\mathrm{\Lambda }/f^41`$ \[the exact expression is $`1/r_0=\sqrt{\lambda U_{}}T`$, in terms of quantities to be defined below, in eqs. (62) and (80)\]. Once the bounce solution is known, it must be substituted back into the action, which can be written as
$$S=\frac{4\pi }{T}_0^{\mathrm{}}𝑑rr^2\left(\frac{1}{2}\varphi ^2+V_{\mathrm{eff}}(\varphi )\right).$$
(60)
The nucleation rate is given by
$$\frac{\mathrm{\Gamma }}{V}=\frac{|\omega _{}|}{2\pi }\left(\frac{S}{2\pi }\right)^{3/2}|𝒟|^{1/2}e^S$$
(61)
where $`\omega _{}`$ is the imaginary frequency of the unstable mode of fluctuations around the the bounce solution, and $`𝒟`$ is the fluctuation determinant factor, to be described below. A typical profile for the bounce solution is shown in figure 2.
For the numerical determination of the bounce solution and action, as well as understanding their parametric dependences, it is convenient to rescale the radius and the field by
$`r`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{r}}{\sqrt{\lambda }T}};\lambda ={\displaystyle \frac{\mathrm{\Lambda }^{}}{f^4}}c_{}^2,`$ (62)
$`\varphi `$ $`=`$ $`ZT\stackrel{~}{\varphi };Z={\displaystyle \frac{fc_{}^{1/ϵ}}{T}}.`$ (63)
Then the action takes the form
$`S`$ $`=`$ $`{\displaystyle \frac{Z^2}{\sqrt{\lambda }}}\stackrel{~}{S}(ϵ,\widehat{m}_1,\lambda ,Z);`$ (64)
$`\stackrel{~}{S}`$ $`=`$ $`4\pi {\displaystyle _0^{\mathrm{}}}d\stackrel{~}{r}\stackrel{~}{r}^2\{\frac{1}{2}(\stackrel{~}{\varphi }^2+\stackrel{~}{\varphi }^2f_2)+Z^2\stackrel{~}{\varphi }^4f_0`$ (65)
$``$ $`{\displaystyle \frac{1}{12\pi }}{\displaystyle \frac{\sqrt{\lambda }}{Z^2}}\left[\left(\frac{c_+}{c_{}}+12Z^2\stackrel{~}{\varphi }^2f_2\right)^{3/2}\left(\frac{c_+}{c_{}}\right)^{3/2}\right]`$ (66)
$`+`$ $`{\displaystyle \frac{3c_b}{8\pi ^2}}\lambda f_2\stackrel{~}{\varphi }^2(\frac{c_+}{c_{}}+6Z^2\stackrel{~}{\varphi }^2f_2)\}`$ (67)
where
$`f_0(\stackrel{~}{\varphi })`$ $`=`$ $`(\stackrel{~}{\varphi }^ϵ1)(\stackrel{~}{\varphi }^ϵ\frac{c_+}{c_{}})`$ (68)
$`f_{n+1}(\stackrel{~}{\varphi })`$ $`=`$ $`\left(1+\frac{1}{4n}\varphi _{\stackrel{~}{\varphi }}\right)f_n(\stackrel{~}{\varphi })`$ (69)
In the following, it will be helpful to keep in mind that $`Z`$ can be extremely small, of order $`\mathrm{\Omega }`$ in (42) when $`ϵ`$ is small, whereas $`\sqrt{\lambda }`$ tends to be closer to unity, depending on the mass of the radion and the definition of the TeV scale (35):
$$\sqrt{\lambda }=\frac{ϵ^{3/4}}{2\sqrt{6}}\frac{m_\varphi }{M_{\mathrm{TeV}}}\frac{(1\sqrt{ϵ_1+ϵ_4})}{\left(\sqrt{}1+\frac{4}{\widehat{m}_1}+2\frac{\sqrt{ϵ}}{\widehat{m}_1}\right)^{1/2}}$$
(70)
In the rescaled variables, the value $`\stackrel{~}{\varphi }=1`$ corresponds to the first zero of the potential, which would be the starting point of the bounce if energy was conserved in the mechanical analog problem, i.e., if there was no viscous damping term $`\varphi ^{}/r`$ in the equation of motion. The actual starting point turns out to have a value in the range $`\stackrel{~}{\varphi }_01.53`$ because of this. The rescaled action $`\stackrel{~}{S}`$ depends mainly on the model parameters $`ϵ`$ and $`\widehat{m}_1`$, for the parameter values which are of interest to us. All the sensitive exponential dependence on $`ϵ`$, namely the factor $`c_{}^{1/ϵ}`$, is removed from $`\stackrel{~}{S}`$. Numerically we find that
$$\stackrel{~}{S}\frac{2}{3}\widehat{m}_1(ϵ\widehat{m}_1)^{3/4},$$
(71)
except when $`\widehat{m}_1`$ becomes close to $`ϵ`$. For $`\widehat{m}_1`$ slightly smaller than $`ϵ`$, $`ϵ_1`$ starts to exceed 1, and $`c_{}`$ becomes negative, signaling the onset of an instability in the radion potential toward coincidence of the two brane positions. Figure 3 shows the dependence of $`\mathrm{log}(\stackrel{~}{S}/\widehat{m}_1)`$ versus $`\mathrm{log}(ϵ\widehat{m}_1)`$.
Figure 3: $`\mathrm{log}_{10}(\stackrel{~}{S}/\widehat{m}_1)`$ versus $`\mathrm{log}_{10}(ϵ\widehat{m}_1)`$, where $`\stackrel{~}{S}`$ is the rescaled bounce action, eq. (65). The other parameters are $`m_\varphi =T=100`$ GeV and $`M_{\mathrm{TeV}}=1`$ TeV.
We have computed the bounce solution and the corresponding action for a range of parameters $`ϵ`$ and $`\widehat{m}_1`$ which can be consistent with the solution to the hierarchy problem (i.e., that $`\widehat{\varphi }10^{16}`$ at the global minimum). The size of the bounce in position space, measured as the width at half-maximum, is small near $`ϵ=\widehat{m}_1`$, and reaches a larger constant value as $`\widehat{m}_1\mathrm{}`$. Using the rescaled radial variable $`\stackrel{~}{r}=r\sqrt{\lambda }T`$, the dependence of the width on $`ϵ`$ and $`\widehat{m}_1`$ is shown in figure 4.
Figure 4: Half-width $`\stackrel{~}{w}`$ of the bounce solution, in terms of the rescaled radial distance $`\stackrel{~}{r}=r\sqrt{\lambda }T`$, versus $`\mathrm{log}_{10}(\widehat{m}_1)`$, for the same parameters as in figure 3.
The action of the bounce can be much greater than or much less than 1, depending on the parameters: for $`ϵ\widehat{m}_11`$, $`S1`$, while for larger values of $`ϵ`$ and $`\widehat{m}_1`$, $`S1`$. Where the crossover occurs ($`S1`$) depends on $`m_\varphi `$, $`T`$ and $`M_{\mathrm{TeV}}`$. This behavior can be inferred from figure 3 (showing $`\stackrel{~}{S}`$) and the dependences of the coefficient in the relation $`S=(Z^2/\sqrt{\lambda })\stackrel{~}{S}`$. Rather than presenting further results for $`S`$ directly however, we will turn to the more relevant quantity, the rate of bubble nucleation. For this we need to determine the prefactor of $`e^S`$ in the rate.
### B Prefactor of Bubble Nucleation Rate
The bounce action is the most important quantity determining the rate of tunneling, since it appears in the exponent of the rate (61). Since we do not have a model for the inflation and reheating of the universe which must occur prior to the bubble nucleation, hence an exact prediction for the reheating temperature which enters the rate, it would not be worthwhile to compute the prefactors in eq. (61) very accurately; however we can estimate their size.
First, consider the frequency $`\omega _{}`$ of the unstable mode. $`\omega _{}^2`$ is the negative eigenvalue of the Schrödinger-like equation for small fluctuations $`\delta \varphi `$ around the bounce solution, which we will denote by $`\varphi _b(r)`$:
$$\left(\delta \varphi ^{\prime \prime }+\frac{2}{r}\delta \varphi ^{}\right)+\frac{^2V_{\mathrm{eff}}}{\varphi ^2}|_{\varphi _b(r)}\delta \varphi =\omega _{}^2\delta \varphi $$
(72)
Rescaling the radius and background field exactly as in eqs. (62-63), eq. (72) becomes
$`(\delta \varphi ^{\prime \prime }`$ $`+`$ $`{\displaystyle \frac{2}{\stackrel{~}{r}}}\delta \varphi ^{})+U(\stackrel{~}{r})\delta \varphi ={\displaystyle \frac{\omega _{}^2}{\lambda T^2}}\delta \varphi ;`$ (73)
\[
Figure 5: The potential for small fluctuations around the bounce solution, $`U(\stackrel{~}{r})=X^2V^{\prime \prime }(\varphi _b(\stackrel{~}{r}))`$, as a function of $`\stackrel{~}{r}`$, for the parameters $`ϵ=0.02,`$ $`\widehat{m}_1=5`$ (left) and $`ϵ=0.1`$, $`\widehat{m}_1=25`$ (right). \]
$`U(\stackrel{~}{r})`$ $`=`$ $`{\displaystyle \frac{1}{\lambda T^2}}{\displaystyle \frac{^2V_{\mathrm{eff}}}{\varphi ^2}}(\stackrel{~}{\varphi }_b(r))`$ (74)
$`=`$ $`f_4+X{\displaystyle \frac{3\sqrt{\lambda }}{\pi }}\left(f_4\sqrt{\frac{c_+}{c_{}}+X}+{\displaystyle \frac{12(Z\stackrel{~}{\varphi }f_3)^2}{\sqrt{\frac{c_+}{c_{}}+X}}}\right)`$ (75)
$`+`$ $`{\displaystyle \frac{6c_b}{8\pi ^2}}\lambda \left(\frac{c_+}{c_{}}f_4+X(f_4+2f_3^2/f_2)\right)`$ (76)
where primes now denote $`\frac{d}{d\stackrel{~}{r}}`$, $`X=12Z^2\stackrel{~}{\varphi }^2f_2`$, and the $`f_i`$ are defined in eq. (68). Except when the radion mass is signficantly larger than 100 GeV, $`\lambda `$ is much smaller than 1, and $`U`$ is dominated by the first two terms in (73). Of these, the first ($`f_4`$) always dominates if $`Z1`$, while the second ($`X`$) can be important near $`r=0`$ if $`Z\text{ }\stackrel{>}{}\text{ }1`$. The two different cases are illustrated in figure 5. Since $`1<\stackrel{~}{\varphi }_0^ϵ<\frac{c_+}{c_{}}`$, both $`f_4`$ and $`X`$ are negative at $`r=0`$, so that
$`U_0U(0)`$ $``$ $`ϵ^{3/2}\sqrt{}1+\frac{4}{\widehat{m}_1}`$ (77)
$`\times `$ $`\left(2+\mathrm{ln}\stackrel{~}{\varphi }_0+12Z^2\stackrel{~}{\varphi }_0^2(\frac{1}{2}+\mathrm{ln}\stackrel{~}{\varphi }_0)\right)`$ (78)
and thus the smallest eigenvalue of eq. (72) is negative. This is the unstable mode of the saddle point solution, with imaginary frequency of order
$$\omega _{}^2U_0\lambda T^2.$$
(79)
Recall that $`|\omega _{}|`$ appears in the prefactor of the nucleation rate $`\mathrm{\Gamma }/V`$.
As $`\stackrel{~}{r}\mathrm{}`$, $`U(\stackrel{~}{r})`$ approaches a maximum value
$$U_{}U(\mathrm{})=\frac{c_+}{c_{}}\left(1\frac{3}{\pi }\sqrt{\lambda \frac{c_+}{c_{}}}+\frac{3c_b}{4\pi ^2}\lambda \frac{c_+}{c_{}}\right),$$
(80)
which determines the asymptotic behavior of the fluctuations at large $`\stackrel{~}{r}`$: $`\delta \varphi e^{\sqrt{U_0}\stackrel{~}{r}}`$. The fluctuations around the false vacuum state ($`\varphi =0`$) thus have a mass given by
$$m^2=V_{\varphi =0}^{\prime \prime }=U_{}\lambda T^2,$$
(81)
which will be relevant for the following.
Next we must estimate the functional determinant factor,
$$𝒟=\frac{\stackrel{}{det}(_\tau ^2^2+V^{\prime \prime }(\varphi _b))}{det(_\tau ^2^2+V^{\prime \prime }(0))},$$
(82)
where $`\tau `$ is imaginary time ($`\tau [0,1/T]`$), $`^2`$ is the three-dimensional Laplacian, $`\varphi _b`$ is the bounce solution, and the prime on $`det^{}`$ means that the three translational zero-mode eigenvalues must be omitted from the determinant for fluctuations around the bounce. These zero modes correspond to spatial translations of the bubble solution. Because of their removal, $`𝒟`$ has dimensions of (energy)<sup>-6</sup>, as is required to get a rate per unit volume in eq. (61).
Ref. has given a thorough account of how to compute $`𝒟`$ by a method which was discussed for one-dimensional systems in . In 3-D one should classify the eigenvalues of the fluctuation operators by the quantum numbers $`n`$ and $`l`$, denoting Matsubara and angular momentum excitations, respectively. Then $`𝒟`$ can be written as a product, $`𝒟=_{n,l}𝒟_{n,l}`$.
Ref. shows that the contribution to $`𝒟`$ from the $`l`$th partial wave can be expressed, to leading order in a perturbative expansion in the potential $`U(r)`$, as
$$𝒟_{n,l}\left(1+h_l^{(1)}\right)^{2l+1}.$$
(83)
The quantity $`h_l^{(1)}`$ has the Green’s function solution
$`h_l^{(1)}`$ $`=`$ $`2{\displaystyle _0^{\mathrm{}}}𝑑rrI_{l+1/2}(\kappa r)K_l(\kappa r)\left(V^{\prime \prime }(\varphi _b(r))m^2\right)`$ (84)
$`=`$ $`2{\displaystyle _0^{\mathrm{}}}𝑑\stackrel{~}{r}\stackrel{~}{r}I_{l+1/2}(\frac{\kappa \stackrel{~}{r}}{\sqrt{\lambda }T})K_l(\frac{\kappa \stackrel{~}{r}}{\sqrt{\lambda }T})\left(U(\stackrel{~}{r})U_{}\right)`$ (85)
using the modified Bessel functions $`I`$ and $`K`$, and the mass $`m`$ of the field in the false vacuum, eq. (81). For general Matsubara frequencies, $`\nu =2\pi nT`$, one defines $`\kappa =\sqrt{m^2+\nu ^2}`$.
The subdeterminant for the $`n=0`$ (zero-temperature) sector of the theory has the usual ultraviolet divergences of quantum field theory, namely the vacuum diagram $``$ (the dot represents one insertion of $`V(\varphi )`$), which should be absorbed by renormalization of the zero of energy for the radion potential. Since we are not attempting to solve the cosmological constant problem here, we are going to ignore all of this and compute only the factor $`𝒟_{0,1}`$, which contains the translational zero modes—or more precisely, which has the zero modes removed. This removal is accomplished by replacing
$$1+h_1^{(1)}\frac{dh_1^{(1)}}{d\kappa ^2}$$
(86)
Notice that this quantity has dimensions of (mass)<sup>-2</sup>, and there are $`2l+1=3`$ such factors, so that $`|𝒟|^{1/2}`$ has dimensions of (mass)<sup>3</sup>, as required. From eq. (84) one can show that
$`{\displaystyle \frac{dh_1^{(1)}}{d\kappa ^2}}={\displaystyle \frac{1}{\lambda T^2U_{}^2}}I_U;`$ (87)
$`I_U{\displaystyle _0^{\mathrm{}}}𝑑yy^2\left(I_{3/2}(y)K_1(y)\right)^{}\left(U(\frac{y}{\sqrt{U_{}}})U_{}\right).`$ (88)
We have numerically evaluated the integral $`I_U`$ for each set of parameters. Our estimate for the fluctuation determinant factor in the nucleation rate can then be written as
$$|𝒟|^{1/2}\left(\frac{\lambda T^2U_{}^2}{I_U}\right)^{3/2}$$
(89)
### C Results for Nucleation Rate
Putting the above ingredients together to find the rate of bubble nucleation per unit volume, $`\mathrm{\Gamma }/V`$, we see that the latter depends on five undetermined parameters: $`ϵ`$, $`\widehat{m}_1`$, $`m_\varphi `$, $`M_{\mathrm{TeV}}`$ and the temperature $`T`$. Ref. showed that, as long as the energy density on the TeV brane is much less than $`M_{\mathrm{TeV}}^4`$, the usual 4-D effective theory governs the Hubble rate:
$$H^2=\frac{\rho }{3M_p^2},$$
(90)
where $`\rho `$ is the total energy density,
$`\rho `$ $`=`$ $`g_{}{\displaystyle \frac{\pi ^2}{30}}T^4+\rho _\varphi ;`$ (91)
$`\rho _\varphi `$ $`=`$ $`\frac{1}{2}\dot{\varphi }^2+V(\varphi )V(\varphi _+)`$ (92)
$``$ $`\frac{3}{8}M_{\mathrm{TeV}}^2m_\varphi ^2{\displaystyle \frac{\sqrt{}1+\frac{4}{\widehat{m}_1}(1+\sqrt{ϵ_1})}{\sqrt{}1+\frac{4}{\widehat{m}_1}+2\frac{\sqrt{ϵ}}{\widehat{m}_1}}}`$ (93)
We take the number of relativistic degrees of freedom, $`g_{}`$, to be 100. The kinetic energy of the radion is zero since $`\varphi =0`$ in the false vacuum, so $`\rho _\varphi `$ is essentially the potential energy of the radion in the false vacuum, assuming the 4-D cosmological constant is zero: $`V(\varphi )V(\varphi _+)=|V(\widehat{\varphi }_+)|`$, which is given by eq. (40). Depending on the parameters, this can be comparable in size or dominate over the energy density of radiation. Using our estimates for the prefactor of the tunneling rate, the logarithm of the ratio of $`\mathrm{\Gamma }/V`$ to $`H^4`$ can be written as
$$\mathrm{ln}\frac{\mathrm{\Gamma }}{VH^4}\mathrm{ln}\left(\frac{\lambda ^2\sqrt{U_0}U_{}^3}{(2\pi )^{5/2}}T^4\left(\frac{3M_p^2}{\rho }\right)^2\left(\frac{S}{I_U}\right)^{3/2}\right)S$$
(94)
where $`S`$ is the action of the bounce solution. The criterion for completion of the phase transition to the true vacuum state is that $`\mathrm{ln}(\mathrm{\Gamma }/VH^4)>0`$. The saddle point approximation leading to eq. (61) is only valid if the action $`S`$ is not much less than 1. Otherwise, the barrier is not effective for preventing the field from rolling to the true minimum, as in a second order phase transition. This situation occurs in the vicinity of $`\mathrm{ln}(\mathrm{\Gamma }/VH^4)150`$ in the following results; thus the transition region where $`\mathrm{\Gamma }/VH^4=1`$ is well within the realm of validity of the approximation.
In figure 6 we show the contours of constant $`\mathrm{ln}(\mathrm{\Gamma }/VH^4)`$ in the plane of $`\mathrm{log}_{10}(\widehat{m}_1)`$ and $`\mathrm{log}_{10}(ϵ)`$, starting with the fiducial values $`T=m_\varphi =100`$ GeV, $`M_{\mathrm{TeV}}=1`$ TeV for the other parameters, and showing how the results change when any one of these is increased. The dependences can be understood from the prefactor $`Z^2/\sqrt{\lambda }`$ in the action, eq. (64):
$$\frac{Z^2}{\sqrt{\lambda }}ϵ^{3/4}\frac{M_{\mathrm{TeV}}^3\mathrm{\Omega }^2}{T^2m_\varphi },$$
(95)
where we recall that $`Z`$ and $`\lambda `$ are given by (62-63) and $`\mathrm{\Omega }`$ by (42). The factor $`\mathrm{\Omega }`$ is responsible for suppressing the bounce action when $`ϵ1`$ or $`ϵ\widehat{m}_11`$, explaining the shape of the allowed regions in each graph. Nucleation of bubbles containing the true minimum becomes faster when the temperature or the radion mass is increased, but slower if the definition of the TeV scale in increased. These dependences are dictated not only by the size of the barrier between the two minima in the effective potential, but also by the size of the bubbles.
Interestingly, the borderline between allowed and forbidden regions of parameter space falls within the range which is relevant from the point of view of building a model of radion stabilization. That is, some choices which would otherwise have been natural and acceptable are ruled out by our considerations. We see furthermore that the choice of $`\widehat{m}_1\mathrm{}`$, as was effectively focused on
\[
Figure 6: Contours of $`\mathrm{ln}(\mathrm{\Gamma }/VH^4)`$ in the plane of $`\mathrm{log}_{10}(\widehat{m}_1)`$ versus $`\mathrm{log}_{10}(ϵ)`$. The shaded regions are where the tunneling rate is too small for the phase transition to complete. Figure (a) has $`T=m_\varphi =100`$ GeV, $`M_{\mathrm{TeV}}=1`$ TeV. The other figures are the same except for the following changes: (b) $`T=400`$ GeV; (c) $`m_\varphi =400`$ GeV; (d) $`M_{\mathrm{TeV}}=2`$ TeV. \]
in ref. , is not the optimal one for achieving a large nucleation rate.
It might be thought that our analysis is rendered less important by the fact that one can always obtain fast nucleation simply by going to high enough temperatures. However it must be remembered that the TeV scale functions as the high-energy cutoff in the Randall-Sundrum scenario: the whole semiclassical description breaks down at super-TeV scales, where quantum gravity effects start to become important. From this point of view, the temperatures of $`100300`$ GeV which we are discussing are already rather high, and a fairly efficient mechanism of reheating at the end of inflation will be needed to generate them.
## V Discussion
In this paper we have presented a somewhat simpler model of radion stabilization by a bulk field ($`\psi `$) than that of Goldberger and Wise ; although the physics is qualitatively identical, we are able to write the radion potential exactly, and thus explore the effect of letting the stabilizing field’s VEV’s on the branes be pinned more or less strongly to their minimum energy values. One such effect is that the mass of the radion can be significantly increased for small values of the parameter $`m_1`$, which is the coefficient of the potential for $`\psi `$ on the TeV brane. Moreover if $`m_1/k\widehat{m}_1`$ is accidentally close to $`ϵ`$, approximately the minimum value consistent with a stable potential, the radion mass can start to diverge, by the factor $`(1\frac{ϵ}{4}\frac{ϵ}{\widehat{m}_1})^{1/2}`$. This modifies somewhat the expectation expressed in ref. that the radion mass will be small relative to the TeV scale, due to a factor of $`ϵ^{3/4}`$.
Our main focus was on the problem that the radion potential has a local minimum at infinite brane separation, and that the barrier between the true and false minima is so small that for generic initial conditions, one would expect the true minimum to be bypassed as the radion field rolls through it. We showed that for a large range of parameters, the high-temperature phase transition to the true minimum is able to complete, thus overcoming the problem. There are however significant constraints on the model parameters, and the initial temperature after inflation, to insure this successful outcome.
There remain some outstanding issues. The form of the radion effective potential is such that the field is able to reach $`\varphi =0`$ in a finite amount of time; yet $`\varphi =0`$ represents infinite brane separation in the extra dimension. This paradoxical situation may be due to the assumption that the stabilizing field, $`\psi `$, is always in its minimum energy configuration at any given moment. In reality $`\psi `$ must require a finite amount of time to respond to changes in the radion. Thus one should solve the coupled problem for time-varying $`\varphi `$ and $`\psi `$ to do better. This is probably a difficult problem, which we leave to future study.
A related question is whether it is correct to treat thermal fluctuations of the radion field $`\varphi `$ analogously to a normal scalar field with values in the range $`(\mathrm{},\mathrm{})`$. Since $`\varphi `$ is related to the size of the extra dimension by $`\varphi =fe^{kb}`$, its range is $`[0,f]`$. We have not studied what effect this might have on the thermal part of the effective potential; instead we assumed that the usual treatment suffices.
Another approximation we made was to ignore the back-reaction of the stabilizing field on the geometry. Ref. has given a method of finding exact solutions to the coupled equations for the warp factor $`a(y)`$ and the stabilizing field $`\psi (y)`$. This method cannot be applied in the present case because it works only for bulk scalar potentials with a special form that, among other things, requires them to be unbounded from below.In this is asserted to be allowed because of special properties of anti-de Sitter space; however we believe that the real reason the bulk potential can be unbounded from below is that the potentials on the branes have the correct sign to prevent an instability. Moreover, since the method of generates only static solutions to the equation of motion, it cannot be used to deduce the radion potential, which is a probe of the response of the geometry when it is perturbed away from a static solution. On the other hand, does show that the neglect of the back reaction is justified for the parameter values which most closely resemble the Goldberger-Wise model.
One might at first feel uneasy about using a 4-D effective description of the problem when in reality our initial condition is a universe with an infinitely large extra dimension. In the Randall-Sundrum scenario, however, this is justified because the graviton is trapped on the transverse length scale of $`1/k1/M_p`$, rather than the size of the 5th dimension (see also ). Moreover the 4-D Friedmann equation (90) was shown by ref. to be valid without actually assuming the radion to be stabilized.<sup>§</sup><sup>§</sup>§The extra dimension is free to expand in this case, and the kinetic energy of the radion simply appears as an additional contribution to the energy density of the universe, as in eq. (90-91). Difficulties with the “wrong” rate of expansion ($`H\rho `$ instead of $`H\sqrt{\rho }`$ ) arise only when one fine-tunes the brane energy densities to prevent radion motion even in the absence of stabilization, which we are not doing here. In any case, changing the form of the expansion rate would have a small effect on our results since this alters only the logarithmic term in (94), not the overwhelmingly dominant term $`S`$.
The problem of shallow barriers in moduli potentials is not unique to the Randall-Sundrum scenario, and a new idea for addressing it was recently presented in ref. . The coupling of the kinetic terms of matter fields to the modulus can give the damping of the modulus motion needed to make it settle in the true minimum in some cases. This effect might provide an alternative to the thermal mechanism we have discussed here.
We thank Mark Wise and Rob Myers for enlightening discussions, Guy Moore for helpful correspondence and Geneviève Boisvert for perceptive criticisms of the manuscript.
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# Hyperbolic character of the angular moment equations of radiative transfer and numerical methods.
## 1 Introduction
In standard problems of radiation<sup>1</sup><sup>1</sup>1 We will use the term radiation for both photons and neutrinos hydrodynamics (RH) where radiation contains a large fraction of the energy and momentum density, the Boltzmann Equation (BE) must be coupled to the hydrodynamic equations in order to obtain the evolution of the system as well as the spectrum and angular distribution of the radiation field. However, an algorithm built to solve the BE numerically (a Boltzmann solver) in a non-stationary case is too time-consuming, from a computational point of view, to allow for the extension to more than one dimension of the existing numerical codes \[Mezzacappa & Bruenn 1993, Yamada, Janka & Suzuki 1999\]. In many cases, instead of the BE, its angular moments are considered, obtaining then the multigroup (or multi-frequency) equations or the even simpler energy averaged equations.
Standard RH methods \[Mihalas & Mihalas 1984\] have found their way into the literature, and one can find several RH codes with different approaches for the radiation transport part, from single-energy-group (or two temperature) such as VISPHOT \[Ensman & Burrows 1992\] or TITAN \[Gehmyer & Mihalas 1994\], to multigroup radiative transfer, such as STELLA \[Blinnikov & Bartunov 1993\]. Other authors built codes devoted more specifically to the radiation transport, but at the expense of detailed hydrodynamics. An example is the code EDDINGTON \[Eastman & Pinto 1993\] in which free expansion is assumed. The hydrodynamics part of most of the existing codes is a one dimensional implicit finite difference scheme, including artificial viscosity terms for problems requiring an accurate treatment of shock waves and discontinuities. During the last decade, a new subclass of numerical methods, the so-called Godunov-type methods, has been gradually substituting the schemes based on numerical viscosity due to their easier extension to multidimensional cases and their greater capabilities in the treatment of strong shocks. The lack of a radiation transport method fully compatible with hydrodynamical Godunov-type schemes is one of the motivations of these paper.
Turning back to the radiation transport part, another important issue related to the angular moment equations is the closure relation. By taking angular moments of the distribution function the complexity of BE is highly reduced, but the information about the angular dependence of the radiation field is partially lost. Each $`n^{th}`$ angular moment equation of the BE contains angular moments of higher order, thus any truncated hierarchy of moment equations contains more unknowns than equations and must be supplied with additional equations or closure relations \[Cernohorsky & Bludman 1994, Groth et al. 1996\]. Theoretically, if the correct closure for a given problem is known, the solution obtained for the first moments of the distribution function by solving the moment equations should be the same as the solution obtained by solving the BE.
Although the moment equations are much simpler to solve than the BE, in many situations an additional simplification is made by neglecting some other terms, leading to the diffusion approximation (DA). In this approximation, however, the resulting energy flux can be higher than the limit predicted by causality, specially in the regions where the radiation mean free path becomes comparable to the characteristic length. Different extensions of DA, like flux-limited diffusion \[Bowers & Wilson 1982\] or artificial opacities \[Dgani & Janka 1992\] have being used to overcome this problem. These extensions only partially solve the breakdown of the DA in the semi-transparent and transparent layers, and all of them are based on the same idea: to reintroduce some terms which had been dropped in the original assumptions of the DA. The term which is usually kept out of the equations is the time derivative of the fluxes. By neglecting this term, the character of the system of equations is changed to parabolic, and causality is therefore violated (disturbances propagate at infinite velocity). In order to develop a method which is valid in all regimes (from diffusion to free streaming) while preserving causality, one must solve the full set of equations keeping its hyperbolic character.
In this paper, we address the problem of establishing a well-defined hyperbolic system of equations for the first angular moments of the BE in the non-stationary spherically symmetric case. We will see that, besides the important question of its hyperbolicity, some constraints on the mathematical properties of the closure relations can be derived. They might help to disregard among some of the closures previously proposed. Moreover, the behaviour of the characteristic fields of the hyperbolic system, which gives information on the velocity of perturbations, can also be used to establish additional constraints to the closures.
A direct consequence of considering the moment equations in hyperbolic form is that it permits the application of powerful numerical techniques that have been developed in recent years for hyperbolic systems of conservation laws (the equations of classical and relativistic fluid dynamics for perfect fluids, for example). Among the different numerical techniques, the so-called high resolution shock capturing (HRSC) methods have a number of interesting features such that stability, being conservative, convergence to physical solutions and high accuracy in regions where the solution is smooth. HRSC methods are based on the resolution of local Riemann problems (an initial value problem with discontinuous data) at the interfaces of numerical shells, ensuring a consistent treatment of discontinuities and steep gradients \[Godunov 1959\]. Their special relativistic extension \[Martí, Ibáñez & Miralles 1991\] has shown its potential in simulations of heavy ion collisions and extra-galactic jets, and different attempts to extend the method to General Relativity have been done \[Banyuls et al. 1997, Pons et al. 1998\]. We refer the interested reader to the recent reviews by Martí & Müller \[Martí & Müller 1999\] and Ibáñez & Martí \[Ibáñez & Martí 1999\] for a detailed description of the current status of HRSC techniques in numerical relativistic hydrodynamics.
Although HRSC are specially designed for hyperbolic systems without source terms which corresponds to the transparent regime, we will show in numerical experiments that, by appropriate treatment of the sources, accurate solutions can be obtained also in the diffusion regime, where the source terms are dominant.
The structure of the paper is the following: In §2 we summarise the deduction of the angular moment equations of the BE in a spherically symmetric case, and for a static background. We also describe the general form of the collision terms when emission-absorption and iso-energetic scattering processes are included. In §3 we discuss the most common techniques used to solve the transport equations, such as diffusion or flux-limited diffusion, remarking their main features and limitations. In §4 the hyperbolic character of the equations and its implications for the closure relations are discussed. In §5, the numerical techniques employed to solve the transport equations as a hyperbolic system of conservation laws are discussed. A number of numerical experiments solving the transport equations in several test problems are displayed in §6, and the feasibility of the hyperbolic treatment in all kinds of regimes is demonstrated. Finally, main conclusions and advantages of our proposal are summarised in §7.
## 2 Transport equations in spherical coordinates
We shall start our discussion from the radiative transfer equations in a static medium, deferring to a future work the inclusion of fluid velocity terms in the equations. Although the restriction to zero velocity seems to be specific, there are some astrophysical scenarios where the assumption that matter is at rest is a reasonable approach. Moreover, the inclusion of velocity terms, in some cases of interest does not change the essentials of our conclusions.
In a static background, the Boltzmann Equation for massless particles in spherical coordinates can be written as follows<sup>2</sup><sup>2</sup>2We work in units where $`c=\mathrm{}=1`$.
$$\omega \left[\frac{}{t}+\mu \frac{}{r}+\frac{(1\mu ^2)}{r}\frac{}{\mu }\right]=\left(\frac{d}{d\tau }\right)_{coll}$$
(1)
where $`=(t,r,\omega ,\mu )`$ is the invariant distribution function, $`\mu `$ is the cosine of the angle of the particle momentum with respect to the radial direction, $`\omega `$ the energy of the particle and the right hand side term is the invariant source or collision term coming from the interaction between the radiation field and matter.
As stated in the introduction, a common method of solving the BE is the method of moments \[Thorne 1981\], which involves taking angular moments of the equation by applying the operator
$$\frac{1}{4\pi }_1^{+1}_0^{2\pi }\mu ^i𝑑\mu 𝑑\mathrm{\Phi }i=0,1,2,\mathrm{}$$
(2)
Denoting by $`E`$, $`F`$ and $`P`$ the angular moments of the specific intensity $`g\omega ^3`$, $`g`$ being the statistical weight ($`g=2`$ for photons and $`g=1`$ for neutrinos)
$`E=E(t,r,\omega )=g{\displaystyle \frac{\omega ^3}{2}}{\displaystyle _1^{+1}}𝑑\mu ,`$
$`F=F(t,r,\omega )=g{\displaystyle \frac{\omega ^3}{2}}{\displaystyle _1^{+1}}𝑑\mu \mu ,`$
$`P=P(t,r,\omega )=g{\displaystyle \frac{\omega ^3}{2}}{\displaystyle _1^{+1}}𝑑\mu \mu ^2,`$ (3)
the equations corresponding to the first two moments (0 and 1) of the BE are:
$`_tE+_rF+{\displaystyle \frac{2F}{r}}=s^0`$ (4)
$`_tF+_rP+{\displaystyle \frac{3PE}{r}}=s^1`$ (5)
where the moments of the collision term are defined as
$$s^0=g\frac{\omega ^2}{2}_1^{+1}\left(\frac{d}{d\tau }\right)_{coll}𝑑\mu $$
(6)
$$s^1=g\frac{\omega ^2}{2}_1^{+1}\left(\frac{d}{d\tau }\right)_{coll}\mu 𝑑\mu $$
(7)
Notice that, in the absence of velocity fields or strong gravitational fields, the different energy groups in a multifrequancy scheme are only coupled through the source terms, rendering the mathematical character of the left hand side of equations (4)-(5) formally identical to the energy-averaged problem.
### 2.1 Collision terms
The most general form for the collision term including emission, absorption and isoenergetic scattering is the following:
$$\left(\frac{d}{d\tau }\right)_{coll}=\omega \left(j\kappa _a+\kappa _s[]\right)$$
(8)
where $`j`$ is the emissivity, $`\kappa _a`$ the absorption opacity including final states blocking (for fermions) or stimulated emission (for bosons) and $`\kappa _s`$ the scattering opacity related to the scattering reaction rate ($`R^s`$) through
$$\kappa _s(\omega )=\left(\frac{\omega }{2\pi }\right)^3_1^{+1}𝑑\mu ^{}_0^{2\pi }𝑑\phi R^s(\omega ,\mathrm{cos}\mathrm{\Theta })((\omega ,\mu ^{})(\omega ,\mu ))$$
(9)
where $`\mathrm{\Theta }`$ is the angle between the in-going and outgoing particle and $`\phi `$ is the azimuthal angle. The opacities and the emissivity are expressed in units of inverse length.
For isotropic scattering the source terms appearing in equations (4) and (5) can be written as
$$s^0=\kappa _a\left(E^{eq}E\right)$$
(10)
$$s^1=\kappa F$$
(11)
where $`\kappa =\kappa _a+\kappa _s`$ and $`E^{eq}`$ is the value of $`E`$ in equilibrium with matter. In the above, we have assumed matter in local thermodynamic equilibrium (LTE), and only the radiation field (photons or neutrinos) is allowed to deviate from equilibrium. The conclusions of our study might be affected by large velocity fields or for frequencies close to line discontinuities, which are present in some astrophysical scenarios, and deserve a more careful study \[Mihalas et al. 1976, Kunasz 1983, Mihalas & Kunasz 1986\].
## 3 Flux-limited diffusion and artificial opacities
One of the approaches most widely used to solve the transport equations numerically is the diffusion approximation, in which the invariant distribution function is assumed to be nearly isotropic in the comoving frame. Thus an expansion in terms of Legendre polynomials to the order $`O(\mu ^2)`$ is enough to maintain the main features of the radiation field
$$(\mu )=_0+3_1\mu $$
(12)
where $`_1_0`$. Consistent with this assumption, the time derivative of the flux is neglected in equation (5), and together with equation (7) gives the following relation for the flux in terms of the energy gradient
$$F(\omega )=\frac{1}{3\kappa (\omega )}\frac{E(\omega )}{r}$$
(13)
As stated before, the DA breaks down when the mean free path is large compared to the typical scale of the problem, and the fluxes calculated with the former formula may give non-causal values, in the sense that the flux can be greater than the energy density. This pathology comes from the fact that we have neglected some terms in the momentum equation in order to obtain a simple formula for the fluxes.
Flux limiters have been introduced ad hoc in the literature to avoid this non-causal behaviour; see, e.g., Minerbo \[Minerbo 1978\], Levermore & Pomraining \[Levermore & Pomraining 1981\] or Cernohorsky & Bludman \[Cernohorsky & Bludman 1994\]. In order to illustrate where the problems stem from, we begin by deriving the flux-limited diffusion equations. Let us define the flux factor, $`f`$, and the Eddington factor, $`p`$:
$$f=\frac{F}{E}$$
(14)
$$p=\frac{P}{E}$$
(15)
By subtracting equation (4) multiplied by $`f`$ from (5), and after some algebra (now keeping all terms), one obtains:
$$F=\frac{(pf^2)}{(\kappa _{ph}+\kappa _J)}_rE$$
(16)
where the quantities, including the flux and Eddington factors, are now energy dependent. The explicit expression for the different opacities are
$$\kappa _{ph}=\kappa _s+\kappa _a\frac{E^{eq}}{E}$$
(17)
$`f\kappa _J=_tf+_rpf_rf+{\displaystyle \frac{(3p12f^2)}{r}}`$ (18)
The physical opacity, $`\kappa _{ph}`$, or effective albedo only depends on the interaction of radiation with matter and includes an out–of–equilibrium correction to absorption opacity. The second term in the denominator, $`\kappa _J`$, is the multigroup extension of the so-called artificial opacity \[Dgani & Janka 1992\] and takes into account geometrical corrections and deviations from the diffusion limit (it vanishes in the limit $`f1`$).
The term in the numerator of equation (16), $`(pf^2)`$, plays the role of a flux–limiter. It has the value $`\frac{1}{3}`$ in the diffusion limit and goes to zero as the free streaming limit ($`f1`$ , $`p1`$) is reached. To close the set of equations, a closure relation for $`p`$ is needed, and it is directly related to flux limiter $`(pf^2)`$. We will re-address the closure relation problem in next section.
With this description of the transport equations, the simple structure of the diffusion equation and its parabolic character is maintained when the first term in $`\kappa _J`$ is neglected ($`_tf`$ = 0). The geometric and velocity field corrections can also be included through artificial opacities as corrections to the physical opacity. Although the flux-limiter helps to avoid non-causal behaviour in the sense $`F>E`$, causality is also violated when the time derivative of $`f`$ is neglected, because the character of the equations change to parabolic, which results in an infinite speed of propagation of the information. By keeping this term, the diffusion scheme is, in some applications, numerically unstable and further numerical tools are needed to handle these instabilities. This is the subject of the next section.
## 4 Hyperbolic Formulation of Transport Equations
The system of equations (4) and (5) can be written in a compact form as follows
$$\frac{\stackrel{}{𝒰}}{t}+\frac{1}{r^2}\frac{\left(r^2\stackrel{}{}\right)}{r}=\stackrel{}{𝒮}$$
(19)
where the vectors $`\stackrel{}{𝒰}`$ and $`\stackrel{}{}`$ are
$$\stackrel{}{𝒰}=(E,F)$$
(20)
$$\stackrel{}{}=(F,pE)$$
(21)
and the sources $`\stackrel{}{𝒮}`$, according to equations (10-11), are given by
$$\stackrel{}{𝒮}=(\kappa _a(E^{eq}E),\kappa F+\frac{EP}{r}).$$
(22)
Notice that the sources are free of partial derivative operators.
The above system of equations (19) is said to be a hyperbolic system of conservation laws (with source terms) when the Jacobian matrix associated to the fluxes ,$`_\stackrel{}{𝒰}\stackrel{}{}`$, has real eigenvalues and there exists a complete set of corresponding eigenvectors. To analyse the Jacobian matrix, we will assume that the Eddington factor is only a function of the flux factor $`p=p(f)`$. This assumption is commonly used by a number of authors who have proposed different closure relations or flux-limiters in the context of flux-limited diffusion. In this case, the Jacobian matrix is
$$_\stackrel{}{𝒰}\stackrel{}{}=\left(\begin{array}{ccc}& 0& 1\\ & pfp^{}& p^{}\end{array}\right)$$
(23)
where $`p^{}={\displaystyle \frac{dp}{df}}`$ .
The eigenvalues of the above matrix are
$$\lambda _\pm =\frac{p^{}\pm \sqrt{p_{}^{}{}_{}{}^{2}+4(pfp^{})}}{2}$$
(24)
and the right eigenvectors are
$$\stackrel{}{r}_+=(1,\lambda _+);\stackrel{}{r}_{}=(1,\lambda _{})$$
(25)
The mathematical character of the system of equations (19) depends on the value of the discriminant in equation (24), which can be written as follows
$$\mathrm{\Delta }=(p^{}2f)^2+4(pf^2).$$
(26)
Since $`p`$ and $`f`$ are normalised moments of a non–negative weight function, on the unit sphere, there is an extra relation between $`f`$ and $`p`$, the Schwarz inequality
$$f^2p1,$$
(27)
that ensures that for any closure satisfying (27) the eigenvalues are always real and hence the hyperbolic character is guaranteed. In addition, causality imposes that the velocity of propagation of perturbations must be smaller than the speed of light, which gives another constraint on the eigenvalues
$$\left|\lambda _\pm \right|1,$$
(28)
that can be expressed in terms of $`f`$, $`p`$ and $`p^{}`$ as follows
$$\frac{1p}{1+f}p^{}\frac{1p}{1f}.$$
(29)
This restriction, together with equation (27) and the condition $`|f|1`$, defines the region in the $`[f,p,p^{}]`$ space allowed for causal closures.
In many astrophysical problems, radiation diffuses out from a central opaque region to an outer optically thiner region. In such a scenario, the radiation angular distribution is nearly isotropic in the central region ($`|f|1`$) and forward peaked in the transparent region, far from the opaque centre ($`f1`$). Thus, the closure in that problem must satisfy
$$p(f=0)=\frac{1}{3}\mathrm{and}p(f=1)=1.$$
(30)
The first limit comes from the diffusion approximation while the second limit comes from the Schwarz inequality (27), and both are satisfied by all the closures found in the literature (see Table I). Another general property is that, in the diffusion regime the Eddington factor has the constant value of $`1/3`$, thus another restriction on the derivative can be imposed
$$p^{}(f=0)=0,$$
(31)
that leads to the characteristic speeds $`\lambda _\pm =\pm \frac{1}{\sqrt{3}}`$. In the limit $`f=1`$, equations (27) and (29) gives the bounds $`0p^{}2`$ and the eigenvalues in this limit are $`\lambda _+=1`$ and $`\lambda _{}=p^{}1`$. According to the theory of hyperbolic systems, this result can be interpreted as a first wave propagating outwards at the speed of light and a second wave with a characteristic speed that depends on the value of $`p^{}`$. Although $`\lambda _{}`$ depends on the form of the closure, for $`f=1`$ the amplitude of this second wave is zero and the value of $`p^{}`$ does not affect the evolution. However, for $`f`$ close but not equal to one, the choice of $`p^{}`$ can be important, as we will show in section §5.2. Similar conclusions can be obtained for the case $`f=1`$.
In table 1 we summarise the main properties of some of the most widely used closures: Minerbo \[Minerbo 1978\], LP \[Levermore & Pomraining 1981\], the linear formula by Auer \[Auer 1984\], Kershaw’s parabolic approach \[Kershaw 1976\], the flux-limiters used by Bruenn \[Bruenn 1985\] and Bowers & Wilson \[Bowers & Wilson 1982\], Janka’s fit to Montecarlo simulations of neutrino transport in Supernovae \[Janka 1991\] and that corresponding to an isotropically emitting sphere \[Cooperstein, van den Horn & Baron 1986\], named vacuum.
There are several comments about Table I that we would like to address. First, one closure relation (BW) leads to non-causal values of one eigenvalue in the diffusion regime. Secondly, four closures do not satisfy the condition $`p^{}(f=0)=0`$ (Auer, Bruenn, BW, Vacuum). Auer’s formula failure is due to an excessive simplification (it is a linear approach) while ’Vacuum’ has been obtained from the assumption of an isotropically radiating sphere, which restricts their validity to $`f0.5`$. The closure relations proposed by Kershaw, Janka, Minerbo and LP have the common features of preserving causality leading to reasonable values of the characteristic velocities ($`\lambda _\pm =\pm \frac{1}{\sqrt{3}}`$) in the diffusion regime. Thirdly, it is worthy to point out that all the closures except Auer’s give $`p^{}1`$ when approaching the radially-streaming limit $`(f=1)`$, which makes both eigenvalues be positive. It ensures that no information can be propagated inwards for sufficiently high values of $`f`$. It is clear that these closures cannot be used to treat problems in which a small perturbation is propagating inwards from the transparent region, as we will show in next section.
In a more general situation, the Eddington factor depends on both, the flux factor and the occupation density corresponding to the zeroth angular moment of the distribution function, $`e=E/\omega ^3`$. The Jacobian matrix associated to the fluxes is
$$\frac{\stackrel{}{}}{\stackrel{}{U}}=\left(\begin{array}{cccc}& 0& & 1\\ & & & \\ & pf\frac{p}{f}|_e+e\frac{p}{e}|_f& & \frac{p}{f}|_e\end{array}\right)$$
and now the eigenvalues are
$$\lambda _\pm =\frac{1}{2}\left(\frac{p}{f}|_e\pm \sqrt{\mathrm{\Delta }}\right)$$
(32)
where the discriminant is
$$\mathrm{\Delta }=\left(\frac{p}{f}|_e2f\right)^2+4\left(pf^2+e\frac{p}{e}|_f\right)$$
(33)
In contrast to (26), the sign of the discriminant might be negative for negative values of $`\frac{p}{e}|_f`$, and the system would become elliptic instead of hyperbolic. As stated before, hyperbolic equations are related to finite speeds of propagation of the perturbations, while elliptic equations are related with boundary problems and do not allow for solutions with propagating signals. This fact indicates that, in order to preserve hyperbolicity, only closures which give positive values of the discriminant should be considered.
A closure of this kind has been proposed by Cernohorsky & Bludman \[Cernohorsky & Bludman 1994\] using arguments of maximum entropy. Minerbo’s, LP’s and the Vacuum closures can be obtained as limiting cases of CB closure in the low occupation limit (classical statistics), large ocuppation limit for Bose-Einstein statistics and maximal forward packing limit for Fermi-Dirac statistics, respectively. We have checked that CB closure gives real and causal eigenvalues for both, Bose-Einstein and Fermi-Dirac statistics, in any regime.
## 5 The numerical approach.
The numerical code employed for the calculations presented in §6 is a finite finite difference scheme in which the left (L) and right (R) states for the Riemann problems set up at each interface have been obtained with a monotonic, piecewise linear reconstruction procedure \[van Leer 1979\]. The numerical fluxes are evaluated following the idea of Harten, Lax & van Leer \[Harten, Lax & van Leer 1983\], that is based on approximate solution of the original Riemann problem with a single intermediate state. The resulting numerical flux is given by
$$\widehat{}^{HLL}=\frac{\psi ^+_L\psi ^{}_R+\psi ^+\psi ^{}(𝒰_R𝒰_L)}{\psi ^+\psi ^{}}$$
(34)
where $`\psi ^+=\text{max}(0,\lambda _+^R,\lambda _+^L)`$ and $`\psi ^{}=\text{min}(0,\lambda _{}^R,\lambda _{}^L)`$.
The time-advance algorithm is a TVD-preserving, third order Runge-Kutta that takes into account the influence of the source terms in intermediate steps. The source term in this problem deserves a special mention. It acts as a relaxation term which leads to a long-time behaviour governed by the reduced parabolic system described in section §3. Although high resolution methods for hyperbolic conservation laws usually fail to capture this feature (unless a very fine grid is used), the use of a modified flux \[Jin & Levermore 1996\] and an implicit linearization of the source term in each Runge-Kutta step allows us to succeed in the treatment of such stiff relaxation term with coarse grids and time-steps longer than the source time-scale ($`1/\kappa `$), as we will show in subsection §6.1. Although we have constrained the numerical experiments to situations when only scattering processes are permitted, the inclusion of absorption-emission processes can be treated in a similar way. However it introduces energy and momentum exchange with the matter and makes necessary to solve the equations of radiation transport coupled with the equations of hydrodynamics or the equations of hydrostatic equilibrium with number and heat transfer. Such a problem is out of the scope of this paper, in which we focussed on the radiation part, and will be addressed in forthcoming works.
## 6 Numerical experiments and discussion.
To check the applicability of HRSC techniques to radiation transport problems, we have performed a number of numerical tests considering several simplified models in which the analytical solution is known. In the remainder of this section we describe the different tests and display the results. Our intention is to prove that our method can solve radiation transport in all kinds of regimes, from diffusion (TEST 1) to radially streaming (TEST 2), going through the semi-transparent region (TEST 5) and in situations where strong gradients are developed (TEST 4).
Ideally, in a multi-group scheme we have as many systems of two equations as energy bins in which the radiation density is splitted. In a radiation hydrodynamics problem all those pairs of equations are coupled to the matter equations and must be solved together. Since our objective is to study the applicability of HRSC techniques for the radiation part, in what follows we assume to have only one energy group or, equivalently, we consider the energy integrated equations. Notice, though, that the conclusions are not affected since, when velocity fields are not taken into account, the equations for different energies are only coupled through the source terms. A different problem will arise when velocity fields are important or gravitational and Doppler red-shift are considered. In those situations either the anisotropy of the source term (in the observer frame) or the advection and aberration terms (in the comoving frame) cause the energy redistribution along the photon (or neutrino) path and a more careful treatment needs to be done.
### 6.1 TEST 1: Diffusion limit
In the diffusion limit, the closure is obviously $`p=1/3`$ and the time derivative of the flux can be neglected, leading to equation (13). Setting the absorption opacity to zero and the scattering opacity to a constant value, the analytical solution of the system formed by (4) and (13), corresponding to an initial value problem consisting in a Dirac delta located in $`r=0`$ at $`t=0`$, is the following
$$E(r,t)=\left(\frac{\kappa }{t}\right)^{3/2}\mathrm{exp}\left\{\frac{3\kappa r^2}{4t}\right\}$$
(35)
$$F(r,t)=\frac{r}{2t}E(r,t)$$
(36)
We imposed $`F(r=0)=0`$ as inner boundary conditions and outflow boundary conditions in the outermost shell.
We have performed two simulations. First, we set a value of $`\kappa =100`$ on a equally spaced grid of 100 shells in the domain $`[0,1]`$. It corresponds to a Peclet number ($`Pe=\kappa \mathrm{\Delta }x`$) of order unity. The cell Peclet number is a measure of the spatial resolution relative to the relaxation scale and a grid is coarse when $`Pe1`$. This first simulation starts at $`t=1`$ to avoid the numerical problems produced by an infinite value in $`t=0`$. Results are shown in Figure 1. The energy density (left panel) and the flux (right panel) are displayed for different times in dimension-less units, $`t=1,2,3,5`$ from top to bottom. The solid line corresponds to the analytical solution and the crosses correspond to the numerical solution. Notice that, although we are finding numerical solutions of system (19), the analytical solution of the diffusion limit is well reproduced by the numerical code, since the opacity is large enough to be close to the diffusion limit. The second simulation has been performed in a grid of 100 shells in the domain $`[0,0.5]`$, and using a value of $`\kappa =10^5`$, which corresponds to $`Pe=5\times 10^4`$. The initial time in this case is 200 and the Courant time-step is of the same order ($`10^2`$) as in the previous simulation, three orders of magnitude larger than the relaxation scale. Results are displayed in Fig. 2. The agreement between the analytical and the numerical solutions clearly proofs the capability of the code to handle situations with very stiff source terms.
### 6.2 TEST 2: Radially Streaming limit
In the radially streaming limit ($`|f|=1`$, $`p=1`$), and assuming there is no interaction between radiation and matter ($`\kappa _a=0`$, $`\kappa =0`$) equations (19) are reduced to
$`{\displaystyle \frac{E(\omega )}{t}}+{\displaystyle \frac{1}{r^2}}{\displaystyle \frac{(r^2E(\omega ))}{r}}=0`$ (37)
The solution to this equation is any linear combination of ingoing and outgoing spherical waves, i.e., any linear combination of functions of the form $`E(t,r)=\frac{1}{r^2}g(r\pm t)`$. For the first numerical experiment we choose a Gaussian spherical wave propagating outwards,
$$E(r,t)=\frac{1}{r^2}\mathrm{exp}\left\{(rt)^2\right\}$$
(38)
$$F(r,t)=E(r,t).$$
(39)
Our numerical grid consists on 200 equally spaced shells going from $`r_1=0.2`$ to $`r_2=10.2`$ and the boundary conditions consist in imposing the analytical solution in both boundaries. We show the results of this test in Figure 3. In the figure we displayed the analytical solution (solid line) and the numerical (crosses) for the energy density. In order to analyse the effect of the particular choice for the closure, we repeated the experiment using different values for $`p^{}(f=1)`$, which leads to different speeds of propagation of the waves, and no difference was observed. The final reason is that all the closures displayed in Table 1 give $`\lambda _+(f=1)=1`$. This is a general property of any closure, since $`pf^2`$ is a positive defined function that goes to zero as $`f1`$ and, therefore, its derivative has to be negative in a vicinity of $`f=1`$ and $`p^{}2f`$. Thus taking the limit $`f1`$ in (24), one obtains
$$\lambda _+=\frac{p^{}+\sqrt{(2p^{})^2}}{2}=1$$
(40)
To understand why the value of $`\lambda _{}`$ is not important in this problem, we study the diagonalized system obtained by adding and subtracting the two equations in (19) after taking $`p=1`$,
$`{\displaystyle \frac{(E+F)}{t}}+{\displaystyle \frac{1}{r^2}}{\displaystyle \frac{(r^2(E+F))}{r}}=0`$ (41)
$`{\displaystyle \frac{(EF)}{t}}{\displaystyle \frac{1}{r^2}}{\displaystyle \frac{(r^2(EF))}{r}}=0.`$ (42)
In this equations it is clearly seen that the quantities propagated with velocities $`\lambda _\pm `$ (Riemann invariants) are $`E\pm F`$, respectively. Since in our particular problem we took $`F=E`$, the quantity propagated by $`\lambda _{}`$ is identically null and the value of the speed of propagation does not affect the solution.
A different situation is found when we add to the initial conditions a second wave propagating in the opposite direction. Figure 4 shows the results from such experiment, in which a second wave with smaller amplitude and opposite direction is located initially at $`r=7`$. The closure used was $`p=1`$, $`p^{}=0`$, which gives $`\lambda _+=1`$, $`\lambda _{}=1`$. In the figure, the analytical solution is denoted by the solid lines and the numerical solution by crosses and each panel corresponds to a different time. It can be seen how the code is able to solve the interaction of the two waves until they completely cross over. The difference between the analytical and numerical solution near the center for the last panel are due to the boundary conditions (initially, it was an ingoing wave).
To illustrate the effect of a wrong speed of propagation, the experiment is repeated with Kershaw’s closure, that gives a correct speed of propagation when $`f=\pm 1`$ but a wrong value otherwise. Results are shown in Figure 5 for the same time steps as in the previous case. At the beginning both waves seem to propagate correctly but, as they cross each other and the flux factor departs from unity, the numerical result differs from the analytical solution. A similar behaviour is obtained with other closures. This simple example manifests the importance of the choice for the closure, since it fixes the speeds of propagation of the waves.
### 6.3 TEST 3: Sphere radiating in the vacuum.
In the two previous tests we studied the limiting cases, diffusive and free-streaming regimes. Another academic problem in the limit of vanishing opacities consists in a sphere of radius $`R`$ radiating isotropically into vacuum (Figure 6). Given the explicit form of the distribution function at the surface of the sphere, $`_0(t)=(t=0,r=R)`$, the general solution can be easily obtained
$$(t,r,\mu )=\{\begin{array}{cc}_0(ts)& \mu x(r,t)\\ 0& \mu <x(r,t)\end{array}$$
(43)
where
$$s=r\left(\mu \sqrt{\mu ^21+(R/r)^2}\right)$$
(44)
$$x(r,t)=\{\begin{array}{cc}\sqrt{1(R/r)^2}& t\sqrt{r^2R^2}\\ \frac{t^2+r^2R^2}{2tr}& rR<t<\sqrt{r^2R^2}\\ 1& t<rR\end{array}$$
(45)
where $`\mu =\mathrm{cos}\alpha `$.
We take a linear time dependence $`_0(t)=Ct`$, being $`C`$ an arbitrary constant, so that the angular integrals can be done analytically. The angular moments of the distribution function are, then, given by
$$E(t,r)=\frac{C}{4}\left[2t(1x)r(1x^2)+r\sqrt{1x^2}+rx^2\mathrm{log}\left(\frac{x}{1+\sqrt{1x^2}}\right)\right]$$
(46)
$$F(t,r)=\frac{C}{12}\left[3t(1x^2)2r(1x^3)+2r(1x^2)^{3/2}\right]$$
(47)
$$P(t,r)=\frac{C}{48}\left[8t(1x^3)6r(1x^4)+3r(2x^2)\sqrt{1x^2}+3rx^4\mathrm{log}\left(\frac{x}{1+\sqrt{1x^2}}\right)\right].$$
(48)
We set up a numerical grid of 200 points, uniformly distributed, from $`r=1`$ to $`r=200`$. The initial conditions are $`F=0`$ and $`E=\alpha /r^2`$ with $`\alpha `$ being a small number, different from zero for numerical purposes. We have taken $`\alpha =10^4`$. The radiation field is introduced through the inner boundary conditions, $`E(t,r=R)=2t`$ and $`f(t,r=R)=0.5`$. Outflow is assumed in the outer boundary.
In Figure 7 the evolution of the radiation energy density as a function of the radial coordinate is shown for three different evolution times. The test has been repeat for three different closure relations from those described in table 1: Vacuum (dashed lines), Kershaw (dotted lines) and Minerbo (dashed-dotted lines). As in the previous figures the analytical solution is represented by the solid line. Although the appropriate closure for this problem (vacuum) gives a better approximation to the real solution, it is remarkable that other closures also give similar results. The differences between different closures stem from the slightly different propagation speeds of the simple waves in the regime of interest. Unlike in the two-wave problem described in previous section, for this specific problem any reasonable closure might be used without major deviations from the real solution.
### 6.4 TEST 4: Radiative cooling and heating
With this test we intend to simulate a more realistic scenario in which, initially, radiation produced in the center of a star diffuses out with constant luminosity. We begin by setting up a sphere which is in radiative equilibrium at a given luminosity $`L`$. For $`p=\frac{1}{3}`$, the stationary solution is given by
$$E(r)=\frac{L}{4\pi }\left[\frac{1}{R^2}+3\kappa \left(\frac{1}{r}\frac{1}{R}\right)\right]$$
(49)
$$F(r)=\frac{L}{4\pi r^2}$$
(50)
where $`R`$ is the radius of the outer boundary, at which $`f=1`$. Starting with the initial model corresponding to $`L=4\pi `$, we let it evolve in three different situations: 1) we take $`L=0`$ as boundary conditions at the inner boundary (located at $`r=1`$), and the radiation field should diffuse out continuously; 2) we change the inner boundary condition to $`L=40\pi `$, ten times larger than the initial model. The behaviour is now the opposite, the energy density must increase until the new stationary solution for the new luminosity is reached; 3) we again take $`L=0`$ at the inner boundary but we switch on a central point–source in the middle of the star ($`r=6`$). The radiation field must evolve towards the corresponding stationary solution, which is different in the two parts of the star defined by the source. Inside, the luminosity is zero and the energy density is constant. In the region external to the energy source, the stationary solution is defined by $`L=\dot{s}`$, where $`\dot{s}`$ is the energy per unit time injected by the source. The results for $`\kappa =100`$ are shown in Figures 8-10, for cases 1 to 3, respectively. All three simulations have been done on the same grid of 200 uniform zones and outflow boundary conditions in the outer edge have been imposed.
Left panels display the energy density as a function of radius. The thick solid line in the figures corresponds to the initial model and the thick dotted line in Figure 9 represent the stationary solution with $`L=40\pi `$. In Figure 9, we can see how the heat wave lasts about 1000 time-units (the diffusion time $`R\kappa `$) to arrive to the surface and a new stationary solution is successfully obtained after $`15`$ times the diffusion time scale.
### 6.5 TEST 5: The semi-transparent regime.
Testing a method in the semi-transparent regime becomes a tough issue due to a lack of analytical solutions to compare with. To overcome this problem, we have performed Monte Carlo simulations, which is the closest solution to an exact one, and results in an excellent test of the performance of the method in the semi-transparent regime.
We have used a Monte Carlo code that simulates the transfer of a radiation field in a spherically symmetric geometry. We consider a region with inner boundary at $`r=R_{in}`$ and outer boundary at $`r=R_{out}`$, divided in 200 spherical shells. Within each shell a constant value of the scattering opacity (elastic and isotropic) is assumed. We do not consider emission or absorption processes. Our Monte Carlo procedure starts by injecting particles at the inner boundary outwardly-directed. Then, we compute the trajectory of each particle as it scatters off matter, according to the standard random walk laws \[Lucy 1999\]. The trajectory ends when the particle escapes by crossing the outer boundary or it re-enters the inner boundary. Once the path of a particle is computed we proceed to calculate the contribution of this particle to the moments of the radiation field. The contribution of each particle to the energy density at a given numerical shell is calculated by adding the path length between two successive scattering events, within the shell, divided by the volume of the shell. Consequently, the total energy density is obtained just by adding the contribution of each particle multiplied by an arbitrary factor Analogously, to obtain the contribution of each particle to the flux, the path length is weighted using the cosine of the angle with respect to the radial direction. Finally, pressure is calculated by means of the same procedure, but using the square of the cosine as weight of the path length.
In our numerical experiment we have taken: $`R_{in}=10`$ and $`R_{out}=100`$. The opacity law is an exponentially decreasing function in radius, taking the value $`\kappa =2`$ at the inner boundary and $`\kappa =2\times 10^5`$ at the outer boundary. The corresponding values of the optical depth are such that an important part of the region is semi-transparent. We have run the Monte Carlo code using $`\mathrm{500\hspace{0.17em}000}`$ particles in order to obtain the stationary solution.
The set up for our transport code is the following: i) the above opacity law, ii) the inner boundary conditions given by the stationary solution (obtained from Monte Carlo), and iii) the closure relation. We have performed two different calculations using Kershaw’s closure relation and a cubic polynomial fit to Monte Carlo results, explicitly
$$p=0.3228+0.1902f0.0476f^2+0.5131f^3$$
(51)
Starting with arbitrary initial profiles, we let the radiation field evolve until the stationary solution is reached and results can be compared to those obtained from the Monte Carlo simulations.
In Figure 12 we show the stationary solutions calculated from Monte Carlo simulations (solid line) and from our code with the two different closures, Kershaw’s (dotted line) and the above fit (dashed line). From the results displayed in the figure we can conclude that the evolutionary code describes the semi-transparent regime accurately. The differences in the values of the energy density (left panel), normalized to its inner boundary value, amounts less than a few percentage in any case. The slight discrepancies in the flux factor are mainly due to the closure relation employed, being less than 6% for Kershaw’s closure and less than 3% for the fit. The closure, therefore, sets the maximum accuracy that can be obtained, although qualitative results at a level of 10% uncertainty can easily be obtained with any reasonable closure in the semi-transparent regime.
### 6.6 A toy model.
Our last experiment intends to illustrate how the numerical method we proposed can succeed in situations when standard flux-limited diffusion techniques fail. We set up a sphere of radius $`r=5`$ with a radiation energy density of $`E=1`$, embedded in a background where the opacity varies as follows:
$$\kappa =\{\begin{array}{cc}\kappa _0& r1\\ \kappa _0/r^2& 1<r10\\ 0& 10<r15\\ \kappa _0& 15<r17\\ 0& r17\end{array}$$
(52)
where we have taken $`\kappa _0=10^3`$. Outside the sphere we take $`E=0`$ (in practice, $`E=10^6`$ for numerical reasons), and initially the flux is zero everywhere. It tries to mimic a situation in which radiation scatters out in a central star, where the diffusion approximation remains valid, through an opacity-varying atmosphere up to the surface ($`r=10`$), where the opacity suddenly vanishes. Then, it crosses a transparent region to find a very opaque layer, where diffusion operates again. This sort of problem is obviously a tough challenge for diffusion-based codes, since they cannot work properly in the transparent region.
In Figure 12 we show results from the evolution of this problem on a grid of 200 uniform shells and using Kershaw’s closure. In the left panel we show the energy density for different times. It can be seen how a gradient is developed in the diffusive regions $`(r<10,15<r<17)`$ while radiative equilibrium $`(dE/dr=0)`$ in the transparent region between the two opaque layers is rapidly obtained. Notice also that only for the last snapshot ($`t=1000`$, dashed curve) there is a flux of energy in the outer transparent region, because it takes some time to the radiation to diffuse through the opaque layer and develop a gradient that keeps the outgoing flow.
## 7 Conclusions
We have studied the mathematical character of the angular moment equations of radiative transfer, considering the implications of using different closure relations. The role played by the closure and their derivatives is more apparent in the hyperbolic formulation. After this analysis, we conclude that the hyperbolic character is assured for a given variety of closures widely used in the literature: those given by $`p=p(f)`$. For general closures $`p=p(e,f)`$, hyperbolicity can not be guaranteed, although the closure proposed by Cernohorsky & Bludman \[Cernohorsky & Bludman 1994\], based on maximum entropy arguments, leads to real eigenvalues. Additional constraints on the closures are imposed on the base of the behaviour of the eigenvalues of the Jacobian matrix since they give the velocity of propagation of perturbations. Causality limitation is written in terms of the closure relation helping us to select, from this point of view, among the different closures reviewed. It turns out that only one of the closures analysed does not satisfy causality, allowing for velocities of the waves higher than the speed of light. For the other closures, although derived without taking into account the mathematical issues addressed in this paper, we found that they are consistent with hyperbolicity and causality.
Writing the moment equations of the distribution function as a hyperbolic system of conservation laws allows one to apply numerical techniques specifically designed for such systems, the so-called HRSC methods. We have applied HRSC methods to solve the transport equations in a static background showing, for a number of test problems, the feasibility of the method. From the results, the main conclusion is that the numerical method can be used in any regime, from optically thick to transparent regions, obtaining numerical solutions with high accuracy. It turns out that the closure relation plays an important role, more evident than in traditional methods for radiation transfer, that can be useful to choose the more appropriate closure for a given problem.
When the radiation transport equations are written as a system of conservation laws, the coupling with hydrodynamics is straightforward \[Lowrie, Morel & Hittinger 1999\]. This feature make us to be confident in the possibility of applying the method to problems involving radiative flows. Moreover, we are optimistic in the applicability to multidimensional problems, since the extension of the method is straightforward.
## Acknowledgements
This work has been supported by the Spanish DGICYT (grant PB97-1432). J.A.P. also acknowledges previous financial support from the Ministerio de Educación y Cultura.
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# Probing High Redshift Radiation Fields with Gamma-Ray Absorption
## 1 Introduction
Our knowledge of UV radiation fields and energy injection into the IGM at $`z>5`$ is fairly tenuous. There are two main constraints: observations of the integrated background light (Madau & Pozzetti 2000, Bernstein et al 1999), and the fact that no Gunn-Peterson trough is observed in the spectra of the highest-redshift quasar to date (Fan et al 2000), implying that the universe must be reionized by $`z=5.8`$. The upcoming Next Generation Space Telescope (NGST) will be able to image high-redshift star clusters or AGNs in rest frame UV continuum emission (Haiman & Loeb 1997,1998), and their redshifts may be obtained via H$`\alpha `$ observations (Oh 1999). Nonetheless, the redshift-binned number counts will be fairly sparse, and one is unlikely to probe sufficiently far down the luminosity function to get a good measure of the comoving emissivity as a function of redshift.
Observations of gamma-ray blazars (“grazars”) probe extragalactic IR and UV radiation fields, by observing the pair production opacity to $`\gamma `$ rays at the high energy end (Gould & Schreder 1967, Stecker, De Jager & Salamon 1992, Madau & Phinney 1996, Primack et al 1999). All theoretical models have confined their predictions to low redshift grazars, with the exception of Salamon & Stecker (1998), who computed the $`\gamma `$-ray opacity up to z=3. They concluded that because the stellar emissivity peaks between z=1 and z=2, the $`\gamma `$-ray opacity shows little increase at high redshift, and thus is not dependent on the initial epoch of galaxy formation.
To date, EGRET has detected 66 gamma-ray loud blazars (Hartman et al 1999), out to redshifts $`z>2`$. The next generation of gamma-ray telescopes (GLAST, CELESTE, STACEE, MAGIC, HESS, VERITAS, and Milagro) should greatly enlarge this sample. If the low redshift correlation between black hole mass and bulge mass (Magorrian et al 1998) continues to high redshift, then it is possible that high-redshift halos could host mini-quasars (Haiman & Loeb 1998, Haehnelt, Natarajan & Rees 1998), which should be detectable in rest frame UV emission by NGST and X-ray emission by Chandra (Haiman & Loeb 1998, 1999) in the redshift range $`z515`$. This raises the exciting possibility that grazars could be detected at similarly high redshifts. It is worth noting that EGRET has detected $`56`$ sources at high Galactic latitudes $`b>10^{}`$ (Mukherjee, Grenier & Thompson 1997), with no known counterparts at other wavelengths. Their spatial distribution and log N- log S plot can be well fit by a Galactic component plus an isotropic, extragalactic contribution. Some of these may well be unidentified high-redshift blazars.
In this paper, I point out that if grazars are detected at high redshifts $`z>3`$, the pair production opacity to gamma ray photons can be used to probe the comoving emissivity longward of the Lyman break at these extremely high redshifts, independent of the star formation rate at lower redshifts. Due to the small escape fraction of ionizing photons $`f_{esc}<5\%`$ from host galaxies, as well as the high photoelectric opacity of the IGM at these wavelengths, the comoving number density of UV photons exhibits a sharp drop at the Lyman edge at all redshifts. Thus, there is only a limited pathlength over which a gamma-ray photon can pair produce against UV photons, before it redshifts to energies which require UV photons above the Lyman edge for pair production to take place. For $`z<z_{break}`$, the universe becomes optically thin to the gamma-ray photon. Thus, the detection of an absorption edge in a high-redshift grazar places an immediate constraint on the mean radiation field over the redshifts $`z_{break}<z<z_s`$. Furthermore, measurement of the different absorption at a given observed energy between blazars at redshifts $`z_1,z_2`$ places an immediate constraint on the radiation field in the redshift range $`z_1<z<z_2`$. Detection of grazars at a number of redshifts would then enable one to probe the UV emissivity history of the universe.
In all numerical estimates, we assume a background cosmology given by the ’concordance’ values of Ostriker & Steinhardt (1995): $`(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda },\mathrm{\Omega }_b,h,\sigma _{8h^1},n)=(0.35,0.65,0.04,0.65,0.87,0.96)`$.
## 2 Gamma-Ray Blazars
A detailed study of the detectability of high-redshift blazars is beyond the scope of this paper. In this section, I merely show that it is plausible that GLAST will be able to detect high redshift blazars.
With a point source sensitivity of $`S(E>100\mathrm{M}\mathrm{e}\mathrm{V})2\times 10^7`$photons $`\mathrm{s}^1\mathrm{cm}^2`$, EGRET has detected $`66`$ high-redshift blazars out to $`z>2`$ (Hartman et al 1999). The associated gamma-ray luminosities correspond to $`L_\gamma =10^{46}10^{49}\mathrm{erg}\mathrm{s}^1`$, and typically dominate the bolometric luminosity of the AGN, with $`L_\gamma /L_B11000`$. The upcoming gamma-ray telescope GLAST (see http://glast.gsfc.nasa.gov) will be 2 orders of magnitude more sensitive, with a detection threshold of $`S(E>100\mathrm{M}\mathrm{e}\mathrm{V})2\times 10^9`$photons $`\mathrm{s}^1\mathrm{cm}^2`$ for a 5 $`\sigma `$ detection with a 50 hour integration, and $`S(E>1\mathrm{G}\mathrm{e}\mathrm{V})10^{10}`$ photons $`\mathrm{s}^1\mathrm{cm}^2`$ (these thresholds correspond to the same detection limit for a $`L_\nu \nu ^\alpha `$ source spectrum where $`\alpha =1`$). Goals for GLAST include a broad energy coverage from $`10\mathrm{MeV}>300\mathrm{GeV}`$, with a spectral resolution of $`2\%`$ in the $`>`$ 10 GeV range, a field of view of $`>`$ 3 sr, and a source location determination accuracy of 30 arcsec- 5 arcmin. During its lifetime, it will perform an all-sky survey similar to that conducted by EGRET. Sources of the same or somewhat fainter luminosity as those detected by EGRET may be seen by GLAST out to extremely high redshifts, $`z10`$.
Will such luminous sources will be present at high redshift? If the AGN is assumed to emit all its energy at gamma-ray wavelengths at the Eddington luminosity, the inferred black hole mass is extremely high, $`M_{bh}=10^{10}M_{}(L_{edd}/10^{48}\mathrm{ergs}^1)`$. However, there are two reasons why gamma ray sources of high apparent luminosity do not require such massive black holes: (i) even if all photons are radiated isotropically, if most of the radiation emerges at high energies (as appears to be the case in gamma-ray blazars), then Klein-Nishina effects must be taken into account (Dermer & Gehrels 1995). The inferred black hole mass, given by $`M_8^{KN}\frac{3\pi d_L^2(m_ec^2)}{21.26\times 10^{46}\mathrm{erg}\mathrm{s}^1}\frac{F(ϵ_l,ϵ_u)}{(1+z)}\mathrm{ln}[2ϵ_l(1+z)]`$ (where $`F(ϵ_l,ϵ_u)`$ is the observed flux between the lower and upper bandpass limits $`ϵ_l,ϵ_u`$) typically drops by 3 orders of magnitude, so a $`10^{48}\mathrm{erg}\mathrm{s}^1`$ source only requires a $`10^7M_{}`$ black hole. (ii) There is strong evidence for relativistic beaming in blazars (e.g., through the observation of superluminal jets (von Montigny et al 1995)). In fact, if blazars were not beamed, we would not be able to see any gamma-rays from them due to the high pair production opacity at the source; beaming reduces the luminosity/radius ratio by a factor $`\delta ^{p+1}`$, allowing photons to escape (Maraschi et al 1992). Here, if $`\alpha `$ is the spectral index of the source, then $`p=3+\alpha `$ for a moving sphere in the Synchrotron Self-Compton (SSC) model, while $`p=4+2\alpha `$ in the External Radiation Compton (ERC) model; $`\delta =[\gamma (1\beta \mathrm{cos}\theta ]^1`$ is the relativistic Doppler factor, and $`\gamma `$ is the Lorentz factor. If $``$ is the initial intrinsic luminosity of the jet in gamma-ray emission, beaming boosts the observed luminosity of the jet to $`L=\delta ^p`$. For $`\theta 0^{}`$, then $`\delta 2\gamma `$, and the observed luminosity is amplified by a factor of thousands. The strong relativistic beaming reduces the fraction of sources which are visible, since they can only be seen when viewed along the jet axis. For instance, for $`\gamma =6`$, and $`\alpha =1`$, in the SSC model the observed luminosity is reduced by an order of magnitude from its maximum if the jet is pointing $`8.5^{}`$ from our line of sight, and two orders of magnitude if the jet is pointing $`14.2^{}`$ from our line of sight. Note that the black hole masses derived for a number of low redshift blazars from variability timescale and transparency arguments lie in the range $`10^710^8M_{}`$ (Cheng et al 1999, Hartman et al 1996, Becker & Kafatos 1995, Romero et al 2000). The luminosity of these blazars is so high they could be seen at high redshift with GLAST, and black hole masses of $`10^710^8M_{}`$ are reasonably abundant at high redshift in certain models of AGN formation (Haiman & Loeb 1998). Note that for a set of blazars of constant intrinsic luminosity $``$ and comoving number density, the redshift distribution of detected sources in a flux-limited survey flattens considerably and extends to higher redshifts as the Lorentz factor increases (Dermer & Gehrels 1995).
At present, the modelling of even the low-redshift population of gamma-ray blazars is a matter of considerable debate. Models which attempt to account for the unresolved gamma ray background with faint blazars either extrapolate the observed $`\gamma `$-ray luminosity function obtained with EGRET (Chiang & Mukherjee 1998) or use an assumed conversion between the observed radio loud AGN luminosity function and the blazar luminosity function (Stecker & Salamon 1996). In this paper, I use a highly simplified model to estimate the detectability of high-redshift blazars. I assume that the intrinsic luminosity $``$ of the jet in gamma-ray emission (prior to beaming) scales with the luminosity of the accretion disk $`=fL_{disk}^\beta `$, where the optical B-band luminosity is taken to be an accurate reflection of $`L_{disk}`$ (in particular, assuming the median quasar spectrum of Elvis et al (1994), a 1 $`M_{}`$ black hole shining at the Eddington luminosity has a B-band luminosity of $`5.7\times 10^3\mathrm{L}_{\mathrm{B},}`$). Such a jet-disk correlation is observed in the ratio of observed radio (i.e., after beaming) to optical luminosities (Falcke , Malkan & Biermann 1995). I also assume relativistic beaming with $`L=\delta ^p`$, where $`p=3+\alpha `$ (and $`\alpha =1`$, the average spectral index observed in the EGRET blazars). The change in the observed luminosity function due to the effects of relativistic beaming is given by (Urry & Padovani 1995):
$$\mathrm{\Phi }_{obs}=𝑑P(L|)\mathrm{\Phi }_{intr}()$$
(1)
where the probability of observing luminosity $`L`$ given the intrinsic luminosity $``$ is given by:
$$P(L|)=P(\delta )\frac{d\delta }{dL}=\frac{1}{\beta \gamma \delta }^{1/p}L^{(p+1)/p}$$
(2)
I use the fit to the observed B-band luminosity function $`\varphi (L_B,z)`$ from Pei (1995). I adjust the relation $`=fL_{disk}^\beta `$ and the Lorentz factor $`\gamma `$ (note that since the distribution of blazar Lorentz factors is unknown, I assume they all have the same Lorentz factor) to fit the number of blazars detected by EGRET and their redshift distribution. I ignore the effects of blazar flaring, which increases the luminosity by some factor $`A`$ (where typically $`A5`$) some fraction $`\xi `$ of the time (where typically $`\xi 0.03`$), which results in a second term $`\xi \varphi (L/A)`$, since these two degrees of freedom, the normalization (chosen by selecting $`\gamma `$ and thus the beaming angle) and luminosity boost (choosen by the combination $`L=\delta ^pfL_{disk}^\beta `$), are already present in our model. I find that $`_{13}=3.2\times 10^2L_{B,13}^{0.9}`$ (where $`L_{13}=(L/10^{13}L_{})`$ and $`\gamma =6`$ provides a good fit (see Fig 1, top panel). Note that since $`L_{\gamma ,13}<(2\gamma )^pL_{B,13}^{0.9}600L_{B,13}^{0.9}`$ (corresponding to $`\theta =0`$), these relations result in a SED which is in reasonable agreement with the observed SEDs of gamma-ray blazars (see Fig 1 in Ghisellini et al 1998), where $`L_\gamma /L_B10100`$ typically, although it can range from $`11000`$. Furthermore, the Lorentz factor $`\gamma 6`$ is in reasonable agreement with the relativistic Doppler factor $`\delta <2\gamma `$ from models which take into account the SED, time variability, and gamma-ray transparency of blazars; the derived $`\delta 1020`$ (Ghisellini et al 1998). The somewhat lower Doppler factors I have adopted conservatively underestimate the number of high-redshift blazars.
I then extrapolate this model to high redshift by applying it to the Press-Schechter based model of Haiman & Loeb (1998) for quasars. In this model, which is calibrated to the observed luminosity function of Pei (1995) at lower redshifts, each halo with $`T_{vir}>10^4`$K hosts a black hole with mass $`M_{bh}=10^{3.2}M_{halo}`$ which shines at the Eddington luminosity for $`t_o10^6`$ years. The result is shown in the bottom panel of Fig 1, which shows that gamma-ray blazars may be detected out to $`z7`$. By the time GLAST is launched, a large database of quasars with known redshifts will be available (e.g. from the SLOAN digital sky survey, SDSS (York et al 2000)), and many high-redshift blazars can be selected simply by identifying their optical counterparts. I emphasize once again that this highly simplified model is only intended to serve as a plausibility argument. The main point is that while the Press-Schechter formalism predicts that massive halos $`M_{halo}>10^{11}M_{}`$ expected to host supermassive black holes of the requisite luminosity become exponentially rare at high redshift, processes which increase the luminosity of a lower luminosity population (beaming $`L=\delta ^p`$, flaring $`L=AL`$) create a power-law tail of bright sources. I have neglected a tail to the distribution of Lorentz factors, or flaring, which could further flatten the redshift distribution of detectable sources, increasing the maximum redshift out to which sources can be seen. Finally, gravitational lensing could bring otherwise undetectable sources into view, although the low optical depths for strong lensing (e.g. $`\tau 6\times 10^3`$ for $`z_s=7`$, Porciani & Madau 2000) imply that this should only have a small impact on number counts.
## 3 Calculating pair production opacity
The pair-production optical depth for a photon observed at energy $`E_o`$ and emitted from a source at redshift $`z_s`$ is given by (e.g., Madau & Phinney 1996):
$$\tau (E_o,z_s)=_0^{z_s}𝑑z\frac{dl}{dz}_1^1d(cos\theta )(1\mathrm{cos}\theta )_{ϵ_{th}}^{\mathrm{}}𝑑ϵn(ϵ,z)\sigma (E,ϵ,\theta )$$
(3)
where $`E=(1+z)E_o`$, and $`ϵ_{th}`$ is given by the criterion that pair production can take place, which requires that $`Eϵ(1\mathrm{cos}\theta )2(m_ec^2)^2`$. For a given energy E, the pair production cross-section $`\sigma (E,ϵ,\theta )`$ rises sharply from the threshold energy $`ϵ_{th}`$, reaches a peak of $`0.26\sigma _T`$ at $`ϵ=2ϵ_{th}`$, and finally falls off as $`ϵ^1`$ for $`ϵϵ_{th}`$. To calculate the optical depth of the universe to gamma-ray photons emitted at high redshift, we need to know the number density of photons as a function of energy and redshift. Before doing so in detail, I make some simple estimates.
One can perform a simple order of magnitude estimate to show that the minimal comoving emissivity required to reionize the universe implies a high pair-production opacity for $`E>E_{th}`$, but a low opacity for $`E<E_{th}`$, where we define the threshold energy in the rest frame of the source for pair production against UV photons at the Lyman edge as $`E_{th}(m_ec^2)^2/ϵ_L18`$GeV. This is because the number density of UV photons plummets at wavelengths shortward of the Lyman limit; the pair production opacity is thus dominated by soft photons $`ϵ<ϵ_L=13.6`$eV. For the universe to be reionized, a minimum of 1 ionizing photon per baryon must be emitted. Thus, the comoving number density of ionizing photons is $`n_\gamma (E>E_L)>n_b(1+n_{rec})10^7(1+n_{rec})\mathrm{cm}^3`$, where $`(1+n_{rec})`$ is the average number of times each baryon recombines during the reionization epoch. The pair production cross section peaks when $`Eϵ(m_ec^2)^2`$, with a value of $`\sigma _{pp}0.26\sigma _T`$. Thus, across a Hubble volume, the optical depth for photons with $`E<E_{th}`$ is $`\tau n_{\gamma ,proper}(E>E_L)\sigma _{pp}l_H5\times 10^3(1+z/10)^{1.5}(1+n_{rec})`$, which is undetectably small. What is the pair production opacity for photons with $`E>E_{th}`$? We can scale the number density of photons longward of the Lyman limit with respect to the number of ionizing photons produced. There are 3 factors involved: (i) $`f_{esc}`$. Observations of our own Galaxy (Dove, Shull, & Ferrara 2000, Bland-Hawthorn & Maloney 1999) and local starbursts (Leitherer et al 1995) find the escape fraction of ionizing photons from galaxies into the IGM is of order $`f_{esc}36\%`$, and radiative transfer calculations find that the escape fraction should decrease strongly with redshift (Woods & Loeb 1999, Ricotti & Shull 2000). The actual comoving number density of ionizing photons produced in starbursts is $`f_{esc}^1n_b(1+n_{rec})`$, where $`f_{esc}5\%`$ is the escape fraction of ionizing photons from the halo. (ii) $`f_{break}`$. Photoelectric absorption in stellar atmospheres produces a sharp drop in the flux at the Lyman edge. For a Salpeter IMF with metallicity $`10^2Z_{}`$, there are $`f_{break}5`$ times as many photons emitted longward of the Lyman limit as there are shortward of it, integrated over the history of the starburst. (iii) $`f_{opacity}`$, due to intergalactic absorption. The energy density $`U_\gamma \frac{4\pi }{c}J_\nu `$ where $`J_\nu ϵ_\nu \lambda _{mpf}`$ where $`\lambda _{mfp}`$ is the mean free path of an photon of frequency $`\nu `$. While the universe is optically thin to photons longward of the Lyman limit, it is optically thick to ionizing photons, which have a much shorter mean free path. The energy density of ionizing photons is lower by a factor $`f_{opacity}=\frac{\lambda _{mfp}(E<E_L)}{\lambda _{mfp}(E>E_L)}>10`$ due to their high absorption rate. The mean free path of ionizing photons decreases strongly with redshift and for $`z>2`$ the radiation field is largely local, due to the large increase in the number of adsorbers (e.g., Madau, Haardt & Rees 1999). Thus, we have $`\frac{n_\gamma (E<E_L)}{n_\gamma (E>E_L)}f_{esc}^1f_{break}f_{opacity}10^3\left(\frac{f_{esc}}{0.05}\right)^1\left(\frac{f_{break}}{5}\right)\left(\frac{f_{opacity}}{10}\right)`$. This implies that the pair production opacity longward of $`ϵ_L`$ over a Hubble length is much greater, $`\tau 5(1+z/10)^{1.5}(1+n_{rec})`$. A similar estimate may be obtained by normalizing to the observed metallicity of the IGM at $`z3`$, $`Z10^2Z_{}`$; this can be shown to correspond to $`10`$ ionizing photons per IGM baryon (Miralda-Escude & Rees 1997), or a comoving number density of $`50n_b`$ for photons longward of the Lyman break, compared to the previous estimate of $`100\left(\frac{f_{esc}}{0.05}\right)^1\left(\frac{f_{break}}{5}\right)n_b`$.
The expectation of a sharp drop in the intensity of the ambient intergalactic radiation field at the Lyman edge is fairly robust. A recent composite spectrum of 29 Lyman break galaxies (LBGs) exhibit an observed flux ratio L(1500)/L(900)=$`4.6\pm 1.0`$, which is consistent with no internal photoelectric absorption, i.e. $`f_{esc}1`$ (Steidel, Pettini & Adelberger 2000). Note that the authors themselves stress this result should be treated as preliminary; the result could be due to a large number of uncertainties or selection effects, among them the fact that these galaxies were selected from the bluest quartile of LBGs. Even so, if this is typical of all high redshift galaxies, the observed Lyman break (likely due to absorption by stellar atmospheres) and expected processing by intergalactic absorption at high redshift imply a Lyman edge $`\frac{n_\gamma (E<E_L)}{n_\gamma (E>E_L)}>50`$. Similar considerations of internal and intergalactic photoelectric absorption apply if quasars dominate the ionizing radiation field.
As a photon redshifts, it has to interact with higher energy photons to pair produce. A photon is able to pair produce until it redshifts below $`E_{th}`$, i.e., the universe is optically thin to a photon once it redshifts below $`(1+z_{th})\frac{E_{th}}{E_s}(1+z_s)`$, where $`E_s`$ is the original energy of the photon at source redshift $`z_s`$. Thus, there is a fairly well-defined pathlength $`\delta l=\frac{c}{H_o}((1+z_s)^{3/2}(1+z_{th})^{3/2})`$ over which a photon may pair produce. This has the fortunate consequence that low redshift photons do not affect photons emitted near $`E_{th}`$; thus, near the threshold energies one is always probing the average radiation field at redshifts comparable to the source redshift. The flux decrement at a given frequency measures the average number density of photons redward of the Lyman break over the associated redshift interval, $`n_\gamma \tau /(\mathrm{\Delta }l\sigma _{pp})`$. Gamma ray absorption measurements thus provide a fairly clean measurement of radiation fields at high redshift which are uncomplicated by radiative transfer effects since (apart from unimportant $`\mathrm{H}_2`$ opacity effects) the universe is optically thin to photons redward of the Lyman limit. In particular, photons longward of the Lyman limit establish a homogeneous, isotropic radiation field early in the history of the universe. As the pathlengths for pair production opacity to become significant are typically of order a Hubble volume, opacity fluctuations due to Poisson fluctuations or source clustering are insignificant. The measured background radiation field may be compared directly against the expected background from measurements of the source luminosity function by direct imaging of high redshift sources in mid-IR with NGST, where fast photometric redshifts may be obtained using the Gunn-Peterson break. A comparison of the two should in principle allow one to check the completeness of a survey at NGST flux limits.
Let us now use a specific model of high-redshift star formation to calculation the abundance of UV photons at high redshift. The solution of the cosmological radiative transfer equation yields the mean specific intensity of the radiation background at the observed frequency $`\nu _o`$, as seen by an observer at redshift $`z_o`$ as (Peebles 1993):
$$J(\nu _o,z_o)=\frac{1}{4\pi }_{z_o}^{\mathrm{}}𝑑z\frac{dl}{dz}\frac{(1+z_o)^3}{(1+z)^3}ϵ(\nu ,z)e^{\tau _{\mathrm{eff}}(\nu _o,z_o,z)}$$
(4)
where $`\nu =\nu _o(1+z)/(1+z_o)`$, and $`\tau _{\mathrm{eff}}(\nu _o,z_o,z)`$ is the effective optical depth of the IGM to radiation emitted at $`z`$ and observed at $`z_o`$ at frequency $`\nu _o`$. Redward of the Lyman edge, radiative transfer is particularly simple as the universe is optically thin, and only the redshifting of photons is important (there is one caveat to this statement: the optical depth of the IGM in the Lyman resonance lines prior to reionization is very large; thus whenever a photon redshifts into a Ly$`\beta `$ or higher order Lyman resonance, it is reprocessed into a Ly$`\alpha `$ and Balmer or lower order line photon (see Haiman, Rees & Loeb 1997)). However, this merely causes a modulation in the spectrum in the 11.2–13.6 eV range, redistributing photons to the Ly$`\alpha `$ and Balmer wavelengths. The large typical energy intervals of target photons (see Fig 3, top panel) implies that this redistribution causes the pair-production opacity to remain the same or increase. I therefore ignore this complication). The number density of photons in an energy interval is then given by $`\frac{dn}{dϵ}=\frac{4\pi }{hc}J_\nu `$. I model the star formation history $`\dot{\mathrm{\Omega }}_{}`$ of the high redshift universe with the semi-analytic models of Haiman & Loeb (1997), in which a fixed fraction $`f_{star}1.717\%`$ of the gas (normalised to reproduce the observed IGM metallicity at $`z=3`$ of $`Z=10^310^2Z_{}`$) in halos with $`T_{vir}>10^4`$K (which are able to undergo atomic cooling) fragment to form a starburst lasting for $`10^7`$ years, and the halo collapse rate is given by the Press-Schechter formalism (Sasaki 1994):
$$\frac{d\dot{N}^{form}}{dM}(M,z)=\frac{1}{D}\frac{dD}{dt}\frac{dn_{PS}}{dM}(M,z)\frac{\delta _c^2}{\sigma ^2(M)D^2}$$
(5)
where $`D(z)`$ is the growth factor, and $`\delta _c=1.7`$ is the threshold above which mass fluctuations collapse. In this case the star formation rate is given by:
$$\dot{\mathrm{\Omega }}_{}(z)=\frac{1}{\rho _c(z)}\frac{1}{t_o}\frac{\mathrm{\Omega }_b}{\mathrm{\Omega }_m}f_{star}_{t(z)t_o}^{t(z)}𝑑t_{M(T_{vir}=10^4K,z)}^{\mathrm{}}𝑑M\frac{d\dot{N}^{form}}{dM}(M,z)M$$
(6)
An approximate fit to the comoving star formation rate in the interval $`3<z<10`$ can be given by $`\dot{\rho }_{}=\mathrm{exp}(a0+a1z+a2z^2)\mathrm{M}_{}\mathrm{yr}^1\mathrm{Mpc}^3h^3`$, where $`a0=0.841,a1=0.395,a2=0.0295`$. The comoving emissivity is:
$$ϵ_\nu (t)=\rho _c_0^t𝑑t^{}F_\nu (tt^{})\dot{\mathrm{\Omega }}_{}(t^{})$$
(7)
where $`F_\nu (\mathrm{\Delta }t)`$ is the stellar population spectrum, defined as the power radiated per unit frequency per unit initial mass by a generation of stars with age $`\mathrm{\Delta }t`$. I obtain this spectrum from the Bruzual & Charlot (1993) code for a $`Z=10^2Z_{}`$ population (i.e., extremely low metallicity), assuming a Salpeter IMF with lower and upper mass cutoffs at 0.1 and 100 $`M_{}`$, and only computing the emissivity redward of the Lyman break. I ignore the effects of dust extinction, which should be negligible at these high redshifts when the metallicities are very low.
### 3.1 Recombination radiation
Is the number density of recombination line photons such as Ly$`\alpha `$ sufficiently high to cause significant pair production opacity? It has been emphasized (e.g., Loeb & Rybicki 1999) that apart from possible dust attenuation, Ly$`\alpha `$ photons are not absorbed by the IGM but resonantly scattered until they redshift out of resonance. Indeed, the Ly$`\alpha `$ radiation intensity has been predicted to be particularly strong prior to the epoch of reionization (Haiman, Rees & Loeb 1997, Baltz, Gnedin & Silk 1998). This is because the optical depth in the Lyman series is very high, and all the Lyman series lines except Ly$`\alpha `$ are absorbed immediately and redistributed. Thus, the Ly$`\alpha `$ and Balmer lines are considerably brighter because they receive the energy of the Lyman series. This has spawned suggestions of detecting the epoch of reionization by a sharp drop in intensity at the rest-frame Ly$`\alpha `$ wavelength (Baltz, Gnedin & Silk 1998, Shaver et al 1999). In particular, Baltz, Gnedin & Silk (1998) find that Ly$`\alpha `$ is about 3 times brighter and H$`\alpha `$ is about 30 times brighter immediately prior to the reionization redshift.
If this is indeed the case, Ly$`\alpha `$ photons might contribute significantly to the pair production opacity. Let us define $`f_{jump}=n_{recomb,Ly\alpha }/n_{tot,Ly\alpha }`$, the jump in the number density of photons longward of the Ly$`\alpha `$ wavelength. This is given by the number of photons $`n_{recomb,Ly\alpha }`$ injected at the Ly$`\alpha `$ wavelength , over the total number of photons $`n_{tot,Ly\alpha }`$, include those redshifting into resonance. Most studies calculate the number of Ly$`\alpha `$ recombination photons by summing over the IGM recombination rate in numerical simulations, by direct estimation of gas clumping in the simulations. In fact, this underestimates the number of Ly$`\alpha `$ photons produced: if the escape fraction of ionizing photons is small, most recombinations occur in the HII regions of dense halos where star formation takes place, and the primary source of Ly$`\alpha `$ photons are the ionising sources themselves. If each ionizing photon is converted into a Ly$`\alpha `$ photon (plus lower energy photons) at the source, then $`\dot{n}_{Ly\alpha }\dot{n}_{ion}(1f_{esc})`$, where $`f_{esc}`$ is the average escape fraction of ionizing photons from a souce. Thus, inserting $`ϵ(\nu ,z)=E_{Ly\alpha }\dot{n}_{ion}(1f_{esc})\delta (\nu \nu _{Ly\alpha })`$ into equation (4), I obtain for $`\nu <\nu _{Ly\alpha }`$, the solution
$$J_\nu ^{Ly\alpha }(z)=\frac{1}{4\pi }\frac{c}{H_o}h_P(\frac{\nu }{\nu _{Ly\alpha }})^3\frac{1}{(\mathrm{\Omega }_m(1+z_s)^3+\mathrm{\Omega }_\mathrm{\Lambda })^{1/2}}\dot{n}_{ion}(z_L)(1f_{esc})$$
(8)
where $`n_{ion}(z)`$ is the production rate of ionizing photons in proper coordinates, prior to attenuation by the ionizing photon escape fraction. This is a lower bound on $`J_\nu ^{Ly\alpha }(z)`$ since it does not include the conversion of photons trapped in higher order Lyman resonances to Ly$`\alpha `$ photons. I find that for the adopted stellar spectra, in which the Lyman edge is typically $`f_{break}5`$, then typically $`f_{jump}1`$, which implies that Ly$`\alpha `$ photons are comparable to UV continuum photons as a source of opacity. The jump factor $`f_{jump}`$ may easily be understood as the ratio of the total number of Lyman alpha photons produced (or, the total number of ionizing photons produced) against the total number of photons in the $`10.213.6`$eV range. If the break at the Lyman edge is reduced, then $`f_{jump}`$ is increased and Ly$`\alpha `$ photons are more important in contributing to the pair-production opacity. In particular, for a power law spectrum with no discontinuity at the Lyman edge (e.g. as for quasars), Ly$`\alpha `$ photons are the dominant source of opacity. Ly$`\alpha `$ photons are also the dominant source of opacity in scenarios where the universe is reionized by zero metallicity stars. The higher effective temperature of these Pop III stars imply that their spectrum is much harder, and significantly fewer photons are emitted longward of the Lyman break (Tumlinson & Shull 2000).
Ly$`\alpha `$ photons could thus an important source of pair production opacity provided the escape fraction is small, i.e. $`(1f_{esc})1`$, and dust attenuation while the Ly$`\alpha `$ photons resonantly scatter in the IGM is unimportant. If Ly$`\alpha `$ photons are the dominant source of pair production opacity, this would be extremely interesting, as it would provide a indirect census of the ionizing photon emissivity, prior to attenuation within the host sources. This could prove to be a good measure of the comoving star formation rate. One way to check if this is the case would be to directly image sources in NGST in rest frame UV longward of Ly$`\alpha `$ (sufficiently far away from the Ly$`\alpha `$ damping wing), as well as in rest frame Balmer line emission (which should also be directly proportional to the production rate of ionizing photons, $`\dot{n}_{H\alpha }\dot{n}_{Ly\alpha }\dot{n}_{ion}(1f_{esc})`$). This will give a sense as to whether Ly$`\alpha `$ or UV continuum photons are a greater source of opacity, provided the trend observed in bright sources extrapolates down to lower luminosities.
### 3.2 Results
Fig 2 shows the predicted attenuation factor as a function of observed photon energy, for grazars at redshifts $`z_s=3,6,10`$, for both high ($`f_{star}=17\%`$) and low ($`f_{star}=1.7\%`$) star formation efficiencies. There are several important features to note. Firstly, the shape of the attenuation curve for Ly$`\alpha `$ is similar to that for UV continuum photons. It is not possible to immediately distinguish between scenarios where Ly$`\alpha `$ photons and UV continuum photons are the dominant source of opacity. Secondly, for lower star formation efficiencies, attenuation both sets in at only at higher energies and the attenuation curve is significantly shallower (i.e., it reaches full attenuation after a much longer energy interval). This is easy to understand. Gamma-ray photons of higher energy have both a longer path-length to travel before they redshift to $`E<E_{th}`$, and a higher number density $`n(ϵ_{th}<ϵ<ϵ_L)`$ of photons to pair produce against, since $`ϵ_{th}(m_ec^2)^2/E`$ is lower. Thus, if the overall number density of UV photons is lower, one must go to higher energies to achieve the same attenuation. In Figure 3, I show the relative contribution of different rest-frame target photon energies and redshift intervals for gamma-ray photons with $`\tau (E_o,z_s)=1`$, for a variety of source redshifts. The shape of these curves is largely dependent on the overall number density of target UV photons. For a lower SFR, the energy $`E_o`$ at which $`\tau (E_o,z_s)=1`$ increases, and the target photon interval and redshift interval contributing to the resultant opacity broadens. Note that most of the opacity comes from redshifts comparable to that of the source. Also, the opacity arises from a fairly narrow target photon energy interval $`313.6`$eV; this varies weakly with source redshift. Thus the results do not depend strongly on the assumed spectral slope of the UV sources. I also obtained by direct computation the contribution to the opacity due to ionizing photons shortward of the Lyman edge. Depending on one’s assumptions for $`f_{esc},f_{break},f_{opacity}`$, it is smaller by 2–4 orders of magnitude, and is completely negligible.
An important technique for isolating the contribution of high-redshift radiation fields would be to use measurements of differential absorption between blazars at different redshift. Consider two blazars with at redshifts $`z_1`$ and $`z_2`$, with $`z_2>z_1`$. At a given observed energy $`E_o`$, photons from each blazar will encounter identical optical depths for $`z<z_1`$. Thus, any additional opacity seem in the spectrum of the blazar at $`z_2`$ must arise from the redshift interval $`z_1<z<z_2`$ alone, $`\mathrm{\Delta }\tau (E_o)=_{z_1}^{z_2}𝑑z\frac{dl}{dz}_1^1d(cos\theta )(1\mathrm{cos}\theta )_{ϵ_{th}}^{\mathrm{}}𝑑ϵn(ϵ,z)\sigma (E,ϵ,\theta )`$. Note that this identification is independent of any uncertainties in the spectral slope or redshift evolution of the UV background. As argued above, radiation field intensity fluctuations are unimportant since the pair-production opacity arises on scales of order a Hubble length. Thus, in Figure (2), any difference between the $`z_s=3,6,10`$ curves is solely due to radiation fields between $`3<z<6`$ and $`6<z<10`$; if star formation ceases for $`z>3`$, the curves would lie on top of one another. If there is little difference between two absorption curves from sources at $`z_1`$ and $`z_2`$, one can get an upper bound on the UV emissivity in the range $`z_1<z<z_2`$; likewise, if two curves are so widely separated that meaningful measurements of $`\mathrm{\Delta }\tau (E_o)`$ cannot be obtained (in particular, if absorption saturates in one of the curves), one can get a lower bound on the UV emissivity in the range $`z_1<z<z_2`$. In Figure (4), I show how the attenuation at a fixed observed photon energy $`E_o`$ is expected to rise with redshift, due to the opacity provided by high redshift UV fields. For lower star formation efficiency, the attenuation rises more slowly; if there is no star formation at high redshift the curve would be flat. In Figure (5) I show how the energy $`E_o`$ for which $`\tau (E_o,z_s)=1`$ is expected to fall with increasing redshift; due to the increased opacity provided by high redshift photons, attenuation sets in at lower energies. Again, if there were no star formation at high redshift, this curve would show no evolution.
In fact, differential absorption measurements provide a powerful cross-check on the technique: if $`\tau (E_o,z_2)<\tau (E_o,z_1)`$, then the assumed unabsorbed blazar spectrum must be evolving with redshift, due to a change in internal absorption, or underlying spectral index. Note that if star formation rates are high, the attenuation occurs rapidly over a small energy interval, and uncertainties due to the extrapolation from the unabsorption portion of the spectrum is small; for lower star formation rates the attenuation occurs over a larger energy interval, and uncertainties due to extrapolation are greater.
## 4 Conclusions
In this paper, I have suggested that if blazars can be detected at high redshift, detection of gamma-ray absorption due to pair production against high-redshift UV photons will provide a valuable probe of high-redshift UV radiation fields. This is because the sharp Lyman edge in the intergalactic radiation fields implies that gamma-ray photons have only a limited redshift interval in which to pair-produce. As they redshift to lower energies, they require photons with $`ϵ>13.6`$eV to pair-produce, so the universe becomes optically thin. The shape of the attenuation curve is primarily sensitive to the overall number density of photons longward of the Lyman edge at high redshifts: the higher this number density, the lower the gamma-ray photon energy at which pair-production opacity sets in. This makes it a useful test of the overall level of star formation and ambient UV radiation fields present at high redshift. Ly$`\alpha `$ photons provide an important contribution to this pair production opacity, and indeed may be the dominant source of opacity if sources with a relatively lower fluxes longward of the Lyman edge (such as quasars or low metallicity stars) are abundant. Finally, measurements of differential absorption between blazars at the same observed energies will allow us to cleanly isolate the increase in opacity due to radiation fields at high redshift.
There are two large uncertainties. The first is whether GLAST will be able to see high-redshift blazars at all. However, in the unified model of AGN, the scarcity of ultra-luminous blazars is a geometrical effect (due to relativistic beaming) rather than a requirement of extremely large black hole masses. In fact, the luminosity boost provided by beaming reduces the black hole mass by several orders of magnitude below that demanded by the Eddington limit. So it is at least plausible that high redshift blazars will be detectable. The second uncertainty is whether absorption seen in a blazar will be internal, rather than due to pair production against photons in the IGM. However, at GeV energies we have some physical understanding of the observed EGRET spectra (e.g., Ghisellini et al 1998); opacity to gamma-ray photons due to internal radiation fields can be constrained by time variability arguments and other constraints. Furthermore, GLAST should assemble an extremely large catalog ($`>`$ few thousand) of low redshift blazars, whose spectra can be studied in detail (and the contribution to opacity due to low redshift star formation can be quantified by other means); provided blazar properties do not evolve too strongly with redshift, we should have a firm handle on the intrinsic unabsorbed blazar spectrum.
## 5 Acknowledgements
I thank David Spergel for his advice and encouragement, Youjun Lu for helpful conversations on blazar properties, Bruce Draine for helpful comments, and Bart Pindor for technical assistance. I thank the Institute of Theoretical Physics, Santa Barbara for its hospitality during the completion of this work. This work is supported by the NASA ATP grant NAG5-7154, and by the National Science Foundation, grant number PHY94-07194.
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# Glitches in Southern Pulsars
## 1 Introduction
Timing observations have shown that pulsars are remarkably stable clocks (e.g. Kaspi, Taylor & Ryba 1994). However, they are not always predictable. Two types of instability have been observed. In the first, the pulsar period, $`P`$, or frequency $`\nu =1/P`$, varies in an apparently random fashion, fluctuating on timescales of days, weeks, months and years \[Cordes & Downs 1985\]. The intensity of the fluctuations can be described by an ‘activity parameter’ based on the amplitude of the second frequency-derivative term in a Taylor-series fit to a set of pulse times of arrival (TOAs):
$$\mathrm{\Delta }(t)=\mathrm{log}\left(\frac{|\ddot{\nu }|}{6\nu }t^3\right)$$
(1)
where $`t`$ is the length of the data span in seconds \[Arzoumanian et al. 1994\]. Conventionally $`t=10^8`$ s is adopted, giving the parameter $`\mathrm{\Delta }_8`$. In most young pulsars, the magnitude of $`\ddot{\nu }`$ is far in excess of that expected from secular slowdown. Furthermore, timing noise is normally very ‘red’ \[Cordes & Downs 1985\] and so the second derivative term is a good indicator of the level of noise. Observations of large samples of pulsars \[Cordes & Downs 1985, D’Amico et al. 1998\] show that the activity parameter is positively correlated with the value of $`\dot{P}`$, the first time-derivative of the pulsar period.
In the second type of period instability, known as a ‘glitch’, the pulse frequency has a sudden increase which typically has a fractional amplitude $`\mathrm{\Delta }\nu _g/\nu `$ in the range $`10^8`$$`10^6`$. These glitches are unpredictable but typically occur at intervals of a few years in young pulsars. Coincident with the glitch, there is often an increase in the magnitude of the frequency derivative, typically by $`1`$ per cent, which sometimes decays with a timescale of weeks to years. The time-dependence of the pulsar frequency after a glitch is generally well described by the following function:
$$\nu (t)=\nu _0(t)+\mathrm{\Delta }\nu _g[1Q(1\mathrm{exp}(t/\tau _d)]+\mathrm{\Delta }\dot{\nu }_pt$$
(2)
where $`\nu _0(t)`$ is the value of $`\nu `$ extrapolated from before the glitch, $`\mathrm{\Delta }\nu _g=\mathrm{\Delta }\nu _d+\mathrm{\Delta }\nu _p`$ is the total frequency change at the time of the glitch ($`t=0`$), where $`\mathrm{\Delta }\nu _d`$ is the part of the change which decays exponentially and $`\mathrm{\Delta }\nu _p`$ is the permanent change in pulse frequency, $`Q=\mathrm{\Delta }\nu _d/\mathrm{\Delta }\nu _g`$, $`\tau _d`$ is the decay time constant and $`\mathrm{\Delta }\dot{\nu }_p`$ is the permanent change in $`\dot{\nu }`$ at the time of the glitch. The increment in frequency derivative is given by
$$\mathrm{\Delta }\dot{\nu }(t)=\mathrm{\Delta }\dot{\nu }_d\mathrm{exp}(t/\tau _d)+\mathrm{\Delta }\dot{\nu }_p=\frac{Q\mathrm{\Delta }\nu _g}{\tau _d}\mathrm{exp}(t/\tau _d)+\mathrm{\Delta }\dot{\nu }_p.$$
(3)
The middle term of equation 2 was originally derived on the basis of the two-component model for glitch recovery \[Baym et al. 1969\], where the initial frequency jump was due to a ‘starquake’ or sudden change in moment of inertia of the solid crust. In this model, the parameter $`Q=I_s/I_c`$ where $`I_s`$ and $`I_c`$ are the moments of inertia of the superfluid component and the crust, respectively.
In later versions of the theory, the frequency jump is due to a sudden unpinning of vortex lines in crustal neutron superfluid \[Anderson & Itoh 1975\]. Post-glitch relaxation may be due to ‘vortex creep’, that is, a slow drift of the vortices across the crustal lattice (e.g. Alpar, Cheng & Pines 1989) or drift of the crustal lattice itself with the vortices remaining pinned \[Ruderman 1991, Ruderman, Zhu & Chen 1998\].
Depending on the internal temperature of the star and other factors, the response to a glitch may be linear or non-linear \[Alpar, Cheng & Pines 1989\]. In the linear regime, the frequency jump at $`t=0`$ decays exponentially with time constant $`\tau _d`$. The fractional change in frequency derivative at the jump is
$$\mathrm{\Delta }\dot{\nu }_g/\dot{\nu }=(I_s/I)(\mathrm{\Delta }\omega _s/\tau _d),$$
(4)
where $`\mathrm{\Delta }\omega _s`$ is the change at the time of the glitch in the rotational lag of the superfluid, and $`\omega _s=2\pi (\nu _s\nu )`$, where $`\nu _s`$ is the rotation frequency of the neutron superfluid. For glitches which are large compared to the steady-state lag, $`\mathrm{\Delta }\omega _s2\pi \mathrm{\Delta }\nu `$ and the change in spin-down rate can be large compared to $`I_s/I`$.
In the non-linear regime, there is essentially no relaxation ($`Q0`$) and
$$\mathrm{\Delta }\dot{\nu }_g/\dot{\nu }=I_s/I,$$
(5)
giving a permanent change $`\mathrm{\Delta }\dot{\nu }_p`$ in $`\dot{\nu }`$ (Eqn 3). Analyses of data for the Vela pulsar and other pulsars \[Alpar et al. 1988, Alpar et al. 1993, Ruderman, Zhu & Chen 1998\] are generally consistent with a rather small superfluid fraction, $`I_s/I10^2`$, participating in the glitch activity. In the Ruderman et al. (1998) model, permanent changes in $`\dot{\nu }`$ may result from a change in the magnetic field configuration associated with the crustal cracking at the time of a glitch.
Until the last few years, the number of known glitches and glitching pulsars was modest. However, with relatively high-frequency searches at low Galactic latitudes \[Clifton et al. 1992, Johnston et al. 1992, Kaspi et al. 1992\] discovering a much larger sample of young pulsars, the number has increased significantly, with 21 pulsars having 46 glitches being listed in the recent review by Lyne (1996). In this paper we describe timing observations of a sample of 40 mostly young pulsars using the Parkes radio telescope which show a total of 30 glitches in 11 pulsars.
## 2 Observations and Analysis
Observations were made using the Parkes 64-m radio telescope between 1990 January and 1998 December at frequencies around 430, 660, 1400 and 1650 MHz. The frequencies below 1 GHz were used only occasionally to improve the determination of dispersion measures (DMs). At all frequencies, cryogenic dual-channel receivers were used. Two separate back-end systems were used to provide the frequency resolution necessary for dedispersing. The early observations used filterbanks and a one-bit digitizer system constructed at Jodrell Bank; a $`2\times 256\times 0.125`$ MHz system was used at 430 MHz, a $`2\times 128\times 0.25`$ MHz system was used at 660 MHz and at the higher frequencies a $`2\times 64\times 5`$ MHz system was used. Further details of the filterbank data acquisition system may be found in Manchester et al. (1996). From 1994 July, data were acquired using a correlator system constructed at Caltech \[Navarro 1994\]. This system used two-bit digitization and an autocorrelator with $`2\times 256`$ lags over a maximum bandwidth of 128 MHz. From mid-1995, the lags were split between two frequency bands which gave simultaneous observations at radio frequencies of 1400 and 1650 MHz \[Sandhu et al. 1997\].
For the filterbank systems, data were folded synchronously with the pulsar period using off-line programs to give mean total intensity pulse profiles. In the correlator system this folding was performed in a hardware integrator having 1024 bins across the pulsar period. For both systems, integration start times were established to better than $`1\mu `$s using time signals from the Observatory clock system. This was related to UTC(NIST) using a radio link to the NASA DSN station at Tidbinbilla and clock offsets kindly provided by the Jet Propulsion Laboratory, Pasadena. In subsequent analysis, the data were summed to form 8 or 16 frequency subbands and time subintegrations of 60 or 90 s duration and stored on disk. These files were then summed in frequency and time using the best available period and DM information to form a single profile for each observation. This was then cross-correlated with a standard profile for the pulsar to give a pulse time of arrival (TOA).
A total of 40 pulsars were monitored during the program. The J2000 and B1950 names of these pulsars, their periods and characteristic ages, and the dates spanned by the observations are listed in Table 1. Depending on the strength of the pulsar, observation times for each TOA were between 2 and 12 min. Observations were obtained at intervals of between 2 and 6 weeks for most of the pulsars in Table 1, with the known glitching pulsars being observed more frequently. The TOAs resulting from this program were analysed using version 11.3 of TEMPO<sup>1</sup><sup>1</sup>1See http://pulsar.princeton.edu/tempo/ for a description of TEMPO. which includes provision for fitting the glitch parameters given in Equation 2. The DE200 solar-system ephemeris \[Standish 1990\] was used to convert TOAs to the solar-system barycentre. The sixth and seventh columns of Table 1 give the rms value of the TOA uncertainty (dependent on the pulse period, pulse shape and signal/noise ratio) and the rms residual after fitting for just pulse frequency and its first derivative. For glitching pulsars, the largest glitch-free interval was used for this fit. These residuals are dominated by the effects of timing noise, and so they are an indication of its amplitude. An exception to this is PSR J1513-5908, which has significant frequency second and third derivatives due to secular slowdown \[Kaspi et al. 1994\]; in this case, the fit giving $`\sigma _T`$ included the first three frequency derivatives.
For several pulsars, the position and DM were measured with comparable or better precision than was previously available. Positions were obtained using the so-called ‘pre-whitening’ method (Kaspi et al. 1994). In this method, the position is fitted along with sufficiently many frequency-derivative terms and in some cases, glitch terms, to ‘absorb’ the timing noise and give an approximately Gaussian distribution of timing residuals. DMs were measured from short sections of data where there were multi-frequency observations. This ensures that contamination by long-term timing noise is not a problem and that error estimates are realistic.
For glitches where the amplitude of the glitch and the interval between the last pre-glitch observation and the first post-glitch observation are not too large, the epoch of the glitch can be determined by requiring that the pulse phase be continuous over the glitch. In this situation, TEMPO gives an estimate of the glitch parameters, including the epoch, and their uncertainty. Where post-glitch decay parameters are estimated, these refer to the long-term decay, typically with timescales of hundreds of days. The resolution of our observations is generally not sufficient to detect the more rapid post-glitch recoveries observed in some pulsars, for example, the Vela pulsar \[Flanagan 1990, McCulloch et al. 1990\].
For larger glitches where there may be one or more turns of phase in the residuals between the bounding observations, the glitch epoch is uncertain. Other glitch parameters are affected by this uncertainty. To estimate uncertainties in this case, the following procedure was adopted. The glitch epoch was taken to be halfway between the bounding observations, with an uncertainty of half their separation. The increments $`\mathrm{\Delta }\nu _g`$ and $`\mathrm{\Delta }\dot{\nu }_g`$ were then computed for an assumed glitch epoch $`t_g`$ by extrapolating the pre- and post-glitch fits to $`t_g`$ and taking differences. Uncertainties were similarly extrapolated and quadrature sums taken for the uncertainties in the increments. This was done separately for $`t_g`$ at the two bounding epochs. The final error estimates were then the quadrature sum of the difference between the increments at either end of the data gap and the larger of the uncertainties in the increments.
## 3 Results
Table 2 lists parameters for the 11 pulsars for which glitches were detected. The pulsar J2000 positions and DMs used in or derived from the analysis are given. Uncertainties in the last digit quoted are given in parentheses. These and other quoted uncertainties from TEMPO fits are twice the formal rms error. Position fits were generally to the largest available data span not too strongly affected by post-glitch post-glitch recovery. As described above, the data spans used to determine the positions and their errors were ‘pre-whitened’ to eliminate the effects of period noise from the positions and their errors. The fifth and sixth columns give the data span used when fitting the position and the final rms timing residual. References for the position and DM are given in the final column.
A total of 30 glitches were observed in these 11 pulsars. Independent fits to the pre-glitch, inter-glitch and post-glitch timing data are given in Table 3. Except for short sections of data, the fits include a $`\ddot{\nu }`$ term. In every case, this term is dominated by recovery from a previous glitch; the $`\ddot{\nu }`$ from the long-term secular slow down is negligible by comparison. Higher-order frequency derivative terms were not fitted, so the rms residuals given in right-most column reflect the presence of random period irregularities, especially for longer data spans.
Glitch parameters are listed in Table 4. If marked by $``$, the glitch epoch is determined by requiring phase continuity across the glitch. The next two columns give the fractional steps in $`\nu `$ and $`\dot{\nu }`$ at the glitch determined by extrapolating the pre- and post-glitch solutions (Table 3) to the glitch epoch, with uncertainties determined as described in Section 2. Parameters in the remaining columns were determined using TEMPO to fit phases across the glitch.
Ten of these glitches have been previously reported: single glitches in PSR J0835$``$4510 \[Flanagan 1996\], PSR J1709$``$4428 \[Johnston et al. 1995, Shemar & Lyne 1996\], PSR J1731$``$4744 \[D’Alessandro & McCulloch 1997\] and PSR J1801$``$2451 \[Lyne et al. 1996a\], and three glitches in PSR J1341$``$6220 \[Kaspi et al. 1992, Shemar & Lyne 1996\] and PSR J1801$``$2306 \[Kaspi et al. 1993, Shemar & Lyne 1996\]. We reanalyse these glitches for completeness and consistency with the results for the other pulsars.
In the following sections we discuss each pulsar in turn.
### 3.1 PSR J0835$``$4510, PSR B0833$``$45, the Vela pulsar
The Vela pulsar is well known to suffer many giant glitches \[Cordes, Downs & Krause-Polstorff 1988, McCulloch et al. 1990, Lyne et al. 1996b\] and is being regularly monitored at a number of observatories. At Parkes, we are unable to make such regular and frequent observations. In this paper, we describe observations of the latest glitch which has only been briefly reported \[Flanagan 1996\].
Table 3 gives the results of fitting for $`\nu `$ and its first two time derivatives, both before and after the glitch. The pre-glitch fit was obtained from an approximately 1-year data span and shows a significant $`\ddot{\nu }`$ resulting from previous glitches. The post-glitch fit given in Table 3 is for an 82-day span commencing 18 days after the glitch. A timing model with polynomial terms up to $`\ddot{\nu }`$ is not a good fit to longer data spans.
Fig. 1 shows the time dependence of the frequency residual $`\mathrm{\Delta }\nu `$ and of $`\dot{\nu }`$ around the time of the glitch. The $`\mathrm{\Delta }\nu `$ values plotted are differences between the values of $`\nu `$ obtained from independent fits to short sections of data, typically of span 20 – 30 days, and those determined from the predictions of the pre-glitch model given in Table 3. This plot shows an approximately exponential recovery in $`\dot{\nu }`$ after the glitch, followed by an approximately linear increase in $`\dot{\nu }`$. This behaviour is similar to that seen in previous Vela glitches \[Lyne et al. 1996b\].
Flanagan (1997) gives $`\mathrm{\Delta }\nu _g/\nu =2.15(2)\times 10^6`$ and an epoch for the glitch of 1996, October 13.394 UT, corresponding to MJD 50369.394. Fractional changes in $`\nu `$ and $`\dot{\nu }`$ at the time of the glitch, obtained by extrapolating the pre-glitch and post-glitch fits (Table 3) to the glitch epoch are given in the fifth and sixth columns of Table 4. Results of fitting the exponential model (Equation 2) to the interval from 350 days before the glitch to 200 days after are also given in Table 4. In this fit, $`\nu `$, $`\dot{\nu }`$ and $`\ddot{\nu }`$ were held at their pre-glitch values (Table 3). The small rms residual shows that the exponential model with about 40 per cent of the initial glitch in frequency decaying on a timescale of about 900 d is a very good representation of the post-glitch relaxation, at least over the data span fitted. Fitting of $`\nu `$, $`\dot{\nu }`$ and $`\ddot{\nu }`$ to the post-glitch data instead of the three parameters of the exponential fit gives a substantially worse fit with rms residual 360 $`\mu `$s and a systematic quartic term in the residuals. For the exponential fit, most of the rms residual comes from the first few post-glitch points which are positive, indicating an additional short-term recovery. However, it was not possible to fit for the parameters of this. The fitted glitch epoch, given in the third column of Table 4, is slightly earlier than that given by Flanagan (1997); this is consistent with the presence of more rapid post-glitch relaxation than modelled.
### 3.2 PSR J1048$``$5832, PSR B1046$``$58
This young pulsar was discovered in a high-frequency survey of the southern Galactic plane \[Johnston et al. 1992\]. Evidence for gamma-ray pulsations has been found by Kaspi et al. (2000). The pulsar suffered two large glitches within the observed data span, the first of size $`\mathrm{\Delta }\nu _g/\nu 3\times 10^6`$ in 1993 February (MJD $`49034`$) and the second of size $`0.7\times 10^6`$ in 1997 December (MJD $`50788`$). Observed frequency residuals are shown in Fig. 2. The expanded plot of $`\mathrm{\Delta }\nu `$ (Panel b) shows a third and much smaller glitch about 100 d before the large first glitch. Fits to the inter-glitch intervals (Table 3) give evidence for large fluctuations in the period. Fig. 3 gives timing residuals for the data span following the largest glitch, showing quasi-random fluctuations with a timescale of a few hundred days.
Table 2 gives an improved position for PSR J1048$``$5832 derived from the data following the largest glitch. Most of the systematic oscillation was removed by fitting up to the twelfth pulse frequency derivative at the same time as the position was determined. A fit to the whole data set, including parameters for the glitches, gave the same position within the combined errors. This position differs by 8″ from that given by Johnston et al. (1995); this difference can be attributed to the different sets and the presence of the strong period irregularities. Recently, Stappers et al. (1999) have used the Australia Telescope Compact Array to determine an interferometric position for the pulsar: R.A. (J2000) 10<sup>h</sup> 48<sup>m</sup> $`12\stackrel{s}{.}604\pm 0\stackrel{s}{.}008`$, Dec. (J2000) $`58\mathrm{°}32\mathrm{}03\stackrel{}{.}75\pm 0\stackrel{}{.}05`$. This position agrees with the timing position in R.A., but differs by $`2\stackrel{}{.}45\pm 0\stackrel{}{.}80`$ in declination, probably as a result of unmodelled period irregularities affecting the timing position.
Extrapolation of fits to the data sets on either side of the glitches gives the estimates of glitch parameters listed in columns 5 and 6 of Table 4. Glitch parameters from TOA fits are given in the remaining columns. Because of the contaminating effect of the period fluctuations, the parameters of the large glitches were determined by fitting the data span 100 – 150 days before and after each glitch. The effects of the earlier small glitch were determined separately by fitting the interval between 80 days prior to it and up to the large glitch, and were subtracted from the fit to the large glitch. Parameters from these fits are given in Table 4. Because of the large systematic period variations, it is not possible to reliably measure the post-glitch decay times. Fig. 3 gives some evidence for a relaxation with a timescale of $`100`$ d following the 1993 glitch and somewhat longer for the 1997 glitch. Setting the decay time to 100 d and 400 d respectively for these two glitches gave the parameters listed in Table 4. The derived $`Q`$ for the 1997 glitch appears significantly larger than that for the 1993 glitch, but this result is uncertain given the uncertainty in the relaxation time and the short data span following the glitch.
### 3.3 PSR J1105$``$6107
PSR J1105$``$6107 is a young pulsar with the relatively short period of 63 ms discovered in 1993 using the Parkes radio telescope \[Kaspi et al. 1997\]. It is located close to but outside the boundaries of the supernova remnant G290.1$``$0.8 and within the error box for the EGRET gamma-ray source 2EG J1103$``$6106 \[Kaspi et al. 2000\]. Its association with these objects remains plausible but unproven. Kaspi et al. (1997) reported on timing observations from 1993 July to 1996 July. Here we extend this interval to the end of 1998, and show that there were two glitches, a large one around the end of November, 1996 (MJD $`50417`$) , and a small one about 200 days later.
Fig. 4 gives the variations in pulse frequency and frequency derivative over the observed data span, clearly showing the larger glitch which is of fractional size $`\mathrm{\Delta }\nu _g/\nu 0.28\times 10^6`$. There was a gap between bracketing observations of about 30 days, so the glitch epoch is not well determined. The expanded plot (b) shows that, despite the substantial period irregularities, there is good evidence for a small glitch about 200 days after the large glitch. Fig. 5 shows phase residuals around the time of this small glitch. Although this glitch is by far the smallest discussed in this paper, there is little doubt about its reality. Another small glitch may have occurred near the end of the data set, but there are insufficient post-glitch observations to confirm this. The position given in Table 2 was determined from a fit to the data prior to the large glitch, with seven frequency derivatives to absorb the period irregularities.
Table 3 gives fits to the three data spans delimited by the data set and the two glitches; the position was held at the value quoted in Table 2. The second derivative terms in the first and third fits are dominated by random period irregularities. Glitch parameters obtained by extrapolating these fits and from phase fits across the glitches are given in Table 4. Because of the period noise, the phase fits were restricted to intervals of 150 days before the first glitch, the inter-glitch interval, and 150 days after the second glitch. The fractional increase $`\dot{\nu }`$ after the first glitch is small, implying a similarly small value of the ratio $`Q/\tau _d`$ (Equation 3). Because of the small data span, $`\tau _d`$ is not well determined, but taking 100 d gives $`Q0.035`$; i.e., only a few percent of the glitch is likely to decay. This is confirmed by a fit to the whole data span including both glitches.
### 3.4 PSR J1123$``$6259
PSR J1123$``$6259, discovered in the Parkes southern pulsar survey \[Manchester et al. 1996\], has a period of 271 ms and a modest period derivative, giving a characteristic age of $`8\times 10^5`$ yr. Despite this relatively large characteristic age, a glitch was detected between 1994 December 17 – 22 (MJD $`49706`$). The pulsar position was determined from the post-glitch data (Table 2); no pre-whitening was necessary as random period irregularities are small in this pulsar. Fits to the pre- and post-glitch data with this position are given in Table 3. Extrapolation of these fits to a central epoch (Table 4) shows that the glitch was of magnitude $`\mathrm{\Delta }\nu /\nu 7.5\times 10^7`$. <sup>4</sup><sup>4</sup>4A glitch of this magnitude was reported by D’Amico et al. (1998) to have occurred at MJD $`48650\pm 20`$ (1992 December). This epoch is in error. It was in fact the December 1994 glitch discussed in this paper.
As Table 4 and Fig. 6 show, there was a small but significant increase in $`|\dot{\nu }|`$ at the time of the glitch. A fit to the entire data span including glitch parameters but with the position and pre-glitch frequency parameters held to the values given in Table 2 and Table 3 is an excellent fit to the data and gives an unambiguous value for the epoch of the glitch with an uncertainty of about 15 min on the assumption that there was phase continuity across the glitch. The derived value of $`Q`$ is very small and the decay time is long (Table 4).
### 3.5 PSR J1341$``$6220
PSR J1341$``$6220 is a young pulsar (characteristic age $`12,000`$ y) discovered by Manchester et al. (1985) and associated with the supernova remnant G308$``$0.1 by Kaspi et al. (1992). Kaspi et al. reported two large glitches ($`\mathrm{\Delta }\nu _g/\nu 1.5\times 10^6`$ and $`1.0\times 10^6`$) and one smaller one ($`2.3\times 10^8`$) within the period 1990 January to 1992 May, making this one of the most actively glitching pulsars known. In this paper, we reanalyse the data from 1990 to 1992 and extend the data set 1998 March. Unfortunately, observations after this time were too sparse to permit unambiguous pulse counting.
The position given in Table 2 was determined in a simultaneous fit across the whole data set, including all glitches (see below). It has smaller estimated uncertainties than the position given by Kaspi et al. (1992), and lies 4$`\stackrel{}{.}`$3 north of it. Given the prominence of period irregularities in this object, this difference, while twice the combined uncertainties, is of marginal significance. Determination of the DM is complicated by the strong scattering tail shown by this pulsar, even at frequencies around 1400 MHz. The value given in Table 2 was determined using TEMPO with data recorded between 1998 March and December at frequencies close to 1400 MHz and 1700 MHz. The relative phase of the standard profiles at the two frequencies was adjusted to allow for the effects of scattering. The derived DM is much more accurate than the value given by Kaspi et al. (1992) and just outside the error range of this value.
Fig. 7 shows the very complicated period history of this pulsar. A total of 12 glitches were observed in the 8.2-year interval, making this the most frequently glitching pulsar known, with a mean glitch interval of 250 days. The next most frequent is PSR B1737$``$30, with a mean interval of 345 days \[Shemar & Lyne 1996\]. Four of the glitches were relatively large, with $`\mathrm{\Delta }\nu _g/\nu >0.7\times 10^6`$ and for two the fractional glitch size exceeds $`10^6`$. Most of the other glitches are quite small, with all except one having $`\mathrm{\Delta }\nu _g/\nu <0.04\times 10^6`$. They are, however, unambiguous on phase-time plots.
For most of the glitches, there is little evidence for any post-glitch relaxation. This is partly because of the relatively poorly sampled data, especially in the earlier years, but also because of the short time between successive glitches. For the third and ninth glitches, there is weak evidence in the $`\dot{\nu }`$ plot for some relaxation.
Table 3 gives fits to the inter-glitch intervals. Because of the limited data span prior to the first observed glitch, the value of $`\dot{\nu }`$ was held fixed at a value close to the long-term mean for this fit. For about half of the intervals, there was a clear $`\ddot{\nu }`$ term in the residuals and this term was fitted.
Changes in $`\nu `$ and $`\dot{\nu }`$ at the time of each glitch, obtained by extrapolation of these fits to the glitch epoch, are given in Table 4. This table also gives glitch parameters from a single fit to the whole data set, solving simultaneously for the pulsar position, pulsar frequency ($`\nu `$), mean frequency derivative ($`\dot{\nu }`$), and the parameters for all 12 glitches. As mentioned above, for most of the glitches, there was no significant post-glitch relaxation. For the third and ninth glitches, a post-glitch exponential decay was fitted; in both cases the decay time constant was not fitted for, but was determined by trial to minimise the final residual. In both cases the derived $`Q`$ values are relatively small.
### 3.6 PSR J1614$``$5047, PSR B1610$``$50
This pulsar has a period of 231 ms and a very large period derivative ($`495\times 10^{15}`$), implying a small characteristic age of $`7400`$ yr. In terms of its period irregularities, it is one of the noisiest pulsars known; this makes it difficult at times to keep track of the pulse phase. Since its discovery in late 1989 \[Johnston et al. 1992\], there has been no clear evidence for a glitch although, particularly where there is a significant gap in the timing data, it is often difficult to distinguish between a (small) glitch and more continuous period irregularities. However, in 1995 June (MJD $`49802`$), there was a massive glitch with $`\mathrm{\Delta }\nu _g/\nu 6.5\times 10^6`$, the largest ever observed in any pulsar (cf. Shemar & Lyne 1996). Fig. 8 shows this glitch and also illustrates the more continuous period irregularities.
These irregularities make it difficult to determine the pulsar position from timing data. Previously published positions have differed by much more than the quoted uncertainties \[Taylor, Manchester & Lyne 1993, Johnston et al. 1995\]. Taking the data span from MJD 50269 to to 50778 (which is free of major irregularities – see below) and fitting for position and two frequency derivatives gives the position quoted in Table 2. An independent fit to the MJD range 48732 – 49093 gave a position with uncertainty of about 2″, consistent with the position from the longer data span. Subsequent fits were made keeping the position fixed at the Table 2 value. Stappers et al. (1999) have recently determined an interferometric position for this pulsar: R.A. (J2000) 16<sup>h</sup> 14<sup>m</sup> $`11\stackrel{s}{.}55\pm 0\stackrel{s}{.}01`$, Dec. (J2000) $`50\mathrm{°}48\mathrm{}01\stackrel{}{.}9\pm 0\stackrel{}{.}1`$. As for PSR J1048$``$5832, the differences most probably result from period irregularities affecting the timing position.
The slope of the post-glitch data in the top two plots indicates a large change in frequency derivative at the time of the glitch. Because of the strong irregularities in the period, the pre- and post-glitch fits given in Table 3 are restricted to intervals of less than 300 days. Extrapolation of those fits gives the glitch parameters in columns 5 and 6 of Table 4. Glitch parameters were also obtained from a single timing solution over the same interval as the pre- and post-glitch fits (MJD 49559 – 50117) and are given in columns 7 – 10 of Table 4. This fit did not converge well when fitting for the glitch decay time, so this was held fixed at 2000 d, a value representative of those obtained from the fitting. From the post-fit residuals, it is clear that there is also a more rapid decay of a portion of the glitch with a timescale of 10 – 20 days. There is good agreement between the two methods of deriving the glitch parameters. From Equation 3, the $`Q`$ and $`\tau _d`$ values given in Table 4 imply a fractional change in $`\dot{\nu }`$ of 0.0096, compared to 0.0097 from the extrapolation. Unfortunately, because of the large gap between the observations bracketing the glitch and the large size of the glitch, it is not possible to accurately determine the glitch epoch from phase continuity.
This is the largest glitch yet observed. The previous largest was for PSR B0355+54 for which $`\mathrm{\Delta }\nu /\nu `$ was $`4.37\times 10^6`$ \[Lyne 1987\]. There was very little decay of the frequency step in PSR B0355+54, whereas the $`Q`$ value for the PSR J1614$``$5047 glitch indicates that more than half of it will decay away over a few years. In this respect, the PSR J1614$``$5047 glitch is more similar to the giant glitches in the Vela pulsar, which decay by a similar amount over similiar timescales.
Phase fitting of post-glitch data revealed two further abrupt changes in pulse frequency. Phase plots for the intervals around these events are shown in Fig. 9. These phase plots have the character of glitch events, with a persistent fractional frequency change $`\mathrm{\Delta }\nu _g/\nu 0.03\times 10^6`$ at about MJD 50170 and 50780, respectively. The first event may indeed be a small glitch – it is not possible to tell because of the large data gap – but the second is not. This frequency change is resolved, taking place over about 10 days. In both cases, there appears to be a significant decrease in $`|\dot{\nu }|`$ associated with the event. These rapid frequency changes could be classed as just part of the frequency irregularities which are prominent in this pulsar, but in general these irregularities have a much longer timescale and are smaller in amplitude. This is demonstrated by the small residuals for the fits to data before these events; for the second event, the fit includes only two frequency derivatives and has an rms residual of only 1.2 ms (Table 3). The MJD 50780 event stands out from the general irregularities with a much larger rate of frequency change over the 10 days. Phase fits to the post-glitch data given in Table 3 have been split into three sections to avoid contamination by these sudden frequency changes. The last ends at MJD 50926 (1998 April) since data are sparse after that point and it is not possible to fit across the gaps without ambiguity. The frequency second derivative terms are very significant in these fits, but differ greatly in both sign and magnitude for the different segments. They are clearly related to the period-noise processes occurring in this star and not to the secular slowdown.
### 3.7 PSR J1709$``$4428, PSR B1706$``$44
This pulsar, discovered by Johnston et al. (1992), is of interest for several reasons. It is young ($`\tau _c17`$ kyr) and has the third highest known value of $`\dot{E}/d^2`$, where $`\dot{E}`$ is the spin-down luminosity and $`d`$ is the pulsar distance, after the Crab and Vela pulsars. It has been detected at gamma-ray wavelengths \[Thompson et al. 1992, Thompson et al. 1996\], X-ray wavelengths \[Becker, Brazier & Trümper 1995\] and possibly at TeV energies \[Kifune et al. 1995\]. A supernova remnant association was proposed by McAdam, Osborne & Parkinson (1993), but deemed unlikely by Nicastro, Johnston & Koribalski (1996).
Johnston et al. (1995) presented results from two years of Parkes timing observations showing that the pulsar suffered a giant glitch of magnitude $`\mathrm{\Delta }\nu _g/\nu 2\times 10^6`$ near the end of May, 1992 ($``$ MJD 48778). In this paper we analyse data from 1990 January to 1998 December. To determine the pulsar position, we have taken data from 150 days after the glitch (MJD 48928 – 51155) and fitted nine frequency derivatives to absorb the period irregularities. This fit had an rms residual of only 211 $`\mu `$s and gave the position listed in Table 2. The derived position is about $`12^{\prime \prime }`$ southeast of the position given by Johnston et al. (1995), but agrees better with an interferometric position given by Frail & Scharringhausen (1997): R.A. (J2000) 17<sup>h</sup> 09<sup>m</sup> $`42\stackrel{s}{.}75\pm 0\stackrel{s}{.}01`$, Dec. (J2000) $`44\mathrm{°}29\mathrm{}06\stackrel{}{.}6\pm 0\stackrel{}{.}3`$. The timing fits described below were obtained with the position held at the Table 2 value.
Fig. 10 shows the large jump with a classic exponential recovery on a timescale of $`100`$ days coupled with a longer term relaxation. Fits of a cubic phase polynomial to the pre-glitch data, to the 150 days after the glitch and to data from 150 days post-glitch are given in Table 3. The post-glitch fit has large systematic residuals, dominated by the post-glitch recovery. Glitch parameters obtained by extrapolating the first and second fits to the glitch epoch are given in Table 4 along with those from a phase fit to the whole data set. This latter fit is strongly affected by the strong period irregularities (Fig. 11) and so the derived parameters are only approximate. There is clear evidence in the post-fit residuals for a more rapid exponential decay with time constant of the order of 100 days. However, it was not possible to fit for the parameters of this owing to the contaminating effect of the period irregularities.
### 3.8 PSR J1731$``$4744, PSR B1727$``$47
PSR J1731$``$4744 is a long-period pulsar (0.830 s) which has a large period derivative giving a relatively low characteristic age, $`80,000`$ years. D’Alessandro & McCulloch (1997) observed this pulsar at Mt Pleasant Observatory in Tasmania from 1987 to 1994 and detected a relatively large glitch at MJD 49387.68 (1994 February 2) but found no evidence for any recovery in 10 months of observations after the glitch. The Parkes observations extend to the end of 1998 (Fig. 12) and show a clear recovery from this glitch which is roughly exponential. They also reveal a second smaller glitch in 1997 September which has a similar recovery.
Table 2 gives an improved position and DM for this pulsar. The position was determined from data between the two glitches by fitting for the position and five frequency-derivative terms to absorb the period irregularities. The DM was determined by fitting to data in the MJD range 49400 to 49900 where there were observations around 430 MHz as well as around 1400 MHz. This value is significantly different from the best previously published value ($`121.9\pm 0.1`$ cm<sup>-3</sup> pc; McCulloch et al. 1973), most probably reflecting a changing line-of-sight path through the interstellar medium. The average rate of DM change over the 24 years is $`0.6`$ cm<sup>-3</sup> pc y<sup>-1</sup>. DM changes have been observed in other pulsars, notably the Crab \[Isaacman & Rankin 1977\] and Vela \[Hamilton, Hall & Costa 1985\] pulsars. The average rate of change of DM for PSR J1731$``$4744 is comparable to that observed for the Vela pulsar and larger than that for the Crab pulsar. The large changes observed for the these two pulsars are attributed to ionised gas within the associated supernova remnant, but PSR J1731$``$4744 has no associated supernova remnant.
Fits to the inter-glitch intervals are given in Table 3. These fits represent the data reasonably well, except that there is a significant quartic term in the residuals for the middle interval. This is mostly due to the period irregularities which are evident in Fig. 12.
Table 4 gives parameters for the two glitches. The fitted parameters were determined by a simultaneous fit to both glitches. As mentioned above, a recovery after both glitches is clearly seen in Fig. 12, and these are well fitted by the exponential model. For the first glitch, about 8 per cent of the glitch is recovered with a time constant of about 260 days. This implies a fractional increment in $`\dot{\nu }`$ at the time of the glitch of $`2.25\times 10^3`$. The glitch parameters are in reasonable agreement with those quoted by D’Alessandro & McCulloch (1997); the differences of a few times the combined uncertainties could result from the different models used for the extrapolation to the time of the glitch. The data span following the second glitch is too short to permit solving for the time constant, so 250 days was assumed, giving a $`Q`$ value of about 25 per cent.
### 3.9 PSR J1801$``$2304, PSR B1758$``$23
This pulsar was detected in a search for short-period pulsars associated with supernova remnants \[Manchester, D’Amico & Tuohy 1985\]. Despite its relatively long period of 415 ms, this pulsar has a short characteristic age (58 kyr). It lies close to the supernova remnant W28, but its association with the remnant remains controversial \[Kaspi et al. 1993, Frail, Kulkarni & Vasisht 1993\]. It has the highest known dispersion measure (1074 cm<sup>-3</sup> pc) and even at 1.4 GHz the profile is very scattered, reducing the precision of pulse timing observations at this and lower frequencies. Furthermore, the pulsar lies very close to the ecliptic plane and suffers frequent glitches, so that the pulsar position is not well determined by pulse timing; the position quoted in Table 2 and used for the timing analyses is from the VLA observations of Frail et al. (1993).
Four glitches in this pulsar occurring between 1986 and the end of 1994 were reported by Kaspi et al. (1993) and Shemar & Lyne (1996). In this paper, we present data on the last two of these glitches and two further glitches, both of relatively small amplitude, occurring in 1995 and 1996 respectively. Fig. 13 shows the variations in pulse frequency and frequency derivative around these four glitches relative to the pre-glitch solution given Shemar & Lyne (1996) for the MJD 48454 glitch. This pre-glitch solution was adopted because that given in Table 3 is based on a data span of only 120 days and has a significantly different value of $`\dot{\nu }`$ from the other fits, probably as a result of intrinsic period irregularities. Fits to the inter-glitch intervals given in Table 3 have relatively large residuals because of these irregularities. Only the fit to the data following the fourth glitch required a $`\ddot{\nu }`$ term, indicating a significant relaxation following the glitch.
Parameters for the four glitches are given in Table 4, both from extrapolation of the polynomial fits given in Table 3 and from a simultaneous fit to the TOA data across all four glitches. Values of $`\mathrm{\Delta }\nu /\nu `$ and $`\mathrm{\Delta }\dot{\nu }/\dot{\nu }`$ for the first two jumps have smaller uncertainties but are consistent with the values given by Shemar & Lyne (1992). For the first three jumps the glitch epoch is determined from the requirement of phase continuity over the glitch; for the fourth glitch the data gap is too large for an unambiguous determination of the glitch epoch. Consistent with the observation of a significant $`\ddot{\nu }`$ term, post-fit residuals were significantly reduced by allowing an exponential relaxation following the fourth jump. However, the presence of period irregularities made fitting for the decay time impossible; the value of 100 days was found by trial to give the best representation of the data.
### 3.10 PSR J1801$``$2451, PSR B1757$``$24
PSR J1801$``$2451 is a young pulsar ($`\tau _c15,000`$ y) which is located at the “beak of the Duck”, in the small nebula associated with the much larger supernova remnant G5.4$``$1.2 \[Manchester et al. 1991, Frail & Kulkarni 1991\]. A giant glitch in the pulse period was observed around MJD 49476 by Lyne et al. (1996). About a year of post-glitch data to MJD 49950 (1995 August) suggested a relaxation of a small fraction of the glitch ($`Q0.005`$) with a characteristic timescale of about 40 days.
In Fig. 14 we show observed pulse frequency variations based on Parkes data from 1992 October to 1998 December. These show the glitch observed by Lyne et al. (1996) and also a second smaller, but still large, glitch around 1997 July. Polynomial fits to the inter-glitch intervals given in Table 3 have large residuals owing to the presence of period irregularities. Extrapolation of the first two of these fits to the glitch epoch determined by Lyne et al. (1996) gives results consistent with those obtained by these authors, with $`\mathrm{\Delta }\nu _g/\nu 2.0\times 10^6`$. The second glitch, separated from the first by about 3.2 years, has a fractional amplitude of $`1.2\times 10^6`$, making it also a member of the giant glitch class.
The longer timespan of the present observations following the first glitch show that the post-glitch relaxation is best described by an exponential recovery with a timescale of several hundred days. It is likely that period irregularities were primarily responsible for the apparent quicker relaxation found by Lyne et al. (1996). There is also evidence for a similar long-timescale relaxation following the second glitch. A TEMPO fit across both glitches yielded the parameters given in Table 4, showing that, for both glitches, about 20 per cent of the glitch decayed. For the first glitch, the fitted timescale was 800 days. No fit of the timescale to the second glitch was possible because of the shorter data span; an assumed decay time of 600 days gave a minimum in the post-fit residuals.
### 3.11 PSR J1803$``$2137, PSR B1800$``$21
PSR J1803$``$2137 is a young pulsar ($`\tau _c16,000`$ y) discovered by Clifton & Lyne (1986). Lyne & Shemar (1996) observed a large glitch ($`\mathrm{\Delta }\nu _g/\nu 4.1\times 10^6`$) in 1990 December which showed an exponential recovery with timescale of $`150`$ d, together with an apparently linear decay in $`\dot{\nu }`$ indicating decay from an earlier unobserved glitch.
Our observations of this pulsar are from 1997 August to 1998 December. Fig. 15 shows that another large glitch of magnitude $`\mathrm{\Delta }\nu _g/\nu 3.2\times 10^6`$ occurred in 1997 November. Again, there the signature of an exponential decay following the glitch is seen. There is, however, no evidence of decay in $`\dot{\nu }`$ preceding the glitch (Table 3), although the data span is rather short. Table 4 gives the extrapolated and fitted parameters for the glitch. Fitting of a single exponential decay gives a relatively good fit to the post-glitch data with $`Q0.13`$ and $`\tau _d640`$ d and an rms residual of 1690 $`\mu `$s. However, there is clear evidence in the residuals immediately after the glitch for a shorter-term decay. The glitch parameters in Table 4 are from a simultaneous fit of two exponential decays, one with a short decay time and the other representing the longer-term decay. The final fit is extremely good, with a final rms residual of only 334 $`\mu `$s and the residuals dominated by random noise. Fitting of the short-term decay resulted in an increase in the estimates for $`Q`$ and $`\tau _d`$ for the longer-term decay by about 25 per cent as shown in Table 4.
## 4 Discussion
In this paper we have presented timing observations of 40 pulsars over various intervals up to a maximum of 8.9 years and analysed these data for improved astrometric and pulse frequency parameters and for glitch activity. In total, 30 glitches were detected in 11 pulsars, including the largest known glitch, with $`\mathrm{\Delta }\nu _g/\nu 6.5\times 10^6`$ in PSR J1614$``$5047. Twelve glitches were detected in the period of PSR J1341$``$6220 in a data span of 8.2 years, making this the most frequently glitching pulsar known. Evidence was found in PSR J1614$``$5047 for a new class of irregularity in which the pulse frequency increases markedly over a few-day interval. There appears to be an accompanying decrease in the magnitude of the frequency derivative, implying a corresponding decrease in braking torque.
Table 5 lists all known glitches, 76 in total, giving the fractional glitch amplitude $`\mathrm{\Delta }\nu _g/\nu `$ and, in parentheses, the approximate MJD of the glitch. Fig. 16 shows a histogram of fractional glitch amplitudes. Largely because of the Vela pulsar, the most common glitches are large, with $`\mathrm{\Delta }\nu _g/\nu `$ in the range $`13\times 10^6`$. In most other pulsars, smaller glitches with fractional amplitude down to $`10^8`$ are more common. Glitches with size smaller than $`10^9`$ are difficult to identify, especially in noisy pulsars, and the sample is certainly incomplete at this level. However, there does appear to be a reduced rate of occurence of glitches with $`\mathrm{\Delta }\nu _g/\nu <10^8`$.
In most models for the glitch phenomenon, the sudden spin-up is triggered by the release of stress built up as a result of the steady spin down of the pulsar. For the original star-quake model \[Baym et al. 1969\], the equilibrium shape of the star becomes less oblate as the star spins down. At some point, the crust cracks and relaxes to (or toward) the new equilibrium shape, reducing its moment of inertia and hence spinning it up. However, this model fails to account for frequent giant glitches, as seen, for example, in the Vela pulsar, as the rate of change of oblateness is too slow. In models based on unpinning of internal superfluid vortices \[Alpar, Cheng & Pines 1989, Ruderman, Zhu & Chen 1998\], the stresses on the pinned vortices build up as the crust slows down until finally some fraction of them unpin and then repin at a larger radius, resulting in a transfer of angular momentum to the crust.
In these models, one expects some relation between the size of the glitch and the length of time that the pulsar has been slowing down since the previous glitch or the integrated change in spin rate since the last glitch. Fig. 17(a) shows glitch fractional amplitudes plotted against the length of the preceding inter-glitch interval ($`\mathrm{\Delta }t_g`$), and in Fig. 17(b) against the change in spin frequency since the previous glitch ($`|\dot{\nu }|\mathrm{\Delta }t_g`$). These figures show that, contrary to expectations, there is no general relation between fractional glitch amplitude and either the time since the previous glitch or the total change in spin frequency since the previous glitch. For the Vela pulsar (marked with a $``$ in the Figure), the smaller glitches do tend to have shorter preceding intervals and all but one of the giant glitches have preceding inter-glitch intervals of about 1000 days. However, for PSR J1341$``$6220 there is if anything an inverse relationship, with larger glitches occuring after shorter intervals. For PSR J1740$``$3015 there is no relationship between glitch size and preceding interval. Fig. 17(b) shows that similar relationships for individual pulsars hold when glitch size is plotted against accumulated spin-down since the last glitch. The three points to the right on this plot are for the Crab pulsar; these too show an inverse relationship.
Alternatively, if a small glitch results from a release of only part of the built-up stress, one might expect another glitch to occur soon after, when the breaking strain is again reached. Conversely, if a large glitch releases all or most of the stress, a long time would be required for it to build up again. This suggests a correlation between glitch size and time to the next glitch. Fig. 18 shows glitch size plotted against duration of the following inter-glitch interval and accumulated spin-down in that period. No such correlation is observed for the Vela pulsar, where the small glitches are followed by relatively long intervals, but a weak positive correlation is seen for PSR J1341$``$6220 and PSR J1740$``$3015. Excepting the Crab pulsar, other pulsars (marked with a dot) have a good correlation between glitch size and accumulated spin-down frequency following the glitch.
These results suggest that the triggering of glitches is a local phenomenon, not dependent on global stresses. This lends some support to the ideas of Ruderman et al. (1998) in which migration of magnetic flux tubes determines the stresses on the pinned vortices. This model is also supported by the observations of the “slow” glitches and subsequent decrease in slow-down rate in PSR J1614$``$5047.
McKenna & Lyne (1990) introduced the glitch activity parameter, $`A_g`$, defined to be the accumulated pulse frequency change $`\mathrm{\Delta }\nu _g`$ due to glitches divided by the observation data span. It therefore has the same units as $`\dot{\nu }`$ and represents the portion of $`\dot{\nu }`$ which is overcome by glitches. This is typically small; the largest known value of $`A_g`$, for the Vela pulsar, is $`2.5\times 10^{13}`$ s<sup>-2</sup>, only about 2 per cent of its spin-down rate $`\dot{\nu }1.6\times 10^{11}`$ s<sup>-2</sup>. Fig. 19(a) shows activity parameter plotted against characteristic age, $`\tau _c`$ for most of the pulsars with extensive timing data. As noted by McKenna & Lyne (1990), there is a clear peak in activity for pulsars with ages between 2,000 and 20,000 years.
However, there is a substantial group of young pulsars with low glitch activity; these so far have had no observed glitch and are represented by upper limits in Fig. 19(a) corresponding to a single glitch of fractional size $`10^9`$. These eight pulsars have a mean data span of more than six years and so, at the least, they are infrequent glitchers. A single glitch of amplitude $`10^6`$ would raise them into the same region as the other young pulsars. The dashed line corresponds to a uniform accumulated $`\dot{\nu }`$ over the characteristic spin-down time. This is approximately true for pulsars with age of less than or about $`10^6`$ years, but older pulsars clearly have less glitch activity than predicted by this rule.
In Fig. 19(b) the activity parameter is plotted against the absolute value of $`\dot{\nu }`$. Some theoretical models for glitches (e.g. Ruderman, Zhu & Chen, 1998) predict that the ratio of the effective spin-up rate due to glitches to the spin-down rate is proportional to the ratio of the moment of inertia of the crustal superfluid to the total moment of inertia of the neutron star. A constant ratio corresponds to a line of slope $`1`$ in Fig. 19(b). For the younger pulsars, except PSR B0531+21 (Crab) and PSR B1509$``$58, the points correspond to a fraction of 1 – 2 per cent (cf. Lyne et al. 1999). Pulsars with spin-down rates less than about $`10^{14}`$ s<sup>-1</sup> have a significantly smaller fraction of their spindown recovered by glitches.
When a significant post-glitch decay is observed, it is generally well described by the exponential relation given in Equation 2. In Fig. 20(a) the fraction of the glitch which decays, $`Q=\mathrm{\Delta }\nu _d/\mathrm{\Delta }\nu _g`$ is plotted against characteristic age. With one exception, PSR B0525+21, old pulsars tend to have low values of $`Q`$. Younger pulsars can have larger $`Q`$, but many glitches in young pulsars have little or no decay.
Fig. 20(b) is a plot of glitch decay timescale versus pulsar characteristic age. This shows that, apart from the Crab pulsar (the three points in the lower left of the figure) in which the decay timescale is very short, there is no relation between decay timescale and pulsar age. We emphasise that this result applies to the longest decay timescale present. This may indicate the unreliability of characteristic age as an indicator of true age and hence internal temperature, or it may indicate that decay time, at least for pulsars with age greater than a few thousand years, is dependent on other properties, for example, the core magnetic field.
These observations have demonstrated the great diversity of glitch properties. Glitch activity is clearly greatest in pulsars with ages of between a few times $`10^3`$ and $`10^5`$ years. The Crab pulsar has distinctly different glitch properties from those of the middle-aged pulsars, and other pulsars of similar age show no glitches at all. Apart from these clear trends, glitch properties vary greatly, both for successive glitches from a given pulsar and across different pulsars, with few systematic trends. These properties suggest that the glitch phenomenon, both the event and the following response, depend on quasi-random processes occurring in the pulsar, rather than global properties such as slow-down rate or age. Crustal processes driven by magnetic field evolution as proposed by Ruderman et al. (1998) seem more consistent with this than the vortex creep models of Alpar et al. (1989).
## ACKNOWLEDGEMENTS
We thank the many colleagues who have helped with observing and software development over the course of this project, and the staff of Parkes Observatory for their always cheerful assistance. In particular, we thank John Sarkissian for doing much of the observing during 1998, and Matthew Britton, Stuart Anderson and John Yamasaki for help with correlator hardware and software development. NW thanks the National Nature Science Foundation (NNSF) of China for their support of this work and VMK acknowledges support from an Alfred P. Sloan Research Fellowship. The Parkes radio telescope is part of the Australia Telescope which is funded by the Commonwealth of Australia for operation as a National Facility managed by CSIRO.
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# Practical Reasoning for Expressive Description LogicsThis paper appeared in the Proceedings of the 6th International Conference on Logic for Programming and Automated Reasoning (LPAR’99), number 1704 Lecture Notes in Artificial Intelligence, pages 161-180. Springer-Verlag, September 1999.
## 1 Motivation
Description Logics (DLs) are a well-known family of knowledge representation formalisms \[DLNS96\]. They are based on the notion of concepts (unary predicates, classes) and roles (binary relations), and are mainly characterised by constructors that allow complex concepts and roles to be built from atomic ones. Sound and complete algorithms for the interesting inference problems such as subsumption and satisfiability of concepts are known for a wide variety of DLs \[SS91, DLNdN91, Sat96, DL96, CDL99\].
To be used in a specific application, the expressivity of the DL must be sufficient to describe relevant properties of objects in the application domain. For example, transitive roles (e.g. “ancestor”) and inverse roles (e.g. “successor”/“predecessor”) play an important rôle not only in the adequate representation of complex, aggregated objects \[HS99\], but also for reasoning with conceptual data models \[CLN94\]. Moreover, reasoning with respect to cyclic definitions is crucial for applying DLs to reasoning with database schemata \[CDL98a\].
The relevant inference problems for (extensions of) DLs that allow for transitive and inverse roles are known to be decidable \[DL96\], and appropriate inference algorithms have been described \[DM98\], but their high degree of non-determinism appears to prohibit their use in realistic applications. This is mainly due to the fact that these algorithms can handle not just transitive roles but also the transitive closure of roles. It has been shown \[Sat96\] that restricting a DL to transitive roles can lead to a lower complexity, and that transitive roles (even when combined with role hierarchies) allow for algorithms that behave quite well in realistic applications \[Hor98\]. However, it remained to show that this is still true when inverse roles and qualifying number restrictions are also present.
This paper extends our understanding of these issues in several directions. Firstly, we present an algorithm that decides satisfiability of $`𝒜𝒞`$ \[SS91\] (which can be seen as a notational variant of the multi modal logic $`𝖪_m`$) extended with transitive and inverse roles, role hierarchies, and qualifying number restrictions, i.e., concepts of the form $`(3hasChildFemale)`$ that allow the description of objects by restricting the number of objects of a given type they are related to via a certain role. The algorithm can also be used for checking satisfiability and subsumption with respect to general concept inclusion axioms (and thus cyclic definitions) because these axioms can be “internalised”. The absence of transitive closure leads to a lower degree of non-determinism, and experiments indicate that the algorithm is well-suited for implementation.
Secondly, we show that $`𝒜𝒞`$ extended with both transitive *and* inverse roles is still in Pspace. The algorithm used to prove this rather surprising result introduces an enhanced *blocking* technique. In general, blocking is used to ensure termination of the algorithm in cases where it would otherwise be stuck in a loop. The enhanced blocking technique allows such cases to be detected earlier and should provide useful efficiency gains in implementations of this and more expressive DLs.
Finally, we investigate the limits of decidability for this family of DLs, showing that relaxing the constraints placed on the kinds of roles allowed in number restrictions leads to the undecidability of all inference problems.
Due to a lack of space we can only present selected proofs. For full details please refer to \[HST98, HST99\].
## 2 Preliminaries
In this section, we present the syntax and semantics of the various DLs that are investigated in subsequent sections. This includes the definition of inference problems (concept subsumption and satisfiability, and both of these problems with respect to terminologies) and how they are interrelated.
The logics we will discuss are all based on an extension of the well known DL $`𝒜𝒞`$ \[SS91\] to include transitively closed primitive roles \[Sat96\]; we will call this logic $`𝒮`$ due to its relationship with the proposition (multi) modal logic $`\mathrm{𝐒𝟒}_{(𝐦)}`$ \[Sch91\].<sup>1</sup><sup>1</sup>1The logic $`𝒮`$ has previously been called $`𝒜𝒞_{R^+}`$, but this becomes too cumbersome when adding letters to represent additional features. This basic DL is then extended in a variety of ways—see Figure 1 for an overview.
###### Definition 2.1
Let $`𝐂`$ be a set of *concept names* and $`𝐑`$ a set of *role names* with transitive role names $`𝐑_+𝐑`$. The set of $`𝒮`$-roles is $`𝐑\{R^{}R𝐑\}`$. The set of $`𝒮`$-*concepts* is the smallest set such that every concept name is a concept, and, if $`C`$ and $`D`$ are concepts and $`R`$ is an $`𝒮`$-role, then $`(CD)`$, $`(CD)`$, $`(\neg C)`$, $`(R.C)`$, and $`(R.C)`$ are also concepts.
To avoid considering roles such as $`R^{}`$, we define a function $`𝖨𝗇𝗏`$ on roles such that $`𝖨𝗇𝗏(R)=R^{}`$ if $`R`$ is a role name, and $`𝖨𝗇𝗏(R)=S`$ if $`R=S^{}`$. We also define a function $`𝖳𝗋𝖺𝗇𝗌`$ which returns $`\mathrm{true}`$ iff $`R`$ is a transitive role. More precisely, $`𝖳𝗋𝖺𝗇𝗌(R)=\mathrm{true}`$ iff $`R𝐑_+`$ or $`𝖨𝗇𝗏(R)𝐑_+`$.
$`𝒮`$ is obtained from $`𝒮`$ by allowing, additionally, for a set of *role inclusion axioms* of the form $`RS`$, where $`R`$ and $`S`$ are two roles, each of which can be inverse. For a set of role inclusion axioms $``$,
$$^+:=(\{𝖨𝗇𝗏(R)𝖨𝗇𝗏(S)RS\},\text{*})$$
is called a role hierarchy, where $`\text{*}`$ is the transitive-reflexive closure of $``$ over $`\{𝖨𝗇𝗏(R)𝖨𝗇𝗏(S)RS\}`$.
$`𝒮𝒬`$ is obtained from $`𝒮`$ by allowing, additionally, for *qualifying number restrictions*, i.e., for concepts of the form $`(nRC)`$ and $`(nRC)`$, where $`R`$ is a *simple* (possibly inverse) role and $`n`$ is a non-negative integer. A role is called *simple* iff it is neither transitive nor has transitive sub-roles.
$`𝒮𝒩`$ is the restriction of $`𝒮𝒬`$ where qualifying number restrictions may only be of the form $`(nR)`$ and $`(nR)`$. In this case, we omit the symbol $``$ and write $`(nR)`$ and $`(nR)`$ instead.
An interpretation $`=(\mathrm{\Delta }^{},^{})`$ consists of a set $`\mathrm{\Delta }^{}`$, called the domain of $``$, and a *valuation* $`^{}`$ which maps every concept to a subset of $`\mathrm{\Delta }^{}`$ and every role to a subset of $`\mathrm{\Delta }^{}\times \mathrm{\Delta }^{}`$ such that, for all concepts $`C`$, $`D`$, roles $`R`$, $`S`$, and non-negative integers $`n`$, the properties in Figure 1 are satisfied, where $`\mathrm{}M`$ denotes the cardinality of a set $`M`$. An interpretation satisfies a role hierarchy $`^+`$ iff $`R^{}S^{}`$ for each $`R\text{*}S^+`$; we denote this fact by $`^+`$ and say that $``$ is a model of $`^+`$.
A concept $`C`$ is called satisfiable with respect to a role hierarchy $`^+`$ iff there is some interpretation $``$ such that $`^+`$ and $`C^{}\mathrm{}`$. Such an interpretation is called a model of $`C`$ w.r.t. $`^+`$. A concept $`D`$ subsumes a concept $`C`$ w.r.t. $`^+`$ (written $`C_^+D`$) iff $`C^{}D^{}`$ holds for each model $``$ of $`^+`$. For an interpretation $``$, an individual $`x\mathrm{\Delta }^{}`$ is called an instance of a concept $`C`$ iff $`xC^{}`$.
All DLs considered here are closed under negation, hence subsumption and (un)satisfiability w.r.t. role hierarchies can be reduced to each other: $`C_^+D`$ iff $`C\neg D`$ is unsatisfiable w.r.t. $`^+`$, and $`C`$ is unsatisfiable w.r.t. $`^+`$ iff $`C_^+A\neg A`$ for some concept name $`A`$.
In \[Baa91, Sch91, BBN<sup>+</sup>93\], the *internalisation* of terminological axioms is introduced, a technique that reduces reasoning with respect to a (possibly cyclic) terminology to satisfiability of concepts. In \[Hor98\], we saw how role hierarchies can be used for this reduction. In the presence of inverse roles, this reduction must be slightly modified.
###### Definition 2.2
A *terminology* $`𝒯`$ is a finite set of *general concept inclusion axioms*, $`𝒯=\{C_1D_1,\mathrm{},C_nD_n\}`$, where $`C_i,D_i`$ are arbitrary $`𝒮𝒬`$-concepts. An interpretation $``$ is said to be a *model* of $`𝒯`$ iff $`C_i^{}D_i^{}`$ holds for all $`C_iD_i𝒯`$. $`C`$ is *satisfiable* with respect to $`𝒯`$ iff there is a model $``$ of $`𝒯`$ with $`C^{}\mathrm{}`$. Finally, $`D`$ *subsumes* $`C`$ with respect to $`𝒯`$ iff for each model $``$ of $`𝒯`$ we have $`C^{}D^{}`$.
The following Lemma shows how general concept inclusion axioms can be *internalised* using a “universal” role $`U`$, that is, a transitive super-role of all roles occurring in $`𝒯`$ and their respective inverses.
###### Lemma 2.3
Let $`𝒯`$ be a terminology, $``$ a set of role inclusion axioms and $`C,D`$ $`𝒮𝒬`$-concepts and let
$$C_𝒯:=\underset{C_iD_i𝒯}{\text{}}\neg C_iD_i.$$
Let $`U`$ be a transitive role that does not occur in $`𝒯,C,D`$, or $``$. We set
$$_U:=\{RU,𝖨𝗇𝗏(R)UR\text{ occurs in }𝒯,C,D\text{, or }\}.$$
Then $`C`$ is satisfiable w.r.t. $`𝒯`$ and $`^+`$ iff $`CC_𝒯U.C_𝒯`$ is satisfiable w.r.t. $`_U^+`$. Moreover, $`D`$ subsumes $`C`$ with respect to $`𝒯`$ and $`^+`$ iff $`C\neg DC_𝒯U.C_𝒯`$ is unsatisfiable w.r.t. $`_U^+`$.
The proof of Lemma 2.3 is similar to the ones that can be found in \[Sch91, Baa91\]. Most importantly, it must be shown that, (a) if a $`𝒮𝒬`$-concept $`C`$ is satisfiable with respect to a terminology $`𝒯`$ and a role hierarchy $`^+`$, then $`C,𝒯`$ have a *connected* model, and (b) if $`y`$ is reachable from $`x`$ via a role path (possibly involving inverse roles), then $`x,yU^{}`$. These are easy consequences of the semantics and the definition of $`U`$.
###### Theorem 2.4
Satisfiability and subsumption of $`𝒮𝒬`$-concepts (resp. $`𝒮`$-concepts) w.r.t. terminologies and role hierarchies are polynomially reducible to (un)satisfiability of $`𝒮𝒬`$-concepts (resp. $`𝒮`$-concepts) w.r.t. role hierarchies.
## 3 Reasoning for $`𝒮`$ Logics
In this section, we present two tableaux algorithms: the first decides satisfiability of $`𝒮𝒬`$-concepts, and can be used for all $`𝒮𝒬`$ reasoning problems (see Theorem 2.4); the second decides satisfiability (and hence subsumption) of $`𝒮`$-concepts in Pspace. Please note that $`𝒮𝒩`$ (and hence $`𝒮𝒬`$) no longer has the finite model property: for example, the following concept, where $`R`$ is a transitive super-role of $`F`$, is satisfiable, but each of its models has an infinite domain.
$$\neg CF^{}.(C1F)R^{}.(F^{}.(C1F))$$
This concept requires the existence of an infinite $`F^{}`$-path, where the first element on the path satisfies $`\neg C`$ while all other elements satisfy $`C1F`$. This path cannot collapse into a cycle: (a) it cannot return to the first element because this element cannot satisfy both $`C`$ and $`\neg C`$; (b) it cannot return to any subsequent element on the path because then this node would not satisfy $`1F`$.
The correctness of the algorithms we are presenting can be proved by showing that they create a *tableau* for a concept iff it is satisfiable. For ease of construction, we assume all concepts to be in *negation normal form* (NNF), that is, negation occurs only in front of concept names. Any $`𝒮𝒬`$-concept can easily be transformed to an equivalent one in NNF by pushing negations inwards \[HNS90\]; with $`\mathrm{}C`$ we denote the NNF of $`\neg C`$. For a concept $`C`$ in NNF we define $`\text{clos}(C)`$ as the smallest set of concepts that contains $`C`$ and is closed under subconcepts and $`\mathrm{}`$. Please note that size of $`\text{clos}(C)`$ is linearly bounded by the size of $`C`$.
###### Definition 3.1
Let $`D`$ be a $`𝒮𝒬`$-concept in NNF, $`^+`$ a role hierarchy, and $`𝐑_D`$ the set of roles occurring in $`D`$ and $`^+`$ together with their inverses. Then $`T=(𝐒,,)`$ is a *tableau* for $`D`$ w.r.t. $`^+`$ iff $`𝐒`$ is a set of individuals, $`:𝐒2^{\text{clos}(D)}`$ maps each individual to a set of concepts, $`:𝐑_D2^{𝐒\times 𝐒}`$ maps each role to a set of pairs of individuals, and there is some individual $`s𝐒`$ such that $`D(s)`$. Furthermore, for all $`s,t𝐒`$, $`C,C_1,C_2\text{clos}(D)`$, and $`R,S𝐑_D`$, it holds that:
1. if $`C(s)`$, then $`\neg C(s)`$,
2. if $`C_1C_2(s)`$, then $`C_1(s)`$ and $`C_2(s)`$,
3. if $`C_1C_2(s)`$, then $`C_1(s)`$ or $`C_2(s)`$,
4. if $`S.C(s)`$ and $`s,t(S)`$, then $`C(t)`$,
5. if $`S.C(s)`$, then there is some $`t𝐒`$ such that $`s,t(S)`$ and $`C(t)`$,
6. if $`S.C(s)`$ and $`s,t(R)`$ for some $`R\text{*}S`$ with $`𝖳𝗋𝖺𝗇𝗌(R)`$, then $`R.C(t)`$,
7. $`x,y(R)`$ iff $`y,x(𝖨𝗇𝗏(R))`$,
8. if $`s,t(R)`$ and $`R\text{*}S`$, then $`s,t(S)`$,
9. if $`(nSC)(s)`$, then $`\mathrm{}S^T(s,C)n`$,
10. if $`(nSC)(s)`$, then $`\mathrm{}S^T(s,C)n`$,
11. if $`(nSC)(s)`$ and $`s,t(S)`$ then $`C(t)`$ or $`\mathrm{}C(t)`$,
where we use $``$ as a placeholder for both $``$ and $``$ and we define
$$S^T(s,C):=\{t𝐒s,t(S)\text{and}C(t)\}.$$
Tableaux for $`𝒮`$-concepts are defined analogously and must satisfy Properties 1-7, where, due to the absence of a role hierarchy, $`\text{*}`$ is the identity.
Due to the close relationship between models and tableaux, the following lemma can be easily proved by induction. As a consequence, an algorithm that constructs (if possible) a tableau for an input concept is a decision procedure for satisfiability of concepts.
###### Lemma 3.2
A $`𝒮𝒬`$-concept (resp. $`𝒮`$-concept) $`D`$ is satisfiable w.r.t. a role hierarchy $`^+`$ iff $`D`$ has a tableau w.r.t. $`^+`$.
### 3.1 Reasoning in $`𝒮𝒬`$
In the following, we give an algorithm that, given a $`𝒮𝒬`$-concept $`D`$, decides the existence of a tableaux for $`D`$. We implicitly assume an arbitrary but fixed role hierarchy $`^+`$. The tableaux algorithm works on a finite completion tree (a tree some of whose nodes correspond to individuals in the tableau, each node being labelled with a set of $`𝒮𝒬`$-concepts), and employs a *blocking* technique \[HS99\] to guarantee termination: If a path contains two pairs of successive nodes that have pair-wise identical label and whose connecting edges have identical labels, then the path beyond the second pair is no longer expanded, it is said to be blocked. Blocked paths can be “unravelled” to construct an infinite tableau. The identical labels make sure that copies of the first pair and their descendants can be substituted for the second pair of nodes and their respective descendants.
###### Definition 3.3
A *completion tree* for a $`𝒮𝒬`$-concept $`D`$ is a tree where each node $`x`$ of the tree is labelled with a set $`(x)\text{clos}(D)`$ and each edge $`x,y`$ is labelled with a set $`(x,y)`$ of (possibly inverse) roles occurring in $`\text{clos}(D)`$; explicit inequalities between nodes of the tree are recorded in a binary relation $`\doteq ̸`$ that is implicitly assumed to be symmetric.
Given a completion tree, a node $`y`$ is called an $`R`$-*successor* of a node $`x`$ iff $`y`$ is a successor of $`x`$ and $`S(x,y)`$ for some $`S`$ with $`S\text{*}R`$. A node $`y`$ is called an $`R`$*-neighbour* of $`x`$ iff $`y`$ is an $`R`$-successor of $`x`$, or if $`x`$ is an $`𝖨𝗇𝗏(R)`$-successor of $`y`$. Predecessors and ancestors are defined as usual.
A node is *blocked* iff it is directly or indirectly blocked. A node $`x`$ is *directly blocked* iff none of its ancestors are blocked, and it has ancestors $`x^{}`$, $`y`$ and $`y^{}`$ such that
1. $`x`$ is a successor of $`x^{}`$ and $`y`$ is a successor of $`y^{}`$ *and*
2. $`(x)=(y)`$ and $`(x^{})=(y^{})`$ *and*
3. $`(x^{},x)=(y^{},y)`$.
In this case we will say that $`y`$ *blocks* $`x`$. Since this blocking technique involves pairs of nodes, it is called *pair-wise* blocking.
A node $`y`$ is *indirectly blocked* iff one of its ancestors is blocked, or it is a successor of a node $`x`$ and $`(x,y)=\mathrm{}`$; the latter condition avoids wasted expansions after an application of the $``$-rule.
For a node $`x`$, $`(x)`$ is said to contain a clash iff $`\{A,\neg A\}(x)`$ or if, for some concept $`C`$, some role $`S`$, and some $`n`$: $`(nSC)(x)`$ and there are $`n+1`$ $`S`$-neighbours $`y_0,\mathrm{},y_n`$ of $`x`$ such that $`C(y_i)`$ and $`y_i\doteq ̸y_j`$ for all $`0i<jn`$. A completion tree is called clash-free iff none of its nodes contains a clash; it is called complete iff none of the expansion rules in Figure 2 is applicable.
For a $`𝒮𝒬`$-concept $`D`$, the algorithm starts with a completion tree consisting of a single node $`x`$ with $`(x)=\{D\}`$ and $`\doteq ̸=\mathrm{}`$. It applies the expansion rules in Figure 2, stopping when a clash occurs, and answers “$`D`$ is satisfiable” iff the completion rules can be applied in such a way that they yield a complete and clash-free completion tree.
The soundness and completeness of the tableaux algorithm is an immediate consequence of Lemmas 3.2 and 3.4.
###### Lemma 3.4
Let $`D`$ be an $`𝒮𝒬`$-concept.
1. The tableaux algorithm terminates when started with $`D`$.
2. If the expansion rules can be applied to $`D`$ such that they yield a complete and clash-free completion tree, then $`D`$ has a tableau.
3. If $`D`$ has a tableau, then the expansion rules can be applied to $`D`$ such that they yield a complete and clash-free completion tree.
The proof can be found in the appendix. Here, we will only discuss the intuition behind the expansion rules and their correspondence to the constructors of $`𝒮𝒬`$. Roughly speaking,<sup>2</sup><sup>2</sup>2For the following considerations, we employ a simpler view of the correspondence between completion trees and models, and need not bother with the path construction mentioned above. the completion tree is a partial description of a model whose individuals correspond to nodes, and whose interpretation of roles is taken from the edge labels. Since the completion tree is a tree, this would not yield a correct interpretation of transitive roles, and thus the interpretation of transitive roles is built via the transitive closure of the relations induced by the corresponding edge labels.
The $``$-, $``$-, $``$\- and $``$-rules are the standard tableaux rules for $`𝒜𝒞`$ or the propositional modal logic $`𝖪_m`$. The $`_+`$-rule is the standard rule for $`𝒜𝒞_{R^+}`$ or the propositional modal logic $`\mathrm{𝖲𝟦}_m`$ extended to deal with role-hierarchies as follows. Assume a situation that satisfies the precondition of the $`_+`$-rule, i.e., $`S.C(x)`$, and there is an $`R`$-neighbour $`y`$ of $`x`$ with $`𝖳𝗋𝖺𝗇𝗌(R)`$, $`R\text{*}S`$ and $`R.C(y)`$. If $`y`$ has an $`R`$-successor $`z`$, then, due to the transitivity of $`R`$, $`z`$ is also an $`R`$-successor of $`x`$. Since $`R\text{*}S`$, it is also an $`S`$-successor of $`x`$ and hence must satisfy $`C`$. This is ensured by adding $`R.C`$ to $`(z)`$
The rules dealing with qualifying number restrictions work similarly to the rules given in \[BBH96\]. For a concept $`(nRC)(x)`$, the $``$-rule generates $`n`$ $`R`$-successors $`y_1,\mathrm{},y_n`$ of $`x`$ with $`C(y_i)`$ . To prevent the $``$-rule from indentifying the new nodes, it also sets $`y_i\doteq ̸y_j`$ for each $`1i<jn`$ . Conversely, if $`(nRC)(x)`$ and $`x`$ has more than $`n`$ $`R`$-neighbours that are labelled with $`C`$, then the $``$-rule chooses two of them that are not in $`\doteq ̸`$ and merges them, together with the edges connecting them with $`x`$. The definition of a clash takes care of the situation where the $`\doteq ̸`$ relation makes it impossible to merge any two $`R`$-neighbours of $`x`$, while the choose-rule ensures that all $`R`$-neighbours of $`x`$ are labelled with either $`C`$ or $`\mathrm{}C`$. Without this rule, the unsatisfiability of concepts like $`(3RA)(1RB)(1R\neg B)`$ would go undetected. The relation $`\doteq ̸`$ is used to prevent infinite sequences of rule applications for contradicting number restrictions of the form $`(nRC)`$ and $`((m)RC)`$, with $`n>m`$. Labelling edges with sets of roles allows a single node to be both an $`R`$ and $`S`$-successor of $`x`$ even if $`R`$ and $`S`$ are not comparable with respect to $`\text{*}`$.
The following theorem is an immediate consequence of Lemma 3.2 and 3.4, and Theorem 2.4.
###### Theorem 3.5
The tableaux algorithm is a decision procedure for the satisfiability and subsumption of $`𝒮𝒬`$-concepts with respect to terminologies.
### 3.2 A PSpace-algorithm for $`𝒮`$
To obtain a (worst-case) optimal algorithm for $`𝒮`$, the $`𝒮𝒬`$ algorithm is modified as follows. (a) Since $`𝒮`$ does not allow for qualifying number restrictions the $``$-, $``$-, and choose-rule can be omitted. In the absence of the choose-rule we may assume all concepts appearing in labels to be in NNF from the (smaller) set of all subconcepts of $`D`$ denoted by $`\text{sub}(D)`$, and in the absence of role hierarchies, edge labels can be restricted to roles (instead of sets of roles). Due to the absence of number restrictions the logic still has the finite model property, and blocking no longer need involve two pairs of nodes with identical labels, but only two nodes with (originally) identical labels. (b) To obtain a PSpace algorithm, we employ a refined blocking strategy which further loosens this “identity” condition to a “similarity” condition. This is achieved by using a second label $``$ for each node. In the following, we will describe and motivate this blocking technique; detailed proofs as well as an extension of this result to $`𝒮𝒩`$ can be found in \[HST98\].
Establishing a PSpace-result for $`𝒮`$ is not as straightforward as it might seem at a first glance. One problem is the presence of inverse roles which might lead to constraints propagating upwards in the tree. This is not compatible with the standard trace technique \[SS91\] that keeps only a single path in memory at the same time, because constraints propagating upwards in the tree may have an influence on paths that have already been visited and have been discarded from memory. There are at least two possibilities to overcome this problem: (1) by guessing which constraints might propagate upwards beforehand; (2) by a *reset-restart* extension of the trace technique described later in this section. Unfortunately, this is not the only problem. To apply either of these two techniques, it is also necessary to establish a polynomial bound on the length of paths in the completion tree. This is easily established for logics such as $`𝒜𝒞`$ that do not allow for transitive roles. For $`𝒜𝒞`$ with transitive roles (i.e., $`𝒮`$), this bound is due to the fact that, for a node $`x`$ to block a node $`y`$, it is sufficient that $`(y)(x)`$. In the presence of inverse roles, we use a more sophisticated blocking technique to establish the polynomial bound.
###### Definition 3.6
A *completion tree* for an $`𝒮`$ concept $`D`$ is a tree where each node $`x`$ of the tree is labelled with two sets $`(x)(x)\text{sub}(D)`$, and each edge $`x,y`$ is labelled with a (possibly inverse) role $`(x,y)`$ occurring in $`\text{sub}(D)`$.
$`R`$-neighbours, -successors, and -predecessors are defined as in Definition 3.3 where, in the absence of role hierarchies, $`\text{*}`$ is the identity on $`𝐑`$.
A node $`x`$ is *blocked* iff $`x`$ has a blocked ancestor $`y`$, or $`x`$ has an ancestor $`y`$ and a predecessor $`x^{}`$ with $`(x^{},x)=S`$, and
$$(x)(y)\text{ and }(x)/𝖨𝗇𝗏(S)=(y)/𝖨𝗇𝗏(S),$$
where $`(x)/𝖨𝗇𝗏(S)=\{𝖨𝗇𝗏(S).C(x)\}`$.
For a node $`x`$, $`(x)`$ is said to contain a *clash* iff $`\{A,\neg A\}(x)`$. A completion tree to which none of the expansion rules given in Figure 3 is applicable is called *complete*.
For an $`𝒮`$-concept $`D`$, the algorithm starts with a completion tree consisting of a single node $`x`$ with $`(x)=(x)=\{D\}`$. It applies the expansion rules in Figure 3, stopping when a clash occurs, and answers “$`D`$ is satisfiable” iff the completion rules can be applied in such a way that they yield a complete and clash-free completion tree.
As for $`𝒮𝒬`$, correctness of the algorithm can be proved by first showing that a $`𝒮`$-concept is satisfiable iff it has a tableau, and next proving the $`𝒮`$-analogue of Lemma 3.4, see \[HST98\].
###### Theorem 3.7
The tableaux algorithm is a decision procedure for satisfiability and subsumption of $`𝒮`$-concepts.
Since blocking plays a major rôle both in the proof of Theorem 3.7 and especially in the following complexity considerations, we will discuss it here in more detail. Blocking guarantees the termination of the algorithm. For DLs such as $`𝒜𝒞`$, termination is mainly due to the fact that the expansion rules can only add new concepts that are strictly smaller than the concept that triggered their application.
For $`𝒮`$ this is no longer true: the $`_+`$-rule introduces new concepts that are the same size as the triggering concept. To ensure termination, nodes labelled with a subset of the label of an ancestor are *blocked*. Since rules can be applied “top-down” (successors are only generated if no other rules are applicable, and the labels of inner nodes are never touched again) and subset-blocking is sufficient (i.e., for a node $`x`$ to be blocked by an ancestor $`y`$, it is sufficient that $`(x)(y)`$), it is possible to give a polynomial bound on the length of paths.
For $`𝒮`$, *dynamic blocking* was introduced in \[HS99\], i.e., blocks are not established on a once-and-for-all basis, but established and broken dynamically. Moreover, blocks must be established on the basis of label *equality*, since value restrictions can now constrain predecessors as well as successors. Unfortunately, this may lead to completion trees with exponentially long paths because there are exponentially many possibilities to label sets on such a path. Due to the non-deterministic $``$-rule, these exponentially many sets may actually occur.
This non-determinism is not problematical for $`𝒮`$ because disjunctions need not be completely decomposed to yield a subset-blocking situation. For an optimal $`𝒮`$ algorithm, the additional label $``$ was introduced to enable a sort of subset-blocking which is independent of the $``$-non-determinism. Intuitively, $`(x)`$ is the restriction of $`(x)`$ to those non-decomposed concepts that $`x`$ must satisfy, whereas $`(x)`$ contains boolean decompositions of these concepts as well as those that are imposed by value restrictions in descendants. If $`x`$ is blocked by $`y`$, then all concepts in $`(x)`$ are eventually decomposed in $`(y)`$. However, in order to substitute $`x`$ by $`y`$, $`x`$’s constraints on predecessors must be at least as strong as $`y`$’s; this is taken care of by the second blocking condition.
Let us consider a path $`x_0,x_1,\mathrm{},x_n`$ where all edges are labelled $`R`$ with $`𝖳𝗋𝖺𝗇𝗌(R)`$, the only kind of path along which the length of the longest concept in the labels might not decrease. If no rules can be applied, then we have, for $`1i<n`$,
$$\begin{array}{ccc}\hfill (x_{i+1})/𝖨𝗇𝗏(R)& & (x_i)/𝖨𝗇𝗏(R)\text{ and}\hfill \\ \hfill (x_i)& & (x_{i+1})\{C_i\}\hfill \end{array}$$
(where $`R.C_i(x_i)`$ triggered the generation of $`x_{i+1}`$). This limits the number of different labels and guarantees blocking after a polynomial number of steps.
###### Lemma 3.8
The paths of a completion tree for a concept $`D`$ have a length of at most $`m^4`$ where $`m=|\text{sub}(D)|`$.
Finally, a slight modification of the expansion rules given in Figure 3 yields a PSpace algorithm. This modification is necessary because the original algorithm must keep the whole completion tree in memory—which needs exponential space even though the length of its paths is polynomially bounded. The original algorithm may not forget about branches because restrictions which are pushed *upwards* in the tree might make it necessary to revisit paths which have been considered before. A *reset-restart* mechanism solves this problem as follows:
Whenever the $``$\- or the $`_+`$-rule is applied to a node $`x`$ and its *predecessor* $`y`$ (Case 2’ of these rules), we delete all successors of $`y`$ from the completion tree (*reset*). While this makes it necessary to *restart* the generation of successors for $`y`$, it makes it possible to implement the algorithm in a depth-first manner which facilitates the re-use of space.
This modification does not affect the proof of soundness and completeness for the algorithm, but of course we have to re-prove termination \[HST98\] as it formerly relied on the fact that we never removed any nodes from the completion tree. Summing up we get:
###### Theorem 3.9
The modified algorithm is a PSpace decision procedure for satisfiability and subsumption of $`𝒮`$-concepts.
## 4 The Undecidability of Unrestricted $`𝒮𝒩`$
Like earlier DLs that combine a hierarchy of (transitive and non-transitive) roles with some form of number restrictions \[HS99, HST98\], $`𝒮𝒩`$ only allows *simple* roles in restrictions, i.e. roles that are neither transitive nor have transitive subroles. The justification for this limitation has been partly on the grounds of a doubtful semantics (of transitive functional roles) and partly to simplify decision procedures. In this section, we will show that allowing arbitrary roles in $`𝒮𝒩`$ number restrictions leads to undecidability. For convenience, we denote $`𝒮𝒩`$ with arbitrary roles in number restrictions by $`𝒮𝒩^+`$.
The undecidability proof uses a reduction of the domino problem \[Ber66\] adapted from \[BS96\]. This problem asks whether, for a set of domino types, there exists a *tiling* of an $`^2`$ grid such that each point of the grid is covered with exactly one of the domino types, and adjacent dominoes are “compatible” with respect to some predefined criteria.
###### Definition 4.1
A domino system $`𝒟=(D,H,V)`$ consists of a non-empty set of domino types $`D=\{D_1,\mathrm{},D_n\}`$, and of sets of horizontally and vertically matching pairs $`HD\times D`$ and $`VD\times D`$. The problem is to determine if, for a given $`𝒟`$, there exists a *tiling* of an $`\times `$ grid such that each point of the grid is covered with a domino type in $`D`$ and all horizontally and vertically adjacent pairs of domino types are in $`H`$ and $`V`$ respectively, i.e., a mapping $`t:\times D`$ such that for all $`m,n`$, $`t(m,n),t(m+1,n)H`$ and $`t(m,n),t(m,n+1)V`$.
This problem can be reduced to the satisfiability of $`𝒮𝒩^+`$-concepts, and the undecidability of the domino problem implies undecidability of satisfiability of $`𝒮𝒩^+`$-concepts.
Ensuring that each point is associated with exactly one domino type and that a point and its neighbours satisfy the compatibility conditions induced by $`H`$ and $`V`$ is simple for most logics (via the introduction of concepts $`C_{D_i}`$ for domino types $`D_i`$, and the use of value restrictions and boolean connectives), and applying such conditions throughout the grid is also simple in a logic such as $`𝒮𝒩^+`$ which can deal with arbitrary axioms. The crucial difficulty is representing the $`\times `$ grid using “horizontal” and “vertical” roles $`X`$ and $`Y`$, and in particular forcing the coincidence of $`XY`$\- and $`YX`$-successors. This can be accomplished in $`𝒮𝒩^+`$ using an alternating pattern of two horizontal roles $`X_1`$ and $`X_2`$, and two vertical roles $`Y_1`$ and $`Y_2`$, with disjoint primitive concepts $`A`$, $`B`$, $`C`$, and $`D`$ being used to identify points in the grid with different combinations of successors. The coincidence of $`XY`$ and $`YX`$ successors can then be enforced using number restrictions on transitive super-roles of each of the four possible combinations of $`X`$ and $`Y`$ roles. A visualisation of the resulting grid and a suitable role hierarchy is shown in Figure 4, where $`S_{ij}^{}`$ are transitive roles.
The alternation of $`X`$ and $`Y`$ roles in the grid means that one of the transitive super-roles $`S_{ij}`$ connects each point $`(m,n)`$ to the points $`(m+1,n)`$, $`(m,n+1)`$ and $`(m+1,n+1)`$, and to no other points. A number restriction of the form $`3S_{ij}`$ can thus be used to enforce the necessary coincidence of $`XY`$\- and $`YX`$-successors. A complete specification of the grid is given by the following axioms:
$$\begin{array}{ccc}\hfill A& & \neg B\neg C\neg DX_1.BY_1.C3S_{11}\text{,}\hfill \\ \hfill B& & \neg A\neg C\neg DX_2.AY_1.D3S_{21}\text{,}\hfill \\ \hfill C& & \neg A\neg B\neg DX_1.DY_2.A3S_{12}\text{,}\hfill \\ \hfill D& & \neg A\neg B\neg CX_2.CY_2.B3S_{22}\text{.}\hfill \end{array}$$
It only remains to add axioms which encode the local compatibility conditions (as described in \[BS96\]) and to assert that $`A`$, $`B`$, $`C`$, and $`D`$ are subsumed by the disjunction of all domino types to enforce the placement of a tile on each point of the grid. The concept $`A`$ is now satisfiable w.r.t. the various axioms (which can be internalised as described in Lemma 2.3) iff there is a compatible tiling of the grid.
## 5 Discussion
A new DL system is being implemented based on the $`𝒮𝒬`$ algorithm described in Section 3.1. Pending the completion of this project, the existing FaCT system \[Hor98\] has been modified to deal with inverse roles using the $`𝒮𝒬`$ blocking strategy, giving a DL which is equivalent to $`𝒮`$ extended with functional roles \[HS99\]; we will refer to this DL as $`𝒮`$ and to the modified FaCT system as I-FaCT.
I-FaCT has been used to conduct some initial experiments with a terminology representing (fragments of) database schemata and inter schema assertions from a data warehousing application \[CDL<sup>+</sup>98\] (a slightly simplified version of the proposed encoding was used to generate $`𝒮`$ terminologies). I-FaCT is able to classify this terminology, which contains 19 concepts and 42 axioms, in less than 0.1s of (266MHz Pentium) CPU time. In contrast, eliminating inverse roles using an embedding technique \[CDR98\] gives an equisatisfiable FaCT terminology with an additional 84 axioms, but one which FaCT is unable to classify in 12 hours of CPU time.
An extension of the embedding technique can be used to eliminate number restrictions \[DL95\], but requires a target logic which supports the transitive *closure* of roles, i.e., *converse*-PDL. The even larger number of axioms which this embedding would introduce makes it unlikely that tractable reasoning could be performed on the resulting terminology. Moreover, we are not aware of any algorithm for *converse*-PDL which does not employ a so-called *cut rule* \[DM98\], the application of which introduces considerable additional non-determinism. It seems inevitable that this would lead to a further degradation in empirical tractability.
As far as complexity is concerned, we have already been successful in extending the PSpace-result for $`𝒮`$ to $`𝒮𝒩`$ \[HST98\]. Currently we are working on an extension of this result to $`𝒮𝒬`$ combining the techniques from this paper with those presented in \[Tob99\].
## Appendix
In this appendix we present the proof of Lemma 3.4, which is repeated here for easier reference.
Lemma. Let $`D`$ be an $`𝒮𝒬`$-concept.
1. (Termination) The tableaux algorithm terminates when started with $`D`$.
2. (Soundness) If the expansion rules can be applied to $`D`$ such that they yield a complete and clash-free completion tree, then $`D`$ has a tableau.
3. (Completeness) If $`D`$ has a tableau, then the expansion rules can be applied to $`D`$ such that they yield a complete and clash-free completion tree.
#### (Termination)
Let $`m=|\text{clos}(D)|`$, $`k=|𝐑_D|`$, and $`n_{\text{max}}`$ the maximum $`n`$ that occurs in a concept of the form $`(nSC)\text{clos}(D)`$. Termination is a consequence of the following properties of the expansion rules:
* The expansion rules never remove nodes from the tree or concepts from node labels. Edge labels can only be changed by the $``$-rule which either expands them or sets them to $`\mathrm{}`$; in the latter case the node below the $`\mathrm{}`$-labelled edge is blocked and this block is never broken.
* Each successor of a node $`x`$ is the result of the application of the $``$-rule or the $``$-rule to $`x`$. For a node $`x`$, each concept in $`(x)`$ can trigger the generation of successors at most once.
For the $``$-rule, if a successor $`y`$ of $`x`$ was generated for a concept $`S.C(x)`$ and later $`(x,y)`$ is set to $`\mathrm{}`$ by the $``$-rule, then there is some $`S`$-neighbour $`z`$ of $`x`$ with $`C(z)`$.
For the $``$-rule, if $`y_1,\mathrm{},y_n`$ were generated by the $``$-rule for $`(nSC)(x)`$, then $`y_i\doteq ̸y_j`$ holds for all $`1i<jn`$. This implies that there are always $`n`$ $`S`$-neighbours $`y_1^{},\mathrm{},y_n^{}`$ of $`x`$ with $`C(y_i^{})`$ and $`y_i^{}\doteq ̸y_j^{}`$ for all $`1i<jn`$, since the $``$-rule never merges two nodes $`y_i^{},y_j^{}`$ with $`y_i^{}\doteq ̸y_j^{}`$, and, whenever an application of the $``$-rule sets $`(x,y_i^{})`$ to $`\mathrm{}`$, there is some $`S`$-neighbour $`z`$ of $`x`$ which “inherits” both $`C`$ and all inequalities from $`y_i^{}`$.
Since $`\text{clos}(D)`$ contains a total of at most $`m`$ $`R.C`$ and $`(nSC)`$ concepts, the out-degree of the tree is bounded by $`mn_{\text{max}}`$.
* Nodes are labelled with non-empty subsets of $`\text{clos}(D)`$ and edges with subsets of $`R_D`$, so there are at most $`2^{2mk}`$ different possible labellings for a pair of nodes and an edge. Therefore, if a path $`p`$ is of length at least $`2^{2mk}`$, then from the pair-wise blocking condition there must be two nodes $`x,y`$ on $`p`$ such that $`x`$ is directly blocked by $`y`$. Furthermore, if a node was generated at distance $`\mathrm{}`$ from the root node, it always remains at this distance, and thus paths are not curled up or shortened. Since a path on which nodes are blocked cannot become longer, paths are of length at most $`2^{2mn}`$. ∎
#### (Soundness)
Let $`𝐓`$ be a complete and clash-free completion tree. A path is a sequence of pairs of nodes of $`𝐓`$ of the form $`p=[\frac{x_0}{x_0^{}},\mathrm{},\frac{x_n}{x_n^{}}]`$. For such a path we define $`𝖳𝖺𝗂𝗅(p):=x_n`$ and $`𝖳𝖺𝗂𝗅^{}(p):=x_n^{}`$. With $`[p|\frac{x_{n+1}}{x_{n+1}^{}}]`$ we denote the path $`[\frac{x_0}{x_0^{}},\mathrm{},\frac{x_n}{x_n^{}},\frac{x_{n+1}}{x_{n+1}^{}}]`$. The set $`\mathrm{𝖯𝖺𝗍𝗁𝗌}(𝐓)`$ is defined inductively as follows:
* For the root node $`x_0`$ of $`𝐓`$, $`[\frac{x_0}{x_0}]\mathrm{𝖯𝖺𝗍𝗁𝗌}(𝐓)`$, and
* For a path $`p\mathrm{𝖯𝖺𝗍𝗁𝗌}(𝐓)`$ and a node $`z`$ in $`𝐓`$:
+ if $`z`$ is a successor of $`𝖳𝖺𝗂𝗅(p)`$ and $`z`$ is not blocked, then $`[p|\frac{z}{z}]\mathrm{𝖯𝖺𝗍𝗁𝗌}(𝐓)`$, or
+ if, for some node $`y`$ in $`𝐓`$, $`y`$ is a successor of $`𝖳𝖺𝗂𝗅(p)`$ and $`z`$ blocks $`y`$, then $`[p|\frac{z}{y}]\mathrm{𝖯𝖺𝗍𝗁𝗌}(𝐓)`$.
Please note that, due to the construction of $`\mathrm{𝖯𝖺𝗍𝗁𝗌}`$, for $`p\mathrm{𝖯𝖺𝗍𝗁𝗌}(𝐓)`$ with $`p=[p^{}|\frac{x}{x^{}}]`$, we have that $`x`$ is not blocked, $`x^{}`$ is blocked iff $`xx^{}`$, and $`x^{}`$ is never indirectly blocked. Furthermore, $`(x)=(x^{})`$ holds.
Now we can define a tableau $`T=(𝐒,,)`$ with:
$$\begin{array}{ccc}\hfill 𝐒& =& \mathrm{𝖯𝖺𝗍𝗁𝗌}(𝐓)\hfill \\ \hfill (p)& =& (𝖳𝖺𝗂𝗅(p))\hfill \\ \hfill (R)& =& \{p,q𝐒\times 𝐒\begin{array}{cc}\multicolumn{2}{c}{\text{Either }q=[p|\frac{x}{x^{}}]\text{ and}}\\ x^{}\text{ is an }R\text{-successor of }𝖳𝖺𝗂𝗅(p)\hfill & \\ \multicolumn{2}{c}{\text{or }p=[q|\frac{x}{x^{}}]\text{ and}}\\ x^{}\text{ is an }𝖨𝗇𝗏(R)\text{-successor of }𝖳𝖺𝗂𝗅(q)\}.\hfill & \end{array}\hfill \end{array}$$
Claim: $`T`$ is a tableau for $`D`$ with respect to $`^+`$.
We show that $`T`$ satisfies all the properties from Definition 3.1.
* $`D([\frac{x_0}{x_0}])`$ since $`D(x_0)`$.
* Property 1 holds because $`𝐓`$ is clash-free; Properties 2,3 hold because $`𝖳𝖺𝗂𝗅(p)`$ is not blocked and $`𝐓`$ is complete.
* Property 4: Assume $`S.C(p)`$ and $`p,q(S)`$. If $`q=[p|\frac{x}{x^{}}]`$, then $`x^{}`$ is an $`S`$-successor of $`𝖳𝖺𝗂𝗅(p)`$ and thus $`C(x^{})`$ (because the $``$-rule is not applicable). Since $`(q)=(x)=(x^{})`$, we have $`C(q)`$. If $`p=[q|\frac{x}{x^{}}]`$, then $`x^{}`$ is an $`𝖨𝗇𝗏(S)`$-successor of $`𝖳𝖺𝗂𝗅(q)`$ and thus $`C(𝖳𝖺𝗂𝗅(q))`$ (because $`x^{}`$ is not indirectly blocked and the $``$-rule is not applicable), hence $`C(q)`$.
* Property 5: Assume $`S.C(p)`$. Define $`x:=𝖳𝖺𝗂𝗅(p)`$. In $`𝐓`$ there is an $`S`$-neighbour $`y`$ of $`x`$ with $`C(y)`$, because the $``$-rule is not applicable. There are two possibilities:
+ $`y`$ is a successor of $`x`$ in $`𝐓`$. If $`y`$ is not blocked, then $`q:=[p|\frac{y}{y}]𝐒`$ and $`p,q(S)`$ as well as $`C(q)`$. If $`y`$ is blocked by some node $`z`$ in $`𝐓`$, then $`q:=[p|\frac{z}{y}]𝐒`$.
+ $`y`$ is a predecessor of $`x`$. Again, there are two possibilities:
- $`p`$ is of the form $`p=[q|\frac{x}{x^{}}]`$ with $`𝖳𝖺𝗂𝗅(q)=y`$.
- $`p`$ is of the form $`p=[q|\frac{x}{x^{}}]`$ with $`𝖳𝖺𝗂𝗅(q)=uy`$. $`x`$ only has one predecessor in $`𝐓`$, hence $`u`$ is not the predecessor of $`x`$. This implies $`xx^{}`$, $`x`$ blocks $`x^{}`$ in $`𝐓`$, and $`u`$ is the predecessor of $`x^{}`$ due to the construction of $`\mathrm{𝖯𝖺𝗍𝗁𝗌}`$. Together with the definition of the blocking condition, this implies $`(u,x^{})=(y,x)`$ as well as $`(u)=(y)`$ due to the pair-wise blocking condition.
In all three cases, $`p,q(S)`$ and $`C(q)`$.
* Property 6: Assume $`S.C(p)`$, $`p,q(R)`$ for some $`R\text{*}S`$ with $`𝖳𝗋𝖺𝗇𝗌(R)`$. If $`q=[p|\frac{x}{x^{}}]`$, then $`x^{}`$ is an $`R`$-successor of $`𝖳𝖺𝗂𝗅(p)`$ and thus $`R.C(x^{})`$ (because otherwise the $`_+`$-rule would be applicable). From $`(q)=(x)=(x^{})`$ it follows that $`R.C(q)`$. If $`p=[q|\frac{x}{x^{}}]`$, then $`x^{}`$ is an $`𝖨𝗇𝗏(S)`$-successor of $`𝖳𝖺𝗂𝗅(q)`$ and hence $`𝖳𝖺𝗂𝗅(q)`$ is an $`R`$-neighbour of $`x^{}`$. Because $`x^{}`$ is not indirectly blocked, this implies $`R.C(𝖳𝖺𝗂𝗅(q))`$ and hence $`R.C(q)`$.
* Property 11: Assume $`(nSC)(p)`$, $`p,q(S)`$. If $`q=[p|\frac{x}{x^{}}]`$, then $`x^{}`$ is an $`S`$-successor of $`𝖳𝖺𝗂𝗅(p)`$ and thus $`\{C,\mathrm{}C\}(x^{})\mathrm{}`$ (since the choose-rule is not applicable). Since $`(q)=(x)=(x^{})`$, we have $`\{C,\mathrm{}C\}(q)\mathrm{}`$. If $`p=[q|\frac{x}{x^{}}]`$, then $`x^{}`$ is an $`𝖨𝗇𝗏(S)`$-successor of $`𝖳𝖺𝗂𝗅(q)`$ and thus $`\{C,\mathrm{}C\}(𝖳𝖺𝗂𝗅(q))\mathrm{}`$ (since $`x^{}`$ is not indirectly blocked and the choose-rule is not applicable), hence $`\{C,\mathrm{}C\}(q)\mathrm{}`$.
* Assume Property 9 is violated. Hence there is some $`p𝐒`$ with $`(nSC)(p)`$ and $`\mathrm{}S^T(p,C)>n`$. We show that this implies $`\mathrm{}S^𝐓(𝖳𝖺𝗂𝗅(p),C)>n`$, in contradiction of either the clash-freeness or completeness of $`𝐓`$. Define $`x:=𝖳𝖺𝗂𝗅(p)`$ and $`P:=S^T(p,C)`$. Due to the assumption, we have $`\mathrm{}P>n`$. We distinguish two cases:
+ $`P`$ contains only paths of the form $`q=[p|\frac{y}{y^{}}]`$. We claim that the function $`𝖳𝖺𝗂𝗅^{}`$ is injective on $`P`$. Assume that there are two paths $`q_1,q_1P`$ with $`q_1q_2`$ and $`𝖳𝖺𝗂𝗅^{}(q_1)=𝖳𝖺𝗂𝗅^{}(q_2)=y^{}`$. Then $`q_1`$ is of the form $`q_1=[p|(y_1,y^{})]`$ and $`q_2`$ is of the form $`q_2=[p|\frac{y_2}{y^{}}]`$ with $`y_1y_2`$. If $`y^{}`$ is not blocked in $`𝐓`$, then $`y_1=y^{}=y_2`$, contradicting $`y_1y_2`$. If $`y^{}`$ is blocked in $`𝐓`$, then both $`y_1`$ and $`y_2`$ block $`y^{}`$, which implies $`y_1=y_2`$, again a contradiction.
Since $`𝖳𝖺𝗂𝗅^{}`$ is injective on $`P`$, it holds that $`\mathrm{}P=\mathrm{}𝖳𝖺𝗂𝗅^{}(P)`$. Also for each $`y^{}𝖳𝖺𝗂𝗅^{}(P)`$, $`y^{}`$ is an $`S`$-successor of $`x`$ and $`C(y^{})`$. This implies $`\mathrm{}S^𝐓(x,C)>n`$.
+ $`P`$ contains a path $`q`$ where $`p`$ is of the form $`p=[q|\frac{x}{x^{}}]`$. Obviously, $`P`$ may only contain one such path. As in the previous case, $`𝖳𝖺𝗂𝗅^{}`$ is an injective function on the set $`P^{}:=P\{q\}`$, each $`y^{}𝖳𝖺𝗂𝗅^{}(P^{})`$ is an $`S`$-successor of $`x`$ and $`C(y^{})`$ for each $`y^{}𝖳𝖺𝗂𝗅^{}(P^{})`$. To show that indeed $`\mathrm{}S^𝐓(x,C)>n`$ holds, we have to prove the existence of a further $`S`$-neighbour $`u`$ of $`x`$ with $`C(u)`$ and $`u𝖳𝖺𝗂𝗅^{}(P^{})`$. This will be “supplied” by $`z:=𝖳𝖺𝗂𝗅(q)`$. We distinguish two cases:
- $`x=x^{}`$. Hence $`x`$ is not blocked. This implies that $`x`$ is an $`𝖨𝗇𝗏(S)`$-successor of $`z`$ in $`𝐓`$. Since $`𝖳𝖺𝗂𝗅^{}(P^{})`$ contains only successors of $`x`$, we have that $`z𝖳𝖺𝗂𝗅^{}(P^{})`$ and, by construction, $`z`$ is an $`S`$-neighbour of $`x`$ with $`C(z)`$.
- $`xx^{}`$. This implies that $`x^{}`$ is blocked in $`𝐓`$ by $`x`$ and that $`x^{}`$ is an $`𝖨𝗇𝗏(S)`$-successor of $`z`$ in $`𝐓`$. The definition of pairwise-blocking implies that $`x`$ is an $`𝖨𝗇𝗏(S)`$-successor of some node $`u`$ in $`𝐓`$ with $`(u)=(z)`$. Again, since $`𝖳𝖺𝗂𝗅^{}(P^{})`$ contains only successors of $`x`$ we have that $`u𝖳𝖺𝗂𝗅^{}(P^{})`$ and, by construction, $`u`$ is an $`S`$-neighbour of $`x`$ and $`C(u)`$.
* Property 10: Assume $`(nSC)(p)`$. Completeness of $`𝐓`$ implies that there exist $`n`$ individuals $`y_1,\mathrm{},y_n`$ in $`𝐓`$ such that each $`y_i`$ is an $`S`$-neighbour of $`𝖳𝖺𝗂𝗅(p)`$ and $`C(y_i)`$. We claim that, for each of these individuals, there is a path $`q_i`$ such that $`p,q_i(S)`$, $`C(q_i)`$, and $`q_iq_j`$ for all $`1i<jn`$. Obviously, this implies $`\mathrm{}S^T(p,C)n`$. For each $`y_i`$ there are three possibilities:
+ $`y_i`$ is an $`S`$-successor of $`x`$ and $`y_i`$ is not blocked in $`𝐓`$. Then $`q_i=[p|\frac{y_i}{y_i}]`$ is a path with the desired properties.
+ $`y_i`$ is an $`S`$-successor of $`x`$ and $`y_i`$ is blocked in $`𝐓`$ by some node $`z`$. Then $`q_i=[p|\frac{z}{y_i}]`$ is the path with the desired properties. Since the same $`z`$ may block several of the $`y_j`$s, it is indeed necessary to include $`y_i`$ explicitly into the path to make them distinct.
+ $`x`$ is an $`𝖨𝗇𝗏(S)`$-successor of $`y_i`$. There may be at most one such $`y_i`$. This implies that $`p`$ is of the form $`p=[q|\frac{x}{x^{}}]`$ with $`𝖳𝖺𝗂𝗅(q)=y_i`$. Again, $`q`$ has the desired properties and, obviously, $`q`$ is distinct from all other paths $`q_j`$.
* Property 7 is satisfied due to the symmetric definition of $``$. Property 8 is satisfied due to the definition of $`R`$-successor that takes into account the role hierarchy $`\text{*}`$. ∎
#### (Completeness)
Let $`T=(𝐒,,)`$ be a tableau for $`D`$ w.r.t. $`^+`$. We use this tableau to guide the application of the non-deterministic rules. To do this, we will inductively define a function $`\pi `$, mapping the individuals of the tree $`𝐓`$ to $`𝐒`$ such that, for each $`x,y`$ in $`𝐓`$:
$$\begin{array}{c}(x)(\pi (x))\hfill \\ \text{if }y\text{ is an }S\text{-neighbour of }x\text{, then }\pi (x),\pi (y)(S)\hfill \\ x\doteq ̸y\text{ implies }\pi (x)\pi (y)\hfill \end{array}\}()$$
Claim: Let $`𝐓`$ be a completion-tree and $`\pi `$ a function that satisfies $`()`$. If a rule is applicable to $`𝐓`$ then the rule is applicable to $`𝐓`$ in a way that yields a completion-tree $`𝐓^{}`$ and an extension of $`\pi `$ that satisfy $`()`$.
Let $`𝐓`$ be a completion-tree and $`\pi `$ be a function that satisfies $`()`$. We have to consider the various rules.
* The $``$-rule: If $`C_1C_2(x)`$, then $`C_1C_2(\pi (x))`$. This implies $`C_1,C_2(\pi (x))`$ due to Property 2 from Definition 3.1, and hence the rule can be applied without violating $`()`$.
* The $``$-rule: If $`C_1C_2(x)`$, then $`C_1C_2(\pi (x))`$. Since $`T`$ is a tableau, Property 3 from Definition 3.1 implies $`\{C_1,C_2\}(\pi (x))\mathrm{}`$. Hence the $``$-rule can add a concept $`E\{C_1,C_2\}`$ to $`(x)`$ such that $`(x)(\pi (x))`$ holds.
* The $``$-rule: If $`S.C(x)`$, then $`S.C(\pi (x))`$ and, since $`T`$ is a tableau, Property 5 of Definition 3.1 implies that there is an element $`t𝐒`$ such that $`\pi (x),t(S)`$ and $`C(t)`$. The application of the $``$-rule generates a new variable $`y`$ with $`(x,y=\{S\}`$ and $`(y)=\{C\}`$. Hence we set $`\pi :=\pi [yt]`$ which yields a function that satisfies $`()`$ for the modified tree.
* The $``$-rule: If $`S.C(x)`$, then $`S.C(\pi (x))`$, and if $`y`$ is an $`S`$-neighbour of $`x`$, then also $`\pi (x),\pi (y)(S)`$ due to $`()`$. Since $`T`$ is a tableau, Property 4 of Definition 3.1 implies $`C(\pi (y))`$ and hence the $``$-rule can be applied without violating $`()`$.
* The $`_+`$-rule: If $`S.C(x)`$, then $`S.C(\pi (x))`$, and if there is some $`R\text{*}S`$ with $`𝖳𝗋𝖺𝗇𝗌(R)`$ and $`y`$ is an $`R`$-neighbour of $`x`$, then also $`\pi (x),\pi (y)(R)`$ due to $`()`$. Since $`T`$ is a tableau, Property 6 of Definition 3.1 implies $`R.C(\pi (y))`$ and hence the $`_+`$-rule can be applied without violating $`()`$.
* The choose-rule: If $`(nSC)(x)`$, then $`(nSC)(\pi (x))`$, and, if there is an $`S`$-neighbour $`y`$ of $`x`$, then $`\pi (x),\pi (y)(S)`$ due to $`()`$. Since $`T`$ is a tableau, Property 11 of Definition 3.1 implies $`\{C,\mathrm{}C\}(\pi (y)\mathrm{}`$. Hence the choose-rule can add an appropriate concept $`E\{C,\mathrm{}C\}`$ to $`(x)`$ such that $`(y)(\pi (y))`$ holds.
* The $``$-rule: If $`(nSC)(x)`$, then $`(nSC)(\pi (x))`$. Since $`T`$ is a tableau, Property 10 of Definition 3.1 implies $`\mathrm{}S^T(\pi (x),C)n`$. Hence there are individuals $`t_1,\mathrm{},t_n𝐒`$ such that $`\pi (x),t_i(S)`$, $`C(t_i)`$, and $`t_it_j`$ for $`1i<jn`$. The $``$-rule generates $`n`$ new nodes $`y_1,\mathrm{},y_n`$. By setting $`\pi :=\pi [y_1t_1,\mathrm{}y_nt_n]`$, one obtains a function $`\pi `$ that satisfies $`()`$ for the modified tree.
* The $``$-rule: If $`(nSC)(x)`$, then $`(nSC)(\pi (x))`$. Since $`T`$ is a tableau, Property 9 of Definition 3.1 implies $`\mathrm{}S^T(\pi (x),C)n`$. If the $``$-rule is applicable, we have $`\mathrm{}S^𝐓(x,C)>n`$, which implies that there are at least $`n+1`$ $`S`$-neighbours $`y_0,\mathrm{},y_n`$ of $`x`$ such that $`C(y_i)`$. Thus, there must be two nodes $`y,z\{y_0,\mathrm{},y_n\}`$ such that $`\pi (y)=\pi (z)`$ (because otherwise $`\mathrm{}S^T(\pi (x),C)>n`$ would hold). From $`\pi (y)=\pi (z)`$ we have that $`y\doteq ̸z`$ cannot hold because of $`()`$, and $`y,z`$ can be chosen such that $`y`$ is not an ancestor of $`z`$. Hence the $``$-rule can be applied without violating $`()`$.
Why does this claim yield the completeness of the tableaux algorithm? For the initial completion-tree consisting of a single node $`x_0`$ with $`(x_0)=\{D\}`$ and $`\doteq ̸=\mathrm{}`$ we can give a function $`\pi `$ that satisfies $`()`$ by setting $`\pi (x_0):=s_0`$ for some $`s_0𝐒`$ with $`D(s_0)`$ (such an $`s_0`$ exists since $`T`$ is a tableau for $`D`$). Whenever a rule is applicable to $`𝐓`$, it can be applied in a way that maintains $`()`$, and, since the algorithm terminates, we have that any sequence of rule applications must terminate. Properties $`()`$ imply that any tree $`𝐓`$ generated by these rule-applications must be clash-free as there are only two possibilities for a clash, and it is easy to see that neither of these can hold in $`𝐓`$:
* $`𝐓`$ cannot contain a node $`x`$ such that $`\{C,\neg C\}(x)`$ because $`(x)(\pi (x))`$ and hence Property 1 of Definition 3.1 would be violated for $`\pi (x)`$.
* $`𝐓`$ cannot contain a node $`x`$ with $`(nSC)(x)`$ and $`n+1`$ $`S`$-neighbours $`y_0,\mathrm{}y_n`$ of $`x`$ with $`C(y_i)`$ and $`y_i\doteq ̸y_j`$ for $`0i<jn`$ because $`(nSC)(\pi (x))`$, and, since $`y_i\doteq ̸y_j`$ implies $`\pi (y_i)\pi (y_j)`$, $`\mathrm{}S^T(\pi (x),C)>n`$, in contradiction to Property 9 of Definition 3.1. ∎
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# FUSE OBSERVATIONS OF DIFFUSE INTERSTELLAR MOLECULAR HYDROGEN
## 1. INTRODUCTION
As the most abundant molecule in the Universe, molecular hydrogen (H<sub>2</sub>) comprises the bulk of the mass of dense interstellar molecular clouds and is the main ingredient of star formation (Shull & Beckwith 1982). The first steps of star formation assemble dense interstellar clouds from diffuse gas, but the formation mechanism for giant molecular clouds from diffuse diffuse gas is still unknown. Little is known about the distribution of diffuse H<sub>2</sub> in the interstellar medium (ISM) of our Galaxy, and major uncertainties remain about its formation, destruction, and recycling into dense clouds and stars. H<sub>2</sub> also plays a central role in our understanding of interstellar chemistry.
Here, we describe early results on H<sub>2</sub> in the Milky Way and the Magellanic Clouds obtained with the NASA Far Ultraviolet Spectroscopic Explorer (FUSE) satellite. The FUSE mission and the capabilities of its spectrograph are described by Moos et al. (2000) and Sahnow et al. (2000). Over its mission lifetime, FUSE will probe hundreds of sight lines through the Galactic disk and halo. We expect a large fraction of these sight lines to exhibit absorption from the Lyman and Werner rotational-vibrational bands of H<sub>2</sub>, some 400 absorption lines between 912 and 1120 Å, arising from rotational levels $`J`$ = 0 – 7 in the ground vibrational state. With FUSE, our team will map the distribution of H<sub>2</sub> in the diffuse ISM, measure its abundance, state of excitation, and rates of formation and destruction. In § 2 we summarize our initial observations and describe our analysis methods. In § 3 we describe results from 8 targets. In § 4 we summarize our results and discuss future directions of FUSE H<sub>2</sub> research.
## 2. OBSERVATIONS AND ANALYSIS
Our FUSE observations were obtained from 1999 September to 1999 November during the commissioning phase of satellite operations. A summary of the observations appears in Table 1. The resolution of the spectrograph across the band was R $``$ 12,000 for all observations, based on studies of sharp interstellar lines of H<sub>2</sub>, Ar I, and Fe II. The data were prepared for analysis by the FUSE data pipeline as described in Moos et al. (2000). Once the two-dimensional spectra were extracted to one dimension, a wavelength solution was applied and the oversampled data were smoothed by a 15-pixel running boxcar to 30 km s<sup>-1</sup> resolution.
In typical interstellar conditions, excited H<sub>2</sub> quickly decays to the ground vibrational state of the ground electronic state (Black & Dalgarno 1973). Observations of H<sub>2</sub> with the Copernicus satellite occasionally detected rotational lines up to $`J=7`$, and we search for $`J10`$. However, lines above $`J=4`$ are difficult to detect in diffuse interstellar clouds, with our typical $`4\sigma `$ limiting equivalent width of 30–40 mÅ and the corresponding column density limit of $`\mathrm{N}(\mathrm{H}_2)10^{14}`$ cm<sup>-2</sup>. We use all available H<sub>2</sub> lines, except where they are blended with other interstellar absorption or airglow lines. Typically we neglect the Lyman (6-0) band, which is lost in interstellar Ly$`\beta `$, and the Lyman (5-0) band, which is coincident with resonance absorption lines C II $`\lambda 1036.34`$ and C II$`{}_{}{}^{}\lambda 1037.02`$ and with O I $`\lambda 1039.64`$ airglow.
We illustrate our data in Figure 1, which displays the FUSE spectrum of ESO 141-G55, a Seyfert 1 galaxy at Galactic latitude $`b=26.71^{}`$. We detect one H<sub>2</sub> Galactic component with N(H$`{}_{2}{}^{})=`$ $`2\times 10^{19}`$ cm<sup>-2</sup>. This spectrum is typical of our FUSE data both in data quality and the abundance of H<sub>2</sub> lines.
We employ a complement of techniques to analyze H<sub>2</sub> absorption. The final products are column densities in rotational levels, N($`J`$), from which we can infer the gas density, UV radiation field, and H<sub>2</sub> formation and destruction rates. For high column density absorbers with damping wings in the R(0) and R(1) lines, we fit the line profiles to derive N(0) and N(1) and use a curve of growth to derive column densities for $`J2`$. For absorbers with $`\mathrm{N}(\mathrm{H}_2)10^{19}`$ cm<sup>-2</sup>, we measure equivalent widths of all H<sub>2</sub> lines and produce a curve of growth to infer a doppler $`b`$-parameter and N($`J)`$ – see Fig. 2. We produce initial estimates of the column densities and excitation temperatures by fitting all lines simultaneously with parameters $`\mathrm{N}(\mathrm{H}_2)`$, $`b`$, and $`T`$, assuming a Boltzmann distribution at rotational temperature $`T_{01}`$ for $`J=0,1`$ and an excitation temperature $`T_{\mathrm{ex}}`$ for $`J2`$.
The Copernicus H<sub>2</sub> survey (Spitzer & Jenkins 1975; Savage et al. 1977) showed that the molecular fraction, $`f_{\mathrm{H2}}=`$2N(H<sub>2</sub>)/\[N(H I) + 2N(H<sub>2</sub>)\], is correlated with E(B-V) and with total H column density. They found that $`f_{\mathrm{H2}}`$ undergoes a transition from low values ($`0.01`$) to high values ($`>0.01`$) at E(B-V) $`0.08`$. Other Copernicus studies indicate that the populations in $`J`$ = 0, 1 and $`J2`$ cannot always be described by a single excitation temperature (Spitzer, Cochran, & Hirshfeld 1974). $`T_{\mathrm{ex}}`$ is generally higher than $`T_{01}`$ and reflects the fluorescent pumping of the rotational levels by incident UV radiation (Jura 1975a,b). The ratio, N(1)/N(0), is usually set by collisional mixing of the ortho- and para-H<sub>2</sub> states (Shull & Beckwith 1982), but these rotational levels may not yet have reached equilibrium in lower density clouds.
## 3. RESULTS FOR INDIVIDUAL SIGHT LINES
### 3.1. Magellanic Cloud Targets
We report observations of one star in the LMC and three stars in the SMC. In the cases of AV 232 and HD 5980, we detect H<sub>2</sub> in both the Galaxy and the SMC (at +119 km s<sup>-1</sup> and +124 km s<sup>-1</sup>; see Table 2). Our survey contains two Magellanic Cloud targets with no detected H<sub>2</sub> at LMC or SMC velocities. Sk -67 111 (LMC) and Sk 108 (SMC) show Galactic H<sub>2</sub> absorption near $`v_{\mathrm{LSR}}0`$ but have no detected H<sub>2</sub> at the MC velocities. Our limiting equivalent width of 60 mÅ ($`4\sigma `$) in the Sk -67 111 spectrum sets limits N(0) $``$ $`2.2\times 10^{14}`$ cm<sup>-2</sup> and N(1) $``$ $`3.2\times 10^{14}`$ cm<sup>-2</sup> from the (7-0) R(0) and R(1) lines. For Sk 108, the limit of 30 mÅ yields N(0) $``$ $`1.1\times 10^{14}`$ cm<sup>-2</sup> and N(1) $``$ $`1.6\times 10^{14}`$ cm<sup>-2</sup>. The absence of H<sub>2</sub> in these sight lines may be due to low H I columns, to intense UV radiation fields that destroy H<sub>2</sub>, or to reduced rates of H<sub>2</sub> formation on grain surfaces in warm low-metallicity gas. Abundances for such gas are given by Welty et al. (1999).
The SMC stars HD 5980 and AV 232 are separated by 75<sup>′′</sup> on the sky. Thus, these sight lines may probe the same interstellar clouds in the Galaxy and SMC. With similar column densities and rotational temperatures, the Galactic components appear to arise in the same gas. However, the AV 232 SMC component exhibits a high rotational temperature $`T_{01}`$ = $`300\pm 60`$ K and a higher excitation temperature $`T_{\mathrm{ex}}`$= $`520\pm 90`$ K, compared with $`T_{01}`$ = $`98\pm 13`$ K and $`T_{\mathrm{ex}}`$= $`300\pm 60`$ K for the HD 5980 absorber. The level populations likely arise from differing excitation rates in separate clouds, $`20`$ pc apart at the SMC.
The only detections of Magellanic Cloud H<sub>2</sub> absorption prior to FUSE were observations with ORFEUS (Richter et al. 1998; de Boer et al. 1998), at low signal to noise but comparable ($`R10^4`$) resolution. With a longer mission lifetime and higher effective area, FUSE offers significant improvement for H<sub>2</sub> studies. For example, with higher S/N, we do not confirm the claimed ORFEUS detections of H<sub>2</sub> (J = 5, 6, 7) at +160 km s<sup>-1</sup> toward HD 5980 (Richter et al. 1998). Our observations of H<sub>2</sub> column densities in the +120 km s<sup>-1</sup> component differ considerably, owing to a better curve of growth (Fig. 2) with $`b=10`$ km s<sup>-1</sup> compared to $`b=6`$ km s<sup>-1</sup> with ORFEUS. Based on 28 FUSE lines ($`J=03`$) compared to 8 ORFEUS lines, we derive log N(0) = 15.10 (compared to 16.57) and log N(1) = 15.30 (compared to 15.90). Within error bars, both studies agree on $`N(2)`$ and $`N(3)`$. Our data yield $`T_{01}=98\pm 13`$ K for the SMC gas.
### 3.2. Extragalactic Targets
Our four extragalactic targets at Galactic latitude $`|b|25^{}`$ probe a large volume of the Galaxy and its halo and allow us to measure the abundance and physical properties of molecular gas outside the disk for the first time. The Copernicus disk result, $`T_{01}=77\pm 17`$ K (Savage et al. 1977), was drawn from 61 stars with $`|b|65^{}`$. From seven FUSE sightlines, we derive $``$$`T_{01}`$$``$ = 107 $`\pm `$ 17 K. In the three optically thick extragalactic sightlines, $``$$`T_{01}`$$``$ = 120 $`\pm `$ 13 K. These results may indicate a higher level of rotational excitation in the upper disk and low halo, which could arise from more intense UV pumping or photoelectric heating from grains in the infrared cirrus. However, these interpretations are still speculative, since the the $`J=0,1`$ levels may not be in equilibrium with the kinetic temperature as a result of UV radiative pumping.
In general, the H<sub>2</sub> rotational populations can be modeled (Jura 1975a,b) to provide physical diagnostics of H<sub>2</sub>-bearing clouds, such as $`n_H`$, $`T_{01}`$, thermal pressure, and UV radiation field. We also derive internally consistent values of N(H<sub>2</sub>) and $`T_{\mathrm{ex}}`$ for these clouds, which could indicate the presence of warm H<sub>2</sub>-bearing gas in the lower Galactic halo. For the three AGN sightlines with detectable H<sub>2</sub>, we find $`n_H`$ = 50, 50, and 36 cm<sup>-3</sup> for Mrk 876, ESO 141-G55, and PG 0804+761 respectively, which yields cloud thicknesses of 1.9, 2.4, and 3.3 pc. The line of sight to ESO 141-G55 intersects a region of enhanced IRAS 100 $`\mu `$m emission (Sembach, Savage, & Hurwitz 1999). The other extragalactic sightlines lie behind high-velocity cloud Complex C (Mrk 876) and weak infrared cirrus (PG 0804+761). Our H<sub>2</sub> results suggest that these clouds are also compressed.
## 4. SUMMARY AND FUTURE DIRECTIONS
The major results of these early FUSE observations are: (1) the ubiquity of H<sub>2</sub> in nearly all sightlines; (2) the generally warmer rotational temperatures, $`T_{01}`$ toward high-latitude gas; and (3) the presence of one low-H<sub>2</sub> extragalactic sightline (PKS 2155-304). We have not detected H<sub>2</sub> in the high velocity clouds (HVCs) toward Mrk 876 or PKS 2155-304, although we intend to continue our HVC searches toward other targets. The HD molecule was not detected in any diffuse cloud, but results from a complementary program (Snow et al. 2000; Ferlet et al. 2000) detect H<sub>2</sub> and HD toward the reddened star HD 73882. Detections of H<sub>2</sub> in other FUSE sight lines are analyzed by Friedman et al. (2000), Mallouris et al. (2000), and Oegerle et al. (2000).
The lack of detectable H<sub>2</sub> toward PKS 2155-304 is quite interesting, in contrast to the general presence of H<sub>2</sub> at high Galactic latitudes. Extragalactic sight lines with low H<sub>2</sub> are valuable for studies of QSO absorption lines (Shull et al. 2000; Sembach et al. 2000), since H<sub>2</sub> line blanketing can add confusion to line identifications (Fig. 1). The observed low N(H<sub>2</sub>) may reflect the intrinsic patchiness of the ISM, but can also be influenced by H<sub>2</sub> formation/destruction in regions of high UV radiation fields and low hydrogen column density. PKS 2155-304 has the lowest column density, log N(H I) $`20.15`$, of any of our targets. The Sk 108 sight line shows no H<sub>2</sub> at SMC velocities, and the Sk -67 111 sight line shows no H<sub>2</sub> at LMC velocities. Although H<sub>2</sub> formation may not reach equilibrium with UV destruction in these diffuse clouds, the H<sub>2</sub> formation rate could depend on the metallicity through the content of dust grains, whose surfaces catalyze H<sub>2</sub> formation. In the Galactic halo, some clouds may be exposed to strong UV dissociating radiation from OB associations in the disk. With more sight lines, we should be able to correlate regions lacking H<sub>2</sub> with regions of low N(H I) and low metallicity (LMC, SMC). Close sight lines such as those to HD 5980 and AV 232 present an opportunity to probe the size of diffuse H<sub>2</sub>-bearing interstellar clouds.
The primary product of future FUSE diffuse H<sub>2</sub> studies will be a map of the distribution of diffuse clouds in the disk and halo of the Milky Way, using a large sample of sight lines toward Galactic and extragalactic targets. Multiple sight lines to the Magellanic Clouds will permit the study of CO/H<sub>2</sub> and possibly HD/H<sub>2</sub> in lower-metallicity environs, including high-velocity clouds in the Galactic halo and Magellanic Stream (Gibson et al. 2000). Observations of multiply-intersected absorbers and H<sub>2</sub> associated with infrared cirrus and high velocity clouds can constrain the size of these diffuse interstellar clouds.
This work is based on data obtained for the Guaranteed Time Team by the NASA-CNES-CSA FUSE mission operated by the Johns Hopkins University. Financial support to U.S. participants has been provided by NASA contract NAS5-32985. J.M.S. also acknowledges the support of NASA astrophysical theory grant NAG5-4063.
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# Optical geometry for gravitational collapse and Hawking radiation
## I Introduction
A conformally static spacetime<sup>*</sup><sup>*</sup>*We adopt the notations and conventions of Ref. . $`(,g_{ab})`$ admits a privileged congruence of timelike curves, corresponding to the flow lines of conformal Killing time $`t`$. Consequently, one can define a family of privileged observers with four-velocity $`n^a=\eta ^a/(\eta _b\eta ^b)^{1/2}`$, where $`\eta ^a`$ is the conformal Killing vector field. The set of these observers can be thought of as a generalization of the Newtonian concept of a rest frame. Their acceleration can be expressed as the projection of a gradient,
$$n^b_bn_a=h_{a}^{}{}_{}{}^{b}_b\mathrm{\Phi }$$
(1)
(see Appendix A for a proof), where $`h_{a}^{}{}_{}{}^{b}=\delta _{a}^{}{}_{}{}^{b}+n_an^b`$ and
$$\mathrm{\Phi }=\frac{1}{2}\mathrm{ln}\left(\eta _a\eta ^a\right);$$
(2)
thus, $`\mathrm{\Phi }`$ is a suitable general-relativistic counterpart of the gravitational potential . One can define the ultrastatic metric $`\stackrel{~}{g}_{ab}=(\eta _c\eta ^c)^1g_{ab}=\mathrm{e}^{2\mathrm{\Phi }}g_{ab}`$, which can be written as $`\stackrel{~}{g}_{ab}=_at_bt+\stackrel{~}{h}_{ab}`$, where $`\stackrel{~}{h}_{ab}=\mathrm{e}^{2\mathrm{\Phi }}h_{ab}`$. The hypersurfaces $`t=\text{const}`$ of $``$ are all diffeomorphic to some three-dimensional manifold $`𝒮`$. If the spacetime is static, it follows from Fermat’s principle that light rays coincide with the geodesics on $`𝒮`$ according to $`\stackrel{~}{h}_{ab}`$ . For this reason, $`\stackrel{~}{g}_{ab}`$ is called the optical metric , and $`(𝒮,\stackrel{~}{h}_{ab})`$ the optical space. We shall also refer to the family of preferred observers $`n^a`$ as the optical frame.
There is a simple operational definition of the optical metric. Suppose that all the observers $`n^a`$ agree to construct a set of synchronized devices that measure the Killing time $`t`$. (Of course, these “clocks” will not agree with those based on local physical processes — e.g., on atomic transitions — but this is totally irrelevant for the following argument.) Then, they use light signals according to a radar procedure, and define the distance between two points $`P,Q𝒮`$ as $`t_{PQP}/2`$, where $`t_{PQP}=t_{QPQ}`$ is the lapse of Killing time corresponding to the round trip of the signal between the observers based at $`P`$ and $`Q`$.There is a one-to-one correspondence between conformally static observers and points of $`𝒮`$. In this way, they attribute the metric $`\stackrel{~}{h}_{ab}`$ to $`𝒮`$.
The notion of optical geometry has recently received considerable attention as a powerful tool in general relativity . It is thus important to investigate to which extent it can be generalized to spacetimes that are not conformally static. One proposal in this direction appears mainly formal, and is probably not sufficient in order to determine $`n^a`$ and $`\stackrel{~}{g}_{ab}`$ in a unique way for an arbitrary spacetime . It is perhaps more helpful to focus on specific situations, that may provide one with additional, physically motivated, hints. In the present article we study the case of a spacetime that describes the gravitational collapse of a spherically symmetric configuration of matter. This problem is interesting for two reasons. First, it represents one of the simplest cases in which the property of conformal staticity does not hold. This is particularly evident if we consider a situation in which the collapsing matter is concentrated on an infinitely thin shell. In this case, the spacetime is composed of two regions, corresponding to the interior and the exterior of the shell, joined through a timelike hypersurface which represents the history of the shell. Both these regions are static when considered separately, their metrics being the Minkowski and the Schwarzschild ones, respectively. However, the fields $`n^a`$ associated with these two metrics do not match in a satisfactory way across the surface of the shell (see Fig. 1). In particular, the horizon is a singular locus for the Schwarzschild frame, but is perfectly regular for the Minkowski observers. This very different behaviour prevents one from considering a single frame that reduces to the Schwarzschild and the Minkowski one, respectively, outside and inside the shell. The failure can be seen as a consequence of the fact that the spacetime is not conformally static in any region containing the shell.See Ref. for considerations related to this point. Indeed, independently of its specific properties, the shell represents a non-stationary boundary between two static regions. A second motivation for studying this class of spacetimes is that they lead to the Hawking effect . Given the success of optical geometry in discussing complicated physical phenomena, one expects that it might give new insight about the process of black hole evaporation. Indeed, this appears to be the case, as we shall see.
The structure of the paper is the following. In the next section we present a general construction of the optical geometry for an arbitrary matter configuration undergoing spherical collapse. Section III is devoted to the analysis of some features that become universal (i.e., model-independent) at late times. In Sec. IV we consider a very simple particular case, the Oppenheimer-Snyder dust model. In Sec. V we argue that the optical geometry picture of collapse is the natural framework for derivations of the Hawking effect based on the “moving mirror analogy,” and that it gives useful insight about the issue of black hole evaporation and the information paradox. Section VI contains a summary of the results, together with some final comments and outlines for future investigations.
## II General construction
The metric of any spherically symmetric spacetime can be written as
$$g=\alpha (t,r)\mathrm{d}t^2+\beta (t,r)\mathrm{d}r^2+r^2\left(\mathrm{d}\theta ^2+\mathrm{sin}^2\theta \mathrm{d}\phi ^2\right),$$
(3)
with $`\alpha `$ and $`\beta `$ positive functions (see, e.g., Ref. , pp. 616–617). In the following we consider situations where matter is confined to a region $`rR(t)`$, with $`R(t)`$ a known function (the “radius of the star”). For $`r>R(t)`$, we assume that the spacetime is empty. However, the treatment can be easily extended to include more general types of collapse — e.g., of electrically charged configurations . According to Birkhoff’s theorem, the metric in the external region is the Schwarzschild one, thus we have $`\alpha (t,r)=\beta (t,r)^1=C(r):=12M/r`$ for $`r>R(t)`$.
In this case, the “rest frame” $`n^a`$ outside the star is just made of the Schwarzschild static observers, $`n^\mu =C(r)^{1/2}\delta _t^\mu `$, and the optical geometry is $`\stackrel{~}{g}_{ab}=C(r)^1g_{ab}`$. Introducing the Regge-Wheeler “tortoise” coordinate $`x`$, such that $`\mathrm{d}x:=C(r)^1\mathrm{d}r`$, we have
$$\stackrel{~}{g}=\mathrm{d}t^2+\mathrm{d}x^2+\stackrel{~}{r}(x)^2\left(\mathrm{d}\theta ^2+\mathrm{sin}^2\theta \mathrm{d}\phi ^2\right),$$
(4)
where $`\stackrel{~}{r}:=C(r)^{1/2}r`$. The Regge-Wheeler coordinate has therefore a very simple geometrical meaning in the optical space: It expresses directly the value of radial distances on $`(𝒮,\stackrel{~}{h}_{ab})`$. Notice that, as far as purely radial motions are concerned, the optical metric (4) gives the same line element as Minkowski spacetime. In particular, no event horizon is present, because the conformal transformation from $`g_{ab}`$ to $`\stackrel{~}{g}_{ab}`$ “sends” the Schwarzschild horizon $`r=2M`$ to infinity. In fact, for a spacetime with metric $`\stackrel{~}{g}_{ab}`$ the points with $`r=2M`$ belong to the null infinity, and the conformal rescaling that carries $`g_{ab}`$ into $`\stackrel{~}{g}_{ab}`$ can be compared to the “decompactification” of a Penrose-Carter diagram, as it is evident from Figs. 24.
To define “natural” observers inside the star is not so easy. In general, the metric in the internal region is not conformally static, and one cannot thus apply the construction based on the timelike conformal Killing vector field, outlined at the beginning of Sec. I. However, even when such a field exists it does not necessarily produce a satisfactory family of internal observers. This point can be clarified by considering again the example of a collapsing shell of matter. Inside the shell the spacetime is flat, by Birkhoff’s theorem; therefore, it would seem obvious to choose inertial observers at fixed distances with respect to the centre of the shell, in order to define a “rest frame.” But such observers are not the natural continuation inside the shell of the Schwarzschild static ones, defined outside. This can be seen by noticing that the horizon $`^+`$ is infinitely ahead in the future for the Schwarzschild observers, but not for those at rest with respect to the centre of the shell. Similarly, the Schwarzschild observers “crowd” near $`^+`$, unlike the internal ones (see Fig. 1). Thus, using the Schwarzschild and the inertial frames would lead to ill-behaved optical metric $`\stackrel{~}{h}_{ab}`$ and potential $`\mathrm{\Phi }`$.
Before constructing explicitly an extension of the Schwarzschild frame that does not suffer from these problems, let us present a graphical discussion of some of its properties. Basically, we are looking for a continuation of the coordinates $`t`$ and $`x`$ inside matter, such that $`\mathrm{d}t/\mathrm{d}x=\pm 1`$ for light signals and fundamental observers are located at $`x=\text{const}`$. In a $`(t,x)`$ diagram, the surface of the star is represented by a line like b in Fig. 4, so that we have still only to establish how the centre $`r=0`$ looks like. To this end, it is convenient to consider the Kruskal diagram in Fig. 5, which shows three incoming radial light rays (or spherical wavefronts). The ray at $`v=v_0`$ simply passes through the centre of the star and is then converted into an outgoing signal, $`u=u_0`$. The ray $`v=v_H`$ reaches $`r=0`$ just on the horizon and then turns into a null generator of $`^+`$. For $`v>v_H`$, all incoming signals enter the black hole region; in particular, the ray $`v=v_P`$ does so exactly when the surface of the star crosses the horizon (event $`P`$ in Fig. 5). Since light signals are still represented by straight lines at $`\pm 45^{}`$ in the $`(t,x)`$ diagram, and $`v_0`$, $`v_H`$, $`v_P`$ all have finite values, it follows that the centre $`r=0`$ must correspond to a line that becomes asymptotically parallel to the one representing the surface of the star, as shown in Fig. 6. Since fundamental observers are represented by vertical lines $`x=\text{const}`$ in the $`(t,x)`$ plane, it is easy to see that their qualitative behaviour in a Kruskal diagram is the one shown in Fig. 7. Their worldlines now accumulate along the whole $`^+`$ and match regularly across the surface of the star. These conditions guarantee that when the metric $`\stackrel{~}{g}_{ab}`$ is used, $`^+`$ becomes a regular portion of the future asymptotic null infinity.
Let us now proceed to construct $`n^a`$ analytically. It is convenient to introduce new coordinates $`(\tau ,\xi )`$ in the internal region $`r<R(t)`$, such that
$$\alpha (t,r)\mathrm{d}t^2+\beta (t,r)\mathrm{d}r^2=\gamma (\tau ,\xi )\left(\mathrm{d}\tau ^2+\mathrm{d}\xi ^2\right),$$
(5)
with $`\gamma `$ a positive function. (Such coordinates always exist, because all two-dimensional spacetimes are conformally flat.) Of course, $`t`$ and $`r`$ become now functions $`t(\tau ,\xi )`$ and $`r(\tau ,\xi )`$ of the new coordinates. Thus, the internal metric reads
$$g=\gamma (\tau ,\xi )\left(\mathrm{d}\tau ^2+\mathrm{d}\xi ^2\right)+r(\tau ,\xi )^2\left(\mathrm{d}\theta ^2+\mathrm{sin}^2\theta \mathrm{d}\phi ^2\right).$$
(6)
Since both $`t`$ and $`\tau `$ are timelike coordinates, the history of a point with $`\theta =\text{const}`$, $`\phi =\text{const}`$ on the surface of the star consists of a sequence of events ordered in $`\tau `$. Therefore, in terms of the internal coordinates, the equation of the surface can be written as $`\xi =\mathrm{\Xi }(\tau )`$, obtained by solving the equation $`r(\tau ,\xi )=R(t(\tau ,\xi ))`$ with respect to $`\xi `$.
The spacetime metric must be continuous across the surface of the star. Thus, the external metric (3) and the internal one, given by (6), must agree in the evaluation of the spacetime interval between two events that occur on the star’s surface. Let us consider two such events, labeled, in the internal coordinates, by $`(t,R(t),\theta ,\phi )`$ and $`(t+\delta t,R(t+\delta t),\theta ,\phi )(t+\delta t,R(t)+R^{}(t)\delta t,\theta ,\phi )`$, where a prime denotes the derivative of a function with respect to its argument. Similarly, in external coordinates we have, for the same events, $`(\tau ,\mathrm{\Xi }(\tau ),\theta ,\phi )`$ and $`(\tau +\delta \tau ,\mathrm{\Xi }(\tau +\delta \tau ),\theta ,\phi )(\tau +\delta \tau ,\mathrm{\Xi }(\tau )+\mathrm{\Xi }^{}(\tau )\delta \tau ,\theta ,\phi )`$. Replacing in Eqs. (3) and (6), and equating the outcomes by continuity, we obtain a differential relation between $`t`$ and $`\tau `$ at the surface of the star,
$`\gamma (\tau ,\mathrm{\Xi }(\tau ))`$ $`\left[1\mathrm{\Xi }^{}(\tau )^2\right]\mathrm{d}\tau ^2`$ (8)
$`=\left[C(R(t))R^{}(t)^2/C(R(t))\right]\mathrm{d}t^2.`$
Integrating Eq. (8) gives a relationship $`\tau =f(t)`$ between the values of $`t`$ and $`\tau `$ at the surface.
The form (6) of the internal metric is convenient because it allows one to readily define null coordinates $`(U,V)`$,
$$U=\tau \xi +U_0,$$
(9)
$$V=\tau +\xi +V_0,$$
(10)
where $`U_0`$ and $`V_0`$ are arbitrary constants. The coordinates $`U`$ and $`V`$ have the usual physical meaning: The locus $`U=\text{const}`$ in spacetime is the history of an outgoing spherical wavefront of light, while $`V=\text{const}`$ represents an incoming one. If we introduce null coordinates $`(u,v)`$ in the outside region as
$$u=tx,$$
(11)
$$v=t+x,$$
(12)
we have that an outgoing spherical wavefront is described, inside the star, by the equation $`U=\text{const}`$ and, outside the star, by $`u=\text{const}`$. Therefore, one can establish a one-to-one correspondence $`U(u)`$ between the values of $`U`$ and $`u`$, defining $`U(u)`$ as the internal $`U`$-label of the wavefront which, outside, is labeled by $`u`$. Similarly, one can define a function $`V(v)`$. The explicit form of $`U(u)`$ can be obtained by solving with respect to $`t`$ the equation
$$tx(R(t))=u$$
(13)
and then substituting the result into
$$U=f(t)\mathrm{\Xi }(f(t))+U_0.$$
(14)
Analogously, $`V(v)`$ is obtained by replacing the solution of
$$t+x(R(t))=v$$
(15)
into
$$V=f(t)+\mathrm{\Xi }(f(t))+V_0.$$
(16)
The functions $`U(u)`$ and $`V(v)`$ can be used to extend the coordinates $`(u,v)`$, hence $`(t,x)`$, also inside the star. It is sufficient to invert them, getting $`u`$ and $`v`$ as functions of $`U`$ and $`V`$, respectively, and then define $`t`$ and $`x`$ as
$$t:=\frac{1}{2}\left(v(V)+u(U)\right),$$
(17)
$$x:=\frac{1}{2}\left(v(V)u(U)\right).$$
(18)
In terms of $`t`$ and $`x`$, the internal metric (6) takes the form
$$g=\gamma U^{}V^{}\left(\mathrm{d}t^2+\mathrm{d}x^2\right)+r^2\left(\mathrm{d}\theta ^2+\mathrm{sin}^2\theta \mathrm{d}\phi ^2\right),$$
(19)
where $`U^{}:=\mathrm{d}U/\mathrm{d}u`$, $`V^{}:=\mathrm{d}V/\mathrm{d}v`$, and all the functions are implicitly supposed to be expressed in terms of $`t`$ and $`x`$. Now we choose
$$\mathrm{\Phi }=\frac{1}{2}\mathrm{ln}\left(\gamma U^{}V^{}\right),$$
(20)
so that the internal optical metric reads
$$\stackrel{~}{g}=\mathrm{d}t^2+\mathrm{d}x^2+\frac{r^2}{\gamma U^{}V^{}}\left(\mathrm{d}\theta ^2+\mathrm{sin}^2\theta \mathrm{d}\phi ^2\right).$$
(21)
In $`(t,r)`$ coordinates, the optical frame $`n^a`$ has components $`n^\mu =\left(\gamma U^{}V^{}\right)^{1/2}\delta _t^\mu `$, i.e., in $`(\tau ,\xi )`$ coordinates,
$$n^\mu =\frac{1}{2(\gamma U^{}V^{})^{1/2}}\left[\left(V^{}+U^{}\right)\delta _\tau ^\mu +\left(V^{}U^{}\right)\delta _\xi ^\mu \right].$$
(22)
This vector field is a satisfactory extension inside the star of the static Schwarzschild frame, and it is not difficult to check that, although the metric (19) is not conformally static, Eq. (1) is still satisfied. Thus, we can continue to interpret $`\mathrm{\Phi }`$ as the gravitational potential.
The factor $`U^{}V^{}`$ can be computed as follows. Let us consider two outgoing light rays, corresponding to the values $`U`$ and $`U+\delta U`$, $`u`$ and $`u+\delta u`$ of the coordinates, with $`\delta U`$ and $`\delta u`$ very small. Equations (13) and (14) give us the coordinate times $`t`$ and $`t+\delta t`$ at which these rays cross the surface of the star, expressed as functions of $`u`$, $`u+\delta u`$, $`U`$, and $`U+\delta U`$. This allows us to find the coefficients that link $`\delta u`$ and $`\delta U`$ to $`\delta t`$; eliminating $`\delta t`$ and taking the limit, we get
$$\frac{\mathrm{d}U}{\mathrm{d}u}=\frac{1\mathrm{\Xi }^{}(f(t(u)))}{1R^{}(t(u))/C(R(t(u)))}f^{}(t(u)),$$
(23)
where the function $`t(u)`$ is implicitly defined by Eq. (13). Similarly, one obtains from Eqs. (15) and (16), considering two incoming light rays,
$$\frac{\mathrm{d}V}{\mathrm{d}v}=\frac{1+\mathrm{\Xi }^{}(f(t(v)))}{1+R^{}(t(v))/C(R(t(v)))}f^{}(t(v)),$$
(24)
where now $`t(v)`$ is given by Eq. (15). Using Eq. (8) we get at the end:
$`U^{}(u)V^{}(v)`$ $`=[{\displaystyle \frac{1\mathrm{\Xi }^{}(f(t(u)))}{1+\mathrm{\Xi }^{}(f(t(u)))}}{\displaystyle \frac{1+\mathrm{\Xi }^{}(f(t(v)))}{1\mathrm{\Xi }^{}(f(t(v)))}}`$ (29)
$`\times {\displaystyle \frac{1+R^{}(t(u))/C(R(t(u)))}{1R^{}(t(u))/C(R(t(u)))}}`$
$`\times {\displaystyle \frac{1R^{}(t(v))/C(R(t(v)))}{1+R^{}(t(v))/C(R(t(v)))}}`$
$`\times {\displaystyle \frac{C(R(t(u)))}{\gamma (f(t(u)),\mathrm{\Xi }(f(t(u))))}}`$
$`\times {\displaystyle \frac{C(R(t(v)))}{\gamma (f(t(v)),\mathrm{\Xi }(f(t(v))))}}]^{1/2}.`$
## III Asymptotic behaviour
Equations (20)–(29) give a complete characterization of the optical geometry inside a collapsing spherically symmetric star. Of course, since they require an explicit knowledge of the functions $`R(t)`$, $`\mathrm{\Xi }(\tau )`$, $`f(t)`$, and $`\gamma (\tau ,\xi )`$, one can use them only within a specific model of collapse — and even in that case their integration will usually require numerical methods. It is therefore remarkable that, at late times (i.e., for $`t,u+\mathrm{}`$), one could establish analytically some features that are universal, in the sense that they do not depend on the details of the model.
In the limit $`u+\mathrm{}`$, it is easy to circumvent the nasty expression on the right hand side of Eq. (23) by noticing that, although the relationship between the functions $`U`$ and $`u`$ is singular on $`^+`$, $`U`$ is regularly connected to the Kruskal retarded null coordinate
$$𝒰=\mathrm{exp}(u/4M)$$
(30)
for all values of $`𝒰`$; this follows from the fact that both $`(U,V)`$ and the Kruskal coordinates are regular at the surface of the star. Then, two outgoing light rays, labeled by $`U`$ and $`U+\delta U`$ inside the star, are labeled by $`𝒰`$ and $`𝒰+\delta 𝒰`$ outside, with $`\delta U=a(𝒰)\delta 𝒰`$, where $`a`$ is a regular positive function that depends on the details of collapse (i.e., $`R`$, $`R^{}`$, $`\mathrm{\Xi }`$, and $`\mathrm{\Xi }^{}`$) at the moment when the light rays cross the surface of the star. For rays near the horizon, i.e., in the limit $`u+\mathrm{}`$, one has $`𝒰0`$ and thus $`\delta Ua(0)\delta 𝒰`$, which gives the desired asymptotic relation
$$\frac{\mathrm{d}U}{\mathrm{d}u}\frac{a(0)}{4M}\mathrm{exp}(u/4M).$$
(31)
(Hereafter, we use the notation $`f_1f_2`$ to express the fact that two functions have the same asymptotic expression in some limit, i.e., $`limf_1/f_2=1`$.) The constant positive factor $`a(0)`$ is all that remains of the details of collapse in the limit $`u+\mathrm{}`$.
The situation is rather different as far as $`\mathrm{d}V/\mathrm{d}v`$ is concerned. Of course, considering two incoming light rays labeled by $`V`$ and $`V+\delta V`$ inside the star, and by $`𝒱`$ and $`𝒱+\delta 𝒱`$ outside, where
$$𝒱=\mathrm{exp}(v/4M)$$
(32)
is the Kruskal advanced null coordinate, one can still claim that $`\delta V=b(𝒱)\delta 𝒱`$, where $`b`$ is a regular positive function depending on the dynamics of the star’s surface when it is crossed by the light rays. However, since the interior of the star at $`^+`$ corresponds to an entire range of values for $`V`$ and $`v`$ (the interval $`[v_H,v_P]`$ in Fig. 5), the function $`\mathrm{d}V/\mathrm{d}v`$, although regular everywhere, has not a universal dependence on $`v`$.
The asymptotic form (31) of $`U^{}(u)`$ and the regular dependence of $`V^{}`$ on $`v`$ are nevertheless sufficient in order to establish the main properties of optical geometry during the late stages of collapse. Since $`\gamma `$, $`U^{}`$, and $`V^{}`$ are regular positive functions, the product $`\gamma U^{}V^{}`$ can simply be written as
$$\gamma U^{}V^{}F(v)\mathrm{exp}(u/4M),$$
(33)
where $`F(v)`$ is a non-vanishing positive function which depends on the details of collapse. Of course, outside the star, $`U𝒰`$, $`V𝒱`$, and $`\gamma =(32M^3/r)\mathrm{exp}(r/2M)`$, so $`\gamma U^{}V^{}=12M/r`$ and $`F(v)\mathrm{exp}(1+v/4M)`$.
From Eq. (22) it is evident that, since $`U^{}0`$ when $`u+\mathrm{}`$, the optical frame behaves like the one in Fig. 7, because the components $`n^\tau `$ and $`n^\xi `$ tend to become equal near $`^+`$. The optical metric (21) and the potential $`\mathrm{\Phi }`$ have the following asymptotic forms:
$$\stackrel{~}{g}\mathrm{d}t^2+\mathrm{d}x^2+\frac{4M^2\mathrm{e}^{(tx)/4M}}{F(t+x)}\left(\mathrm{d}\theta ^2+\mathrm{sin}^2\theta \mathrm{d}\phi ^2\right);$$
(34)
$$\mathrm{\Phi }\frac{xt}{8M}+\frac{1}{2}\mathrm{ln}F(t+x).$$
(35)
Using Eq. (1), one can easily compute the acceleration of the fundamental observers. The only nonvanishing component is
$$n^a_an_x=\frac{\mathrm{\Phi }}{x}=\frac{1}{8M}+\frac{F^{}(t+x)}{2F(t+x)},$$
(36)
which gives, up to the sign, the value of the “gravitational field” at late times. It is interesting to notice that, outside the star, $`n^a_an_x=1/4M`$, which coincides with the surface gravity of the black hole (see Appendix B for a general proof). On the other hand, inside the star, the term $`F^{}/2F`$ gives a correction to the surface gravity that varies from place to place on the horizon and depends on the model of collapse.
Let us now consider a timelike hypersurface with equation $`\xi =F(\tau )`$, such that near $`^+`$ one can write $`\xi \nu \tau +\text{const}`$, where $`\nu `$ is a constant (from the expression (6) of the metric, it follows that $`1<\nu <1`$). This is the case, for example, of the centre of the star, $`\xi 0`$, or of the star’s surface, for which we have $`\nu =\mathrm{d}\mathrm{\Xi }/\mathrm{d}\tau `$ evaluated at the horizon. Equations (9) and (10) give then $`(1+\nu )\mathrm{d}U(1\nu )\mathrm{d}V`$. Near $`^+`$, the coordinate $`v`$ is approximately constant on the submanifold identified by $`\xi =F(\tau )`$ (for example, in the cases of the centre and of the surface of the star, it is equal to $`v_H`$ and $`v_P`$, in the notations of Fig. 5), and this relation can be rewritten as
$$\mathrm{d}vA\mathrm{exp}(u/4M)\mathrm{d}u,$$
(37)
where $`A`$ is a cumulative positive constant and we have used Eq. (31). Integrating and using Eq. (12), we get
$$t+x\overline{v}K\mathrm{exp}(u/4M),$$
(38)
where $`K:=4MA>0`$ and the integration constant $`\overline{v}`$ is the advanced time at which the surface $`\xi =F(\tau )`$ crosses $`^+`$. Equation (38) can be rewritten using Eq. (11) as
$$xt+\overline{v}K\mathrm{exp}(t/2M)$$
(39)
(see also Ref. , p. 869, and Ref. ). This equation expresses the asymptotic behaviour, in $`(t,x)`$ coordinates, of any worldline that crosses $`^+`$. Of course, it is valid both inside and outside the star.
## IV Example: The Oppenheimer-Snyder model
In order to apply the general techniques developed so far to a specific case, let us consider the simplest model of a collapsing star, in which matter is a ball of dust with uniform density . In this case, the internal solution is part of a spatially closed Friedman spacetime (see Ref. , pp. 851–856), so we have $`\gamma (\tau ,\xi )=a(\tau )^2/a_0^2`$ and
$$r(\tau ,\xi )=a(\tau )\mathrm{sin}(\xi /a_0),$$
(40)
where $`0\tau <\pi a_0`$ ($`\tau =0`$ corresponds to the beginning of collapse), $`a_0`$ is a constant, $`0\xi \mathrm{\Xi }_0`$, with $`\mathrm{\Xi }_0<\pi a_0/2`$ corresponding to the surface of the star, and
$$a(\tau )=a_0\mathrm{cos}^2(\tau /2a_0).$$
(41)
It is convenient to introduce the dimensionless variable $`\sigma :=\tau /2a_0`$. Then, the function $`R(t)`$ is defined implicitly by the two equations
$$R=R_i\mathrm{cos}^2\sigma $$
(42)
and
$`t`$ $`=`$ $`2M\mathrm{ln}\left[{\displaystyle \frac{\left({\displaystyle \frac{R_i}{2M}}1\right)^{1/2}+\mathrm{tan}\sigma }{\left({\displaystyle \frac{R_i}{2M}}1\right)^{1/2}\mathrm{tan}\sigma }}\right]`$ (43)
$`+`$ $`4M\left({\displaystyle \frac{R_i}{2M}}1\right)^{1/2}\left[\sigma +{\displaystyle \frac{R_i}{4M}}\left(\sigma +\mathrm{sin}\sigma \mathrm{cos}\sigma \right)\right],`$ (44)
where $`R_i=a_0\mathrm{sin}(\mathrm{\Xi }_0/a_0)`$ is the radius of the star at the beginning of collapse. Note that for $`\sigma =\sigma ^{}:=\mathrm{arccos}\sqrt{2M/R_i}`$, which corresponds to the event horizon, one has $`t=+\mathrm{}`$, as it should be.
In this model one has $`\mathrm{\Xi }(\tau )=\mathrm{\Xi }_0=\text{const}`$. The functions $`U(u)`$ and $`V(v)`$ can, in principle, be obtained by the systems
$$\begin{array}{ccc}u& =& tX\\ U& =& \tau \mathrm{\Xi }_0+U_0\end{array}\}$$
(46)
and
$$\begin{array}{ccc}v& =& t+X\\ V& =& \tau +\mathrm{\Xi }_0+V_0\end{array}\},$$
(47)
respectively, where $`X=R+2M\mathrm{ln}\left(R/2M1\right)`$ and both $`R`$ and $`t`$ are expressed in terms of $`\tau `$, according to Eqs. (42) and (LABEL:T). From now on, we shall exploit the freedom in the constants $`U_0`$ and $`V_0`$, choosing $`U_0=V_0=\mathrm{\Xi }_0`$, so that $`U=V=\tau `$ at the surface of the star. In practice, Eq. (LABEL:T) is a transcendental one for $`\tau `$ and cannot be inverted. However, we can easily find the inverse functions $`u(U)`$ and $`v(V)`$. Consider an outgoing spherical wavefront of light, characterized by $`U=\overline{U}=\text{const}`$ and $`u=\overline{u}=\text{const}`$, respectively inside and outside the star. In particular, we must have $`\overline{u}=t(\overline{\tau })X(\overline{\tau })`$, where $`\overline{\tau }`$ is the value of $`\tau `$ at the moment when the wavefront crosses the surface of the star. But from Eq. (46) we have also $`\overline{U}=\overline{\tau }`$, so Eqs. (42) and (LABEL:T) allow us to conclude that, inside the star, the relationship between $`U`$ and $`u`$ is
$`u(U)`$ $`=`$ $`t(\tau =U)X(\tau =U)`$ (48)
$`=`$ $`4M\mathrm{ln}\left({\displaystyle \frac{R_i}{2M}}\mathrm{cos}^2\sigma _U1\right)`$ (49)
$`+`$ $`4M\mathrm{ln}\left[\left({\displaystyle \frac{R_i}{2M}}1\right)^{1/2}+\mathrm{tan}\sigma _U\right]R_i\mathrm{cos}^2\sigma _U`$ (50)
$`+`$ $`2M\mathrm{ln}\mathrm{cos}^2\sigma _U+4M\left({\displaystyle \frac{R_i}{2M}}1\right)^{1/2}\sigma _U`$ (51)
$`+`$ $`R_i\left({\displaystyle \frac{R_i}{2M}}1\right)^{1/2}\left(\sigma _U+\mathrm{sin}\sigma _U\mathrm{cos}\sigma _U\right),`$ (52)
with $`\sigma _U=U/2a_0`$. Analogously, we find
$`v(V)`$ $`=`$ $`t(\tau =V)+X(\tau =V)`$ (53)
$`=`$ $`4M\mathrm{ln}\left[\left({\displaystyle \frac{R_i}{2M}}1\right)^{1/2}+\mathrm{tan}\sigma _V\right]+R_i\mathrm{cos}^2\sigma _V`$ (54)
$`+`$ $`2M\mathrm{ln}\mathrm{cos}^2\sigma _V+4M\left({\displaystyle \frac{R_i}{2M}}1\right)^{1/2}\sigma _V`$ (55)
$`+`$ $`R_i\left({\displaystyle \frac{R_i}{2M}}1\right)^{1/2}\left(\sigma _V+\mathrm{sin}\sigma _V\mathrm{cos}\sigma _V\right),`$ (56)
where $`\sigma _V=V/2a_0`$. Substituting Eqs. (52) and (56) into Eqs. (17) and (18), we can express $`t`$ and $`x`$ inside the star as functions of $`U`$ and $`V`$.
Along the worldlines of the observers belonging to the optical frame one has $`x=\text{const}`$, i.e., $`\mathrm{d}u/\mathrm{d}v=1`$. Writing Eqs. (52) and (56) in differential form,
$$\mathrm{d}u=\frac{R_i^2\mathrm{cos}^4\sigma _U}{2a_0M}\frac{\left({\displaystyle \frac{R_i}{2M}}1\right)^{1/2}+\mathrm{tan}\sigma _U}{{\displaystyle \frac{R_i}{2M}}\mathrm{cos}^2\sigma _U1}\mathrm{d}U,$$
(57)
$$\mathrm{d}v=\frac{R_i^2\mathrm{cos}^4\sigma _V}{2a_0M}\frac{\left({\displaystyle \frac{R_i}{2M}}1\right)^{1/2}\mathrm{tan}\sigma _V}{{\displaystyle \frac{R_i}{2M}}\mathrm{cos}^2\sigma _V1}\mathrm{d}V,$$
(58)
we find the slope of these worldlines in $`(U,V)`$ coordinates:
$`{\displaystyle \frac{\mathrm{d}U}{\mathrm{d}V}}`$ $`={\displaystyle \frac{\left(\mathrm{cos}^2\sigma _U{\displaystyle \frac{2M}{R_i}}\right)\mathrm{cos}^4\sigma _V}{\left(\mathrm{cos}^2\sigma _V{\displaystyle \frac{2M}{R_i}}\right)\mathrm{cos}^4\sigma _U}}`$ (60)
$`\times {\displaystyle \frac{\left({\displaystyle \frac{R_i}{2M}}1\right)^{1/2}\mathrm{tan}\sigma _V}{\left({\displaystyle \frac{R_i}{2M}}1\right)^{1/2}+\mathrm{tan}\sigma _U}}.`$
This equation can be integrated numerically for suitable values of $`M`$ and $`R_i`$, and produces diagrams that agree with our qualitative sketch of Fig. 7.
The function $`u(U)`$ assumes a particularly simple asymptotic form near the event horizon, in agreement with the general analysis in Sec. III. The surface of the star crosses the horizon when $`R=2M`$, i.e., at a time $`\tau =\tau ^{}`$ such that $`\mathrm{cos}^2\sigma ^{}=2M/R_i`$. Since the horizon is a null hypersurface with $`U=U^{}=\text{const}`$, we have also that $`\mathrm{cos}^2\sigma _U^{}=2M/R_i`$. Then, for $`UU^{}`$ the first term on the right hand side of Eq. (52) dominates and one can write, asymptotically,
$$u(U)4M\mathrm{ln}\left(\sigma _U^{}\sigma _U\right).$$
(61)
Consistently, one finds from Eq. (57) that $`\mathrm{d}u/\mathrm{d}U4M/(UU^{})`$; in fact, this relationship is not restricted to the present model and holds for a generic collapse, as it follows from Eq. (31). Since $`v(V)`$ remains finite, we have also
$$t(U,V)2M\mathrm{ln}\left(\sigma _U^{}\sigma _U\right)$$
(62)
and
$$x(U,V)2M\mathrm{ln}\left(\sigma _U^{}\sigma _U\right).$$
(63)
Notice that $`t+\mathrm{}`$ and $`x\mathrm{}`$ for $`UU^{}`$. The worldlines of the observers belonging to the optical frame, $`x=\text{const}`$, tend to lie along the event horizon, as one can also see from Eq. (60), which implies $`\mathrm{d}U/\mathrm{d}V0`$ as $`UU^{}`$. All these features are very satisfactory from the point of view of a smooth extension of the Schwarzschild rest frame inside the star, and are in full agreement with the general results obtained in Sec. III.
## V Hawking effect
We now want to figure out how collapse looks like in the optical geometry. For this purpose it is convenient to consider the embedding diagrams of an equatorial plane at different values of $`t`$ (see Ref. for a general discussion of embedding diagrams). When $`R>3M`$, there are no qualitative differences with respect to the conventional description (see Fig. 8). However, as $`R`$ becomes smaller than $`3M`$, a “throat” develops in the external space, in correspondence with $`r=3M`$, while the surface of the star expands progressively (Figs. 9 and 10), escaping to $`x\mathrm{}`$ with the asymptotic law (39). Thus, “collapse” actually corresponds, in the optical geometry, to a sort of expansion into a space that is created by the process itself. This picture is less absurd than it may seem, if one remembers the operational meaning of optical distance outlined in Sec. I. Consider an observer standing at a large value of $`r`$, who sends light signals on a mirror that has been previously placed at the centre of the star, and then defines his distance from the centre simply in terms of the lapse of Killing time $`t`$ taken by the round trip. Since the time delay becomes progressively larger, he will deduce that a collapsing star recedes from him at an increasing speed. Because of spherical symmetry, the only geometrical picture consistent with this description is the one of Figs. 810, where the star expands into a “lower space” that is created in the course of collapse.
If a quantum field is present, this dynamical process disturbs its modes analogously to what happens when there is a moving boundary. In fact, in the two-dimensional section shown in Fig. 6, the situation is exactly the same as if there were a moving boundary, that accelerates asymptotically towards the speed of light, with the law (39). In Minkowski spacetime, this is known to lead, at late times $`t`$, to a thermal flux of radiation with temperature $`(8\pi M)^1`$ (see also Ref. , pp. 102–109 and 229). Given the identity between the $`(t,x)`$ part of the line element of $`\stackrel{~}{g}_{ab}`$ with the one of a two-dimensional Minkowski spacetime, one expects therefore that even in the late stages of collapse there must be a flux of radiation at temperature $`T_H=(8\pi M)^1`$. This is precisely the quantum emission found by Hawking , whose existence follows therefore in a natural way from the picture of collapse based on the optical geometry. It is remarkable that, although the moving boundary analogy has often been used in the literature , optical geometry gives a very simple physical explanation of why it works: Essentially, because gravitational collapse is a particular case of a moving boundary.<sup>§</sup><sup>§</sup>§Although the optical space $`𝒮`$ has not a boundary in the technical sense of the topology of manifolds, the centre of the star works in the same way as far as fields are concerned, because the boundary conditions at $`r=0`$ coincide with those of perfect reflection.
Actually, the analogy with a moving boundary in Minkowski spacetime fails to hold exactly when the angular dimensions are taken into account, because the field propagates on a non-Euclidean geometry. However, this has the only effect of distorting the spectrum of the emitted radiation, according to the way waves are scattered by the curvature of space. Such a correction is well-known in the theory of the Hawking effect .
In this picture the event horizon does not play any role in deriving Hawking’s radiation. In this respect, it is worth reminding that in the optical geometry the horizon is very much alike to a portion of the asymptotic null infinity. To claim that it has anything to do with the thermal flux from black holes would be exactly analogous to saying that the future null infinity is relevant for the radiation produced by a moving boundary in Minkowski spacetime. Thus, explanations of the effect based on quantum tunneling mechanisms appear rather implausible from the perspective of optical geometry.
A similar situation occurs for black hole entropy. In the optical geometry, the celebrated relation $`S=A/4`$ has no direct meaning, because the horizon has an infinite area according to the metric $`\stackrel{~}{h}`$. Indeed, the horizon is located at $`x\mathrm{}`$ in the optical space. If some notion of entropy can be introduced, it can therefore be associated only with the region $`r>2M`$. One possibility is to attribute the entropy entirely to the Hawking radiation, in the following way. Let us consider Fig. 10 again, where the space around a collapsing star at sufficiently late times is visualized as made of a vast “lower” region connected to the usual “upper” one through a mouth at $`r=3M`$. To observers living at large values of $`r`$, the $`r=3M`$ surface looks like the boundary of a three-dimensional cavity containing thermal black body radiation at the temperature $`T_H`$. (One may even consider the $`r<3M`$ region as analogous to a Kirchhoff cavity in ordinary thermodynamics.) Whenever a small amount of energy $`\mathrm{\Delta }E`$ escapes to infinity as Hawking’s radiation, or is added to the collapsing star in the form of accreted matter, the total entropy of the radiation contained in the cavity is modified by the amount $`\mathrm{\Delta }S=\mathrm{\Delta }E/T_H`$. Since $`\mathrm{\Delta }E`$ coincides with the change in the mass, as measured from infinity, we recover the Beckenstein-Hawking expression $`S=4\pi M^2`$ in its differential form. Thus, optical geometry suggests that one should regard the so-called black hole entropy as being actually associated with the Hawking radiation surrounding the collapsing star.That there must be a trapping effect on waves (usually attributed to the reflection off an effective potential barrier) is rather evident from the embedding diagram of Fig. 10. This phenomenon, and its relevance for the study of long-lived gravitational-wave modes, has been discussed in some detail in Ref. . In the present context, it is responsible for the persistence of a consistent amount of radiation in the region with $`r<3M`$. This interpretation contrasts with a viewpoint often expressed, according to which $`S`$ is a property of spacetime, and is similar to others that attribute entropy to a “thermal atmosphere” of the black hole .
Until now, we have assumed that the quantum field is a test one, i.e., that it does not affect the background spacetime. Clearly, this approximation becomes invalid at late times, when the amount of energy carried away by the Hawking radiation that leaks through the neck at $`r=3M`$ becomes a non-negligible fraction of the mass $`M`$ of the star. Let us then sketch qualitatively how back-reaction modifies the picture of collapse suggested by the optical geometry. Roughly, the main effect of Hawking’s emission is to decrease the value of $`M`$. This has essentially three consequences on the diagram of Fig. 10: First, it decreases the size of the mouth at $`r=3M`$; second, it increases the curvature of the optical space at the mouth and in the lower region; third, it increases the value of the Hawking temperature $`T_H`$. Thus, as the process goes on and more radiation is able to escape to $`r+\mathrm{}`$, the geometry in the region $`r>3M`$ of space becomes closer and closer to the Euclidean one, while the throat at $`r=3M`$ shrinks down, becoming progressively sharper. The region $`r<3M`$ has larger and larger negative curvature and contains, at large negative values of $`x`$, matter that escapes to $`x\mathrm{}`$ with increasing acceleration.
Although this description is purely qualitative and is not based on an explicitly constructed model, it gives us important clues about the issue of the final state of a collapsing star, when Hawking’s radiation is taken into account. First of all, it implies that the black hole, in the strict sense of the region beyond the future event horizon $`^+`$, does not form. This is evident in the optical geometry point of view, where $`^+`$ corresponds to $`x\mathrm{}`$ and $`t+\mathrm{}`$. In the usual language, one would say that radiation makes the value of $`M`$ decrease, which in turn makes the horizon shrink, and eventually reduce to a point before the star could cross it. Thus, the very existence of the Hawking effect would prevent the formation of black holes by collapse.This was suggested already in the early years following Hawking’s discovery , but the idea has never been pursued further, at least to the authors’ knowledge.
If this is the case, what are then we left with in the limit $`t+\mathrm{}`$? A straightforward extrapolation of the process that we have just described suggests that, for $`M0`$, the throat at $`r=3M`$ pinches off, leaving two spaces — one Euclidean, the other with infinite negative curvature — with just one point in common. This is clearly a degenerate, and highly implausible, situation; one would rather think that the process stops when the throat reaches a critical size (perhaps at the Planckian scale), because of as yet unknown physical processes. Such a situation corresponds to the hypothesis of remnants in the common description .
Remnants have originally been introduced as a possible resolution of the information paradox . However, the viability of this hypothesis has sometimes been questioned, because the small scale and mass of remnants would allow them to store very little information . Roughly, the argument is the following. A physical system with size $`l`$, total energy $`E`$, and Hamiltonian bounded from below has a number of states of the order of $`El`$; thus, it can contain a maximum information which is also of order $`El`$. The status of this claim is rather controversial . However, even assuming that there is indeed such an information bound, it would hardly represent a difficulty in optical geometry, where the “remnant” is actually an enormously vast region, because ordinary distances are rescaled by a huge factor.<sup>\**</sup><sup>\**</sup>\**This is similar to what happens in the “cornucopion” scenario of dilaton gravity coupled to electromagnetism (see Ref. for criticisms of this model). However, the internal geometries and the physical contexts are very different in the two cases. The possibility that information be stored inside a large region delimited by a small neck has been discussed in general by Giddings . Even if $`l`$ should represent the size of the remnant as seen from the exterior, as suggested by Bekenstein , the bound would still be circumvented thanks to the existence, in the lower space, of an enormous supply of negative energy associated with the field $`\mathrm{\Phi }`$.
## VI Conclusions and outlooks
We have constructed the optical geometry for a spherically symmetric collapsing body. The procedure adopted is a simple extension of the technique used when defining optical distance operationally in a static spacetime. Essentially, for any spacetime event $`P`$, one considers the in- and out-going spherical light-fronts that cross at $`P`$. These correspond to well-defined values of advanced and retarded time $`v`$ and $`u`$, that can be read at infinity as the affine parameters along the null generators of $`^{}`$ and $`^+`$. The event $`P`$ is then labeled by $`v`$ and $`u`$, in addition to the angular coordinates, and can also be identified by a timelike coordinate $`t:=(v+u)/2`$ and a spacelike one, $`x:=(vu)/2`$. The optical metric is then defined as the only one which is conformal to $`g_{ab}`$ — the metric of spacetime in general relativity — and which reproduces the two-dimensional Minkowskian line element on $`\theta =\text{const}`$, $`\phi =\text{const}`$ sections. Outside the collapsing star, $`g_{ab}`$ has the Schwarzschild form and $`x`$ coincides with the Regge-Wheeler “tortoise” coordinate, thus giving the well-known optical metric of an empty, spherically symmetric spacetime. However, our construction allows one to extend the optical geometry smoothly inside the star.
Associated with the optical space are the optical frame $`n^a`$, made of the observers at $`x=\text{const}`$, and the scalar potential $`\mathrm{\Phi }`$, essentially the logarithm of the conformal factor that links $`g_{ab}`$ and $`\stackrel{~}{g}_{ab}`$. These two concepts are related to each other in the following sense. The optical frame defines a notion of “rest in the gravitational field,” thus the four-acceleration $`n^b_bn_a`$ of the observers can be identified with the gravitational acceleration, changed of sign. Since $`n^b_bn_a`$ is equal to the spatial gradient of $`\mathrm{\Phi }`$ in the optical frame, see Eq. (1), $`\mathrm{\Phi }`$ can be thought of as a covariant generalization of the gravitational potential. This allows one to give a precise meaning to the notion of gravitational field inside a collapsing object. It must be noted that a gravitational field so defined, although analogous to the same Newtonian concept, nevertheless differs from it in some essential details. As an example, consider a collapsing spherical shell. It is obvious from Fig. 7 that $`n^b_bn_a0`$ everywhere, thus leading to the conclusion that the gravitational field is nonvanishing inside the shell, as well as outside. This contrasts with intuition shaped after Newton’s theory, according to which the gravitational field inside the shell should be zero. However, this difference is not surprising if one considers that $`\mathrm{\Phi }`$ does not obey the Poisson equation, but a nonlinear generalization of it, as can be easily checked on replacing $`g_{ab}=\mathrm{e}^{2\mathrm{\Phi }}\stackrel{~}{g}_{ab}`$ in the trace of the Einstein equation.
Although we have focused our treatment on the case of a star collapsing in empty space, generalizations to, e.g., the collapse of electrically charged bodies are straightforward (see, e.g., Ref. for a configuration with extremal charge, $`Q^2=M^2`$). In fact, the construction presented in Sec. II makes a heavy use of the null structure on $`^{}`$ in defining the coordinates $`v`$ and $`u`$, but seems rather general otherwise. It would be important to understand whether a similar technique could be used to define optical geometry in still more general situations — for example, for non-spherically symmetric collapse.
We have seen that, in optical geometry, gravitational collapse corresponds to a very fast expansion of the star into a “lower space,” with the centre receding from the optical observers $`n^a`$ at a speed that approaches exponentially the speed of light. This process excites the modes of a quantum field, as it happens in the presence of a moving boundary, leading to the production of Hawking radiation. Thus, optical geometry provides one with a physical origin of the formal “moving mirror analogy” of the Hawking effect.
The issue of the final state of black hole evaporation is also clarified by the use of optical geometry. There are essentially two possibilities, in both of which the black hole does not form, strictly speaking. Either the evaporation process continues until one remains with flat space and a (physically unaccessible) infinitely warped “lower space” or, perhaps more likely, the process stops at some scale leaving an enormously large remnant. In both cases there is no information paradox. Of course, since the issue involves distances of the order of the Planck length, a definitive answer is beyond the limits of applicability of present-day physics. Nevertheless, it would be useful to get a more detailed insight into the evaporation process by reformulating simple models that include back-reaction (see, e.g., Ref. ) in the language of optical geometry.
Of course, all our conclusions apply to black holes deriving from collapse. However, it is not difficult to extend them to eternal black holes, simply studying the quantum field theory on the optical spacetime associated with the Schwarzschild solution. In this case there is no moving boundary, and the optical metric looks (as far as the $`t`$ and $`x`$ coordinates are concerned) exactly like the Minkowski one. In particular, the observers at $`x=\text{const}`$, belonging to the optical frame, correspond to the Minkowskian inertial observers. Thus, one expects that they should register no particles, provided that a condition of “no incoming radiation” is imposed for $`t\mathrm{}`$. This is precisely what happens in the so-called Boulware (or Schwarzschild) state $`|0_B`$ . The optical geometry viewpoint provides therefore support for regarding $`|0_B`$, rather than the Hartle-Hawking (or Kruskal) and the Unruh states $`|0_H`$ and $`|0_U`$ (see Ref. , pp. 281–282), as describing the quantum vacuum around an eternal black hole.
The Boulware state is often considered pathological because the expectation value of several physical quantities diverges at the horizon . For example, for the stress-energy-momentum tensor operator of a massless scalar field $`\varphi `$ one has, after renormalization, that $`0_B|T_{ab}|0_B_{\mathrm{ren}}`$ depends on $`r`$ as $`(12M/r)^1`$ for $`r2M`$. It is remarkable that, in the optical spacetime $`(,\stackrel{~}{g}_{ab})`$, such a factor is exactly canceled out, because $`\varphi `$ is conformally transformed in $`\stackrel{~}{\varphi }=\mathrm{e}^\mathrm{\Phi }\varphi `$, and $`0_B|\stackrel{~}{T}_{ab}|0_B_{\mathrm{ren}}`$ turns out to be finite everywhere, representing just the vacuum polarization due to curvature. A similar situation occurs when considering the response function $`\mathrm{\Pi }(\omega |r)`$ for an ideal static detector in the Boulware state. In the spacetime $`(,g_{ab})`$ one has, near the horizon ,
$$\mathrm{\Pi }(\omega |r)\frac{1}{12M/r}\frac{\omega }{2\pi }\mathrm{\Theta }(\omega ),$$
(64)
where $`\mathrm{\Theta }`$ is the step function. However, after the conformal rescaling the response is given by
$$\stackrel{~}{\mathrm{\Pi }}(\omega |r)\frac{\omega }{2\pi }\mathrm{\Theta }(\omega ),$$
(65)
which not only is finite, but is also the answer one would expect in a proper “vacuum.” Thus, the pathologies of the Boulware state are removed by the conformal transformation to the optical spacetime, in which $`|0_B`$ behaves as a satisfactory quantum vacuum.
###### Acknowledgements.
It is a pleasure to thank Stefano Liberati for many helpful discussions. Part of this work was done while SS and JA were at the Department of Astronomy and Astrophysics of Chalmers University. SS was also partially supported by the Interdisciplinary Laboratory of SISSA and by ICTP.
## Appendix A: Proof of Eq. (1)
¿From the definition (2) of $`\mathrm{\Phi }`$ it follows that the vector field $`n^a`$ can be expressed as $`n^a=\mathrm{e}^\mathrm{\Phi }\eta ^a`$. Thus,
$$n^b_bn_a=\mathrm{e}^{2\mathrm{\Phi }}\left(\eta ^b_b\eta _a\eta _a\eta ^b_b\mathrm{\Phi }\right).$$
(A.1)
Using now the relation
$$_a\eta _b+_b\eta _a=\frac{1}{2}g_{ab}_c\eta ^c,$$
(A.2)
valid for any conformal Killing vector field $`\eta ^a`$ (see e.g. Ref. , pp. 443–444), we can write
$$\eta ^b_b\eta _a=\mathrm{e}^{2\mathrm{\Phi }}_a\mathrm{\Phi }+\frac{1}{2}\eta _a_b\eta ^b.$$
(A.3)
Furthermore,
$$\eta ^b_b\mathrm{\Phi }=\frac{\eta ^b\eta ^c_b\eta _c}{\eta _a\eta ^a}=\frac{\eta ^b\eta ^c_{(b}\eta _{c)}}{\eta _a\eta ^a}=\frac{1}{4}_b\eta ^b.$$
(A.4)
Substituting Eq. (A.4) into Eq. (A.3) and then the latter into Eq. (A.1) we obtain, after some trivial algebra, Eq. (1).
## Appendix B: Physical meaning of the surface gravity in optical geometry
Here we show that, for a static black hole, the surface gravity coincides with the magnitude of the gravitational pull on the natural observers for $`t+\mathrm{}`$, as measured with respect to the optical metric. A convenient expression for the surface gravity $`\kappa `$ is
$$\kappa =lim\left[\frac{\left(\eta ^a_a\eta _c\right)\left(\eta ^b_b\eta ^c\right)}{\eta _d\eta ^d}\right]^{1/2},$$
(B.1)
where “$`lim`$” stands for the limit as the horizon is approached (see Ref. , p. 332). From the Killing equation and the definition of $`\mathrm{\Phi }`$ it follows that $`\eta ^a_a\eta _c=_c\left(\eta ^a\eta _a\right)/2=\mathrm{e}^{2\mathrm{\Phi }}_c\mathrm{\Phi }`$. Substituting into Eq. (B.1) we obtain
$$\kappa =lim\left(\stackrel{~}{g}^{ab}_a\mathrm{\Phi }_b\mathrm{\Phi }\right)^{1/2}=lim\left(\stackrel{~}{h}^{ab}\stackrel{~}{D}_a\mathrm{\Phi }\stackrel{~}{D}_b\mathrm{\Phi }\right)^{1/2},$$
(B.2)
where $`\stackrel{~}{D}`$ is the Riemannian connection associated with $`\stackrel{~}{h}_{ab}`$, and we have used the stationarity property in the form $`\eta ^a_a\mathrm{\Phi }=0`$. By Eq. (1), the last term in Eq. (B.2) is just the value of the gravitational pull on the observers of the optical frame.
This proof can be generalized without any difficulty to the case of a stationary black hole, simply by replacing $`\eta ^a`$ with a Killing vector field which is normal to the horizon.
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# Globular Clusters in the dE,N galaxy NGC 3115 DW1: New Insights from Spectroscopy and HST Photometry
## 1 Introduction
The study of globular cluster systems of dwarf galaxies complements the numerous studies of such systems in larger elliptical and spiral galaxies. Few globular cluster systems (GCSs) around dwarf galaxies beyond the Local group have been studied to date with respect to their cluster system (see Ashman & Zepf 1998). This is mostly due to the low numbers of globular clusters present in such galaxies. However, their properties are relevant for a number of globular cluster system formation scenarios. Dwarf galaxies are expected to provide insight into how the smallest galaxies build up a system of globular clusters. Further, their properties must be known in order to verify scenarios in which larger globular cluster systems are predicted to build up by the accretion of proto-galactic fragments or dwarf galaxies (Kissler-Patig et al. 1998a, Côté et al. 1998, Hilker et al. 1999). These scenarios relate to the older idea that galaxy halos might have formed through the assembly of such small stellar systems (e.g. Searle & Zinn, 1978).
Photometric studies of several globular cluster systems in dwarf galaxies were carried out by Durrell et al. (1996a) and Miller et al. (1998). Durrell et al. (1996a) studied the systems of 11 dwarf galaxies in the Virgo cluster. All were found to host globular cluster candidates and have specific frequencies ranging from 3 to 8, similar to Local Group dwarfs and giant elliptical galaxies. Miller et al. (1998) studied 24 dwarf ellipticals in the Virgo and Fornax clusters as well as in the Leo group. They found that dE,N galaxies had higher specific frequencies than dE galaxies, with values around $`S_\mathrm{N}=6.5\pm 1.2`$, increasing with increasing $`M_V`$ (decreasing luminosity). Not much is known yet about the metallicities of globular clusters in dwarf galaxies. Minniti et al. (1996) constructed a metallicity distribution for all Local Group dwarf galaxies and noticed that the distribution was peaked around \[Fe/H\]$`1.7`$ dex with no clusters more metal-rich than \[Fe/H\]$`=1.0`$ dex. Durrell et al. (1996a) derived metallicities from Washington colors for two of their Virgo dwarf ellipticals and obtained a mean metallicity of \[Fe/H\]$`=1.45\pm 0.2`$ dex.
Durrell et al. (1996b) studied the GCS of the dE NGC 3115 DW1 in more detail and found it to be relatively rich (see below for a more detailed description of their results). This motivated us to carry out spectroscopy for some of the globular cluster candidates in this galaxy to get a more detailed picture of their chemical and kinematical properties. Further, HST/WFPC2 data were available from the archive, allowing us to study the sizes of some of the clusters.
NGC 3115 DW1 is a dE1,N galaxy in the vicinity of the giant S0 galaxy NGC 3115. It is located RA: 10h 05m 41.6s; Dec: $`07^\mathrm{o}`$ 58′ 53.5″ ($`l=248.12^\mathrm{o}`$; $`b=36.69^\mathrm{o}`$). We will assume a distance of $`11_{2.3}^{+5.0}`$ Mpc throughout this paper following Durrell et al. (1996b). Additional properties will be given in the text where they are relevant. In Section 2 we describe the new data. In Section 3 we analyze the spectroscopic data giving first a brief kinematical study of the globular cluster system before discussing the abundances and the overall metallicity of the system. In Section 4 we revisit the previous photometry (number of clusters, colors, specific frequency), compare the photometric metallicities with the spectroscopic ones and add the HST/WFPC2 imaging to derive sizes for two clusters. In Section 5 we discuss whether NGC 3115 DW1 could have suffered stripping by its giant companion. We summarize our results in Section 6.
## 2 Data
### 2.1 Spectroscopic Observation
46 spectra of candidate globular clusters in NGC 3115 DW1 were obtained with the LRIS-A spectrograph (Oke et al., 1995) on the Keck-II telescope during the nights of 1996 December 18th and 19th. The observations were performed in multi-slit mode using two slit masks containing 24 (mask A) and 22 slits (mask B), respectively. Mask A was exposed for 4800 seconds while mask B was exposed 6600 seconds. The multi-object-spectra images were taken with a Tektronix 2048$`\times `$2048 pix<sup>2</sup> chip. The seeing was around 0.8″ during both nights. For all observations the first order of the 600/5000 grism (i.e. 600 lines mm<sup>-1</sup>) was used with slitlets of 1″ which provided a dispersion of 1.3 Å pix<sup>-1</sup> (at 5000 Å) and a spectral resolution of 4.1 Å FWHM<sup>-1</sup>.
All frames were binned (1$`\times 2`$) perpendicular to the dispersion axis during the read-out. For each night the images were reduced individually and eventually combined. A master bias was created from 5 zero-second images, taken at the end of each night. A master flat was obtained from 5 twilight-flat images obtained each night. Each object image was bias subtracted and flat-fielded in a standard way. With the IRAF<sup>1</sup><sup>1</sup>1IRAF is distributed by the National Optical Astronomy Observatories, which are operated by the Association of Universities for Research in Astronomy, Inc., under cooperative agreement with the National Science Foundation. package apall we traced, extracted and sky subtracted all object spectra. The sky was “optimally” subtracted, i.e. modeled with variance weighting perpendicular to the object spectra before being subtracted (Horne, 1986). HgKrNe-calibration-lamp spectra were obtained for wavelength calibration. The calibration spectra were traced, extracted, and sky subtracted exactly in the same way as the object spectra. The wavelength calibration was verified on sky spectra included in each slit spectrum. The overall standard deviation of the wavelength calibration was determined to be $`\sigma _{\mathrm{cal}}0.12`$Å for both nights.
The field-of-view of LRIS-A is 6′ $`\times `$ 8′. We placed the multi-slit masks on the central/north and south-eastern part of NGC 3115 DW1 as can be seen in Fig. 1. The slit masks were aligned such that the slits were pointing towards the center of the giant S0 galaxy NGC 3115. Due to the restrictions in making the masks (long enough slits, wavelength coverage, etc.) about 50% of the targets could be selected from the previous photometry of Durrell et al. (1996b) (see Fig. 1). The other 50% was ‘blindly’ selected from an LRIS acquisition image. Only 6 of the former group lay within 48″ of the galaxy center, i.e. within the region where Durrell et al. (1996b) detected a clear overabundance of objects. We therefore expected $``$ 6+ objects in our total sample to be true globular clusters.
### 2.2 Additional Photometry
In order to measure sizes of globular cluster candidates that were confirmed by our spectroscopy, additional data for the globular cluster system were obtained from the HST archive<sup>2</sup><sup>2</sup>2 Based on observations made with the NASA/ESA Hubble Space Telescope, obtained from the data archive at the Space Telescope Science Institute. STScI is operated by the Association of Universities for Research in Astronomy, Inc. under NASA contract NAS 5-26555.. Short exposures of 160 seconds in F555W and 320 seconds in F814W filter were taken from program GO:5999 (PI: A.Phillips). A simple reduction procedure was applied to all HST images: The images were biased, flat-fielded and calibrated as described in Holtzman et al. (1995), including corrections accounting for CTE and field distortions. The brightest clusters were matched on both the F555W and F814W images and their sizes were studied. In addition, we re-analyzed the ground-based B,V data<sup>3</sup><sup>3</sup>3The data were kindly provided in electronic form by Patrick Durrell. from Durrell et al. (1996b). The field-of-views of both photometric datasets are over-plotted in Fig. 1.
## 3 Analysis of Spectroscopic Data
In this section we extract heliocentric radial velocities $`v_{\mathrm{rad}}`$ of globular clusters and perform a mass estimate of the host galaxy, NGC 3115 DW1. Subsequently, we measure abundances using the Lick/IDS passband definitions and infer a mean metallicity of the GCS.
### 3.1 Kinematics
We derived radial velocities by two different methods. Velocities were derived by cross-correlation with high S/N template spectra of two bright GCs in M31: 225-280 and 158-213 (for nomenclature see Huchra et al., 1982). The cross-correlation was carried out with the IRAF task fxcor. All measurements are summarized in Table 1. The “internal” errors which are given by the cross-correlation code lie about $`\sigma 80100`$ km s<sup>-1</sup> for well defined high-S/N spectra and increase rapidly as the object spectrum becomes less defined.
For 7 high-S/N spectra we also estimated the radial velocity by measuring redshifts of individual absorption lines (see footnote in Table 1). The mean radial velocity from cross-correlation served as a first guess to the mean wavelength-shift determination. We used the IRAF package rvidlines which employs a center1d code to match the center of each individual absorption feature (see manual of rvidlines). We obtained an “internal” mean error from the averaging process which is $``$ 100 km s<sup>-1</sup>. The results and their errors are included in Table 1.
#### 3.1.1 Selection of Globular Clusters
After obtaining the radial velocities from cross-correlation and in some cases from wavelength shifts of individual absorption lines, we combined all available measurements for further selection of globular clusters by radial velocity. Velocities derived from cross-correlation were given twice the weight of velocities from absorption lines. The average radial velocity for each object can be found in the last but one column of Table 1.
Figure 2 shows a radial-velocity histogram of objects with $`500v_{\mathrm{rad}}1000`$ km s<sup>-1</sup>. The figure shows a hint of bimodality due to the extended tail of velocities $`>`$ 500 km s<sup>-1</sup>. Based on the measured radial velocity of NGC 3115 DW1 of $`700`$ km s<sup>-1</sup> ($`715\pm 62`$ km s<sup>-1</sup> de Vaucouleurs et al. 1991, $`716\pm 19`$ km s<sup>-1</sup> Peterson & Caldwell 1993, and $`698\pm 42`$ km s<sup>-1</sup> Capaccioli et al. 1993) and an adopted velocity dispersion of the GCS in NGC 3115 DW1 of $`100`$ km s<sup>-1</sup> (see below), globular clusters should have radial velocities in the range $`4001000`$ km s<sup>-1</sup>. Seven candidates lie within this range. This cut enables a reliable differentiation between foreground objects of low radial velocity and globular clusters in NGC 3115 DW1.
We used the KMM code (Ashman et al., 1994) to obtain the statistical significance of bimodality. The KMM code fits two Gaussians to the data using Maximum-Likelihood techniques. The data can be fit with two Gaussians of identical dispersion or two Gaussians of independent dispersion. We applied both techniques. The mean radial velocities of the two Gaussians are $`48\pm 26`$ km s<sup>-1</sup> and $`567\pm 46`$ km s<sup>-1</sup> for identical-dispersion and $`52\pm 28`$ km s<sup>-1</sup> and $`575\pm 39`$ km s<sup>-1</sup> for independent-dispersion fitting. The error is the error of the mean calculated from the variance of each distribution. KMM estimates also the fraction of data points which are part of each sub-distribution. Of 28 objects in the histogram (one additional object, D26, is not included due to $`v_{\mathrm{rad}}<500`$ km s<sup>-1</sup>) in Figure 2, KMM assigns 21 objects to the sub-population with the lower mean radial velocity and 7 to the sub-population with the higher mean radial velocity which is in good agreement with the expectations (see 2.1). The confidence level for bimodality is $`>99`$%.
The (weighted) mean radial velocity of all 7 detected globular clusters is $`v_{\mathrm{rad}}=572\pm 30`$ km s<sup>-1</sup>. This value deviates from the measured radial velocity of NGC 3115 DW by $`1.6\sigma `$ which is a hint that we are biased towards lower velocities by both the small sample and the choice of slit-mask position on the sky. Since we detected globular clusters predominantly in the central and northern field of NGC 3115 DW1 our measurements could be influenced by a systematic rotation of the GCS (see also Sect. 3.1.4).
All 22 objects with low radial velocities have a weighted mean velocity of $`v_{\mathrm{rad}}=50\pm 19`$ km s<sup>-1</sup>. Assuming a simple stellar rotation model for the Milky Way we expect the foreground stars in the direction of NGC 3115 DW1 ($`l=248.12^\mathrm{o}`$ and $`b=36.69^\mathrm{o}`$) to have $`v_{\mathrm{rad}}=220\mathrm{sin}2l\mathrm{cos}^2b=98`$ km s<sup>-1</sup> (van de Kamp, 1967) which is in rough agreement with our measurement.
In summary, within our data of 46 spectroscopically-analyzed objects we found 7 globular clusters 22 foreground stars, and 15 background galaxies (9 of them are significantly clumped about $`v_{\mathrm{rad}}24000`$ km s<sup>-1</sup> or $`z0.08`$; the velocity dispersion of this potential galaxy cluster is $`\sigma =1300`$ km s<sup>-1</sup>). 2 objects (D21,L6) could not be identified reliably and were therefore dropped. The radial-velocity data for all objects are summarized in Table 1.
#### 3.1.2 Mass Estimate of NGC 3115 DW1
Our data sample of 7 globular clusters is only sufficient for a first rough mass estimate of NGC 3115 DW1 and its M/L ratio. Furthermore, we rely in this section on the assumption that the system is not influenced by the nearby giant S0 galaxy, an assumption that we will question in Sec. 5.
We used two mass estimators which are extensively described and tested by Bahcall & Tremaine (1981) and Heisler et al. (1985). We applied the Virial Mass Estimator (VME) and the Projected Mass Estimator (PME) accounting for different orbit characteristics of the globular clusters. Since the alignment of our multi-slit masks is biased towards clusters of the central and northern quadrant of the galaxy, our data is subject to unknown systematic east-west rotation of the entire globular cluster system. If the angular momentum vector points along the north-south axis we will have under- or overestimated the mass depending on the direction of rotation. Although, it is possible to provide a lower-mass limit from our spatially constrained cluster sample by varying the mean systemic velocity (see below), it is not possible to correct completely for the unknown total rotation of the globular cluster system (see also 3.1.4).
We let the radial velocity of NGC 3115 DW1 (which is assumed to be the mean velocity of the GCS as well) vary over a wide range while the measured radial velocities of the globular clusters remained fixed. During each step of 10 km s<sup>-1</sup> we calculated the mass of the galaxy with each mass estimator. The resulting plot (mass vs. $`v_{\mathrm{rad}}`$ of the GCS) is shown in Figure 3. Table 2 shows the lower mass limits calculated with all mass estimators.
At the highest measured radial velocity of NGC 3115 DW1 ($`v_{\mathrm{rad}}=716\pm 19`$ km s<sup>-1</sup>, Peterson & Caldwell, 1993) we obtain total masses in the range $`6.310^{10}`$M to $`3.610^{11}`$M. The estimated masses for the measured mean radial velocity of the globular cluster system ($`v_{\mathrm{rad}}=572\pm 30`$ km s<sup>-1</sup>) are in the range $`2.110^{10}`$M to $`1.010^{11}`$M. The errors in Table 2 are the intrinsic, statistic, and systematic uncertainties of the mass-estimate. The first is the uncertainty of the code itself (Bahcall & Tremaine, 1981) while the remaining are due to the limited sample size and the uncertain distance. All masses are the total mass estimates within a galactocentric radius of $`r189.4`$″ or $`R10.1`$ kpc, respectively. This is the projected radial distance of the outermost globular cluster (L63).
Assuming an isotropic orbit distribution for globular clusters in NGC 3115 DW1 the lower-mass limit for the galaxy of $`M_{\mathrm{PME}}=(4.8\pm 2.3)10^{10}`$M seems rather large for a dwarf elliptical. The absolute magnitude of $`M_V=17.7`$ mag (Durrell et al., 1996b) is high as well (more than a magnitude brighter than M32 for example, and similar to NGC 4486B, although NGC 3115 DW1 has a dissimilar structure to that of these low-L Es). The mean absolute V-magnitude for nearby dwarf ellipticals is $`<M_V>16.9`$ mag (Ferguson & Binggeli, 1994). Considering its mass and luminosity we address the fact that NGC 3115 DW1 appears to be a transition-type galaxy between luminous dEs and low-luminosity ellipticals in the discussion section.
#### 3.1.3 Radial Dependencies and Mass-to-Light Ratios
Limiting the data set of radial velocities to smaller radii we can probe the radial mass dependencies in NGC 3115 DW1. Clearly, statistical errors become important when we reduce the already small data set by removing the outermost globular clusters. Nonetheless, we calculate mass estimates for different radii (using the PME and assuming isotropic orbits) since the kinematics of the two outer clusters might be influenced by the nearby giant S0. The results are summarized in Table 3.
Combining the former findings with the photometry of Durrell et al. (1996b), who found total magnitudes $`V_\mathrm{T}=12.63\pm 0.06`$ mag and $`B_\mathrm{T}=13.57\pm 0.09`$ mag, we estimate rough mass-to-light ratios for different radii. Applying a King profile (King, 1962) with the parameters $`r_c=`$14.4″ and $`c`$1.4 (Durrell et al., 1996b) we obtain at the outermost globular-cluster projected radius $`r=189.4`$″ ($`R=10.1`$ kpc) M/L$`{}_{V}{}^{}=52\pm 25`$ with a $`1\sigma `$ uncertainty due to statistical sample-size uncertainties and photometric errors. The systematic error due to distance uncertainties of NGC 3115 DW1 is $`{}_{32}{}^{}{}_{}{}^{+67}`$. Going inwards, the M/L<sub>V</sub> ratio drops. At a projected radius of $`r=56.3`$″ ($`R=3.0`$ kpc) the mass-to-light ratio is M/L$`{}_{V}{}^{}=22\pm 13`$.
Our analysis can be expanded by combining our results with the M/L<sub>V</sub> measurements for the innermost part of NGC 3115 DW1. Peterson & Caldwell (1993) measure within $`r3`$″ (or $`R160`$ pc) of NGC 3115 DW1 a central velocity dispersion of $`\sigma =2030`$ km s<sup>-1</sup>. They derive a M/L<sub>V</sub> of $`3\pm 2`$ within $`r3`$″. A gradient of M/L<sub>V</sub> can therefore be traced outward from the center of the galaxy, although the uncertainties are quite large. However, all values fit well into the range given by Ferguson & Binggeli (1994) for dwarf elliptical galaxies (M/L$`{}_{V}{}^{}5`$ for Fornax up to M/L$`{}_{V}{}^{}100`$ for Ursa Minor).
#### 3.1.4 Rotation and Velocity Dispersion
Using the Maximum-Likelihood method of Pryor & Meylan (1993) we measure a marginal net rotation of $`v_{\mathrm{rot}}=75\pm 70`$ km s<sup>-1</sup> for the globular cluster system of NGC 3115 DW1. Note that to date no dE was reported to show significant rotation of its stellar body (Ferguson & Binggeli, 1994). The position angle of the rotation axis is $`\theta =90^\mathrm{o}\pm 60^\mathrm{o}`$ (poorly defined given the weak rotation). This result may be of course biased by the incomplete spatial coverage of our small sample. Correcting for the net rotation we obtain a line-of-sight velocity dispersion of $`\sigma =130\pm 15`$ km s<sup>-1</sup> for the full sample of 7 clusters, and $`\sigma =74\pm 36`$ km s<sup>-1</sup> for the inner 5 clusters.
The major axis of the dE1,N galaxy of $`\theta _{\mathrm{gal}}=100^\mathrm{o}\pm 10^\mathrm{o}`$ (Durrell et al., 1996b) appears at face value nearly parallel to the rotation axis of the GCS although the latter is not well defined as mentioned above. Spectroscopy for more clusters with better spatial coverage is certainly needed to establish the axis alignment.
### 3.2 Metallicity
#### 3.2.1 Mean Metallicity of the Globular Cluster System
The S/N of our spectra are insufficient to reliably establish individual cluster metallicities. To measure a mean abundance for the GCS we combined all the individual spectra into a high-S/N ‘mean’ spectrum. All individual globular cluster spectra as well as the combined spectrum are shown in Figure 4.
For abundance measurements we used the passband definitions of Brodie & Huchra (1990) and the new Lick/IDS passband definitions of Trager et al. (1998). All abundances (both of single spectra and the mean spectrum) are given in Table 4. Brodie & Huchra (1990) calibrate single element abundances with \[Fe/H\] using a large sample of Milky Way and M31 globular clusters. We used their calibration to estimate the mean \[Fe/H\] for the entire GCS, based on the composite spectrum. The mean metallicity of our globular cluster sample is $``$\[Fe/H\]$`_{\mathrm{GCS}}=0.97\pm 0.11`$ dex. Exactly the same value is obtained from the weighted mean of the individual measurements (see Table 5).
Overall, the mean metallicity does not follow the empirical GCS-metallicity — galaxy-luminosity relation (see Fig. 5.7 in Ashman & Zepf, 1998). NGC 3115 DW1 ($`M_V=17.7`$ mag) falls in the transition region between dwarfs and elliptical galaxies, whereas the mean metallicity of the GCS falls in the range of metallicities found in giant elliptical galaxies. According to this empirical relation NGC 3115 DW1 appears to be slightly too metal-rich for its luminosity.
#### 3.2.2 Abundance Ratios
We compare the mean abundance ratios of the GCS with abundance ratios of globular clusters in other galaxies. Trager et al. (1998) provide a compilation of abundances of Milky Way and M31 globular clusters. Abundances of NGC 1399 globular clusters were measured by Kissler-Patig et al. (1998b). Both data sets use the definitions of the Lick/IDS system. In order to minimize the statistical noise of abundance measurements we calculate a mean iron abundance $`\text{Fe}`$ and a mean metal abundance $`[\text{MgFe}]`$ (see González, 1993, for a detailed discussion).
Figure 5 shows abundance ratios for several dominant elements. The upper four panels show abundance ratios relative to the \[MgFe\] index. In the upper left panel the age-sensitive H$`\beta `$ abundance is seen to be in good agreement with Milky Way and M31 data, although possibly at the lower (older) edge. The upper right panel shows the G-band index compared with the \[MgFe\] index. The G-band is a primary metallicity indicator (Brodie & Huchra, 1990). The data show no abundance anomalies.
The two middle panels of Figure 5 could in principal be used to examine the $`\alpha `$-element content of these globular clusters. For the stellar light of (mostly brighter) ellipticals Worthey et al. (1992) found an $`\alpha `$-element enhancement. The \[$`\alpha `$/Fe\] enhancement is a very sensitive indicator of the star formation rate in a galaxy. As $`\alpha `$-elements are preferentially created in SNe type II, their enhancement indicates a violent star formation and/or a top-heavy IMF. A depression, or normal values, of \[$`\alpha `$/Fe\] would result from quiet star formation in which SNe Ia dominate the enrichment processes. Given the relatively high mean metallicity ($``$\[Fe/H\]$`_{\mathrm{GCS}}=0.97\pm 0.11`$ dex), a normal $`\alpha `$-element ratio would suggest that these clusters formed from enriched material either during epochs of quiet star formation, or at the very beginning of a burst. Better data could lead to interesting insights on this topic.
## 4 Photometry
### 4.1 Spectroscopic check on photometry
The main result from the photometry of Durrell et al. (1996b) was that the GCS in NGC 3115 DW1 is rich with a specific frequency of $`S_\mathrm{N}=4.9\pm 1.9`$ and a total globular cluster population $`N_{\mathrm{GC}}=59\pm 23`$. Our spectroscopy and the photometric data set of Durrell et al. have 22 objects in common (see Figure 1 and Table 6). 6 of these 22 objects have projected radial distances of $`r48`$″ ($`R2.6`$ kpc). Durrell et al. consider the GCS to lie within this radius (mainly because the surface over-density of objects disappears beyond it). We can confirm 4 of the 6 objects as bona-fide globular clusters (D7, D14, D15, and D46). One object was found to be a background galaxy and another cannot be identified either by radial velocity or by its spectrum.
Assuming that this sample of 6 objects is a statistically representative sample of objects in the projected vicinity of NGC 3115 DW1 the upper limit of the probability of finding a globular cluster within a radius of $`r48`$″ ($`R2.6`$ kpc) around the center of NGC 3115 DW1 is $`f_{\mathrm{GC}}=5/60.83`$ (the lower limit is $`f_{\mathrm{GC}}=4/60.67`$; if we exclude the non-identified object). This is in good agreement with the findings of Durrell et al. (1996b), although the statistical significance is very low. Durrell et al. measure the contaminating surface density of background objects to be $`\sigma =6.4\pm 1.9`$ arcmin<sup>-2</sup>. Within their radial limit of $`r48`$″ there are a total of $`13`$ background objects. For a total population of $`N_{\mathrm{GC}}=59\pm 23`$, the probability of picking a globular cluster within $`r48`$″ (and at the magnitude limit of the photometry of Durrell et al.) is $`f_{\mathrm{GC}}=59/(13+59)0.82`$. Our spectroscopic results suggest that there is no need to make any correction to the values for specific frequency and total globular cluster population size derived from photometry.
Miller et al. (1998) measure specific frequencies for dwarf elliptical galaxies in the Virgo and Fornax cluster and find a $`\mathrm{log}(S_\mathrm{N})M_V`$ relations for the nucleated dwarfs. NGC 3115 DW1 has a higher specific frequency ($`S_\mathrm{N}=4.9\pm 1.9`$) than the $`S_\mathrm{N}2.2`$ derived from Miller’s et al. relation for group and cluster dE,N galaxies.
Figure 6 shows the CMD of objects with $`r48`$″ ($`R2.6`$ kpc) which have been marked by open squares. Objects with photometry by Durrell et al., spectroscopically identified foreground stars, background galaxies, and globular clusters are indicated. We determined the mean color of 5 globular clusters (2 spectroscopically confirmed globular clusters are not included in the photometric sample) in this CMD with Maximum-Likelihood techniques. We obtained $`(BV)_{\mathrm{GCS}}=0.82\pm 0.04`$ mag with a dispersion of $`\sigma (BV)_{\mathrm{GCS}}=0.06\pm 0.04`$ mag. Durrell et al. found $`(BV)=0.74\pm 0.03`$ mag, $`\sigma (BV)=0.13`$ mag for the total sample i.e. corresponding to a lower mean metallicity. Our subset seems to be slightly biased towards metal-rich objects.
### 4.2 Comparison of photometrically- and spectroscopically-derived metallicities
In order to constrain the significance of photometrically-derived metallicities we transform the $`(BV)`$ color into a \[Fe/H\]-metallicity and compare it with the findings of our abundance measurements. For this purpose we use the relation of Couture et al. (1990)
$$\text{[Fe/H]}=5.0(BV)_\mathrm{o}4.86$$
(1)
with \[Fe/H\] being the independent parameter during the calibration of the equation. The application of equation 1 to all de-reddened (E$`{}_{(BV)}{}^{}=0.052`$ mag, Schlegel et al., 1998) globular-cluster colors in our NGC 3115 DW1 sample leads to photometrically derived metallicities which can be compared with the \[Fe/H\] values from spectroscopy. All data are summarized in Table 5. The resulting weighted mean metallicity is $``$\[Fe/H\]$`_{\mathrm{GCS}}=0.93\pm 0.11`$ dex (for the 5 globular clusters) with a dispersion of $`\sigma (`$\[Fe/H\]$`)_{\mathrm{GCS}}=0.41\pm 0.20`$ dex in good agreement with the values derived from spectroscopy. The uncertainty results from the photometric error of the color only. No transformation uncertainty was included.
### 4.3 Globular Cluster Sizes
We matched two spectroscopically-confirmed globular clusters, D15 and D25 (see Figure 1 and Table 6), in the HST images taken from the archive (both on WF chips). At the distance of NGC 3115 DW1 the WF chips resolve globular clusters with core radii of $`r_c`$5 pc. Only 14% of Milky-Way globular clusters show core radii larger that 5 pc (Harris, 1999). However, we can use the HST data to derive upper limits for the globular cluster sizes. Only images taken through the F814W filter were used because of their higher S/N.
The radial source profile $`I(r)`$ (i.e. the PSF of the final image) is a convolution of the object profile $`O(r)`$ with the telescope PSF $`T(r)`$ and an additive noise term $`R(r)`$ ($`r`$ being the radial distance from the center of the profile);
$$I(r)=\underset{r_{\mathrm{min}}}{\overset{r_{\mathrm{max}}}{}}O(s)T(r,s)𝑑s+R(r).$$
(2)
To calculate the telescope-PSF profile $`T(r)`$ we used the TinyTim v4.4 code by Krist & Hook (1997) which gives a semi-analytic estimation of the HST-PSF for each chip, each chip position, and each filter. We adopted a King profile (King, 1962) for $`O(r)`$ which appears to be a good fit to globular-cluster radial profiles in Milky Way (e.g. Trager et al., 1995) and extragalactic systems (Grillmair et al. 1996 in M31, Elson & Freeman 1985 in LMC, Kundu & Whitmore 1998 in NGC 3115, Kundu et al. 1999 in M87, and Puzia et al. 1999 in NGC 4472). From an analysis of Milky Way globular clusters (Harris, 1999) we chose $`c\mathrm{log}(r_t/r_c)=1.5`$ as a concentration parameter.
In equation 2 we neglect the additive noise term since our size-estimation errors are dominated by the convolution of the poorly-defined charge-diffusion matrix with the optical HST-PSF (see Krist & Hook, 1997, for a detailed discussion). Note that the charge-diffusion smears 25% of the infalling light of the central pixel among its neighbors. For consistency with other work we continue to use this convolution throughout our analysis despite the fact that the diffusion correction has been derived only for the F555W filter and is thought to be wavelength dependent.
Five King profiles were generated with core radii in the range $`r_c=0.10.5`$ pix ($`R_c=0.52.7`$ pc). In addition, we generated HST-PSFs for both our identified globular clusters using individual specifications (e.g. filter, chip, chip position). Both HST-PSFs were convolved with all the King profiles. Aperture photometry was applied to all generated profiles and both globular clusters on the HST images. For this purpose we used SExtractor (Bertin & Arnouts, 1996) and measured magnitudes in 30 apertures with diameters in the range 1–30 pix. All magnitudes were normalized to the average aperture magnitudes in the range 10–30 pix, i.e. $`I(r)=I_o(r)I_{1030\mathrm{pix}}`$. Figure 7 shows the profiles of both globular clusters, a raw HST-PSF profile, and two of the convolution profiles with core radii of $`r_c=0.2`$ and 0.4 pix.
Both globular cluster profiles deviate significantly from the raw HST-PSF profile which indicates that both D15 and D25 are resolved. The S/N of the F814W image drops to 1 at an aperture diameter of 6 pix. For larger apertures there is insufficient signal to detect any deviations from a raw HST-PSF. For smaller apertures both globular cluster profiles lie between King profiles of core radii 0.2 pix and 0.4 pix. We deduce an upper limit for both globular cluster sizes of $`r_c=2.1_{0.4}^{+0.9}`$ pc at an adopted distance of NGC 3115 DW1 of $`d=11_{2.3}^{+5.0}`$ Mpc (Durrell et al., 1996b). 27% of Milky-Way globular clusters have core radii larger than 2.1 pc (Harris, 1999) and 10% have sizes in the range defined by the errors of the NGC 3115 DW1 clusters.
This upper limit compares well with the results of Kundu & Whitmore (1998) and Kundu et al. (1999). Using HST photometry, these authors find typical half-light radii<sup>4</sup><sup>4</sup>4The half-light radius is comparable with the core radius of a King profile. of $`r_h=2.0\pm 0.1`$ pc and $`r_h2.5`$ pc for globular clusters in NGC 3115 and M87, respectively.
## 5 Discussion
In section 3.1.4, we derived a high globular cluster velocity dispersion, and thus a high galaxy mass, when we included the two outermost globular clusters. The high mass is not unexpected given the bright absolute magnitude of NGC 3115 DW1. Based on its $`M_B=16.8`$ mag (Durrell et al., 1996b), NGC 3115 DW1 falls in the transition region between dwarfs and ellipticals in the mass-luminosity relation of Dekel & Silk (1986) (see Fig. 3 therein). Its high mass ($`M_{\mathrm{PME}}=(4.8\pm 2.3)10^{10}`$M) and the high velocity dispersion ($`\sigma =130\pm 15`$ km s<sup>-1</sup>, see below) are more consistent with an elliptical galaxy. We therefore discuss whether the two outermost clusters could be in the process of being stripped by the nearby giant S0 galaxy NGC 3115.
### 5.1 Possible Stripping?
Figure 1 and Table 6 show that 2 (L1 and L63) of the 7 globular clusters have significantly larger projected radii, i.e. 161.4″ ($`8.6`$ kpc, L1) and 189.4″ ($`10.1`$ kpc, L63), than the “inner” ($`r56.3`$$`=3`$ kpc) globular clusters. These large projected distances from NGC 3115 DW1 could be due to stripping by the nearby S0 galaxy NGC 3115. Figure 8 shows the relative positions of NGC 3115 DW1 and NGC 3115. The projected distance between the two galaxies is 17.3′ which corresponds to 55 kpc at the distance of $`d`$11 Mpc. The mean radial velocities of L1 and L63 are $`v_{\mathrm{rad}}=420\pm 29`$ km s<sup>-1</sup> and $`v_{\mathrm{rad}}=605\pm 74`$ km s<sup>-1</sup>, respectively. Only L1 shows a significant deviation from the systemic velocity of NGC 3115 DW1 ($`v_{\mathrm{rad}}=698\pm 74`$ km s<sup>-1</sup>, Capaccioli et al., 1993) and NGC 3115 ($`v_{\mathrm{rad}}=663\pm 6`$ km s<sup>-1</sup>, Capaccioli et al., 1993).
We expect no contamination from globular clusters of the nearby galaxy NGC 3115. Kavelaars (1997) found the surface density over-abundance of globular clusters around NGC 3115 (power-law index of radial distribution $`\alpha =1.8\pm 0.5`$) disappearing at 6′ radius from the center of NGC 3115 (at a photometric limit of $`V=23.5`$ mag). The globular clusters L1 and L63 have a radial distance to NGC 3115 of $`14`$′. The projected GC surface density of the GCS of NGC 3115 at the position of these two clusters is $`<0.01`$ arcmin<sup>-2</sup>. The extrapolated GC surface density of NGC 3115 DW1 at this position lies between 0.2 and 6.9 GCs arcmin<sup>-2</sup>, given the large uncertainties on the density profile. As the numbers are too small (we only found 2 clusters to the north and 0 to the south) it cannot statistically be concluded whether the two globular clusters found in the northern field are chance detections or a statistically significant overabundance.
Assuming that both galaxies are roughly at the same distance, we can estimate the dwarf galaxy’s gravitational potential and the ratio of potentials of NGC 3115 DW1 and NGC 3115. Both globular clusters are at about 1/5 of the distance separating NGC 3115 DW1 and NGC 3115. As a rough estimate, we assume that NGC 3115 DW1 and NGC 3115 have similar M/L<sub>V</sub>. In this case, the ratio of their $`M_V`$’s would imply that NGC 3115 has a mass 10 times larger than NGC 3115 DW1. Hence, the gravitational potentials are comparable at the projected position of the distant globular clusters (L1 and L63). Since the mass of NGC 3115 is likely to be higher than the adopted value (assuming an extended dark matter halo) the motion of both globular clusters is no longer dominated by the gravitational potential of NGC 3115 DW1 alone. Both clusters could then be considered as intergalactic globular clusters.
Note that stripping of globular clusters appears to be common among interacting galaxies. Da Costa & Armandroff (1995) show that four globular clusters of the Sagittarius dSph are in the process of being stripped by the Milky Way and are being added to its globular cluster system. Other studies have indicated that stripping may be important in galaxy clusters (e.g. in the Fornax cluster Kissler-Patig et al. 1999, Hilker et al. 1999). However, there are no other (optical) hints of interaction from NGC 3115 DW1’s stellar light. Durrell et al. (1996b) found the isophotes to be consistent with little or no tidal disruption out to a projected radius of 60″ (corresponding to 3.2 kpc) where their photometric errors start to dominate.
A simple test for the stripping hypothesis would be a wide-field study of the system in order to rule out (spectroscopically) the presence of any similar clusters around NGC 3115 DW1.
### 5.2 The Expected Velocity Dispersion
A look at the fundamental plane of dwarf elliptical galaxies (e.g. Peterson & Caldwell, 1993) shows that NGC 3115 DW1 fits reasonably well into the relation for dwarf and giant elliptical galaxies, under their assumption of $`M_V=16.7`$ mag. Adopting the absolute magnitude of $`M_V=17.7`$ mag (Durrell et al., 1996a) the galaxy falls slightly off the relation and would imply a higher velocity dispersion than measured in the central 3″. With the measured velocity dispersion of $`\sigma =74\pm 36`$ km s<sup>-1</sup> for the 5 globular clusters inside $`r<56.3`$″ ($`R<3`$ kpc) we obtain from the fundamental-plane relation of Peterson & Caldwell (1993) an absolute magnitude of $`M_V=18.0\pm 0.5`$ mag. The measured velocity dispersion for the total sample of 7 clusters inside $`r<189.4`$″ ($`R<10.1`$ kpc) of $`\sigma =130\pm 15`$ km s<sup>-1</sup>, would correspond to far brighter absolute magnitude ($`M_V=19.5\pm 0.5`$ mag) than the measured $`M_V=17.7`$ mag. This discrepancy can be explained by a close encounter and subsequent stripping of the dwarf galaxy’s halo by the nearby S0 galaxy NGC 3115. Stripping of outer halo regions might well have introduced violent perturbations and led to an enhanced velocity dispersion of the halo region (which is traced by the globular clusters). The relaxation time of such a system far exceeds the Hubble time (Binney & Tremaine, 1994) and therefore it is not possible to reject this scenario just from considerations of dynamical timescales.
Alternatively, a high velocity dispersion in the outskirts of a galaxy could be due to a dark-matter dominated massive halo. The outer parts of a number of lower-luminosity Local Group galaxies are known to be dominated by dark matter (e.g. Mateo, 1998). This picture could explain the fact that we measure an uncommonly high mass for a dwarf elliptical (see sec. 3.1.2) at a projected radius of 189.4″ ($`10.1`$ kpc).
We cannot discriminate between the above possibilities at this point.
## 6 Summary
Using LRIS multi-slit spectra we confirm, on the basis of their radial velocities, 7 of the 46 objects in our spectroscopic sample as bona-fide globular clusters associated with the bright ($`M_V=17.7`$ mag) dE1,N galaxy, NGC 3115 DW1.
We verify the findings of Durrell et al. (1996b) (within a projected radius of $`r48`$″ corresponding to $`R2.6`$ kpc) who derived the specific frequency $`S_\mathrm{N}=4.9\pm 1.9`$ and a total globular cluster population size $`N_{\mathrm{GC}}=59\pm 23`$. The spectroscopic verification of foreground and background contamination indicates that no revision of these results is necessary. NGC 3115 DW1 remains a dE,N galaxy with a somewhat high $`S_\mathrm{N}`$ value.
A mass estimate using the Projected Mass Estimator (PME) for isotropic globular cluster orbits yields a total galaxy mass of $`M_{\mathrm{gal}}=(4.8\pm 2.3)10^{10}M_{}`$ (with the error being the internal uncertainty of the mass estimation code) inside a radius $`r189.4`$″ ($`R10.1`$ kpc) and $`M_{\mathrm{gal}}=(1.8\pm 1.0)10^{10}M_{}`$ inside $`r56.3`$″ ($`R3.0`$ kpc). This estimate is a lower mass limit (see 3.1.2) and assumes that the outer globular clusters are not influenced by the nearby giant S0 NGC 3115. Using two mass estimators (i.e. PME and VME, see Sect. 3.1.2) and various assumptions for the systemic velocity and the cluster orbits, we derive masses between $`210^{10}`$M and $`410^{11}`$ M. The mass increases with radius (see Table 3). Inside R$`<160`$ pc the mass-to-light ratio was found to be M/L$`{}_{V}{}^{}=3\pm 2`$ (Peterson & Caldwell, 1993) and increases with radius, leading to M/L$`{}_{V}{}^{}=52\pm 25`$ at $``$ 10 kpc (using the PME and assuming isotropic orbits), compatible with dark matter dominated outer regions. However, we cannot at present exclude the possibility that the high velocity dispersion is due to stripping of the two outer clusters by the nearby giant companion.
A kinematic analysis shows that the globular cluster system has a marginal net rotation of $`v_{\mathrm{rot}}=75\pm 70`$ km s<sup>-1</sup> with a position angle of the rotation axis $`\theta =90^\mathrm{o}\pm 60^\mathrm{o}`$. Subtracting the net rotation we find a line-of-sight velocity dispersion of the globular cluster system of $`\sigma =130\pm 15`$ km s<sup>-1</sup> for the total sample of 7 globular clusters and $`\sigma =74\pm 36`$ km s<sup>-1</sup> for the inner 5 clusters (see Sect. 3.1.4).
We measure mean abundances (using Lick/IDS passband definitions) from a combined mean spectrum of all 7 globular clusters and derive a mean GCS metallicity of $``$\[Fe/H\]$`_{\mathrm{GCS}}=0.97\pm 0.11`$ dex. All abundance ratios appear similar to the ones measured in Milky Way, M31 and NGC 1399 globular clusters.
The mean color of the spectroscopically confirmed globular clusters is $`(BV)_{\mathrm{GCS}}=0.82\pm 0.04`$ mag with a dispersion $`\sigma (BV)_{\mathrm{GCS}}=0.06\pm 0.04`$ mag.
Applying the color-metallicity calibration of Couture et al. (1990) we obtain a photometric mean metallicity $``$\[Fe/H\]$`_{\mathrm{GCS}}=0.93\pm 0.11`$ dex with a dispersion of $`\sigma (\text{[Fe/H]})_{\mathrm{GCS}}=0.41\pm 0.20`$ dex (the error being the photometric uncertainty).
For two globular clusters (L1 and L63) with HST photometry we derive upper limits for their core radii. These were found to be $`r_c=2.1_{0.4}^{+0.9}`$ pc.
We would like to thank Patrick Durrell for providing his ground-based photometry in electronic form. We also thank Duncan Forbes for his help during the observation. Some of the data presented herein were obtained at the W.M. Keck Observatory, which is operated as a scientific partnership among the California Institute of Technology, the University of California and the National Aeronautics and Space Administration. The Observatory was made possible by the generous financial support of the W.M. Keck Foundation. This work was supported by National Science Foundation grant number AST990732 and Faculty Research funds of the University of California, Santa Cruz.
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# Two photon background for Higgs boson searches at the LHC
## 1 Introduction
The understanding of electroweak symmetry breaking is a major motivation for present and future collider physics. Inside the Standard Model (SM) the LHC is expected to detect the predicted Higgs boson, if its mass is below say 800 GeV. On one hand the direct search at the LEP experiments place a stringent lower bound whereas on the other hand the precision measurements indicate an upper bound for the mass of the Higgs boson $`^\mathrm{?}`$:
$$107.7\text{GeV}<M_{Higgs}<188\text{GeV}\text{(95\% c.l.)}$$
(1)
Also for the light Higgs boson of the minimal supersymmetric standard model a lower bound of 88 GeV exists, where on the contrary the upper bound of around 130 GeV is dictated by the structure of the interaction terms. The quartic couplings are related to the gauge couplings here, not a free parameter as in the SM. The most promising signal for such neutral Higgs bosons in the mass range of $`80\text{GeV}<M_{Higgs}<140\text{GeV}`$ comes from the gluon fusion process with a subsequent decay of the Higgs boson into a photon pair. Though the branching fraction is small, $`B(H\gamma \gamma )2\times 10^3`$ in the SM, one expects a sharp peak in the invariant mass distribution of the photon pair, because of the narrow Higgs width of a few MeV. Though the signal is very sharp one is confronted with a huge background. To get a quantitative understanding of the signal to background ratio one has to know the background as best as one can.
The background can be split into three contributions:
Both photons originating from the hard partonic interaction.
At least one of the photons is created in the hadronization of a QCD parton.
The photons are decay products of mesons like $`\pi ^0`$, $`\eta `$, and so on.
Whereas the latter is essentially reducible the former are commonly called irreducible backgrounds though one can reduce the fragmentation contribution substantially by imposing isolation criteria, i.e. one restricts the hadronic energy inside a cone (defined in rapidity and azimuthal angle space) around the photon momentum to be less than some value.
In the following the calculation of this background at full next-to-leading order is shortly sketched. The inclusion of the NLO fragmentation contributions is a step beyond the calculations of $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$. A prediction of the background $`M_{\gamma \gamma }`$ distribution at LHC is presented and its stability under variation of the different scales entering the computation is discussed.
## 2 Photon pair production in hadronic collisions at NLO
The Born term for photon pair production, $`q\overline{q}\gamma \gamma `$, is of order $`\alpha ^2`$. The radiative corrections consist out of all virtual and real emission corrections which lead to an order $`\alpha ^2\alpha _s`$ correction. One observes that in subprocesses like $`qg\gamma \gamma q`$ a final state singularity is present if the photon is collinear to the quark. To absorb these singularities into a photon fragmentation function $`D_{\gamma /q}`$ one has subsequently also to take into account other partonic reactions like e.g. $`qg\gamma q`$. As $`D_{\gamma /q}`$ behaves like $`\alpha /\alpha _s`$ in the investigated kinematic regime, the power counting of the couplings is the same as the one of the Born term. As both photons can be collinear to external partons also both photons should be allowed to be created in hadronization. Two–fragmentation processes are needed for consistency together with the respective higher order corrections, see $`^\mathrm{?}`$ for more details. Finally there is another sizable contribution, $`gg\gamma \gamma `$, where the gluons/photons are attached to a fermion loop. Formally a higher order correction, this contribution is not negligible due to the large gluon flux at the LHC.
The diagrammatic calculation depends on three unphysical scales:
Renormalization scale due to ultraviolet divergence,
Factorization scale due to QCD initial state collinear singularities,
Fragmentation scale due to final state collinear singularities.
The relation between the hadronic and the partonic cross sections for a fragmentation process is done by folding the partonic cross sections with the $`\mu _{FACT}`$-dependent parton distribution functions and the $`\mu _{FRAG}`$-dependent photon fragmentation functions $`D_{\gamma /j\{q,\overline{q},g\}}`$.
## 3 The $`M_{\gamma \gamma }`$ distribution at the LHC
The comparison of our calculation with recent Tevatron data is showing a good agreement $`^\mathrm{?}`$. This is especially true for infrared save quantities as the $`M_{\gamma \gamma }`$ distribution <sup>a</sup><sup>a</sup>aInfrared sensitive observables are showing also good agreement for the tails of the distributions. In the infrared sensitive region resummation of large logarithms is needed $`^\mathrm{?}`$.. In Fig. 1 we plot the invariant mass distribution of a photon pair at the LHC with and without applying isolation cuts. As an example we demand the transversal hadronic energy to be less than $`E_{Tmax}=5`$ GeV in a cone $`R=0.4`$ around the photon defined in rapidity and azimuthal angle space through $`R=\sqrt{(yy_\gamma )^2+(\varphi \varphi _\gamma )^2}`$. Apart from these, standard cuts were taken as indicated in the figure. The scale choice is $`\mu =\mu _{FACT}=\mu _{FRAG}=M_{\gamma \gamma }/2`$. For all figures the MRST2 set of parton distribution functions $`^\mathrm{?}`$ and the BFG photon fragmentation functions $`^\mathrm{?}`$ were used.
To get an idea of the residual scheme dependence of our result for the lower curve in Fig. 1 (with isolation), we varied the renormalization and factorization scales and compared to that curve by plotting $`[d\sigma /dM_{\gamma \gamma }(\mu ,\mu _{FACT})d\sigma /dM_{\gamma \gamma }(M_{\gamma \gamma }/2,M_{\gamma \gamma }/2)]/[d\sigma /dM_{\gamma \gamma }(M_{\gamma \gamma }/2,M_{\gamma \gamma }/2)]`$ in Fig. 2.
Due to isolation the fragmentation scale dependence is marginal <sup>b</sup><sup>b</sup>bThis is because we consider the fragmentation contributions in the calculation. Thus the result is of beyond leading logarithmic accuracy. and is kept equal to $`\mu _{FACT}`$ for simplicity. In Fig. 2 one observes that there is an accidental stability if $`\mu ,\mu _{FACT}`$ are changed in the same directions from $`\mu =\mu _{FACT}=M_{\gamma \gamma }/2`$ to $`\mu =\mu _{FACT}=2M_{\gamma \gamma }`$. This is due to the fact that at LHC for the given parton energies the effect of larger (smaller) $`\mu `$ – means smaller (larger) NLO corrections – is compensated by the increased (decreased) gluon flux if $`\mu _{FRAG}`$ is also increased (decreased). By varying in an anti-diagonal way one gets a more conservative estimation of the uncertainty. One observes variations between 10 and 20 per cent.
## 4 Conclusion
A calculation for the photon pair production at hadron colliders at full NLO exists and is implemented into a partonic event generator, called DIPHOX $`^\mathrm{?}`$. Tevatron data are well described by the code and predictions for the LHC show a residual scale dependence of the order of 10 to 20 per cent <sup>c</sup><sup>c</sup>cTo get a further stabilization at least the dominant processes would have to be calculated one order higher. E.g. the corrections to the box contribution have to be included but also the corrections to the gluon induced subprocesses.. For infrared sensitive observables, like e.g. $`P_T`$ distributions, resummation has to be included in the calculation. In contrary the $`M_{\gamma \gamma }`$ distribution is now known at NLO not only in the direct but also in the fragmentation part which allows more reliable quantitative studies of effects due to isolation criteria.
## Acknowledgments
I would like to thank my collaborators J. Ph. Guillet, E. Pilon and M. Werlen for giving me the opportunity to present our results at the Moriond conference.
This work was supported by the EU Fourth Training Programme ”Training and Mobility of Researchers”, Network ”Quantum Chromodynamics and the Deep Structure of Elementary Particles”, contract FMRX–CT98–0194 (DG 12 - MIHT).
## References
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# Effect of order-parameter fluctuations on the Halperin-Lubenski-Ma first-order transition in superconductors and liquid crystals
## Abstract
We show that order-parameter fluctuations in a good type-I superconductor or a liquid crystal always increase the size of the first-order transition. This behavior is eventually changed when the system crosses over to inverted-XY critical behavior, with the size of the first-order transition vanishing as a power law with a crossover exponent. We find good agreement between our theory and a recent experiment on the nematic-smectic-A first-order transition in 8CB-10CB mixtures of liquid crystals.
More than 25 years ago, Halperin, Lubensky and Ma (HLM) and Coleman and Weinberg demonstrated that when a scalar field is coupled to a gauge field, the fluctuations of the gauge field can change the nature of the phase transition in the theory from continuous to first order. The coupling of a scalar field (or an order parameter (OP)) to a massless gauge field arises often in physics: the Meissner transition in superconductors , the Higgs mechanism in particle physics , and the nematic-smectic-A transition in liquid crystals are among the best-known examples , . The analysis of HLM initially centered on the fluctuations of the gauge field and neglected the OP fluctuations, which is justifiable for good type-I superconductors (with Ginzburg-Landau parameter $`\kappa 1`$), where the size (to be precisely defined shortly) of the first-order transition is larger. This approximation, however, inevitably breaks down for strong type-II materials ($`\kappa 1`$), where neglecting the OP fluctuations would yield a first-order transition well into the critical region. Close to four dimensions, the effect of the OP fluctuations can be studied using the Wilson-Fisher renormalization group , which, however, in this case leads only to “run-away” flows and no stable critical points. This is usually interpreted as a sign of a first-order transition , . Today, based on accumulated analytical and numerical evidence, it is generally believed that the transition for $`\kappa 1`$ is again second order, in the so-called inverted-XY universality class . The corresponding topology of the flow of the coupling constants under scaling transformation is depicted in Fig. 1. It may thus seem natural to expect that the OP fluctuations should decrease the size of the first-order transition, finally reducing it to zero at the crossover to inverted-XY critical behavior.
The fluctuation effects in question are unfortunately too fine to be observable even in high-$`T_c`$ superconductors on account of the smallness of the fine-structure constant . But in liquid crystals, the coupling of the smectic OP to the director fluctuations is stronger, and the fluctuations at the nematic-smectic-A transition become an issue of central importance. A recent experiment on two-component liquid-crystal mixtures found a surprising result: the size of the first-order transition is larger than the prediction of the HLM theory. Motivated by this result, in this Letter we consider theoretically the effect of OP fluctuations on the first-order transition in a type-I material. We show quite generally that for a small enough Ginzburg-Landau parameter $`\kappa `$, the size of the first-order transition is indeed always larger than the HLM result. Crudely, the reason is that for a good type-I material, the effective Ginzburg-Landau parameter $`\kappa `$ always decreases at large scales, making the material only more type I and the transition more strongly first order. The type-I region is defined here as left of the separatrix in Fig. 1, where the flow is qualitatively the same as that of HLM near four dimensions. This behavior is eventually changed as the separatrix between type-I (first-order) and type-II (inverted-XY) regimes is approached, where the size of the transition goes to zero as a power law, with a crossover exponent that characterizes the tricritical point. Using the simple one-loop recursion relations for the coupling constants in an isotropic version of the
Ginzburg-Landau-Wilson action, we show that the experimental data of Yethiraj and Bechhoefer are fit rather well by our theory.
We are interested in a general phase transition described by the three-dimensional Ginzburg-Landau-Wilson theory for a fluctuating complex OP $`\mathrm{\Psi }(\stackrel{}{r})`$ minimally coupled to a fluctuating vector potential (Higgs scalar electrodynamics):
$`H={\displaystyle }d^3\stackrel{}{r}[|(ie\stackrel{}{A}(\stackrel{}{r}))\mathrm{\Psi }(\stackrel{}{r})|^2+{\displaystyle \frac{a(T)}{2}}|\mathrm{\Psi }(\stackrel{}{r})|^2`$ (1)
$`+{\displaystyle \frac{b}{2}}|\mathrm{\Psi }(\stackrel{}{r})|^4+{\displaystyle \frac{c}{6}}|\mathrm{\Psi }(\stackrel{}{r})|^6+{\displaystyle \frac{1}{2}}(\times \stackrel{}{A}(\stackrel{}{r}))^2],`$ (2)
where $`a(T)=\alpha (TT^{})/T^{}`$, and we assume the superconducting gauge, $`\stackrel{}{A}=0`$. For simplicity, we neglect the anisotropy inherent to liquid crystals but retain the sixth-order term in the Ginzburg-Landau expansion, which will be needed later for comparison with the experiment. Fields and lengths have been chosen so that the number of couplings in the theory is minimal. HLM showed that neglecting the fluctuations of the OP and integrating out the vector field $`\stackrel{}{A}`$ around the uniform OP configuration $`\mathrm{\Psi }_0`$, to the lowest order in charge $`e`$ one obtains the corrected mean-field Ginzburg-Landau free-energy per unit volume:
$$F=\frac{a(T)}{2}\mathrm{\Psi }_0^2\frac{e^3\sqrt{2}}{3\pi }\mathrm{\Psi }_0^3+\frac{b}{2}\mathrm{\Psi }_0^4+\frac{c}{6}\mathrm{\Psi }_0^6,$$
(3)
in units where $`k_BT^{}=1`$. The negative cubic term in the free-energy implies a first-order transition. A useful measure of the size of the first-order transition is the parameter $`t=(T_cT^{})/T^{}`$, where $`T_c`$ is the first-order transition temperature. To keep the algebra simple, we will assume $`b>0`$ and temporarily set $`c=0`$ in Eq. 3, and turn $`c`$ back to a finite value only later when we compare our results with the experiment. From Eq. 3, the size of the transition is then
$$\alpha t_{HLM}=\frac{e^4}{(6\pi )^2}F(\kappa ),$$
(4)
where $`\kappa ^2=b/2e^2`$ is the dimensionless Ginzburg-Landau parameter, and $`F(\kappa )=1/\kappa ^2`$ . For $`\kappa 1`$, the first-order transition thus occurs at higher temperatures, where the OP fluctuations may indeed be neglected. The HLM result becomes asymptotically correct in this limit.
To include the OP fluctuations, we assume that the first-order transition in the theory (Eq. 2) with the coupling constants $`e`$ and $`b`$ occurs at some $`tt_{HLM}`$. The renormalizability of the theory (Eq. 2) implies that rescaling the cutoff $`\mathrm{\Lambda }\mathrm{\Lambda }/s`$ by an arbitrary factor $`s>1`$ is equivalent to changing the coupling constants into $`e(s)`$ and $`b(s)`$ , with the first-order transition for these renormalized couplings occurring at some new temperature $`t_0`$. The variation of the renormalized couplings $`e(s)`$, $`b(s)`$, and the temperature $`t(s)`$ with $`s`$ is described by differential recursion relations of the form
$$\frac{d\lambda (s)}{d\mathrm{ln}(s)}=\beta _\lambda [e^2(s),b(s),t(s)],$$
(5)
where $`\lambda =\{e^2,b,t\}`$. If $`t_0`$ is large enough, the OP fluctuations at the first-order transition in the rescaled theory indeed become negligible, and $`t_0`$ may be approximated by the mean-field (HLM) expression:
$$\alpha t_0\frac{e^4(s)}{(6\pi )^2}F(\kappa (s)).$$
(6)
Together with Eqs. 4 and 5, and the boundary conditions $`t(1)=t`$ and $`t(s)=t_0`$, Eq. 6 determines implicitly the actual size of the transition $`t`$. Clearly, the above idea is quite general and applicable to other weakly first-order transitions. First, let us demonstrate that neglecting the interactions between fluctuations by renormalizing the coupling constants only according to dimensional analysis gives just the HLM result. Power counting in Eq. 2 implies that $`t(s)=ts^2`$, $`e^2(s)=e^2s`$, $`b(s)=bs`$. Thus, $`\kappa (s)=\kappa `$, and dividing Eqs. 4 and 6 gives $`t=t_{HLM}`$. More generally, for small $`e^2`$ and $`b`$, and for $`te^2,b`$, the renormalized couplings obey the differential equations (Eq. 5) with :
$$\beta _t=t(2+ue^2vb+𝒪(e^4,b^2,be^2)),$$
(7)
$$\beta _e=e^2xe^4+𝒪(e^6),$$
(8)
$$\beta _b=byb^2+zbe^2we^4+𝒪(b^3,e^6,be^4,b^2e^2).$$
(9)
The signs in the above equations are chosen so that the numerical coefficients $`u`$, $`v`$, $`x`$, $`y`$, $`z`$ and $`w`$ are positive. Their values in principle are non-universal and weakly dependent on renormalization procedure. The differential recursion relation for the Ginzburg-Landau parameter $`\kappa `$, for $`\kappa 1`$, may be easily obtained from Eqs. 8 and 9:
$$\frac{d\kappa ^2}{d\mathrm{ln}(s)}=\frac{w}{2}e^2+𝒪(e^4,e^2\kappa ^2).$$
(10)
Note that the right-hand side in the last equation is negative, and $`\kappa `$, if initially small enough, always decreases under renormalization. This is a consequence of the $`e^4`$ term in Eq. 9, which generates a negative quartic coupling and tends to drive the transition first order. The same term is responsible for the runaway flows near four dimensions in the original HLM analysis. For small charge $`e`$, one may neglect the nonlinear terms and solve Eqs. 7 and 8 to find, for small $`\kappa `$,
$$\frac{t}{t_{HLM}}=\frac{\kappa ^2}{\kappa ^2(\sqrt{t_0/t})}.$$
(11)
Since $`\kappa (s)<\kappa `$, the solution to the above equation always gives $`t>t_{HLM}`$. In sum, while near the Gaussian fixed point $`t`$ and $`e^2`$ approximately scale according to dimensional analysis, a small Ginzburg-Landau parameter $`\kappa `$ acquires a small negative dimension and renormalizes downwards, thus increasing the size of the first-order transition.
Assuming the general topology of the flow diagram as shown in Fig. 1 the size of the first-order transition will go to zero as the coupling constants approach the separatrix between type-I and type-II regimes. If the separatrix at the critical surface $`t=0`$ lies at some $`b_c(e)`$, close to the separatrix one finds
$$t\left[b_c(e)b\right]^{1/\varphi },$$
(12)
where $`\varphi =r\nu `$ is the crossover exponent, with $`\nu ^1=\beta _t/t`$, taken at the tricritical fixed point. Here, $`r`$ is the (positive) scaling dimension of the second relevant scaling variable at the tricritical point, which in general is a linear combination of $`e^2`$, $`b`$, and $`t`$. Near four dimensions, our result agrees with that of HLM: the transition is then always first order, unless $`e=0`$. The role of the tricritical point near four dimensions is played by the $`XY`$ fixed point, which has two relevant directions, $`t`$ and $`e^2`$, with $`\nu =\nu _{xy}`$ and $`r=\epsilon `$, where $`\epsilon =4d`$. Instead of Eq. 12, one then has $`te^{2/(\epsilon \nu _{xy})}`$, which is just Eq. 17 in .
By tuning the concentration $`x`$ of 10CB liquid crystals in 8CB, one can vary the temperature range of the nematic phase, and, in effect, tune the parameter $`b`$ in Eq. 1. For larger concentrations, $`b`$ is negative and the transition is more strongly first order. The latent-heat data for $`x>0.42`$ are well-fit by mean-field theory (Eq. 3) with the parameters $`c=1`$, $`\alpha =3.35`$, $`b=b_0(x^{}x)`$, with $`b_0=0.395`$ and $`x^{}=0.42`$, and $`e^2=0.0421`$ . For $`x<x^{}`$, the quartic term in Eq. (3) becomes positive, but the transition remains first order, in agreement with the HLM theory. At smaller concentrations, $`t`$ continues to decrease, as expected, but there is a clear deviation from the HLM result . To attempt to fit the data in the whole concentration range $`0<x<0.65`$, we take the above parameters to set the initial values of the couplings, and then evolve $`t`$, $`e`$, and $`b`$ according to Eqs. 7-9, with $`u=v=1/4`$, $`x=1/16`$, $`y=(2\sqrt{2}+1)/8`$, $`z=1/2`$, and $`w=1/(2\sqrt{2})`$ . We neglect the change of $`c`$ under scaling, since the results are quite insensitive to its precise value. The only free parameter left is the final value of the temperature $`t_0`$ at which it becomes safe to neglect the OP fluctuations. We chose $`t_0=0.01`$, an order of magnitude larger than the largest measured $`t`$ in the experiment. The quality of the fit turns out not to be critically dependent on this choice. The fit to our theory and the comparison with HLM is shown on Fig. 2. Note that for $`x>0.35`$, there is very little change from the HLM result , but as $`b`$ increases at smaller concentrations, deviations from HLM become significant, in agreement with the experiment . One expects that if one could increase $`b`$ further, eventually $`t`$ would go to zero in accord with Eq. (12). The full non-perturbative structure of the flow diagram in Fig. 1, with the tricritical and the inverted-XY fixed points, however, is beyond the simple one-loop $`\beta `$-functions we used . The data suggest that, at small charge, the tricritical point is likely to lie at some $`\kappa _c>\sqrt{2}`$, which corresponds approximately to the smallest concentration used in the experiment ($`x=0`$).
Although anisotropy is known to be important in liquid crystals , its inclusion, besides introducing two new couplings, does not essentially change the structure of Eqs. 7\- 9 . In particular, the sign of the $`e^4`$ term in Eq. 9 stays negative, and our main point that the size of the transition in good type-I materials is larger than the mean-field prediction remains valid. Some small quantitative differences from our result would be expected, however.
In conclusion, we have shown that the naive expectation that order-parameter fluctuations would always decrease the strength of the first-order transition predicted by Halperin, Lubensky, and Ma is incorrect for good type-I materials. While the transition strength does decrease monotonically as the material becomes more type II, it becomes smaller than the HLM prediction only near the tricritical point that separates type-I and type-II regimes. By using one-loop renormalization group to calculate the evolution of Landau coefficients in the free energy, we can account for the deviations from HLM predictions that were observed in a recent experiment.
Our arguments also lead one to expect a non-trivial exponent near the tricritical point where the transition eventually should become second order. One would expect this point to be accessible experimentally. For example, one can view the variation of concentration $`x`$ in the 8CB-10CB experiments as a way of continuously adjusting the effective molecular length (between 8 and 10 CH<sub>2</sub> lengths). Mixing 8CB with a shorter molecule (for example, 6CB ) might then give access to the tricritical point and allow at least a crude measurement of the crossover exponent $`\varphi `$.
This work has been supported by NSERC (Canada). IFH was also supported by an award from Research Corporation.
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# Bose–Einstein condensation in trapped dipolar gases
## Abstract
We discuss Bose–Einstein condensation in a trapped gas of bosonic particles interacting dominantly via dipole–dipole forces. We find that in this case the mean–field interparticle interaction and, hence, the stability diagram are governed by the trapping geometry. Possible physical realisations include ultracold heteronuclear molecules, or atoms with laser induced electric dipole moments.
Bose–Einstein condensation (BEC) of trapped atomic gases offers unique possibilities to highlight a general physical problem of how the nature and stability of a Bose–condensed state is influenced by the character of interparticle interaction. In this respect, especially interesting are ultra–cold gases with attractive interaction between particles (scattering length $`a<0`$). As known , spatially homogeneous condensates with $`a<0`$ are absolutely unstable with regard to local collapses. The presence of the trapping field changes the situation drastically. This has been revealed in the successful experiments at Rice on BEC of magnetically trapped atomic <sup>7</sup>Li ($`a=14`$ Å). As found in theoretical studies , if the number of Bose–condensed particles is sufficiently small (of order $`10^3`$ in the conditions of the Rice experiments) and the spacing between the trap levels exceeds the mean–field interparticle interaction $`n_0|g|`$ ($`n_0`$ is the condensate density, $`g=4\pi \mathrm{}^2a/M`$, where $`M`$ is the atom mass), there will be a metastable Bose–condensed state. In other words, the condensate is stabilized if the negative pressure caused by the interparticle attraction is compensated by the quantum pressure imposed by the trapping potential. In some sense, this is similar to the gas–liquid phase transition in a classical system with interparticle attraction: The gas phase is stable as long as the thermal pressure exceeds the (negative) interaction–induced pressure (see ).
The recent success in creating ultra–cold molecular clouds opens fascinating prospects to achieve quantum degeneracy in trapped gases of heteronuclear molecules. In a sufficiently high electric field ”freezing” their rotational motion, these molecules interact via the dipole–dipole forces. This interaction is long–range and anisotropic (partially attractive), and there is a non–trivial question of achieving BEC and manipulating condensates in trapped gases of dipolar particles.
Thus far, only the interaction between (small) atomic dipoles has been included in the discussion of the condensate properties. Góral et al. considered the effect of magnetic dipole interaction in a trapped spin–polarized atomic condensate. Magnetic dipoles are small (of the order the Bohr magneton $`\mu _B`$), and even for atoms like Chromium ($`6\mu _B`$) the magnetic interactions are dominated by the Van der Waals forces. Nevertheless, for a relatively small scattering length $`a`$ the condensate wave function may develop novel structures reflecting the interplay between the two types of forces. These effects can be amplified by modyfing $`a`$, which hopefully will soon become a standard technique , and could eventually appear in other systems, such as polar molecules, as pointed out in Ref. . Similar effects have been discussed by Yi and You for ground–state atoms with electric dipole moments induced by a high dc field (of the order of 10<sup>6</sup>V/cm). These authors have demonstrated the validity of the Gross–Pitaevskii equation (GPE) for this system, constructed the corresponding pseudopotential, and determined an effective scattering length.
In this Letter we discuss BEC in a trapped gas of dipolar particles, where the interparticle interaction is dominated by the dipole–dipole forces. Possible realizations include the (electrically polarized) gas of heteronuclear molecules as they have large permanent electric dipoles. We also propose a method of creating a polarized atomic dipolar gas by laser coupling of the atomic ground state to an electrically polarized Rydberg state. Similarly to the condensates with $`a<0`$, dipolar condensates are unstable in the spatially homogeneous case, and can be stabilized by confinement in a trap. However, we find a striking difference from common atomic condensates: In the BEC regime the sign and the value of the dipole–dipole interaction energy in the system is strongly influenced by the trapping geometry and, hence, the stability diagram depends crucially on the trap anisotropy. This offers new possibilities for controlling and engineering macroscopic quantum states. Remarkably, for dipoles oriented along the axis of a cylindrical trap we have found a critical value $`l_{}=0.4`$ for the ratio of the radial to axial frequency $`l=(\omega _\rho /\omega _z)^{1/2}`$: Pancake traps with $`l<l_{}`$ mostly provide a repulsive mean field of the dipole–dipole interaction, and thus the dipolar condensate in these traps will be stable at any number of particles $`N`$. For $`l>l_{}`$ the stability of the condensate requires $`N<N_c`$, where the critical value $`N_c`$ at which the collapse occurs is determined by the condition that (on average) the mean–field interaction is attractive and close to $`\omega _\rho `$.
We consider a condensate of dipolar particles in a cylindrical harmonic trap. All dipoles are assumed to be oriented along the trap axis. Accordingly, the dipole–dipole interaction potential between two dipoles is given by $`V_d(𝐑)=(d^2/R^3)(13\mathrm{cos}^2\theta )`$, where $`d`$ is the dipole moment, $`𝐑`$ the distance between the dipoles, and $`\theta `$ the angle between the vector $`𝐑`$ and the dipole axis. The dipole–dipole interaction is long–range, and one can no longer use the pseudopotential approximation for the mean field. Similar to , we describe the dynamics of the condensate wave function $`\psi (𝐫,t)`$ by using the time–dependent GPE
$`i\mathrm{}{\displaystyle \frac{}{t}}\psi (\stackrel{}{r},t)=\{{\displaystyle \frac{\mathrm{}^2}{2m}}^2+{\displaystyle \frac{m}{2}}(\omega _\rho ^2\rho ^2+\omega _z^2z^2)+\mathrm{\Gamma }0`$ (1)
$`\mathrm{\Gamma }0+g|\psi (\stackrel{}{r},t)|^2+d^2{\displaystyle }d\stackrel{}{r}^{}{\displaystyle \frac{13\mathrm{cos}^2\theta }{|\stackrel{}{r}\stackrel{}{r}^{}|^3}}|\psi (\stackrel{}{r}^{},t)|^2\}\psi (\stackrel{}{r},t).`$ (2)
Here $`\psi (𝐫,t)`$ is normalized to the total number of condensate particles $`N`$. The third term on the rhs corresponds to the mean–field interaction due to (short–range) Van der Waals forces, and the last term to the mean field of the dipole–dipole interaction. Assuming that the interparticle interaction is mostly related to the dipole–dipole forces and $`d^2|g|=4\pi \mathrm{}^2|a|/M`$, we omit the (Van der Waals) term $`g|\psi (𝐫,t)|^2\psi (𝐫,t)`$.
The wave function of the relative motion of a pair of dipoles is influenced by the dipole–dipole interaction at interparticle distances $`|𝐫𝐫^{}|r_{}=Md^2/\mathrm{}^2`$. This influence is ignored in the dipole–dipole term of Eq.(2), as the main contribution to the integral comes from distances $`|𝐫𝐫^{}|`$ of order the spatial size of the condensate, which we assume to be much larger than $`r_{}`$.
As mentioned above, in the spatially homogeneous case the dipolar condensate is unstable. For all dipoles parallel to each other, by using the Bogolyubov method one easily finds the anisotropic dispersion law for elementary excitations: $`\epsilon (𝐤)=[E_k^2+8\pi E_kn_0d^2(1/3\mathrm{cos}^2\theta _k)]^{1/2}`$, where $`E_k=\mathrm{}^2k^2/2M`$, $`n_0`$ is the condensate density, and $`\theta _k`$ the angle between the excitation momentum $`𝐤`$ and the direction of the dipoles. The instability is clearly seen from the fact that at small $`k`$ and $`\mathrm{cos}^2\theta _k>1/3`$ one has imaginary excitation energies $`\epsilon `$.
To understand the influence of the trapping field on the behavior of the dipolar condensate, we have numerically simulated Eq.(2) for various values of the number of particles $`N`$, dipole moment $`d`$, and the trap aspect ratio $`l`$. By evolving Eq.(2) in imaginary time, we have found the condition under which the condensate is stabilized by the trapping field and investigated static properties of this (metastable) Bose–condensed state.
For the stationary condensate the wave function $`\psi (𝐫,t)=\psi _0(𝐫)\mathrm{exp}(i\mu t/\mathrm{})`$, where $`\mu `$ is the chemical potential, and the lhs of Eq.(2) becomes $`\mu \psi _0(𝐫)`$. The important energy scales of the problem are the trap frequencies $`\omega _z`$, $`\omega _\rho `$ and the dipole–dipole interaction energy per particle, defined as $`V=(1/N)V_d(𝐫𝐫^{})\psi _0^2(𝐫)\psi _0^2(𝐫^{})𝑑𝐫𝑑𝐫^{}`$. Accordingly, the quantity $`V/\mathrm{}\omega _\rho `$, the aspect ratio of the trap $`l`$, and the (renormalized) number of particles $`\sigma =Nr_{}/a_{\mathrm{max}}`$, (with $`a_{\mathrm{max}}=(\mathrm{}/2M\omega _{\mathrm{min}})^{1/2}`$ being the maximal oscillator length of the trap) form the necessary set of parameters allowing us to determine the chemical potential and give a full description of the behavior of a trapped dipolar condensate.
We have found that the dipolar condensate is stable either at $`V>0`$, or at $`V<0`$ with $`|V|<\mathrm{}\omega _\rho `$. This requires $`N<N_c`$, where the critical number $`N_c`$ depends on the trap aspect ratio $`l`$. The calculated dependence $`N_c(l)`$ is presented in Fig. 1 and clearly indicates the presence of a critical point $`l_{}=0.4`$. In pancake traps with $`l<l_{}`$ the condensate is stable at any $`N`$, because $`V`$ always remains positive (see Fig. 2). For small $`N`$ the shape of the cloud is Gaussian in all directions. With increasing $`N`$, the quantity $`V`$ increases and the cloud first becomes Thomas–Fermi in the radial direction and then, for a very large $`N`$, also axially. The ratio of the axial to radial size of the cloud, $`L=L_z/L_\rho `$, continuously decreases with increasing number of particles and reaches a limiting value at $`N\mathrm{}`$ (see Fig. 3). In this respect, for a very large $`N`$ we have a pancake Thomas–Fermi condensate.
For $`l1`$ the mean–field dipole–dipole interaction is always attractive. The quantity $`|V|`$ increases with $`N`$ and the shape of the cloud changes (see Fig. 2 and Fig. 3). In spherical traps the cloud becomes more elongated in the axial direction and near $`N=N_c`$ the shape of the cloud is close to Gaussian, with the aspect ratio $`L=2.1`$. In cigar–shaped traps ($`l1`$) especially interesting is the regime where $`\mathrm{}\omega _z|V|\mathrm{}\omega _\rho `$. In this case the radial shape of the cloud remains the same Gaussian as in a non–interacting gas, but the axial behavior of the condensate will be governed by the dipole–dipole interaction which acquires a quasi1D character. Thus, one has a (quasi) 1D gas with attractive interparticle interaction and is dealing with a stable (bright) soliton–like condensate where attractive forces are compensated by the kinetic energy . With increasing $`N`$, $`L_z`$ decreases. Near $`N=N_c`$, where $`|V|`$ is close to $`\mathrm{}\omega _\rho `$, the axial shape of the cloud also becomes Gaussian and the aspect ratio takes the value $`L3.0`$. For $`l_{}ł<1`$, the dipole–dipole interaction energy is positive for small number of particles and increases with $`N`$. The quantity $`V`$ reaches its maximum, and the further increase in $`N`$ reduces $`V`$ and makes the cloud less pancake. At the critical point $`N=N_c`$ the shape of the cloud is close to Gaussian and the aspect ratio $`L<3.0`$, tending to $`1`$ as $`ll_{}`$.
To gain insight in the nature of the stability of dipolar condensates in pancake traps we performed a variational ansatz assuming a Gaussian shape of the cloud (c.f. ):
$$\psi _0=N^{1/2}(2\pi )^{3/4}(L_\rho ^2L_z)^{1/2}e^{\rho ^2/4L_\rho ^2}e^{z^2/4L_z^2}.$$
(3)
Minimizing the energy functional $`H`$ of the system, we found the aspect ratio of the cloud $`L`$ and established that it decreases with increasing $`\sigma `$ for $`l<l_{}`$, and increases otherwise. The point $`l_{}`$ can be estimated by requiring $`dL/d\sigma |_{\sigma =0}=0`$, which provides the value $`l_{}=0.41`$ in good agreement with the numerical calculation.
For understanding the behavior of the dipolar condensate near the critical point $`N=N_c`$, we observe that at this point the local minimum of $`H`$ becomes a saddle point. Hence, at $`N_c`$ one has $`(^2H/L_z^2)(^2H/L_\rho ^2)(^2H/L_zL_\rho )^2=0`$, in addition to $`H/L_z=H/L_\rho =0`$. This gives the relation between $`L`$ and $`l`$ at the criticality:
$$\frac{(2L^2+1)(5+10l^4)}{2L^4+1}\frac{6B(L)(1+2l^4)}{(L^21)^2}1=2\frac{l^4}{L^2},$$
(4)
where $`B(L)=2+L^23L\mathrm{arctan}[\sqrt{1L^2}/L]/\sqrt{1L^2}`$. Similarly, one can find the corresponding expressions for $`\sigma `$, $`L_z`$ and $`L_\rho `$ as a function of $`\omega _{z,\rho }`$. The result of Eq.(4) differs by less than $`15`$% from our numerical calculations.
We have also analyzed the dynamics of the instability, evolving (in real time) the initially stable condensate by a slow increase of $`\sigma `$. In our approach based on the GPE, similarly to the case of $`a<0`$ (see ), the dipolar condensate collapses to a point on a finite time scale. The ratio $`L`$ increases moderately in the course of the collapse. The results for the spherical trap are presented in Fig. 4. The collapse in cigar–shaped traps occurs in a similar way, since the initial shape of the collapsing dipolar cloud is almost the same as in spherical traps (see Fig.1).
As we see, the ground state of a dipolar gas exhibits a very rich behavior and one finds various BEC regimes. In order to electrically polarize an ultra–cold cloud of heteronuclear molecules, and thus create a molecular dipolar gas one should have the electric field which provides the Stark potential $`dE`$ greatly exceeding the spacing between the lowest rotational levels of the molecule. Then the rotational motion of the molecules will be ”frozen” and their dipole moments will be oriented along the direction of the field. For most of the diatomic molecules the rotational level spacing is in the range from $`0.1`$ to $`1`$ K, and hence the required electric field is $`E10^3`$ V/cm.
At present, molecular condensates are not achieved experimentally. We thus propose and analyze an alternative method of inducing electric dipole moments, which can be used in atomic condensates. The idea is to apply a constant electric field and to optically admix the permanent dipole moment of a low–lying Rydberg state to the atomic ground state. Rydberg states of hydrogen and alkali atoms exhibit a linear Stark effect : in hydrogen, for example, an electric field $`E_s`$ splits the manifold of Rydberg states with given principal quantum number $`n`$ and magnetic quantum number $`m`$ into $`2(n|m|1)`$ Stark states. The outermost Stark states have (large) permanent dipole moments $`d_Rn^2ea_B`$ (with $`a_B`$ the Bohr radius), and there will be an associated dipole–dipole force between the atoms.
This dipole–dipole interaction can be controlled with a laser . This is achieved either by admixing the permanent dipole moment of the Stark states to the atomic ground state with an off–resonant cw laser, or by a stroboscopic excitation with a sequence of laser pulses. The pulses tuned to the lowest Stark state of a given Rydberg manifold should be separated by the time $`T`$, have duration $`2\mathrm{\Delta }tT`$ and area $`2\pi `$ . The field $`E_s`$ and the ”dressing” light have to be chosen such that they do not couple the selected lowest Stark state to other states, and the spacing ($`nea_BE_s`$) between the adjacent Stark states should greatly exceed the mean–field dipole–dipole interaction in order to avoid the interaction–induced coupling. Stroboscopic excitation “dresses” the atomic internal states, so that each atom acquires a time averaged dipole moment of the order of $`d_s=n^2ea_Bf`$, oriented in the direction of $`E_s`$. Even though the quantity $`f=\mathrm{\Delta }t/T`$ is assumed to be small, the induced dipole can be rather large for $`n1`$. Taking for example $`\mathrm{\Delta }t=1`$ns, $`T=10\mu `$s, and $`n=20`$, we obtain $`d_s=0.1`$D. In the limit $`f0`$, $`n^2f=\mathrm{const}`$, the resulting time dependent Hamiltonian can be replaced by its time average, leading to Eq.(2) with $`d=d_s`$. A characteristic time scale in Eq.(2) is of order the inverse trap frequency $`\omega ^1`$. Hence, in our case the dynamics of the system is described by Eq.(2) with $`d=d_s`$, if the condition $`\mathrm{\Delta }t,T\omega ^1`$ is satisfied. This has been tested numerically for $`\mathrm{\Delta }t/T=10^4`$: We reproduced our previous results of the static GPE by solving explicitly Eq.(2) for the stroboscopic ”dressing” of atoms, i.e. by setting $`d=d_s/f`$ in the time intervals $`\mathrm{\Delta }t`$ and $`d=0`$ otherwise.
Aside from inducing permanent electric dipoles, the stroboscopic dressing of atoms will somewhat modify the trapping potential and the scattering length related to the Van der Waals interatomic forces. The corresponding corrections will be proportional to the small parameter $`f`$. There will also be losses due to spontaneous emission and black body radiation . The rates of these processes become comparable with each other for $`n=20`$ , where the corresponding decay time is of order $`20\mu `$s. In our scheme the lifetime will be thus $`20\mu \mathrm{s}/f0.2`$ s.
The laser resonant with a bare transition frequency “dresses” only the atoms that are sufficiently separated from their neighbors, since otherwise the dipole–dipole interaction shifts atomic resonances. Atomic pairs are ”shielded” and not dressed at interatomic distances smaller than $`R_{}`$, where the latter follows from the equation $`\mathrm{}\mathrm{\Omega }d_R^2/R_{}^3`$, with $`\mathrm{\Omega }=\pi /\mathrm{\Delta }t`$ being the Rabi frequency associated with the “dressing” laser. For $`n=20`$ and $`\mathrm{\Delta }t=1`$ns we have $`\mathrm{\Omega }=500`$MHz and $`R_{}=0.7\mu `$m.
The most dangerous ”underwater stone” concerns inelastic decay processes. Fortunately, the “shielding” can suppress Penning ionization: a strong suppression is expected if atoms practically do not move during the short time $`\mathrm{\Delta }t`$, and the distances at which the ionization occurs ($`n^2a_B`$) are significantly smaller than $`R_{}`$. For the parameters considered above this should be the case. The dipole–dipole interaction also induces the change of the effective Rabi frequency, and therefore the $`2\pi `$ pulse condition is not strictly satisfied. Hence, a fraction of atoms remains in the Rydberg state between the stroboscopic pulses and decays due to spontaneous emission. This fraction can be reduced by decreasing the quantity $`\stackrel{~}{n}d^2\mathrm{\Delta }t`$, where $`\stackrel{~}{n}`$ is the gas density. A detailed analysis of inelastic processes in the conditions of stroboscopic dressing of atoms requires a separate investigation.
We acknowledge support from Deutsche Forschungsgemeinschaft (SFB 407), TMR ERBXTCT–96–002, the Stichting voor Fundamenteel Ondersoek der Materie (FOM), Alexander von Humboldt Stiftung, the Russian Foundation for Basic studies, and INTAS. We thank K. Góral, H. Knöckel, A. Muryshev, T. Pfau, K. Rza̧żewski, A. Sanpera, and E. Tiemann for fruitful discussions.
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# References
A Numerical Analysis to the $`\pi `$ and $`K`$ Coupled–Channel Scalar Form-factor
Wei Liu, Hanqing Zheng and Xiao-Lin Chen
1) Department of Physics, Peking University, Beijing 100871, P. R. China
## Abstract
A numerical analysis to the scalar form-factor in the $`\pi \pi `$ and $`KK`$ coupled–channel system is made by solving the coupled-channel dispersive integral equations, using the iteration method. The solutions are found not unique. Physical application to the $`\pi \pi `$ central production in the $`pppp\pi \pi `$ process is discussed based upon the numerical solutions we found.
PACS numbers: 14.40.Aq; 11.55.Fv; 13.75.Lb
The I=J=0 channel $`\pi \pi `$ interactions are of great physical interests. Because the interaction of the I=J=0 channel is very strong, the input bare singularities for a given model can be severely renormalized and distorted by the strong attractive force and extra dynamical singularities may be generated . Therefore the I=J=0 channel affords an ideal test ground for models of strong interactions. Also, the lightest glueball is expected to lie in the I=J=0 channel, as predicted by the lattice QCD calculations. Therefore a detailed study to the dynamics in this channel becomes especially interesting. The low energy I=J=0 $`\pi \pi `$ system manifests itself in various production processes intensively measured by experiments, $`i.e.,`$ from $`\pi N\pi \pi N`$, $`\gamma \gamma \pi \pi `$ to $`J/\mathrm{\Psi }\varphi \pi \pi `$, etc., with the center of mass energy, $`\sqrt{s}`$, ranging from the $`\pi \pi `$ threshold to a few GeV. When $`\sqrt{s}`$ exceeds the $`K\overline{K}`$ threshold, a single–channel analysis to the $`\pi \pi `$ system becomes inadequate, and instead, a coupled–channel analysis of $`\pi \pi `$ and $`K\overline{K}`$ system has to be made. Among various $`\pi \pi `$ production processes particularly interesting cases are those I=J=0 final states which are generated weakly. In such circumstances, the final state particles will not scatter back to the initial states, and therefore the complicated dynamics involved is considerably simplified without loss of information on the I=J=0 final state interactions. For example, in such a case the Watson–Migdal’s theorem on final state interactions applies. The $`\pi \pi `$ final states “weakly” produced can be again categorized into two classes: One is that the $`\pi \pi `$ production vertex contains left–hand singularities like in the case $`\gamma \gamma \pi \pi `$ <sup>1</sup><sup>1</sup>1A recent publication on related subject can be found in Ref. . and the another is not, like in the case $`K\pi \pi `$ and in the $`\pi \pi `$ central production process $`pppp\pi \pi `$. Physical situation in the absence of left hand singularities is further simplified, since the dynamical complexity from the production vertex is removed which would otherwise disturb our analysis on the $`\pi \pi `$ (final state) interaction itself. In such a simplified situation, it then becomes reasonable to assume that the $`\pi \pi `$ production amplitude is factorized as a product of a form-factor–like quantity, which we denote as $`A`$, containing all the dynamical singularities from the right–hand cut from $`\pi \pi `$ final state interactions, and the production vertex which is a smooth analytic function of $`s`$ on the complex $`s`$–plane except possibly at infinity since it contains neither the left–hand nor the right–hand singularities.
For a coupled-channel system of $`\pi \pi `$ and $`KK`$, the spectral representation of the form-factor, $`𝐀(A_1,A_2)`$, satisfies the following relation,
$$\begin{array}{ccc}\mathrm{I}mA_1\hfill & =& \hfill A_1\rho _1T_{11}^++A_2\rho _2T_{21}^+,\\ \mathrm{I}mA_2\hfill & =& \hfill A_2\rho _2T_{22}^++A_1\rho _1T_{12}^+,\end{array}$$
(1)
whereas the unitarity relation of the scattering matrix $`𝐓`$ reads,
$$\mathrm{Im}𝐓=𝐓\rho 𝐓^+,$$
(2)
where $`\rho diag(\rho _1,\rho _2)`$ is the matrix of the kinetic phase-space factor. The $`𝐓`$ matrix may contain left–hand cut but $`𝐀`$ does not. Especially, $`𝐀`$ is analytic on the entire physical sheet of the complex $`s`$ plane except on the cut along the real positive axis starting from $`2\pi `$ threshold. The form-factor $`𝐀`$ has the same analytic structure as the scalar form-factor <sup>2</sup><sup>2</sup>2For the definition of the scalar form-factor, see for example Ref. . and the two are different only up to a polynomial. Eqs. (1) and (2) are assumed to be correct down to the lowest threshold, in the absence of anomalous thresholds. That means when $`4m_\pi ^2s4m_K^2`$ Eq. (1) takes the form,
$$\begin{array}{ccc}\mathrm{I}mA_1\hfill & =& \hfill A_1\rho _1T_{11}^+,\\ \mathrm{I}mA_2\hfill & =& \hfill A_1\rho _1T_{12}^+.\end{array}$$
(3)
From previous discussions, it is realized that the physical problem of studying final state interactions in the production process without left–hand singularities is reduced to the mathematical problem of solving Eq. (1) (and Eq. (3)), provided that the T matrix is known. In the single channel case the analytic solution of the form factor can be obtained. The spectral representation of the form-factor is given by the first equation in Eq. (3) from which the classical Omnès solution can be established,
$`A(s)`$ $`=`$ $`P(s)\mathrm{exp}\left({\displaystyle \frac{s}{\pi }}{\displaystyle _{4m_\pi ^2}^{\mathrm{}}}{\displaystyle \frac{\delta _\pi (s^{})}{s^{}(s^{}siϵ)}}𝑑s^{}\right),`$ (4)
where $`P(s)`$ is a polynomial and $`\delta _\pi `$ is the $`\pi \pi `$ scattering phase. The Omnès solution is remarkable in relating the form-factor to $`\delta _\pi `$. The Eq. (4) has been used in Ref. to study the final state interactions in $`K2\pi `$ and $`pppp\pi \pi `$ systems. In the coupled–channel case, however, no analytic solutions of Eq. (1) can be obtained in general. In the following we will instead study the coupled–channel system by numerical method .
Since the coupled–channel form-factor $`𝐀`$ is analytic on the physical sheet of the complex $`s`$ plane except for the cut along the real axis, assuming the $`𝐓`$ matrix is known we can search for solutions of the amplitude $`𝐀`$ in Eq. (1) by solving the following dispersion relation:
$$𝐀=\frac{1}{\pi }_R\frac{𝐀\rho 𝐓^+(s^{})}{s^{}siϵ}𝑑s^{},$$
(5)
where the integration is performed on the unitarity cut $`R`$, starting from $`4m_\pi ^2`$ to $`\mathrm{}`$. The Eq. (5), according to Muskhelishvili , contains two fundamental solutions, $`\varphi _\mathrm{𝟏}`$ and $`\varphi _\mathrm{𝟐}`$. Assuming that $`\varphi _n`$ behaves as $`\varphi _ns^{\chi _n}`$ as $`s\mathrm{}`$ one has
$$\underset{n}{}\chi _n=\frac{1}{2\pi }[argdet𝐒]_R=\frac{1}{\pi }(\delta _\pi (\mathrm{})+\delta _K(\mathrm{})),$$
(6)
where $`𝐒`$ is the coupled–channel $`S`$ matrix. According to , any solution of the integral equation (5) can be written as a linear composition of the two fundamental solutions,
$$\varphi =\underset{n}{}P_n(s)\varphi _n,n=1,2$$
(7)
where $`P_n(s)`$ are polynomials of $`s`$. The polynomials are not determined from analyticity alone. Other physical input has to be implemented to fix their coefficients. For example, chiral perturbation theory can afford an expansion of the scalar form-factor in powers of $`s`$ when $`s`$ is small . Since there can be many solutions of Eq. (5) and we notice that, since the numerical integration in Eq. (5) has to be truncated somewhere (denoted as $`\mathrm{\Lambda }`$ below), all the information on the asymptotic behaviour are lost. As a consequence, it is difficult to distinguish the so–called fundamental solutions from others, since the difference between the two essentially comes from their asymptotic behaviours. Therefore in the present numerical scheme instead of searching for the fundamental solutions, we follow the recipe of Ref. , that is to search for two linearly independent solutions, $`𝐀^\mathrm{𝟏}`$ and $`𝐀^\mathrm{𝟐}`$ which are normalized at $`s=0`$ as,
$$A_1^1(0)=1,A_2^1(0)=0,$$
(8)
and
$$A_1^2(0)=0,A_2^2(0)=1.$$
(9)
In the following we use the $`\pi \pi `$ and $`KK`$ coupled–channel fit of the T matrix from Au, Morgen and Penington as an educative example to solve the coupled–channel dispersive integral equation (5). The integral in Eq. (5) is truncated at $`\mathrm{\Lambda }1.5GeV`$. We find that the influence of the cutoff is rather local, i.e., it has little effects on the behaviour of the solution in a large region of $`s`$, even close to $`\mathrm{\Lambda }`$. One can also use a mild regulation function instead of the truncation at $`\mathrm{\Lambda }`$ and the result is essentially the same except at $`s\mathrm{\Lambda }`$. We use the iteration method to solve Eq. (5). That is,
$$𝐀^{(𝐧+\mathrm{𝟏})}(s)=𝐀(0)+\frac{s}{\pi }_{4m_\pi ^2}^\mathrm{\Lambda }\frac{\mathrm{Real}\left[𝐀^{(𝐧)}(𝐬^{})\rho 𝐓^+(s^{})\right]}{s^{}(s^{}siϵ)}𝑑s^{},$$
(10)
where
$$𝐀(0)=(1,0)\mathrm{or}(0,1).$$
(11)
A once–subtracted form of the dispersive integral in the above Eq. (10) is helpful in incorporating the boundary conditions, Eqs. (8) and (9). The routine converges rather rapidly (after about 20 steps), which confirms the claim in Ref. and the solutions are not unique due to the reason already mentioned above. This means that the iteration may converge to different solutions depending on different initial values for the iteration. In Fig. 1 and Fig. 2 a few examples generated from our numerical recipe are shown. We see from Fig. 1 that all the phases are identical to the phase of $`T_{11}`$ below the $`K\overline{K}`$ threshold as required by the final state theorem. Above the $`K\overline{K}`$ threshold, however, the phases from different solutions of $`A_1`$ can be very different and also deviate from the phase of $`T_{11}`$. In Fig. 2 the corresponding magnitude of the different solutions is also shown. We see that around $`1GeV`$ region there can be zeros or dips to compensate the peak generated by $`f_0(980)`$. The physical discussion on the necessity to introduce the protective zero appearing in the $`\pi \pi `$ production processes can be found in Ref. . In the present numerical approach there is generally no difficulty to pick suitable solutions, from all, to fit experimental data. An example is shown in Fig. 3 where the $`\pi \pi `$ production cross-section in the $`pppp\pi \pi `$ process is reproduced. The cross-section can be expressed by ,
$`{\displaystyle \frac{d\sigma }{d\sqrt{s}}}{\displaystyle \frac{(s4m_\pi ^2)^{1/2}}{s^{3/2}}}|F(s)|^2,`$ (12)
where $`F(s)`$ is a linear combination of solutions $`A_1^1`$ and $`A_1^2`$,
$$F(s)=\alpha _1A_1^1+\alpha _2A_1^2.$$
(13)
It is assumed here that $`\pi \pi `$ produced in the $`pppp\pi \pi `$ process are generated from the fusion of Pomeron pairs and no left–hand singularity is involved in the production vertex. Therefore $`\alpha _1`$ and $`\alpha _2`$ (which contain information about the $`\pi \pi `$ production vertex) appearing in the above equation are expected to be polynomials with a weak dependence on $`s`$.
In fact, the T matrix given in Ref. is obtained by a multi-pole K matrix fit in which the background effects including the left–hand cut one are simulated by a polynomial. Therefore, the T matrix does not contain any left–hand singularity. Under such an approximation it is not really necessary to search for the numerical solutions of $`A_1`$ in order to make use of Eq. (13). Since $`A_1`$ has to be a linear combination of $`T_{11}`$ and $`T_{21}`$ when neglecting the left–hand singularities in $`T`$, one can use instead of Eq. (13) the following form,
$$F(s)=\alpha _1^{}T_{11}+\alpha _2^{}T_{21},$$
(14)
to fit the experimental data, as originally done in Ref. . The method presented in the current note applies to more general cases when the fine details of the left–hand singularities of the scattering $`T`$ matrix are carefully taken into account. Also our program can be easily extended to the situation when the left–hand cut of the production vertex is included. We will investigate these more realistic cases in future.
$`\mathrm{𝐀𝐜𝐤𝐧𝐨𝐰𝐥𝐞𝐝𝐠𝐞𝐦𝐞𝐧𝐭}`$: One of the authors, H.Z. would like to thank M. Locher and V. Markushin for helpful discussions. He is supported in part by National Natural Science Foundation of China under grant No. 19775005.
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# The Solar Core and Solar Neutrinos
## 1 Properties of Solar and Stellar Models
Simple solar models orient us with gross structure: the core $`(r/R_{}\begin{array}{c}<\\ \end{array}`$ 0.3, where nuclear fusion generates the luminosity $`L_{})`$, the radiative zone (0.3 $`\begin{array}{c}<\\ \end{array}r/R_{}<`$ 0.71) and the convective zone (CZ, $`r/R_{}>`$ 0.71) (Figure 1). Less massive stars $`(M<M_{})`$ have, according to stellar models, even deeper convective zones; while for $`M>M_{}`$, the outer convective zone disappears, and a convective inner core appears as the central temperature gradient surpasses a critical value. The Sun might have had a convective inner core, drastically changing the $`\varphi `$ predictions, but there is now decisive evidence against core convection (see Section 3 below). A higher central temperature would also have led to CNO-dominance and higher neutrino fluxes.
A solar model is standard (an SSM) if the model contains all the physics of matter, gravitation, and nuclear fusion needed to obtain a star, but nothing more. Conceptually, stellar structure and evolution divide into three levels by time scales. For the Sun, chemical evolution needs $``$ 10 Gyr; thermal equilibrium, about 10 Myr (Kelvin-Helmholtz time); and hydrostatic equilibrium, about 5 minutes. The hydrodynamic time controls the helioseismic $`p`$\- and $`g`$-modes. (Late in the Sun’s evolution, the chemical time scale will be shortened and the hierarchy blurred.) Structure means only the thermal and mechanical features of a star.
The properties of matter needed are constitutive relations giving pressure $`P`$, opacity $`\kappa `$, and specific luminosity generation $`\epsilon `$ as functions of density $`\rho `$ and temperature $`T`$. If the thermal structure is given, then the equation of state reduces to a barytrope $`P`$ = $`P(\rho )`$, and the mechanical structure alone becomes a closed problem characterized by the stiffness profile $`\mathrm{\Gamma }`$ = $`d\mathrm{ln}P/d\mathrm{ln}\rho `$. The special case of $`\mathrm{\Gamma }`$ = $`\gamma 1+1/n`$ = constant is a polytrope of index $`n`$. All $`n<`$ 5 polytropes have finite mass and radius. $`n`$ = 0 is the constant-density case.
The initial conditions are fixed mass $`M`$ and the element mass abundances $`X_i`$ (the zero-age star assumed chemically homogeneous). The boundary conditions are zero mass and luminosity at the center and (nearly) zero pressure and density at the surface. The $`X_i`$ develop gradients by evolution, as nuclear fusion and heavy element diffusion act, and additional helium accumulates in the core.
Even within the SSM framework, variations are possible. The $`\kappa `$ and $`\epsilon `$ functions and fusion cross sections must be calculated from atomic and nuclear physics and extrapolated into regimes not directly testable. For the dominant luminosity-producing $`ppI`$ reactions (terminating through <sup>3</sup>He–<sup>3</sup>He fusion to <sup>4</sup>He), $`\epsilon `$ and the reaction rates are almost fixed by $`L_{}`$. But reactions not strongly connected with $`L_{}`$ are not well-constrained: the $`ppII`$ and $`ppIII`$ chains (terminating from <sup>7</sup>Be through <sup>7</sup>Li and <sup>8</sup>B to $`2^4`$He, respectively) are sensitive functions of the core temperature and nuclear cross sections, as well as mildly dependent on the core density. $`L_{}`$ and other global properties place only weak constraints on the $`ppII`$ and $`ppIII`$ rates.
## 2 Generalizing the Standard Solar Model
Generic properties of solar structure are restricted by boundary conditions. The SSM can then be generalized in two ways. One is to calibrate with the specific model. This procedure defines a generalized SSM family, although not the most general.
Its most convenient implementation is through homology or power-law scaling, derivable analytically or made evident by numerical solutions. Exploration of model space by varying SSM inputs is actually a special case of homology, which amounts to “small” perturbations of the logarithms of inputs and outputs. Such “perturbative” analysis works over a surprisingly large range, so long as the power law-relations are stable.
The first signs of homological behavior in SSMs were found in the 1000 SSM Monte Carlo study of Bahcall and Ulrich. Subsequent work over a wider model range revealed a much broader validity for homology. The underlying analytic structure was derived by Bludman and Kennedy. Starting with structure alone, one keeps only dimensional and scaling behavior of macroscopic variables, dropping the differential nature of the equations. Assuming multifactor power laws for the equation of state, opacity, and luminosity generation, we found homological relations for the mechanical and thermal structure, assuming fixed powers. This requirement restricts the homology to the radiative and core regions, below the CZ. The dominant luminosity production is by $`ppI`$, carried outwards entirely by radiative diffusion. The constitutive relations are
$$P/\rho =\mathrm{}T/\mu ,\kappa (\rho ,T)=\kappa _o(X_i)\rho ^nT^s,\epsilon (\rho ,T)=\epsilon _o(X)\rho ^\lambda T^\nu .$$
(1)
Expanded about the SSM, the exponents are: $`n`$ = 0.43, $`s`$ = 2.5, $`\lambda `$ = 1.0, $`\nu `$ = 4.2. $`\mu `$, $`\kappa _o`$, and $`\epsilon _o`$ are composition-dependent. The luminosity constraint reads:
$$\varphi (pp)+(0.977)\varphi (\mathrm{Be})+(0.751)\varphi (\mathrm{B})+(0.956)\varphi (\mathrm{CNO})=6.55\times 10^{10}\mathrm{cm}^2\mathrm{sec}^1.$$
(2)
The boundary conditions are imposed in a way appropriate to a single star: $`M_{}`$, $`L_{}`$, and $`R_{}`$ fixed. Homology then gives a family of possible non-convective interiors consistent with observed outer solar features and parametrized by $`\rho _c`$ and $`T_c`$:
$`\rho _c\epsilon _o^{0.34}\kappa _o^{0.40}\mu _c^{0.52}L_{}^{0.085},T_c\epsilon _o^{0.13}\kappa _o^{0.034}\mu _c^{0.22}L_{}^{0.17}.`$ (3)
The resulting $`\varphi `$ scale stably with $`\rho _c`$ and $`T_c`$ over a large range, giving the two-parameter homological mechanical/thermal variations of the SSM: $`\varphi (i)\rho _c^{\alpha _i}T_c^{\beta _i}`$, with $`(\alpha _i,\beta _i)`$ for $`pp`$, Be, and B $`\nu `$’s being $`(0.1,0.7)`$, $`(0.7,9)`$, and $`(0.3,21)`$, respectively. The highest reactions in the $`pp`$ chain have the famous extreme sensitivity to $`T_c`$, while all sensitivities to $`\rho _c`$ are mild and arise from the small luminosity contribution made by the $`ppII`$, $`ppIII`$, and CNO chains. It should be stressed that $`\rho _c`$ and $`T_c`$ are model outputs, like the $`\varphi (i)`$. These exponents reproduce the 1000-SSM Monte Carlo and clarify that the entire homological class of SSMs has the wrong pattern of fluxes to explain the observed energy dependence of $`\varphi :`$ lower energies are more suppressed. (Variation of nuclear cross sections also fail to explain the pattern.) This conclusion depends only on $`ppI`$-dominance in $`\epsilon `$, radiative diffusion in $`\kappa `$, and the ideal gas law.
It is instructive to compare these homology results with the approximations often used to model stars. Perhaps the simplest (after the constant-density case) is the Eddington standard model, based on a constant ratio of radiation to matter pressure throughout the star and equivalent to an $`(n,\gamma )`$ = $`(3,4/3)`$ polytrope. This model, once its free parameters are fit, represents many main sequence stars not badly. Homology applied to the mechanical structure alone automatically leads to a polytrope. But the Bludman-Kennedy homology is more general than a polytrope, as it applies to both mechanical and thermal structure. It also scales correctly in the evolved core, where the molecular weight $`\mu `$ changes substantially, reflecting the chemical evolution that makes the present Sun differ from its zero-age incarnation. No polytrope fits this behavior.
A more complete version of homology is possible if the differential structure is retained and rewritten using scale-invariant homology variables For the entire mechanical, thermal, and chemical structural system, the dimensionless differential equations are not less complex than a full SSM. But if the structure is restricted to the mechanical alone and a barytrope $`P(\rho )`$ assumed, simple dimensionless structure equations follow. The key to the mechanical structure turns out to be the $`\mathrm{\Gamma }`$ profile (Figure 2). Constant $`\mathrm{\Gamma }`$ gives a polytrope again; in fact, $`\mathrm{\Gamma }_{\mathrm{SSM}}`$ 4/3 outside the inner core, up to the CZ, where it rises to 5/3 (the adiabatic value). But within the core, $`\mu `$ rises and $`\mathrm{\Gamma }`$ drops towards to the center, where $`\mathrm{\Gamma }_c^{\mathrm{SSM}}`$ 8/9. Described in terms of a polytrope, the effective index $`n_{\mathrm{eff}}`$ = $`(\mathrm{\Gamma }1)^1`$ rises in the core, diverging at $`\mathrm{\Gamma }`$ = 1. Further towards the center, $`n_{\mathrm{eff}}`$ rises from minus infinity to a finite negative value at the center. Such behavior in the inner core is not even approximately polytropic and explains why attempts to use polytropes to approximate the Sun only work (and only crudely) over the whole Sun and fail badly in the core.
The other approach to generalizing the SSM is to work with a few reasonable assumptions, following these to simple, testable predictions. Some powerful results are available, although restricted to mechanical structure only, to which helioseismology is the key: adiabatic sound waves are mechanical perturbations with information about sound speed, equation of state, and density and pressure profiles. Taking full advantage of these results requires both $`p`$\- (pressure) and $`g`$\- (gravity) modes. Their spectra can be inverted to yield adiabatic sound speed $`c_{\mathrm{ad}}(r)`$ and bouyancy frequency $`N(r)`$ profiles, from which follows the complete mechanical structure. The BP98 $`N(r)`$ profile is shown in Figure 3.
The thermal and chemical structures cannot be directly probed by helioseismology, as their associated time scales are so long. Simple homology implies $`L\mu ^4M^3\kappa _o^1\epsilon _o^0`$, but changes in $`L_{}`$ large enough to explain the solar $`\nu `$ deficit are probably ruled out by paleoclimatology. Direct tests of these aspects of the SSM require comparison to other Sun-like stars. Such stars vary from the Sun somewhat in mass and chemical composition and could lie anywhere on their respective evolutionary tracks. Comparison properties include luminosity, surface temperature, and photospheric radius. With accurate photometry and parallaxes, accurate luminosities and colors are achievable. The intermediary between these observations and stellar structure is stellar atmosphere models, which have advanced considerably in the last 30 years. Although still oversimplified, the models are good enough for solar-type stars to infer ranges for surface $`T`$, $`g`$, abundances, and turbulence. An exciting possibility will be opened by asteroseismology of Sun-like stars, as observation of stellar seismic modes (especially $`g`$-modes) would lead to direct characterization of stellar interiors.
## 3 Observational Issues
$`M_{}`$, $`L_{}`$, $`R_{}`$, and surface $`T`$, as well as surface and proto-solar (meteoric) abundances, are well measured. Helioseismic observations have directly or indirectly captured millions of $`p`$-modes, allowing an accurate inversion of $`c_{\mathrm{ad}}(r)`$ down to $`r/R_{}`$ = 0.05 (Figure 4). A convective core $`(\gamma `$ = 5/3) is ruled out, although circulation of heavy elements not affecting heat transport cannot be at present.
The inferred sound speed peaks off-center at $`r/R_{}`$ 0.07. Hydrostatic equilibrium requires $`dc_{\mathrm{ad}}/dr`$ = 0 at the center, but an off-center peak occurs where $`\mathrm{\Gamma }`$ = 1. This peak and $`dc_{\mathrm{ad}}/dr>`$ 0 for $`r/R_{}<`$ 0.07 indicate $`\mathrm{\Gamma }<`$ 1 there, a crucial confirmation of the core’s chemically evolved state. A complete profile down to the center becomes possible with the lowest $`p`$-modes. But complete inversion for model-independent mechanical structure would be possible only if the higher $`g`$-modes were observed; analogous inversion would yield $`N(r)`$, and together $`N`$ and $`c_{\mathrm{ad}}`$ yield $`\mathrm{\Gamma }`$ and other mechanical profiles, the first truly independent test of the SSM.
Comparing the Sun with other sun-like stars has been possible for many decades, albeit at poor precision. But the recent Hipparcos-Tycho star catalogs (edition I in 1997: 1.1 M stars; edition II in 2000: 2.5 M stars) have revolutionized astrometry, raising the accuracy of nearby stellar parallaxes by up to a factor of 10 or better (1 m-arcsec). Luminosities accurate to $`<`$ 5% for nearby stars (within 25 pc) can now be inferred. With good color measurements and best current stellar atmosphere models, surface $`T`$’s can be limited to 1%. Even more dramatic astrometric improvements could come from the proposed SIM and GAIA orbital systems, to be launched in 2006 and 2009, respectively: 4 $`\mu `$-arcsec parallax errors and luminosity errors limited only by photometry, tenths of a percent.
## Acknowledgments
The author thanks the CSNP for the opportunity to present our results in honor of Frank Avignone’s life and work. This work was done in collaboration with Sidney Bludman (Univ. Pennsylvania and DESY) and Susana Tuzzo (Univ. Florida), who helped sample and interpret the Hipparcos and SIMBAD catalogs, and was supported by the U.S. DOE under grants DE-FG02-97ER41029 (Univ. Florida) and DE-FG06-90ER40561 (Univ. Washington), the Institute for Fundamental Theory (Univ. Florida), and the Eppley Foundation for Research. The Institute for Nuclear Theory (Univ. Washington) and the Aspen Center for Physics extended their hospitality. The author is also grateful for important interactions with John Bahcall (Institute for Advanced Study), Jørgen Christensen-Dalsgaard (Aarhus Univ.), Christopher Essex (Univ. Western Ontario), Wick Haxton (Univ. Washington), Marc Pinsonneault (Ohio State Univ.), and Dimitri Pourbaix (ESA Hipparcos Science Team) concerning solar models, helioseismology, non-equilibrium thermodynamics, and the new classical astronomy of Hipparcos.
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# 1 Introduction
## 1 Introduction
Among variety of all affine Lie (super)algebras<sup>1</sup><sup>1</sup>1We conclude the prefix ”super” in brackets to stress that the Lie (super)algebras include both the Lie algebras and the own Lie superalgebras. (both quantized and non-quantized) the affine (super)algebras of rank 2 play same key role. In the first place, all affine series of the type $`A(n|m)^{(1)}`$, $`B(n|m)^{(1)}`$, $`C(n)^{(1)}`$, $`D(n|m)^{(1)}`$, $`A(2n|2m1)^{(2)}`$, $`A(2n1|2m1)^{(2)}`$,$`C(n)^{(2)}`$, $`D(n|m)^{(2)}`$ and $`A(2n|2m)^{(4)}`$ are started from the affine (super)algebras of rank 2. Secondly, the contragredient Lie (super)algebras of rank 2 are basic structural blocks of any affine (super)algebras of arbitrary rank. This fact permits, for example, to reduce the proof of basic theorems for extremal projector and the universal $`R`$-matrix to the proof of such theorems for the (super)algebras of rank 2 (see Refs. , , ). Moreover the representation theory of the affine (super)algebras (both quantized and non-quantized) contains some typical elements of the representation theory of the affine (super)algebras of rank 2. Besides, in applications of the affine (super)algebras, first of all the affine (super)algebras of rank 2 are used by virtue of their simplicity.
In this paper we give detailed description of the quantum untwisted affine algebra $`U_q(A_1^{(1)})`$ ($`U_q(\widehat{sl}(2))`$) and the quantum twisted affine superalgebra $`U_q(C(2)^{(2)})`$ ($`U_q(\widehat{osp}(2|2))^{(2)}`$). Moreover our goal is to show that these quantum (super)algebras are described in unified way. Namely, we present in unified way their defining relations and actions of the braid group associated with the Weyl group, the construction of the Cartan-Weyl bases, the complete list of all permutation relations of the Cartan-Weyl generators corresponding to all root vectors and finally the unified formula for their extremal projector and universal $`R`$-matrix. We extend also the unified description to so called ’new realizations’ of the algebras. Here we present a unified description of the universal $`R`$-matrices for corresponding Hopf structures in a multiplicative form as well as in the form of contour integrals. Difference between both considered here quantum (super)algebras is only determined by a phase factor which is equal to $`1`$ for $`U_q(A_1^{(1)})`$ and it is equal to $`1`$ for $`U_q(C(2)^{(2)})`$. This situation is similar to the finite-dimensional case. Namely, in the paper it was shown that all quantum (super)algebras $`U_q(g)`$, where $`g`$ are the finite-dimensional contragredient Lie (super)algebras of rank 2, are divided into three classes. Each such class is characterized by the same Dynkin diagram and the same reduced root system, provided that we neglect ’colour’ of the roots, and all (super)algebras of the same class have the unified defining relations, unified construction and properties of the Cartan-Weyl basis and unified formula for the universal $`R`$-matrix. Difference between the (super)algebras of the same class is determined by some phase factors which takes values $`\pm 1`$ depending on the colour of the nodes of their Dynkin diagram. Concerning the Cartan Weyl bases for the quantum affine algebra $`U_q(A_1^{(1)})`$ and quantum affine superalgebra $`U_q(C(2)^{(2)})`$ it should be noted that certain results presented here can be founded in literature separately for the each case (e.g., see Refs. , , , and ).
Basic information about the (super)algebras $`A_1^{(1)}`$ and $`C(2)^{(2)}`$ is presented in the tables 1a and 1b (see Refs. , , ). In the table 1a there are listed the standard and symmetric Cartan matrices $`A`$ and $`A^{sym}`$, the corresponding extended symmetric matrices $`\overline{A}^{sym}`$ and their inverses $`(\overline{A}^{sym})^1`$, as well as the sets of odd simple roots (odd roots), the Dynkin diagrams (diagram), and the dimensions of these (super)algebras (dim). We remind some elementary definitions of the colour of the roots:
* All even roots are called white roots. A white root is pictured by the white node .
* An odd root $`\gamma `$ is called a grey root if $`2\gamma `$ is not any root. This root is pictured by the grey node $``$.
* An odd root $`\gamma `$ is called a dark root if $`2\gamma `$ is a root. This root is pictured by the dark node .
We also remind the definition of the reduced system of the positive root system $`\mathrm{\Delta }_+`$ for any contragredient (super)algebras of finite growth.
* The system $`\underset{¯}{\mathrm{\Delta }}_+`$ is called the reduced system if it is defined by the following way: $`\underset{¯}{\mathrm{\Delta }}_+=\mathrm{\Delta }_+\backslash \{2\gamma \mathrm{\Delta }_+|\gamma \mathrm{is}\mathrm{odd}\}`$. That is the reduced system $`\underset{¯}{\mathrm{\Delta }}_+`$ is obtained from the total system $`\mathrm{\Delta }_+`$ by removing of all doubled roots $`2\gamma `$ where $`\gamma `$ is a dark odd root.
The total and reduced root systems of the (super)algebras $`A_1^{(1)}`$ and $`C(2)^{(2)}`$ are listed in the table 1b. It is convenient to present the total root systems $`\mathrm{\Delta }=\mathrm{\Delta }_+(\mathrm{\Delta }_+)`$ and reduced root systems $`\underset{¯}{\mathrm{\Delta }}=\underset{¯}{\mathrm{\Delta }}_+(\underset{¯}{\mathrm{\Delta }}_+)`$ by the pictures: Figs. 1, 2a, 2b. Comparing Fig. 1 and Fig. 2b we see that the reduced root systems of $`A_1^{(1)}`$ and $`C(2)^{(2)}`$ coincide if we neglect colour of the roots.
Table 1a
| | | | | | |
| --- | --- | --- | --- | --- | --- |
| $`g(A,\mathrm{{\rm Y}})`$ | $`A=A^{sym}`$ | $`\overline{A}^{sym}`$ | $`\left(\overline{A}^{sym}\right)^1`$ | odd | diagram |
| | | | | | |
| $`A_1^{\left(1\right)}`$ | $`\left(\begin{array}{cc}\hfill 2& \hfill 2\\ \hfill 2& \hfill 2\end{array}\right)`$ | $`\left(\begin{array}{ccc}\hfill 0& \hfill 1& \hfill 0\\ \hfill 1& \hfill 2& \hfill 2\\ \hfill 0& \hfill 2& \hfill 2\end{array}\right)`$ | $`\left(\begin{array}{ccc}\hfill 0& \hfill 1& \hfill 1\\ \hfill 1& \hfill 0& \hfill 0\\ \hfill 1& \hfill 0& \hfill \frac{1}{2}\end{array}\right)`$ | $`\mathrm{}`$ | |
| $`C\left(2\right)^{\left(2\right)}`$ | $`\left(\begin{array}{cc}\hfill 2& \hfill 2\\ \hfill 2& \hfill 2\end{array}\right)`$ | $`\left(\begin{array}{ccc}\hfill 0& \hfill 1& \hfill 0\\ \hfill 1& \hfill 2& \hfill 2\\ \hfill 0& \hfill 2& \hfill 2\end{array}\right)`$ | $`\left(\begin{array}{ccc}\hfill 0& \hfill 1& \hfill 1\\ \hfill 1& \hfill 0& \hfill 0\\ \hfill 1& \hfill 0& \hfill \frac{1}{2}\end{array}\right)`$ | $`\{\delta \alpha ,a\}`$ | |
Table 1b
| | | |
| --- | --- | --- |
| $`g(A,\mathrm{{\rm Y}})`$ | $`\mathrm{\Delta }_+`$ | $`\underset{¯}{\mathrm{\Delta }}_+`$ |
| | | |
| $`A_1^{\left(1\right)}`$ | $`\left\{\alpha ,n\delta \pm \alpha ,n\delta \right|n\mathrm{N}\mathrm{N}\}`$ | $`\left\{\alpha ,n\delta \pm \alpha ,n\delta \right|n\mathrm{N}\mathrm{N}\}`$ |
| $`C\left(2\right)^{\left(2\right)}`$ | $`\left\{\alpha ,\mathrm{\hspace{0.17em}2}\alpha ,n\delta \pm \alpha ,\mathrm{\hspace{0.17em}2}n\delta \pm 2\alpha ,n\delta \right|n\mathrm{N}\mathrm{N}\}`$ | $`\left\{\alpha ,n\delta \pm \alpha ,n\delta \right|n\mathrm{N}\mathrm{N}\}`$ |
Fig. 1. The total and reduced root system ($`\mathrm{\Delta }=\underset{¯}{\mathrm{\Delta }}`$) of $`A_1^{(1)}(\widehat{sl}_2)`$
.
Fig. 2a. The total root system $`\mathrm{\Delta }`$ of $`C(2)^{(2)}`$
Fig. 2b. The reduced root system $`\underset{¯}{\mathrm{\Delta }}`$ of $`C(2)^{(2)}`$
.
## 2 Defining relations of $`U_q(A_1^{(1)})`$ and $`U_q(C(2)^{(2)})`$
The quantum (q-deformed) affine (super)algebras $`U_q(A_1^{(1)})`$ and $`U_q(C(2)^{(2)})`$ are generated by the Chevalley elements $`k_\mathrm{d}^{\pm 1}:=q^{\pm h_\mathrm{d}}`$, $`k_\alpha ^{\pm 1}:=q^{\pm h_\alpha }`$, $`k_{\delta \alpha }^{\pm 1}:=q^{\pm h_{\delta \alpha }}`$, $`e_{\pm \alpha }`$, $`e_{\pm (\delta a)}`$ with the defining relations
$`k_\gamma k_\gamma ^1`$ $`=`$ $`k_\gamma ^1k_\gamma =1,[k_\gamma ^{\pm 1},k_\gamma ^{}^{\pm 1}]=0,`$ (2.1)
$`k_\gamma e_{\pm \alpha }k_\gamma ^1`$ $`=`$ $`q^{\pm (\gamma ,\alpha )}e_{\pm \alpha },k_\gamma e_{\pm (\delta \alpha )}k_\gamma ^1=q^{\pm (\gamma ,\delta \alpha )}e_{\pm (\delta \alpha )},`$ (2.2)
$`[e_\alpha ,e_\alpha ]`$ $`=`$ $`[h_\alpha ]_q,[e_{\delta \alpha },e_{\delta +\alpha }]=[h_{\delta \alpha }]_q,`$ (2.3)
$`[e_\alpha ,e_{\delta +\alpha }]`$ $`=`$ $`0,[e_\alpha ,e_{\delta \alpha }]=0,`$ (2.4)
$`[e_{\pm \alpha },[e_{\pm \alpha },[e_{\pm \alpha },e_{\pm (\delta \alpha )}]_q]_q]_q`$ $`=`$ $`0,`$ (2.5)
$`[[[e_{\pm \alpha },e_{\pm (\delta \alpha )}]_q,e_{\pm (\delta \alpha )}]_q,e_{\pm (\delta \alpha )}]_q`$ $`=`$ $`0,`$ (2.6)
where ($`\gamma =\mathrm{d},\alpha ,\delta \alpha `$), $`(\mathrm{d},\alpha )=0`$, $`(\mathrm{d},\delta )=1`$, and $`[h_\beta ]_q:=(k_\beta k_\beta ^1)/(qq^1)`$. The brackets $`[,]`$ and $`[,]_q`$ are the super-, and q-super-commutators:
$$\begin{array}{ccc}\hfill [e_\beta ,e_\beta ^{}]& =& e_\beta e_\beta ^{}(1)^{\vartheta (\beta )\vartheta (\beta ^{})}e_\beta ^{}e_\beta ,\hfill \\ \hfill [e_\beta ,e_\beta ^{}]_q& =& e_\beta e_\beta ^{}(1)^{\vartheta (\beta )\vartheta (\beta ^{})}q^{(\beta ,\beta ^{})}e_\beta ^{}e_\beta .\hfill \end{array}$$
(2.7)
Here the symbol $`\vartheta ()`$ means the parity function: $`\vartheta (\beta )=0`$ for any even root $`\beta `$, and $`\vartheta (\beta )=1`$ for any odd root $`\beta `$.
Remark. The left sides of the relations (2.5) and (2.6) are invariant with respect to the replacement of $`q`$ by $`q^1`$. Indeed, if we remove the q-brackets we see that the left sides of (2.5) and (2.6) contain the symmetric functions of $`q`$ and $`q^1`$. This property permits to write the q-commutators in (2.5) and (2.6) in the inverse order, i.e.
$`[[[e_{\pm (\delta \alpha )},e_{\pm \alpha }]_q,e_{\pm \alpha }]_q,e_{\pm \alpha }]_q`$ $`=`$ $`0,`$ (2.8)
$`[e_{\pm (\delta \alpha )},[e_{\pm (\delta \alpha )},[e_{\pm (\delta \alpha )},e_{\pm \alpha }]_q]_q]_q`$ $`=`$ $`0.`$ (2.9)
Now we prove the useful proposition.
###### Proposition 2.1
(i) In the quantum (super)algebras $`U_q(A_1^{(1)})`$ and $`U_q(C(2)^{(2)})`$ the following relations
$`[[e_{\pm \alpha },[e_{\pm \alpha },e_{\pm (\delta \alpha )}]_q]_q,[[e_{\pm \alpha },[e_{\pm \alpha },e_{\pm (\delta \alpha )}]_q]_q,[e_{\pm \alpha },e_{\pm (\delta \alpha )}]_q]_q]_q`$ $`=`$ $`0,`$ (2.10)
$`[[e_{\pm (\delta \alpha )},[e_{\pm (\delta \alpha )},e_{\pm \alpha }]_q]_q,[[e_{\pm (\delta \alpha )},[e_{\pm (\delta \alpha )},e_{\pm \alpha }]_q]_q,[e_{\pm (\delta \alpha )},e_{\pm \alpha }]_q]_q]_q`$ $`=`$ $`0.`$ (2.11)
are fulfilled.
(ii) On the contrary, if the relations (2.1)–(2.4) and (2.10), (2.11) are satisfied then also the relations (2.5), (2.6) are valid.
Thus, the proposition says that under the conditions (2.1)–(2.4) the relations (2.5), (2.6) and (2.10), (2.11) are equivalent.
Proof. Let us assume that the relations (2.1)–(2.6) are fulfilled. We take the relations (2.6) and apply to them the corresponding q-commutator with fourth power of $`e_{\pm \alpha }`$, i.e. $`[e_{\pm \alpha },[e_{\pm \alpha },[e_{\pm \alpha },[e_{\pm \alpha },a_\pm ]_q\mathrm{}]_q=0`$, where $`a_\pm `$ is the left side of (2.6). After tedious calculation we arrive to the relations (2.10). The relations (2.11) are proved in similar way. Namely, the relations $`[e_{\pm (\delta \alpha )},[e_{\pm (\delta \alpha )},[e_{\pm (\delta \alpha )},[e_{\pm (\delta \alpha )},b_\pm ]_q\mathrm{}]_q=0`$, where $`b_\pm `$ is the right-side of (2.5) in the form (2.8), results in (2.11). On the contrary, if the relations (2.1)–(2.4) and (2.10), (2.11) are realized then from the relations $`[e_\alpha ,[e_\alpha ,[e_\alpha ,[e_\alpha ,a_{}^{}]\mathrm{}]=0`$ and $`[e_{(\delta +\alpha )},[e_{(\delta +\alpha )},[e_{(\delta +\alpha )},[e_{(\delta +\alpha )},b_{}^{}]\mathrm{}]=0`$, where correspondingly $`a_{}^{}`$ and $`b_{}^{}`$ are the left sides of the relations (2.10) and (2.11), follow the relations (2.5) and (2.6). $`\mathrm{}`$
The standard Hopf structure of the quantum (super)algebras $`U_q(A_1^{(1)})`$ and $`U_q(C(2)^{(2)})`$ is given by the following formulas for the comultiplication $`\mathrm{\Delta }_q`$ and antipode $`S_q`$:
$$\begin{array}{ccccc}\hfill \mathrm{\Delta }_q(k_\gamma ^{\pm 1})& =& k_\gamma ^{\pm 1}k_\gamma ^{\pm 1},S_q(k_\gamma ^{\pm 1})& =& k_\gamma ^1,\hfill \\ \hfill \mathrm{\Delta }_q(e_\beta )& =& e_\beta 1+k_\beta ^1e_\beta ,S_q(e_\beta )& =& k_\beta e_\beta ,\hfill \\ \hfill \mathrm{\Delta }_q(e_\beta )& =& e_\beta k_\beta +1e_\beta ,S_q(e_\beta )& =& e_\beta k_\beta ^1,\hfill \end{array}$$
(2.12)
where $`\beta =\alpha ,\delta \alpha `$; $`\gamma =\mathrm{d},\beta `$.
It is not hard to verify by direct calculations for the defining relations (2.1)–(2.6) that the quantum affine (super)algebras $`U_q(A_1^{(1)})`$ and $`U_q(C(2)^{(2)})`$ have the following simple involute (anti)automorphisms.
(i) The non-graded antilinear antiinvolution or conjugation ”:
$$\begin{array}{ccccc}\hfill (q^{\pm 1})^{}& =& q^1,(k_\gamma ^{\pm 1})^{}& =& k_\gamma ^1,\hfill \\ \hfill e_\beta ^{}& =& e_\beta ,e_\beta ^{}& =& e_\beta \hfill \end{array}$$
(2.13)
($`(xy)^{}=y^{}x^{}`$ for $`x,yU_q(g)`$).
(ii) The graded antilinear antiinvolution or graded conjugation ”:
$$\begin{array}{ccccc}\hfill (q^{\pm 1})^{}& =& q^1,(k_\gamma ^{\pm 1})^{}& =& k_\gamma ^1,\hfill \\ \hfill e_\beta ^{}& =& (1)^{\vartheta (\beta )}e_\beta ,e_\beta ^{}& =& e_\beta \hfill \end{array}$$
(2.14)
($`(xy)^{}=(1)^{\mathrm{deg}x\mathrm{deg}y}y^{}x^{}`$ for any homogeneous elements $`x,yU_q(g)`$).
(iii) The Chevalley graded involution $`\omega `$:
$$\begin{array}{ccccc}\hfill \omega (q^{\pm 1})& =& q^1,\omega (k_\gamma ^{\pm 1})& =& k_\gamma ^{\pm 1},\hfill \\ \hfill \omega (e_\beta )& =& e_\beta ,\omega (e_\beta )& =& (1)^{\theta (\beta )}e_\beta .\hfill \end{array}$$
(2.15)
(iv) The Dynkin involution $`\tau `$ which is associated with the automorphism of the Dynkin diagrams of the (super)algebras $`A_1^{(1)}`$ and $`C(2)^{(2)}`$:
$$\begin{array}{ccccc}\hfill \tau (q^{\pm 1})& =& q^{\pm 1},\tau (k_\mathrm{d}^{\pm 1})& =& k_\mathrm{d}^{\pm 1},\hfill \\ \hfill \tau (k_\beta ^{\pm 1})& =& k_{\delta \beta }^{\pm 1},\tau (k_\beta ^{\pm 1})& =& k_{\delta +\beta }^{\pm 1},\hfill \\ \hfill \tau (e_\beta )& =& e_{\delta \beta },\tau (e_\beta )& =& e_{\delta +\beta }.\hfill \end{array}$$
(2.16)
Here in (2.13)–(2.16) $`\beta =\alpha ,\delta \alpha `$; $`\gamma =\mathrm{d},\beta `$.
It should be noted that the graded conjugation ”” and the Chevalley graded involution $`\omega `$ are involute (anti)automorphism of the fourth order, i.e., for example, $`(\omega )^4=id`$. Note also that the Dynkin involution $`\tau `$ commutes with all other three involutions, i.e. $`\tau (x^{})=(\tau (x))^{}`$, $`\tau (x^{})=(\tau (x))^{}`$ and $`\omega \tau (x)=\tau \omega (x)`$ for any element $`xU_q(g)`$ ($`g=A_1^{(1)},C(2,0)^{(2)}`$).
In the next Section we consider a q-analog of automorphisms connected with the Weyl group of the (super)algebras $`A_1^{(1)}`$ and $`C(2,0)^{(2)}`$. This q-analog defines actions of the braid group associated with the Weyl group.
## 3 Braid group actions
We introduce the morphisms $`T_\alpha `$ and $`T_{\delta \alpha }`$ defined by the following formulas:
$$\begin{array}{ccccc}\hfill T_\alpha (q^{\pm 1})& =& q^{\pm 1},T_\alpha (k_\gamma ^{\pm 1})& =& k_\gamma ^{\pm 1}k_\alpha ^{\frac{2(\alpha ,\gamma )}{(\alpha ,\alpha )}},\hfill \\ \hfill T_\alpha (e_\alpha )& =& e_\alpha k_\alpha ,T_\alpha (e_\alpha )& =& (1)^{\theta (\alpha )}k_\alpha ^1e_\alpha ,\hfill \\ \hfill T_\alpha (e_{\delta a})& =& \frac{1}{a}[e_\alpha ,[e_\alpha ,e_{\delta a}]_q]_q,& & \\ \hfill T_\alpha (e_{\delta +a)})& =& \frac{\left(1\right)^{\theta \left(\alpha \right)}}{a}[[e_{\delta +\alpha },e_\alpha ]_{q^1},e_\alpha ]_{q^1},& & \end{array}$$
(3.1)
$$\begin{array}{ccccc}\hfill T_{\delta \alpha }(q^{\pm 1})& =& q^{\pm 1},T_{\delta \alpha }(k_\gamma ^{\pm 1})& =& k_\gamma ^{\pm 1}k_{\delta \alpha }^{\frac{2(\delta \alpha ,\gamma )}{(\alpha ,\alpha )}},\hfill \\ \hfill T_{\delta \alpha }(e_{\delta \alpha })& =& e_{\delta +\alpha }k_{\delta a},T_{\delta \alpha }(e_{\delta +\alpha })& =& (1)^{\theta (\alpha )}k_{\delta a}^1e_{\delta \alpha },\hfill \\ \hfill T_{\delta \alpha }(e_\alpha )& =& \frac{1}{a}[e_{\delta \alpha },[e_{\delta \alpha },e_\alpha ]_q]_q,& & \\ \hfill T_{\delta \alpha }(e_\alpha )& =& \frac{\left(1\right)^{\theta \left(\alpha \right)}}{a}[e_\alpha ,e_{\delta +\alpha }]_{q^1},e_{\delta +\alpha }]_{q^1},& & \end{array}$$
(3.2)
where $`\gamma =\mathrm{d},\alpha ,\delta \alpha `$. It is not difficult to prove by direct verification that the morphisms $`T_\alpha ^1`$ and $`T_{\delta \alpha }^1`$ given by
$$\begin{array}{ccccc}\hfill T_\alpha ^1(q^{\pm 1})& =& q^{\pm 1},T_\alpha ^1(k_\gamma ^{\pm 1})& =& k_\gamma ^{\pm 1}k_\alpha ^{\frac{2(\alpha ,\gamma )}{(\alpha ,\alpha )}},\hfill \\ \hfill T_\alpha ^1(e_\alpha )& =& (1)^{\theta (\alpha )}k_\alpha ^1e_\alpha ,T_\alpha ^1(e_\alpha )& =& e_\alpha k_\alpha ,\hfill \\ \hfill T_\alpha ^1(e_{\delta a})& =& \frac{1}{a}[[e_{\delta \alpha },e_\alpha ]_q,e_\alpha ]_q,& & \\ \hfill T_\alpha ^1(e_{\delta +a)})& =& \frac{\left(1\right)^{\theta \left(\alpha \right)}}{a}[e_\alpha ,[e_\alpha ,e_{\delta +\alpha }]_{q^1}]_{q^1},& & \end{array}$$
(3.3)
$$\begin{array}{ccccc}\hfill T_{\delta \alpha }^1(q^{\pm 1})& =& q^{\pm 1},T_{\delta \alpha }^1(k_\gamma ^{\pm 1})& =& k_\gamma ^{\pm 1}k_{\delta \alpha }^{\frac{2(\delta \alpha ,\gamma )}{(\alpha ,\alpha )}},\hfill \\ \hfill T_{\delta \alpha }^1(e_{\delta \alpha })& =& (1)^{\theta (\alpha )}k_{\delta a}^1e_{\delta +\alpha },T_{\delta \alpha }^1(e_{\delta +\alpha })& =& e_{\delta \alpha }k_{\delta a},\hfill \\ \hfill T_{\delta \alpha }^1(e_\alpha )& =& \frac{1}{a}[[e_\alpha ,e_{\delta \alpha }]_q,e_{\delta \alpha }]_q,& & \\ \hfill T_{\delta \alpha }^1(e_\alpha )& =& \frac{\left(1\right)^{\theta \left(\alpha \right)}}{a}[e_{\delta +\alpha },[e_{\delta +\alpha },e_\alpha ]_{q^1}]_{q^1}& & \end{array}$$
(3.4)
are inverses to $`T_\alpha `$ and $`T_{\delta \alpha }`$, i.e.
$$T_\alpha T_\alpha ^1=T_\alpha ^1T_\alpha =1,T_{\delta \alpha }T_{\delta \alpha }^1=T_{\delta \alpha }^1T_{\delta \alpha }=1.$$
(3.5)
Here in (LABEL:WG1)–(LABEL:WG4) and in what follows we use the notation:
$$a:=[(\alpha ,\alpha )]_q=\frac{q^{(\alpha ,\alpha )}q^{(\alpha ,\alpha )}}{qq^1}.$$
(3.6)
###### Proposition 3.1
(i) The morphisms $`T_\alpha `$ and $`T_{\delta \alpha }`$ (and also $`T_\alpha ^1`$ and $`T_{\delta \alpha }^1`$) commute with the graded conjugation ””, i.e.
$$(T_\alpha (x))^{}=T_\alpha (x^{}),(T_{\delta \alpha }(x))^{}=T_{\delta \alpha }(x^{})$$
(3.7)
for any element $`xU_q(g)`$.
(ii) The morphisms $`T_\alpha ^{\pm 1}`$ and $`T_{\delta \alpha }^{\pm 1}`$ are also compatible with the Chevalley graded involution $`\omega `$, in sense that:
$$T_\alpha \omega =\omega T_\alpha ^1,T_{\delta \alpha }\omega =\omega T_{\delta \alpha }^1.$$
(3.8)
(iii) The morphisms $`T_\alpha ^{\pm 1}`$ and $`T_{\delta \alpha }^{\pm 1}`$ are connected with each other by the Dynkin involution $`\tau `$, in sense that:
$$T_\alpha \tau =\tau T_{\delta \alpha },T_\alpha ^1\tau =\tau T_{\delta \alpha }^1.$$
(3.9)
This proposition can be proved by direct verification for the Chevalley basis.
###### Proposition 3.2
The morphisms $`T_\alpha `$ and $`T_{\delta \alpha }`$ (and also $`T_\alpha ^1`$ and $`T_{\delta \alpha }^1`$) are the automorphisms of the quantum (super)algebras $`U_q(A_1^{(1)})`$ and $`U_q(C(2)^{(2)})`$.
Proof. The proposition is proved by direct verification that the defining relations remain valid under the actions of the given morphisms. To this end we apply Proposition 2.1. Note that under the action of $`T_\alpha `$ the relations (2.4) and (2.5) are transformed into each other, the relation (2.6) is transformed to (2.10). Analogously, under action of $`T_{\delta \alpha }`$ the relations (2.4) and (2.6) are transformed into each other, the relations (2.5) are transformed to (2.11). In addition, it is useful to apply the relations (3.7). $`\mathrm{}`$
It is easy to see that all these automorphisms $`T_\alpha ^{\pm 1}`$ and $`T_{\delta \alpha }^{\pm 1}`$ are not Hopf algebra automorphisms of $`U_q(A_1^{(1)})`$ and $`U_q(C(2)^{(2)})`$, in sense that, e.g., $`T_\alpha T_\alpha \mathrm{\Delta }_q\mathrm{\Delta }_qT_\alpha `$.
In the case of $`U_q(g)`$, where $`g`$ is a finite-dimensional simple Lie algebras, the automorphisms of type $`T_\alpha ^{\pm 1}`$ and $`T_{\delta \alpha }^{\pm 1}`$ are called the Lusztig automorphisms .
Introduce the following root vectors:
$$\begin{array}{ccccc}\hfill e_\delta & :=& [e_\alpha ,e_{\delta \alpha }]_q,e_\delta & :=& [e_{\delta +\alpha },e_\alpha ]_{q^1},\hfill \\ \hfill \stackrel{~}{e}_\delta & :=& [e_{\delta \alpha },e_\alpha ]_q,\stackrel{~}{e}_\delta & :=& [e_\alpha ,e_{\delta +\alpha }]_{q^1}.\hfill \end{array}$$
(3.10)
It is not difficult to verify that under the actions of the automorphisms $`T_\alpha ^{\pm 1}`$ and $`T_{\delta \alpha }^{\pm 1}`$ the elements $`e_{\pm \delta }`$ and $`\stackrel{~}{e}_{\pm \delta }`$ are transformed as follows:
$$\begin{array}{ccccc}\hfill T_\alpha (\stackrel{~}{e}_{\pm \delta })& =& (1)^{\theta (\alpha )}e_{\pm \delta },T_\alpha ^1(e_{\pm \delta })& =& (1)^{\theta (\alpha )}\stackrel{~}{e}_{\pm \delta },\hfill \\ \hfill T_{\delta \alpha }(e_{\pm \delta })& =& (1)^{\theta (\alpha )}\stackrel{~}{e}_{\pm \delta },T_{\delta \alpha }^1(\stackrel{~}{e}_{\pm \delta })& =& (1)^{\theta (\alpha )}e_{\pm \delta }.\hfill \end{array}$$
(3.11)
Therefore
$$\begin{array}{ccccc}\hfill T_{2\delta }(e_{\pm \delta })& =& e_{\pm \delta },T_{2\delta }^1(e_{\pm \delta })& =& e_{\pm \delta },\hfill \\ \hfill \stackrel{~}{T}_{2\delta }(\stackrel{~}{e}_{\pm \delta })& =& \stackrel{~}{e}_{\pm \delta },\stackrel{~}{T}_{2\delta }^1(\stackrel{~}{e}_{\pm \delta })& =& \stackrel{~}{e}_{\pm \delta },\hfill \end{array}$$
(3.12)
where the elements $`T_{2\delta }`$ and $`\stackrel{~}{T}_{2\delta }`$ called the translation operators are given by
$$\begin{array}{ccccc}\hfill T_{2\delta }& =& T_\alpha T_{\delta \alpha },T_{2\delta }^1& =& T_{\delta \alpha }^1T_a^1,\hfill \\ \hfill \stackrel{~}{T}_{2\delta }& =& T_{\delta \alpha }T_\alpha ,\stackrel{~}{T}_{2\delta }^1& =& T_\alpha ^1T_{\delta \alpha }^1.\hfill \end{array}$$
(3.13)
###### Proposition 3.3
The automorphisms $`T_\delta :=T_\alpha \tau `$, $`T_\delta ^1:=\tau T_\alpha ^1`$ and $`\stackrel{~}{T}_\delta :=T_{\delta \alpha }\tau `$, $`\stackrel{~}{T}_\delta ^1:=\tau T_{\delta \alpha }^1`$ are the square roots of the automorphisms of $`T_{2\delta }^{\pm 1}`$ and $`\stackrel{~}{T}_{2\delta }^{\pm 1}`$ correspondingly, i.e.
$$\begin{array}{ccccc}\hfill T_\delta ^2& =& T_{2\delta },T_\delta ^2& =& \stackrel{~}{T}_{2\delta }^1,\hfill \\ \hfill \stackrel{~}{T}_\delta ^2& =& \stackrel{~}{T}_{2\delta },\stackrel{~}{T}_\delta ^2& =& \stackrel{~}{T}_{2\delta }^1.\hfill \end{array}$$
(3.14)
Moreover
$$\begin{array}{ccccc}\hfill T_\delta (e_{\delta \alpha })& =& e_\alpha k_\alpha ,T_\delta (e_{\delta +\alpha })& =& (1)^{\theta (\alpha )}k_\alpha ^1e_\alpha ,\hfill \\ \hfill T_\delta ^1(e_\alpha )& =& (1)^{\theta (\alpha )}k_{\delta \alpha }^1e_{\delta +\alpha },T_\delta ^1(e_\alpha )& =& e_{\delta \alpha }k_{\delta \alpha },\hfill \end{array}$$
(3.15)
$$\begin{array}{ccccc}\hfill \stackrel{~}{T}_\delta (e_\alpha )& =& e_{\delta +\alpha }k_{\delta \alpha },\stackrel{~}{T}_\delta (e_\alpha )& =& (1)^{\theta (\alpha )}k_{\delta \alpha }^1e_{\delta \alpha },\hfill \\ \hfill \stackrel{~}{T}_\delta ^1(e_{\delta \alpha })& =& (1)^{\theta (\alpha )}k_\alpha ^1e_\alpha ,\stackrel{~}{T}_\delta ^1(e_{\delta +\alpha })& =& e_\alpha k_\alpha ,\hfill \end{array}$$
(3.16)
and also
$$\begin{array}{ccccc}\hfill T_\delta (e_{\pm \delta })& =& (1)^{\theta (\alpha )}e_{\pm \delta },T_\delta ^1(e_{\pm \delta })& =& (1)^{\theta (\alpha )}e_{\pm \delta },\hfill \\ \hfill \stackrel{~}{T}_\delta (\stackrel{~}{e}_{\pm \delta })& =& (1)^{\theta (\alpha )}\stackrel{~}{e}_{\pm \delta },\stackrel{~}{T}_\delta ^1(\stackrel{~}{e}_{\pm \delta })& =& (1)^{\theta (\alpha )}\stackrel{~}{e}_{\pm \delta }.\hfill \end{array}$$
(3.17)
Proof. From (3.9) we have that $`T_{\delta \alpha }=\tau T_\alpha \tau `$ and therefore $`T_{2d}=T_\alpha T_{\delta \alpha }=T_\alpha \tau T_\alpha \tau =T_\delta ^2`$. Analogously $`T_\alpha =\tau T_{\delta \alpha }\tau `$ and therefore $`\stackrel{~}{T}_{2d}=T_{\delta \alpha }T_\alpha =T_{\delta \alpha }\tau T_{\delta \alpha }\tau =\stackrel{~}{T}_\delta ^2`$. The formulas (3.15)–(3.17) are trivial.
In the next section we construct the Cartan-Weyl basis and describe its properties in detail.
## 4 Cartan-Weyl basis for $`U_q(A_1^{(1)})`$ and $`U_q(C(2)^{(2)})`$
A general scheme for construction of a Cartan-Weyl basis for quantized Lie algebras and superalgebras was proposed in Ref. . The scheme was applied in detail at first for quantized finite-dimensional Lie (super)algebras and then to quantized non-twisted affine algebras .
This procedure is bases on a notion of “normal ordering” for the reduced positive root system. For affine Lie (super)algebras this notation was formulated in (see also , ). In our case the reduced positive system has only two normal orderings:
$$\alpha ,\delta +\alpha ,2\delta +\alpha ,\mathrm{},\mathrm{}\delta +\alpha ,\delta ,2\delta ,3\delta ,\mathrm{},\mathrm{}\delta ,\mathrm{}\delta \alpha ,\mathrm{},3\delta \alpha ,2\delta \alpha ,\delta \alpha ,$$
(4.1)
$$\delta \alpha ,2\delta \alpha ,3\delta \alpha ,\mathrm{},\mathrm{}\delta \alpha ,\delta ,2\delta ,3\delta ,\mathrm{},\mathrm{}\delta ,\mathrm{}\delta +\alpha ,\mathrm{},2\delta +\alpha ,\delta +\alpha ,\alpha .$$
(4.2)
The first normal ordering (4.1) corresponds to “clockwise” ordering for positive roots in Fig. 1, 2b if we start from root $`\alpha `$ to root $`\delta \alpha `$. The inverse normal ordering (4.2) corresponds to “anticlockwise” ordering for the positive roots when we move from $`\delta \alpha `$ to $`\alpha `$.
In accordance with the normal ordering (4.1) we set
$`e_\delta `$ $`:=`$ $`[e_\alpha ,e_{\delta \alpha }]_q,e_\delta :=[e_{\delta +\alpha },e_\alpha ]_{q^1},`$ (4.3)
$`e_{n\delta +\alpha }`$ $`:=`$ $`\frac{1}{a}[e_{(n1)\delta +\alpha },e_\delta ],e_{n\delta \alpha }:=\frac{1}{a}[e_\delta ,e_{(n1)\delta \alpha }],`$ (4.4)
$`e_{(n+1)\delta \alpha }`$ $`:=`$ $`\frac{1}{a}[e_\delta ,e_{n\delta \alpha }],e_{(n+1)\delta +\alpha }:=\frac{1}{a}[e_{n\delta +\alpha },e_\delta ],`$ (4.5)
$`e_{n\delta }^{}`$ $`:=`$ $`[e_\alpha ,e_{n\delta \alpha }]_q,e_{n\delta }^{}:=[e_{n\delta +\alpha },e_\alpha ]_{q^1},`$ (4.6)
where $`n=1,2,\mathrm{}`$, and $`a`$ is given by the formula (3.6). Analogously for the inverse normal ordering (4.2) we set
$`\stackrel{~}{e}_\delta `$ $`:=`$ $`[e_{\delta \alpha },e_\alpha ]_q,\stackrel{~}{e}_\delta :=[e_\alpha ,e_{\delta +\alpha }]_{q^1},`$ (4.7)
$`\stackrel{~}{e}_{n\delta +\alpha }`$ $`:=`$ $`\frac{1}{a}[\stackrel{~}{e}_\delta ,\stackrel{~}{e}_{(n1)\delta +\alpha }],\stackrel{~}{e}_{n\delta \alpha }:=\frac{1}{a}[\stackrel{~}{e}_{(n1)\delta \alpha },\stackrel{~}{e}_\delta ],`$ (4.8)
$`\stackrel{~}{e}_{(n+1)\delta \alpha }`$ $`:=`$ $`\frac{1}{a}[\stackrel{~}{e}_{n\delta \alpha },\stackrel{~}{e}_\delta ],\stackrel{~}{e}_{(n+1)\delta +\alpha }:=\frac{1}{a}[\stackrel{~}{e}_\delta ,\stackrel{~}{e}_{n\delta +\alpha }],`$ (4.9)
$`\stackrel{~}{e}_{n\delta }^{}`$ $`:=`$ $`[e_{\delta \alpha },\stackrel{~}{e}_{(n1)\delta +\alpha }]_q,\stackrel{~}{e}_{n\delta }^{}:=[e_{\delta +\alpha },\stackrel{~}{e}_{(n1)\delta \alpha }]_{q^1},`$ (4.10)
where $`n=1,2,\mathrm{}`$. Thus, we have two systems of the Cartan-Weyl generators: ’direct’ and ’inverse’. Each such system together with the Cartan generators $`k_\alpha ^{\pm 1}`$, $`k_{\delta \alpha }^{\pm 1}`$, $`e_{\pm \alpha }`$, $`e_{\pm (\delta \alpha )}`$ are called the q-analog of the Cartan-Weyl basis (or simply the Cartan-Weyl basis) for the quantum (super)algebras $`U_q(A_1^{(1)})`$ and $`U_q(C(2)^{(2)})`$.
Now we consider some properties of these bases. First of all, the explicit construction of the Cartan-Weyl generators (4.3)–(4.6) (or (4.7)–(4.10)) permits easy to find their properties with respect to the (anti)involutions (2.13)–(2.15). For example, it is evident that
$$e_{\pm \gamma }^{}=e_\gamma ,\gamma \underset{¯}{\mathrm{\Delta }}_+.$$
(4.11)
and also
$$\begin{array}{ccccc}\hfill e_{n\delta +\alpha }^{}& =& (1)^{(n+1)\theta (\alpha )}e_{n\delta \alpha },e_{n\delta \alpha }^{}& =& (1)^{n\theta (\alpha )}e_{n\delta +\alpha },\hfill \\ \hfill e_{n\delta \alpha }^{}& =& (1)^{n\theta (\alpha )}e_{n\delta +\alpha },e_{n\delta +\alpha }^{}& =& (1)^{(n1)\theta (\alpha )}e_{n\delta \alpha },\hfill \\ \hfill e_{n\delta }^{}& =& (1)^{n\theta (\alpha )}e_{n\delta },e_{n\delta }^{}& =& (1)^{n\theta (\alpha )}e_{n\delta }.\hfill \end{array}$$
(4.12)
Further, it is easy to see that the ’direct’ and ’inverse’ Cartan-Weyl generators (4.3)–(4.6) and (4.7)–(4.10) have very simple connection by the Dynkin involution $`\tau `$:
$$\begin{array}{cccccc}\hfill \tau (e_{n\delta +\alpha })& =& \stackrel{~}{e}_{(n+1)\delta \alpha },\tau (\stackrel{~}{e}_{n\delta +\alpha })& =& e_{(n+1)\delta \alpha }& (nZZ),\hfill \\ \hfill \tau (e_{n\delta \alpha })& =& \stackrel{~}{e}_{(n1)\delta +\alpha },\tau (\stackrel{~}{e}_{n\delta \alpha })& =& e_{(n1)\delta +\alpha }& (nZZ),\hfill \\ \hfill \tau (e_{n\delta })& =& \stackrel{~}{e}_{n\delta },\tau (\stackrel{~}{e}_{n\delta })& =& e_{n\delta }& (n0).\hfill \end{array}$$
(4.13)
The transformation properties with respect to the automorphisms $`T_\alpha `$ and $`T_{\delta \alpha }`$ can be not hard obtained with the help of (LABEL:WG1), (LABEL:WG2), (3.11), and they have the form
$$\begin{array}{ccccc}\hfill T_\alpha (\stackrel{~}{e}_{n\delta +\alpha })& =& (1)^{(n+1)\theta (\alpha )}e_{n\delta \alpha },T_\alpha (\stackrel{~}{e}_{n\delta \alpha })& =& (1)^{n\theta (\alpha )}e_{n\delta +\alpha },\hfill \\ \hfill T_\alpha (\stackrel{~}{e}_{n\delta \alpha })& =& (1)^{(n1)\theta (\alpha )}e_{n\delta +\alpha },T_\alpha (\stackrel{~}{e}_{n\delta +\alpha })& =& (1)^{n\theta (\alpha )}e_{n\delta \alpha },\hfill \\ \hfill T_\alpha (\stackrel{~}{e}_{n\delta })& =& (1)^{n\theta (\alpha )}e_{n\delta },T_\alpha (\stackrel{~}{e}_{n\delta })& =& (1)^{n\theta (\alpha )}e_{n\delta },\hfill \end{array}$$
(4.14)
where $`n>0`$, and
$$\begin{array}{ccccc}\hfill T_{\delta \alpha }(e_{k\delta +\alpha })& =& (1)^{k\theta (\alpha )}\stackrel{~}{e}_{(k+2)\delta \alpha },T_{\delta \alpha }(e_{k\delta \alpha })& =& (1)^{(k+1)\theta (\alpha )}\stackrel{~}{e}_{(k+2)\delta +\alpha },\hfill \\ \hfill T_{\delta \alpha }(e_{l\delta \alpha })& =& (1)^{l\theta (\alpha )}\stackrel{~}{e}_{(l2)\delta +\alpha },T_{\delta \alpha }(e_{l\delta +\alpha })& =& (1)^{(l1)\theta (\alpha )}\stackrel{~}{e}_{(l+2)\delta \alpha },\hfill \\ \hfill T_{\delta \alpha }(e_{m\delta })& =& (1)^{m\theta (\alpha )}\stackrel{~}{e}_{m\delta },T_{\delta \alpha }(e_{m\delta })& =& (1)^{m\theta (\alpha )}\stackrel{~}{e}_{m\delta },\hfill \end{array}$$
(4.15)
for $`k0`$, $`l>1`$, $`m>0`$. As corollary of the formulas (4.13)–(4.15) we easy find the actions of the translation operators $`T_\delta `$ and $`\stackrel{~}{T}_\delta `$:
$$\begin{array}{ccccc}\hfill T_\delta (e_{k\delta +\alpha })& =& (1)^{k\theta (\alpha )}e_{(k+1)\delta +\alpha },T_\delta (e_{k\delta \alpha })& =& (1)^{(k+1)\theta (\alpha )}e_{(k+1)\delta \alpha },\hfill \\ \hfill T_\delta (e_{l\delta \alpha })& =& (1)^{l\theta (\alpha )}e_{(l1)\delta \alpha },T_\delta (e_{l\delta +\alpha })& =& (1)^{(l1)\theta (\alpha )}e_{(l+1)\delta +\alpha },\hfill \\ \hfill T_\delta (e_{m\delta })& =& (1)^{m\theta (\alpha )}e_{m\delta },T_\delta (e_{m\delta })& =& (1)^{m\theta (\alpha )}e_{m\delta }\hfill \end{array}$$
(4.16)
for $`k0`$, $`l>1`$, $`m>0`$, and
$$\begin{array}{ccccc}\hfill \stackrel{~}{T}_\delta (\stackrel{~}{e}_{n\delta +\alpha })& =& (1)^{(n1)\theta (\alpha )}\stackrel{~}{e}_{(n1)\delta +\alpha },\stackrel{~}{T}_\delta (\stackrel{~}{e}_{n\delta \alpha })& =& (1)^{n\theta (\alpha )}\stackrel{~}{e}_{(n+1)\delta +\alpha },\hfill \\ \hfill \stackrel{~}{T}_\delta (\stackrel{~}{e}_{n\delta \alpha })& =& (1)^{(n+1)\theta (\alpha )}\stackrel{~}{e}_{(n+1)\delta \alpha },T_\delta (\stackrel{~}{e}_{n\delta +\alpha })& =& (1)^{n\theta (\alpha )}\stackrel{~}{e}_{(n+1)\delta +\alpha },\hfill \\ \hfill \stackrel{~}{T}_\delta (\stackrel{~}{e}_{n\delta })& =& (1)^{n\theta (\alpha )}\stackrel{~}{e}_{n\delta },\stackrel{~}{T}_\delta (\stackrel{~}{e}_{n\delta })& =& (1)^{n\theta (\alpha )}\stackrel{~}{e}_{n\delta },\hfill \end{array}$$
(4.17)
where $`n>0`$. (Also see (3.15) and (3.17)). Using the formulas (4.16)–(4.17) we can easy find the actions for the inverse translation operators $`T_\delta ^1`$, $`\stackrel{~}{T}_\delta ^1`$ and $`T_{2\delta }^1`$, $`\stackrel{~}{T}_{2\delta }^1`$. These actions are not written here. From the relations (4.16)–(4.17) it is clear that the operators $`T_\delta ^{\pm 1}`$ and $`\stackrel{~}{T}_\delta ^{\pm 1}`$ can be used for construction of the Cartan-Weyl generators (4.3)–(4.6) starting from the Chevalley basis. In the case of the quantum untwisted affine algebras the similar procedure was applied in the paper .
###### Proposition 4.1
The root vectors (4.3)–(4.6) satisfy the following permutation relations:
$$\begin{array}{ccccc}\hfill k_\mathrm{d}e_{n\delta \pm \alpha }k_\mathrm{d}^1& =& q^{n(\mathrm{d},\delta )}e_{n\delta \pm \alpha },k_\mathrm{d}e_{n\delta }^{}k_\mathrm{d}^1& =& q^{n(\mathrm{d},\delta )}e_{n\delta }^{},\hfill \\ \hfill k_\gamma e_{n\delta \pm \alpha }k_\gamma & =& q^{\pm (\gamma ,\alpha )}e_{n\delta \pm \alpha },k_\gamma e_{n\delta }^{}k_\gamma ^1& =& e_{n\delta }\hfill \end{array}$$
(4.18)
for any $`nZZ`$ and any $`\gamma \underset{¯}{\mathrm{\Delta }}_+`$, and also
$`[e_{n\delta +\alpha },e_{n\delta \alpha }]`$ $`=`$ $`(1)^{n\theta (\alpha )}\frac{k_{n\delta +\alpha }k_{n\delta +\alpha }^1}{qq^1}(n0),`$ (4.19)
$`[e_{n\delta \alpha },e_{n\delta +\alpha }]`$ $`=`$ $`(1)^{(n1)\theta (\alpha )}\frac{k_{n\delta \alpha }k_{n\delta \alpha }^1}{qq^1}(n>0);`$ (4.20)
$`[e_{n\delta +\alpha },e_{(n+2m1)\delta +\alpha }]_q`$ $`=`$ $`(q_\alpha ^21){\displaystyle \underset{l=1}{\overset{m1}{}}}q_\alpha ^le_{(n+l)\delta +\alpha }e_{(n+2m1l)\delta +\alpha },`$ (4.21)
$$\begin{array}{ccc}\hfill [e_{n\delta +\alpha },e_{(n+2m)\delta +\alpha }]_q& =& (q_\alpha 1)q_\alpha ^{m+1}e_{(n+m)\delta +\alpha }^2+\hfill \\ & & +(q_\alpha ^21)\underset{l=1}{\overset{m1}{}}q_\alpha ^le_{(n+l)\delta +\alpha }e_{(n+2ml)\delta +\alpha }\hfill \end{array}$$
(4.22)
for any integers $`n0,m>0`$;
$$[e_{(n+2m1)\delta \alpha },e_{n\delta \alpha }]_q=(q_\alpha ^21)\underset{l=1}{\overset{m1}{}}q_\alpha ^le_{(n+2m1l)\delta \alpha }e_{(n+l)\delta \alpha },$$
(4.23)
$$\begin{array}{ccc}\hfill [e_{(n+2m)\delta \alpha },e_{n\delta \alpha }]_q& =& (q_\alpha 1)q_\alpha ^{m+1}e_{(n+m)\delta \alpha }^2\hfill \\ & & (q_\alpha ^21)\underset{l=1}{\overset{m1}{}}q_\alpha ^le_{(n+l)\delta \alpha }e_{(n+2ml)\delta \alpha }\hfill \end{array}$$
(4.24)
for any integers $`n,m>0`$;
$$\begin{array}{ccc}\hfill [e_{n\delta +\alpha },e_{(n+2m1)\delta +\alpha }]& =& (1)^{(n1)\theta (\alpha )}(q_\alpha ^21)\times \hfill \\ & & \times \underset{l=n}{\overset{n+m1}{}}q_\alpha ^lk_{n\delta \alpha }e_{(ln)\delta +\alpha }e_{(n+2m1l)\delta +\alpha }+\hfill \\ & & +(q_\alpha ^21)\underset{l=1}{\overset{n1}{}}(1)^{l\theta (\alpha )}q_\alpha ^lk_\delta ^le_{(n+l)\delta +\alpha }e_{(n+2m1l)\delta +\alpha },\hfill \end{array}$$
(4.25)
$$\begin{array}{ccc}\hfill [e_{n\delta +\alpha },e_{(n+2m)\delta +\alpha }]& =& (1)^{(n1)\theta (\alpha )}(q_\alpha ^21)\times \hfill \\ & & \times \underset{l=n}{\overset{n+m1}{}}q_\alpha ^lk_{n\delta \alpha }e_{(n+l)\delta +\alpha }e_{(n+2m1l)\delta +\alpha }+\hfill \\ & & +(q_\alpha ^21)\underset{l=1}{\overset{n1}{}}(1)^{l\theta (\alpha )}q_\alpha ^lk_\delta ^le_{(n+l)\delta +\alpha }e_{(n+2m1l)\delta +\alpha }\hfill \\ & & (1)^{(n1)\theta (\alpha )}(q_\alpha 1)q_\alpha ^{mn+1}k_{n\delta \alpha }e_{m\delta +\alpha }^2\hfill \end{array}$$
(4.26)
for any integers $`n,m0`$;
$$\begin{array}{ccc}\hfill [e_{(n+2m1)\delta \alpha },e_{n\delta \alpha }]& =& (1)^{(n+1)\theta (\alpha )}(q_\alpha ^21)\times \hfill \\ & & \times \underset{l=n+1}{\overset{n+m1}{}}q_\alpha ^le_{(n+2m1l)\delta +\alpha }e_{(ln)\delta \alpha }k_{n\delta +\alpha }^1\hfill \\ & & (q_\alpha ^21)\underset{l=1}{\overset{n1}{}}(1)^{l\theta (\alpha )}q_\alpha ^le_{(n+2m1l)\delta \alpha }e_{(n+l)\delta \alpha }k_\delta ^l,\hfill \end{array}$$
(4.27)
$$\begin{array}{ccc}\hfill [e_{(n+2m)\delta \alpha },e_{n\delta \alpha }]& =& (1)^{(n+1)\theta (\alpha )}(q_\alpha ^21)\times \hfill \\ & & \times \underset{l=n}{\overset{n+m1}{}}q_\alpha ^le_{(n+2ml)\delta \alpha }e_{(ln)\delta +\alpha }k_{n\delta +\alpha }^1\hfill \\ & & (q_\alpha ^21)\underset{l=1}{\overset{n1}{}}(1)^{l\theta (\alpha )}q_\alpha ^le_{(n+2ml)\delta \alpha }e_{(n+l)\delta \alpha }k_\delta ^l+\hfill \\ & & +(1)^{(n1)\theta (\alpha )}(q_\alpha 1)q_\alpha ^{mn+1}e_{m\delta \alpha }^2k_{n\delta +\alpha }^1\hfill \end{array}$$
(4.28)
for any integers $`n0,m>0`$;
$`[e_{n\delta +\alpha },e_{m\delta \alpha }]_q`$ $`=`$ $`e_{(n+m)\delta }^{}(n0,m>0),`$ (4.29)
$`[e_{n\delta +\alpha },e_{m\delta \alpha }]`$ $`=`$ $`(1)^{(m+1)\theta (\alpha )}e_{(nm)\delta }^{}k_{m\delta +\alpha }^1(n>m0),`$ (4.30)
$`[e_{m\delta +\alpha },e_{n\delta \alpha }]`$ $`=`$ $`(1)^{(m1)\theta (\alpha )}k_{m\delta \alpha }e_{(nm)\delta }^{}(n>m>0),`$ (4.31)
$`[e_{n\delta }^{},e_{m\delta }^{}]`$ $`=`$ $`[e_{n\delta }^{},e_{m\delta }^{}]=0(n>0,m>0),`$ (4.32)
$$[e_{n\delta +\alpha },e_{m\delta }^{}]=q_\alpha ^{m+1}ae_{(n+m)\delta +\alpha }+(q_\alpha ^21)\underset{l=1}{\overset{m1}{}}q_\alpha ^le_{(n+l)\delta +\alpha }e_{(ml)\delta }^{}$$
(4.33)
for any integers $`n0,m>0`$;
$$[e_{m\delta }^{},e_{n\delta \alpha }]=q_\alpha ^{m+1}ae_{(n+m)\delta \alpha }+(q_\alpha ^21)\underset{l=1}{\overset{m1}{}}q_\alpha ^le_{(ml)\delta }^{}e_{(n+l)\delta \alpha }$$
(4.34)
for any integers $`n,m>0`$;
$$\begin{array}{ccc}\hfill [e_{n\delta +\alpha },e_{m\delta }^{}]& =& (1)^{(n1)\theta (\alpha )}q_\alpha ^{m+1}ak_{n\delta \alpha }e_{(mn)\delta +\alpha }\hfill \\ & & (1)^{(n1)\theta (\alpha )}(q_\alpha ^21)k_{n\delta \alpha }\underset{l=n}{\overset{m1}{}}q_\alpha ^le_{(ln)\delta +\alpha }e_{(ml)\delta }^{}+\hfill \\ & & +(q_\alpha ^21)\underset{l=1}{\overset{n1}{}}(1)^{l\theta (\alpha )}q_\alpha ^lk_\delta ^le_{(n+l)\delta +\alpha }e_{(ml)\delta }^{}\hfill \end{array}$$
(4.35)
for any integers $`mn>0`$;
$$\begin{array}{ccc}\hfill [e_{n\delta +\alpha },e_{m\delta }^{}]& =& (1)^{m\theta (\alpha )}q_\alpha ^{m+1}ak_\delta ^me_{(n+m)\delta +\alpha }+\hfill \\ & & +(q_\alpha ^21)\underset{l=1}{\overset{m1}{}}(1)^{l\theta (\alpha )}q_\alpha ^lk_\delta ^le_{(n+l)\delta +\alpha }e_{(ml)\delta }^{}\hfill \end{array}$$
(4.36)
for any integers $`n>m>0`$;
$$\begin{array}{ccc}\hfill [e_{m\delta }^{}e_{n\delta \alpha }]& =& (1)^{(n+1)\theta (\alpha )}q_\alpha ^{m+1}ae_{(mn)\delta \alpha }k_{n\delta +\alpha }^1\hfill \\ & & (1)^{(n+1)\theta (\alpha )}(q_\alpha ^21)\underset{l=n+1}{\overset{m1}{}}q_\alpha ^le_{(ml)\delta }^{}e_{(ln)\delta \alpha }k_{n\delta +\alpha }^1+\hfill \\ & & +(q_\alpha ^21)\underset{l=1}{\overset{n}{}}(1)^{l\theta (\alpha )}q_\alpha ^le_{(ml)\delta }^{}e_{(n+l)\delta \alpha }k_\delta ^l\hfill \end{array}$$
(4.37)
for any integers $`m>n0`$;
$$\begin{array}{ccc}\hfill [e_{m\delta }^{}e_{n\delta \alpha }]& =& (1)^{m\theta (\alpha )}q_\alpha ^{m+1}ae_{(n+m)\delta \alpha }k_\delta ^m+\hfill \\ & & +(q_\alpha ^21)\underset{l=1}{\overset{m1}{}}(1)^{l\theta (\alpha )}q_\alpha ^le_{(ml)\delta }^{}e_{(n+l)\delta \alpha }k_\delta ^l\hfill \end{array}$$
(4.38)
for any integers $`nm>0`$.
Here in the relations (4.21)–(4.38) and in what follows $`q_\alpha :=(1)^{\theta (\alpha )}q^{(\alpha ,\alpha )}`$.
Outline of proof. First of all, the formulas (4.18) are trivial. The relations (4.19) and (4.20) are obtained by application the translation operators $`T_\delta ^n`$ and $`T_\delta ^n`$ to the relations (2.3). Further, in terms of the generators (4.3)–(4.6) the relation (2.5) means that $`[e_\alpha ,e_{\delta +\alpha }]_q=0`$. Applying to it the operator $`T_\delta ^n`$, we obtain the relation (4.21) for $`m=1`$. In the case $`m>1`$ the formulas (4.21) and (4.22) are proved for arbitrary $`m`$ by induction. If we apply the operator $`T_\delta ^k`$ to the relations (4.21) and (4.22) for $`n=0`$, then in the case $`k<m`$ we obtain the relations (4.25) and (4.26), in the case $`m<k<2m`$ we obtain the relations which are obtained from (4.27) and (4.28) by the conjugation ””, and finally for $`k>2m`$ we get the relations which are obtained from (4.23) and (4.24) by the conjugation ””. Further, the relation (4.29) for $`n=0`$ is trivial (see (4.6)). Applying to (4.29) with $`n=0`$ the operators $`T_\delta ^n`$, we can obtain for any $`n>0`$ and $`m>0`$ the relation (4.29) as well as the relation (4.30). The relation (4.31) can be obtained from (4.29) by repeated application of the operator $`T_\delta ^1`$. The relations (4.33) in the case $`n=0`$ and (4.34) in the case $`n=1`$ are proved by direct verification with the help of the previous results. Repeated application of the operators $`T_\delta ^{\pm 1}`$ to these relations results in the general case $`n,m>0`$. The relation (4.32) is proved by direct verification with the help of the relations (4.33) and (4.34). At last, the relations (4.35)–(4.38) can be obtained from (4.33) and (4.34) by repeated application of the operator $`T_\delta ^1`$. $`\mathrm{}`$
The imaginary root vectors $`e_{n\delta }^{}`$ do not satisfy the relations of the type (4.19) and therefore we introduce the new imaginary roots vectors $`e_{\pm n\delta }`$ by the following (Schur) relations:
$$e_{n\delta }^{}=\underset{p_1+2p_2+\mathrm{}+np_n=n}{}\frac{\left(\left(1\right)^{\theta \left(\alpha \right)}\left(qq^1\right)\right)^{{\scriptscriptstyle p_i}1}}{p_1!\mathrm{}p_n!}e_\delta ^{p_1}\mathrm{}e_{n\delta }^{p_n}.$$
(4.39)
In terms of generating functions
$`^{}(u)`$ $`:=`$ $`(1)^{\theta (\alpha )}(qq^1){\displaystyle \underset{n1}{}}e_{n\delta }^{}u^n,`$ (4.40)
$`(u)`$ $`=`$ $`(1)^{\theta (\alpha )}(qq^1){\displaystyle \underset{n1}{}}e_{n\delta }u^n`$ (4.41)
the relation (4.39) may be rewritten in the form
$$^{}(u)=1+\mathrm{exp}(u)$$
(4.42)
or
$$(u)=\mathrm{ln}(1+^{}(u)).$$
(4.43)
This provides a formula inverse to (4.39)
$$e_{n\delta }=\underset{p_1+2p_2+\mathrm{}+np_n=n}{}\frac{\left(\left(1\right)^{\theta \left(\alpha \right)}\left(q^1q\right)\right)^{{\scriptscriptstyle p_i}1}\left(_{i=1}^np_i1\right)!}{p_1!\mathrm{}p_n!}(e_\delta ^{})^{p_1}\mathrm{}(e_{n\delta }^{})^{p_n}.$$
(4.44)
The new root vectors corresponding to negative roots are obtained by the Cartan conjugation $`(^{})`$:
$$e_{n\delta }=(e_{n\delta })^{}.$$
(4.45)
###### Proposition 4.2
The new root vectors $`e_{\pm n\delta }`$ satisfy the following commutation relations:
$`[e_{n\delta +\alpha },e_{m\delta }]`$ $`=`$ $`(1)^{(m1)\theta (\alpha )}a(m)e_{(n+m)\delta +\alpha }(n0,m>0),`$ (4.46)
$`[e_{m\delta },e_{n\delta \alpha }]`$ $`=`$ $`(1)^{(m1)\theta (\alpha )}a(m)e_{(n+m)\delta \alpha }(n,m>0),`$ (4.47)
$`[e_{n\delta +\alpha },e_{m\delta }]`$ $`=`$ $`(1)^{(n+m)\theta (\alpha )}a(m)k_{n\delta \alpha }e_{(mn)\delta +\alpha }(mn>0),`$ (4.48)
$`[e_{n\delta +\alpha },e_{m\delta }]`$ $`=`$ $`(1)^{\theta (\alpha )}a(m)k_\delta ^me_{(n+m)\delta +\alpha }(n>m>0),`$ (4.49)
$`[e_{m\delta },e_{n\delta \alpha }]`$ $`=`$ $`(1)^{(n+m)\theta (\alpha )}a(m)e_{(mn)\delta \alpha }k_{n\delta +\alpha }^1(m>n0),`$ (4.50)
$`[e_{m\delta },e_{n\delta +\alpha }]`$ $`=`$ $`(1)^{\theta (\alpha )}a(m)e_{(n+m)\delta \alpha }k_\delta ^m(nm>0),`$ (4.51)
$`[e_{n\delta },e_{m\delta }]`$ $`=`$ $`\delta _{nm}a(m)\frac{k_\delta ^mk_\delta ^m}{qq^1}(n,m>0),`$ (4.52)
where
$$a(m):=\frac{q^{m(\alpha ,\alpha )}q^{m(\alpha ,\alpha )}}{m\left(qq^1\right)}.$$
(4.53)
This can be proved by direct calculation, applying the relations of Proposition (4.1) and the actions of the translation operators $`T_\delta ^{\pm 1}`$.
All the relations of Propositions (4.1), (4.53) together with the ones obtained from them by the conjugation describe complete list of the permutation relations of the Cartan-Weyl bases corresponding to the ’direct’ normal ordering (4.1). Applying to these relations the Dynkin involution $`\tau `$, it is easy to obtain these results for the ’inverse’ normal ordering (4.2).
## 5 Extremal projector for $`U_q(A_1^{(1)})`$ and $`U_q(C(2)^{(2)})`$
A general formula for the extremal projector for quantized contragredient Lie (super)algebras of finite growth was presented in Refs. , , . Here we specialize this result to our case $`U_q(g)`$, where $`g=A_1^{(1)},C(2)^{(2)}`$.
By definition, the extremal projector for $`U_q(g)`$ is a nonzero element $`p:=p(U_q(g))`$ of the Taylor extension $`T_q(g)`$ of $`U_q(g)`$ (see Refs. , , ), satisfying the equations
$$e_\alpha p=pe_\alpha =0,e_{\delta \alpha }p=pe_{\delta +\alpha }=0,p^2=p.$$
(5.1)
The explicit expression of the extremal projector $`p`$ for our case $`U_q(g)`$ can be presented as follows:
$$p=p_+p_0p_{},$$
(5.2)
where the factors $`p_+`$, $`p_0`$ and $`p_+`$ have the following form
$$p_+=\underset{n0}{\overset{}{}}p_{n\delta +\alpha },p_0=\underset{n1}{}p_{n\delta },p_{}=\underset{n1}{\overset{}{}}p_{n\delta \alpha }.$$
(5.3)
The elements $`p_\gamma `$ are given by the formula
$`p_{n\delta +\alpha }`$ $`=`$ $`{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}\frac{\left(1\right)^m}{\left(m\right)_{\overline{q}_\alpha }!}\phi _{n,m}^+e_{n\delta \alpha }^me_{n\delta +\alpha }^m,`$ (5.4)
$`p_{n\delta }`$ $`=`$ $`{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}\frac{\left(1\right)^m}{m!}\phi _{n,m}^{\mathrm{\hspace{0.17em}0}}e_{n\delta }^me_{n\delta }^m,`$ (5.5)
$`p_{n\delta \alpha }`$ $`=`$ $`{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}\frac{\left(1\right)^m}{\left(m\right)_{\overline{q}_\alpha }!}\phi _{n,m}^{}e_{n\delta +\alpha }^me_{n\delta \alpha }^m,`$ (5.6)
where the coefficients $`\phi _m^+`$, $`\phi _m^0`$ and $`\phi _m^{}`$ are determined as follows:
$`\phi _{n,m}^+`$ $`=`$ $`\frac{\left(1\right)^{mn\theta \left(\alpha \right)}\left(qq^1\right)^mq^{m\left(\frac{m1}{4}+n\right)(\alpha ,\alpha )}}{\underset{r=1}{\overset{m}{}}\left(k_{n\delta +\alpha }q^{\left(n+\frac{1}{2}+\frac{r}{2}\right)(\alpha ,\alpha )}\left(1\right)^{\left(r1\right)\theta \left(\alpha \right)}k_{n\delta +\alpha }^1q^{\left(n+\frac{1}{2}+\frac{r}{2}\right)(\alpha ,\alpha )}\right)},`$ (5.7)
$`\phi _{n,m}^{\mathrm{\hspace{0.17em}0}}`$ $`=`$ $`\frac{n^m\left(qq^1\right)^{n+m}q^{mn(\alpha ,\alpha )}}{\left(q^{n(\alpha ,\alpha )}q^{n(\alpha ,\alpha )}\right)^m\left(k_\delta ^nq^{n(\alpha ,\alpha )}k_\delta ^nq^{n(\alpha ,\alpha )}\right)^m},`$ (5.8)
$`\phi _{n,m}^{}`$ $`=`$ $`\frac{\left(1\right)^{m\left(n1\right)\theta \left(\alpha \right)}\left(qq^1\right)^mq^{m\left(\frac{m5}{4}+n\right)(\alpha ,\alpha )}}{\underset{r=1}{\overset{m}{}}\left(k_{n\delta \alpha }q^{\left(n\frac{1}{2}+\frac{r}{2}\right)(\alpha ,\alpha )}\left(1\right)^{\left(r1\right)\theta \left(\alpha \right)}k_{n\delta \alpha }^1q^{\left(n\frac{1}{2}+\frac{r}{2}\right)(\alpha ,\alpha )}\right)}.`$ (5.9)
Here in the relations (5.4)–(5.6) and in what follows we use the notation $`\overline{q}_\alpha :=(1)^{\theta (\alpha )}q^{(\alpha ,\alpha )}`$, and the symbol $`(m)_{\overline{q}_\alpha }`$ is defined by the formula (6.4).
Acting by the extremal projector $`p`$ on any highest weight $`U_q(g)`$-module $`M`$ we obtain a space $`M^0=pM`$ of highest weight vectors for $`M`$ if $`pM`$ has no singularities. An effective example of application of the extremal projector for the case of the quantum algebra $`U_q(gl(n,𝐂\mathrm{l})`$ can be found in Ref. .
## 6 Universal $`R`$-matrix for $`U_q(A_1^{(1)})`$ and $`U_q(C(2)^{(2)})`$
Any quantum (super)algebra $`U_q(g)`$ is a non-cocommutative Hopf (super)algebra which has the intertwining operator called the universal $`R`$-matrix. By definition , the universal $`R`$-matrix for the Hopf (super)algebra $`U_q(g)`$ is an invertible element $`R`$ of the Tylor extension $`T_q(g)T_q(g)`$ of $`U_q(g)U_q(g)`$ (see Refs. ), satisfying the equations
$`\stackrel{~}{\mathrm{\Delta }}_q(a)`$ $`=`$ $`R\mathrm{\Delta }_q(a)R^1aU_q(g),`$ (6.1)
$`(\mathrm{\Delta }_qid)R`$ $`=`$ $`R^{13}R^{23},(id\mathrm{\Delta }_q)R=R^{13}R^{12},`$ (6.2)
where $`\stackrel{~}{\mathrm{\Delta }}_q`$ is the opposite comultiplication: $`\stackrel{~}{\mathrm{\Delta }}_q=\sigma \mathrm{\Delta }_q`$, $`\sigma (ab)=(1)^{\mathrm{deg}a\mathrm{deg}b}ba`$ for all homogeneous elements $`a,bU_q(g)`$. In the relation (6.2) we use the standard notations $`R^{12}=a_ib_iid`$, $`R^{13}=a_iidb_i`$, $`R^{23}=ida_ib_i`$ if $`R`$ has the form $`R=a_ib_i`$. We employ the following standard notation for the q-exponential:
$$\mathrm{exp}_q(x):=1+x+\frac{x^2}{\left(2\right)_q!}+\mathrm{}+\frac{x^n}{\left(n\right)_q!}+\mathrm{}=\underset{n0}{}\frac{x^n}{\left(n\right)_q!},$$
(6.3)
where
$$(n)_q:=\frac{q^n1}{q1}.$$
(6.4)
A general formula for the universal $`R`$-matrix $`R`$ for quantized contragredient Lie (super)algebras was presented in Refs. . Here we specialize this result to our case $`U_q(g)`$, where $`g=A_1^{(1)},C(2)^{(2)}`$.
The explicit expression of the universal $`R`$-matrix $`R`$ for our case $`U_q(g)`$ can be presented as follows:
$$R=R_+R_0R_{}K.$$
(6.5)
Here the factors $`K`$ and $`R_\pm `$ have the following form
$`K`$ $`=`$ $`q^{\frac{1}{(\alpha ,\alpha )}h_\alpha h_\alpha +h_\delta h_\mathrm{d}+h_\mathrm{d}h_\delta },`$ (6.6)
$`R_+`$ $`=`$ $`{\displaystyle \underset{n0}{\overset{}{}}}_{n\delta +\alpha },R_{}={\displaystyle \underset{n1}{\overset{}{}}}_{n\delta \alpha }.`$ (6.7)
The elements $`_\gamma `$ are given by the formula
$$_\gamma =\mathrm{exp}_{\overline{q}_\gamma }\left(A(\gamma )(qq^1)(e_\gamma e_\gamma )\right),$$
(6.8)
where
$`A(\gamma )`$ $`=`$ $`\{\begin{array}{ccc}& (1)^{n\theta (\alpha )}\hfill & \mathrm{if}\gamma =n\delta +\alpha ,\hfill \\ & (1)^{(n1)\theta (\alpha )}\hfill & \mathrm{if}\gamma =n\delta \alpha .\hfill \end{array}`$ (6.11)
Finally, the factor $`R_0`$ is defined as follows
$$R_0=\mathrm{exp}\left((qq^1)\underset{n>0}{}d(n)e_{n\delta }e_{n\delta }\right),$$
(6.12)
where $`d(n)`$ is the inverse to $`a(n)`$, i.e.
$$d(n)=\frac{n\left(qq^1\right)}{q^{n(\alpha ,\alpha )}q^{n(\alpha ,\alpha )}}.$$
(6.13)
## 7 The ’new realization’
Let us denote by $`d`$ the Cartan element $`h_\mathrm{d}`$ and by $`c`$ the Cartan element $`h_\delta `$, emphasizing that $`d`$ defines homogeneous gradation of the algebra and $`k_\delta =q^{h_\delta }`$ is the central element. It will be convenient in the following to add its square roots $`q^{\pm \frac{c}{2}}=k_\delta ^{\pm \frac{1}{2}}`$. Let us introduce the new notations: $`e_n:=e_{n\delta +\alpha }`$ ($`n0`$), $`e_n:=(1)^{(n1)\theta (\alpha )}k_{n\delta +\alpha }e_{n\delta +\alpha }`$ ($`n>0`$), and $`f_n:=e_{n\delta \alpha }k_{n\delta \alpha }`$ ($`n>0`$), $`f_n:=(1)^{(n+1)\theta (\alpha )}e_{n\delta \alpha }`$ ($`n0`$). We also put $`a_n:=e_{n\delta }q^{\frac{nc}{2}}`$ ($`n1`$), and $`a_n:=(1)^{n\theta (\alpha )}e_{n\delta }q^{\frac{nc}{2}}`$ ($`n1`$). Collect the elements $`e_n`$, $`f_n`$ ($`nZZ`$) and $`a_{\pm n}`$ ($`n1`$) into the generating functions (”fields”)
$$\begin{array}{ccccc}\hfill e(z)& =& \underset{nZZ}{}e_nz^n,\psi _+(z)& =& k_\alpha ^1\mathrm{exp}\left((1)^{\theta (\alpha )}(qq^1)\underset{n=1}{\overset{\mathrm{}}{}}a_nz^n\right),\hfill \\ \hfill f(z)& =& \underset{nZZ}{}f_nz^n,\psi _{}(z)& =& k_\alpha \mathrm{exp}\left((1)^{\theta (\alpha )}(q^1q)\underset{n=1}{\overset{\mathrm{}}{}}a_nz^n\right),\hfill \end{array}$$
(7.1)
such that
$$\mathrm{deg}e(z)=\mathrm{deg}f(z)=\theta (\alpha ),\mathrm{deg}\psi _\pm (z)=0.$$
(7.2)
These fields satisfy the following conjugation conditions with respect to graded conjugation ””:
$$\begin{array}{ccccc}\hfill (e(z))^{}& =& f(z^1),(f(z))^{}& =& (1)^{\theta (\alpha )}e(z^1),\hfill \\ \hfill (\psi _+(z))^{}& =& \psi _{}(z^1),(\psi _{}(z))^{}& =& \psi _+(z^1),\hfill \end{array}$$
(7.3)
and have the following symmetry with respect to the translation operator $`T_\delta `$:
$$\begin{array}{ccccc}\hfill T_\delta (e(z))& =& (1)^{\theta (\alpha )}ze((1)^{\theta (\alpha )}z),T_\delta (f(z))& =& (1)^{\theta (\alpha )}z^1f((1)^{\theta (\alpha )}z),\hfill \\ \hfill T_\delta (\psi _+(z))& =& q^c\psi _+((1)^{\theta (\alpha )}z),T_\delta (\psi _{}(z))& =& q^c\psi _{}((1)^{\theta (\alpha )}z).\hfill \end{array}$$
(7.4)
###### Proposition 7.1
In terms of the fields (7.1) the relations of Section 4 can be rewritten in the following compact form
$$\begin{array}{ccc}\hfill [q^c,\mathrm{everything}]& =& 0,\hfill \\ \hfill u^d\phi (v)u^d& =& \phi (uv)\hfill \end{array}$$
(7.5)
where $`\phi (v)=e(v),f(v),\psi _\pm (v)`$, and also
$`\psi _\pm (u)\psi _\pm (v)`$ $`=`$ $`\psi _\pm (v)\psi _\pm (u),`$ (7.6)
$`(u\overline{q}_\alpha v)e(u)e(v)`$ $`=`$ $`(\overline{q}_\alpha uv)e(v)e(u),`$ (7.7)
$`(uq_\alpha v)f(u)f(v)`$ $`=`$ $`(q_\alpha uv)f(v)f(u),`$ (7.8)
$`\psi _\pm (u)e(v)\left(\psi _\pm (u)\right)^1`$ $`=`$ $`(1)^{\theta (\alpha )}\frac{\overline{q}_\alpha q^{\frac{c}{2}}uv}{q^{\frac{c}{2}}u\overline{q}_\alpha v}e(v),`$ (7.9)
$`\psi _\pm (u)f(v)\left(\psi _\pm (u)\right)^1`$ $`=`$ $`(1)^{\theta (\alpha )}\frac{q_\alpha q^{\pm \frac{c}{2}}uv}{q^{\pm \frac{c}{2}}uq_\alpha v}f(v),`$ (7.10)
$`\left(\psi _+(u)\right)^1\psi _{}(v)\psi _+(u)\left(\psi _{}(v)\right)^1`$ $`=`$ $`\frac{\left(q^cuq_\alpha v\right)\left(q^cu\overline{q}_\alpha v\right)}{\left(q^cv\overline{q}_\alpha u\right)\left(q^cvq_\alpha u\right)},`$ (7.11)
$$[e(u),f(v)]=\frac{1}{qq^1}\left(\delta (\frac{u}{v}q^c)\psi _{}(vq^{c/2})\delta (\frac{u}{v}q^c)\psi _+(uq^{c/2})\right).$$
(7.12)
Here in (7.11) $`\delta (z)=_{nZZ}z^n`$, and the brackets $`[,]`$ in the relation (7.12) mean the supercommutator:
$$[e(u),f(v)]=e(u)f(v)(1)^{\theta (\alpha )}f(v)e(u).$$
(7.13)
Given description is called the ’new realization’, or the current realization of the quantum affine superalgebras $`U_q(A_1^{(1)})`$ and $`U_q(C(2)^{(2)})`$. It should be noted that the relations (7.1) and (7.6)–(7.12) differ from the corresponding relations of Refs. , , by replacement $`q`$ to $`q^1`$. The current realization possesses its own graded comultiplication structure, different from (2.12):
$$\begin{array}{ccc}\hfill \mathrm{\Delta }_q^{(D)}(c)& =& c1+1c,\hfill \\ \hfill \mathrm{\Delta }_q^{(D)}(d)& =& d1+1d,\hfill \\ \hfill \mathrm{\Delta }_q^{(D)}(\psi _\pm (z))& =& \psi _\pm (zq^{\pm \frac{c_2}{2}})\psi _\pm (zq^{\frac{c_1}{2}}),\hfill \\ \hfill \mathrm{\Delta }_q^{(D)}(e(z))& =& e(z)1+\psi _{}(zq^{\frac{c_1}{2}})e(zq^{c_1}),\hfill \\ \hfill \mathrm{\Delta }_q^{(D)}(f(z))& =& f(zq^{c_2})\psi _+(zq^{\frac{c_2}{2}})+1f(z),\hfill \end{array}$$
(7.14)
$$\begin{array}{ccc}\hfill S_q^{(D)}(c)& =& c,S_q^{(D)}(d)=d,\hfill \\ \hfill S_q^{(D)}(\psi _\pm (z))& =& \left(\psi _\pm (z)\right)^1,\hfill \\ \hfill S_q^{(D)}(e(z))& =& \left(\psi _{}(zq^{\frac{c}{2}})\right)^1e(zq^c),\hfill \\ \hfill S_q^{(D)}(f(z)& =& f(zq^c)\left(\psi _+(zq^{\frac{c}{2}})\right)^1,\hfill \end{array}$$
(7.15)
$$\epsilon (c)=\epsilon (d)=\epsilon (e(z))=\epsilon (f(z))=0,\epsilon (\psi ^\pm (z))=1.$$
(7.16)
Here $`\mathrm{\Delta }_q^{(D)}`$, $`S_q^{(D)}`$, and $`\epsilon `$ are the comultiplication, antipode and counite correspondingly. The two comultiplications $`\mathrm{\Delta }_q`$ and $`\mathrm{\Delta }_q^{(D)}`$ are related by the twist :
$$\mathrm{\Delta }_q^{(D)}(x)=F^1\mathrm{\Delta }(x)F,$$
(7.17)
where $`F=R_+^{21}`$, with $`R_+`$ given by (6.7)–(6.11), such that the universal $`R`$-matrix for the comultiplication $`\mathrm{\Delta }_q^{(D)}`$ equals to
$$^{(D)}=R_0R_{}KR_+^{21}$$
(7.18)
with the factors from (6.5). In the generators $`e_n`$, $`f_n`$ and $`a_n`$ it can be rewritten as follows
$$^{(D)}=𝒦\overline{},$$
(7.19)
where
$`𝒦`$ $`=`$ $`q^{\frac{h_\alpha h_\alpha }{(\alpha ,\alpha )}}q^{\frac{1}{2}(cd+dc)}\mathrm{exp}\left((qq^1){\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\overline{d}(n)a_na_n\right)q^{\frac{1}{2}(cd+dc)},`$ (7.20)
$`\overline{}`$ $`=`$ $`{\displaystyle \underset{nZZ}{\overset{}{}}}\mathrm{exp}_{\overline{q}_\alpha }\left((q^1q)f_ne_n\right),`$ (7.21)
and
$$\overline{d}(n)=\frac{n\left(qq^1\right)}{q_\alpha ^nq_\alpha ^n}.$$
(7.22)
It is possible to give another presentation of the element $`\overline{}`$ in the completed algebras $`\overline{U}(g)`$, where $`g`$ is either $`A_1^{(1)}`$ or $`C(2)^{(2)}`$ , . The completion is done with respect to open neighborhoods of zero $`\overline{U}_r=_{s>r}U_s`$, where $`U_s`$ consists of all the elements from $`U(g)`$ of degree $`s`$. The completed algebra acts on (infinite-dimensional) representations of highest weight and admits the series over monomials $`x_{i_1}x_{i_2}\mathrm{}x_{i_n}`$, $`i_1i_2\mathrm{}i_n,`$ with $`x=e,f,a`$ and fixed $`i_k`$. The matrix coefficients of the products of the currents $`e(z_1)e(z_2)\mathrm{}e(z_n)`$ and $`f(z_1)f(z_2)\mathrm{}f(z_n)`$, defined originally as formal series, converge to meromorphic in $`𝐂\mathrm{l}^n`$ functions with the poles at $`z_i=0`$ and $`z_i=q_\alpha ^1z_j`$, $`ij`$.
Let $`t(z)=(qq^1)f(z)e(z)`$. As before, we understand the product $`t(z_1)\mathrm{}t(z_n)`$ as operator-valued meromorphic function in $`\left(𝐂\mathrm{l}^{}\right)^n`$ with simple poles at $`z_i=q_\alpha ^1z_j`$, $`ij`$. Define
$$\overline{}^{}=1+\underset{n>0}{}\frac{1}{n!\left(2\pi i\right)^n}\underset{D_n}{\mathrm{}}\frac{dz_1}{z_1}\mathrm{}\frac{dz_n}{z_n}t(z_1)\mathrm{}t(z_n),$$
(7.23)
and integration region $`D_n`$ is defined as $`D_n=\{|z_i|=1,i=1,\mathrm{},n\}`$ for $`|q|<1`$ and, more generally, by
$$D_n=\left\{\left|z_i\underset{\genfrac{}{}{0pt}{}{j=1,\mathrm{},n,}{ji}}{}(z_iq_\alpha z_j)\right|=1,i=1,\mathrm{},n\right\}$$
(7.24)
for any $`q`$, such that $`q_\alpha ^N1`$, $`NZZ\{0\}`$.
###### Proposition 7.2
The action of the tensor $`^{}=𝒦\overline{}^{}`$ in tensor product of highest weight modules is well defined and coincides with the action of the universal $`R`$-matrix (7.19)
The integrals in (7.20) can be computed explicitly. Let us put by induction
$$t^{(n)}(z)=\underset{z_1=z\overline{q}_\alpha ^{2n2}}{Res}t(z_1)t^{(n1)}(z)\frac{dz_1}{z_1},$$
(7.25)
where $`t^{(1)}(z)=t(z)`$. In the components the fields $`t^{(n)}(z)`$ look as follows:
$$t^{(n)}(z)=C_ne^{(n)}(z)f^{(n)}(z),$$
(7.26)
where
$$\begin{array}{ccc}\hfill C_n& =& (1)^{(n1)\theta (\alpha )}(qq^1)^n\stackrel{~}{q}_\alpha ^{\frac{n(n1)}{2}}(\stackrel{~}{q}_\alpha 1)^{n1}(n1)_{\stackrel{~}{q}_\alpha }!(n)_{\stackrel{~}{q}_\alpha }!,\hfill \\ \hfill e^{(n)}(z)& =& e(z)e(\overline{q}_\alpha z)e(\overline{q}_\alpha ^{\mathrm{\hspace{0.33em}2}}z)\mathrm{}e(\overline{q}_\alpha ^{n2}z)e(\overline{q}_\alpha ^{n1}z),\hfill \\ \hfill f^{(n)}(z)& =& f(\overline{q}_\alpha ^{n1}z)f(\overline{q}_\alpha ^{n2}z)\mathrm{}f(\overline{q}_\alpha ^{\mathrm{\hspace{0.33em}2}}z)f(\overline{q}_\alpha z)f(z),\hfill \end{array}$$
(7.27)
such that
$$\begin{array}{ccc}\hfill e^{(n)}(z)& =& \underset{mZZ}{}\left(z\overline{q}_\alpha ^n\right)^m\underset{\genfrac{}{}{0pt}{}{\lambda _1\mathrm{}\lambda _n,}{\lambda _1+\mathrm{}+\lambda _n=m}}{}\frac{q_\alpha ^{\lambda _1+2\lambda _2+\mathrm{}n\lambda _n}}{\underset{jZZ}{}\left(\lambda _j^{}\lambda _{j+1}^{}\right)_{\overline{q}_\alpha }!}e_{\lambda _n}e_{\lambda _{n1}}\mathrm{}e_{\lambda _1},\hfill \\ \hfill f^{(n)}(z)& =& \underset{mZZ}{}\left(zq_\alpha \right)^m\underset{\genfrac{}{}{0pt}{}{\lambda _1\mathrm{}\lambda _n,}{\lambda _1+\mathrm{}+\lambda _n=m}}{}\frac{q_\alpha ^{\lambda _1+2\lambda _2+\mathrm{}n\lambda _n}}{\underset{jZZ}{}\left(\lambda _j^{}\lambda _{j+1}^{}\right)_{q_\alpha }!}f_{\lambda _n}f_{\lambda _{n1}}\mathrm{}f_{\lambda _1}.\hfill \end{array}$$
(7.28)
Here $`\lambda _j^{}=\mathrm{\#}k`$, such that $`\lambda _kj`$, and $`\stackrel{~}{q}_\alpha :=q^{(\alpha ,\alpha )}`$. The product in denominator is finite, since there are only finitely many distinct $`\lambda _j^{}`$ for a given choice of $`\lambda _k`$. Then, repeating the calculations in , we get vertex type presentation of the element $`\overline{}^{}`$:
$$\overline{}^{}=\mathrm{exp}\left(\underset{n>0}{}\frac{1}{n}I_n\right),$$
(7.29)
where the sequence of operators
$$I_n=\frac{t^{\left(n\right)}\left(z\right)dz}{2\pi iz}$$
(7.30)
commute between themselves:
$$[I_n,I_m]=0,n,m>0.$$
The vertex operator presentation (7.29) is convenient for applications to integrable representations: it is expressed through integrals over the fields, which number is precisely $`k`$ for level $`k`$ integrable representations.
## 8 Final Remarks
The aim of this paper is to describe in unified way with detail the $`q`$-deformed untwisted affine algebra $`U_q(\widehat{sl}(2))=U_q(A_1^{(1)})`$ and twisted superalgebra $`U_q(osp(2|2))^{(2)}=U_q(C(2)^{(2)})`$. In order to describe the complete list of quantum affine (super)algebras of rank 2 one should consider some more three quantum affine (super)algebras: $`U_q(sl(1|3))^{(4)}=U_q(A(0,2))^{(4)}`$, $`U_q(sl(3))^{(2)}=U_q(A_2^{(2)})`$ and $`U_q(\widehat{osp}(1|2))=U_q(B(0,1)^{(1)}`$. The Dynkin diagram of the superalgebra $`A(0,2)^{(4)}`$ has as geometric structure as the (super)algebras $`A_1^{(1)}`$ and $`C(2)^{(2)}`$ but in this case the root $`\alpha `$ is even and $`\delta \alpha `$ is odd one, and the sector of imaginary roots has odd roots. Therefore in the case of the quantum superalgebra $`A(0,2)^{(4)}`$ the relations of the type (4.29)–(4.32) are more complicated and they demand special consideration. The second family of two quantum affine (super)algebras $`U_q(A_2^{(2)})`$ and $`U_q(B(0,1)^{(1)})`$ are described by the same Dynkin diagram with different colors of roots. Preliminary results in this direction are given in , where in particular the Cartan-Weyl basis of basic affine superalgebra $`U_q(\widehat{osp}(1|2))`$ is considered. The unified description of three mentioned above quantum affine (super)algebras, analogous to the one given in the present paper, is in preparation.
## Acknowledgments
This work was supported (S.M. Khoroshkin, V.N. Tolstoy) by the Russian Foundation for Fundamental Research, grant No.98-01-00303, by the program of French-Russian scientific cooperation (CNRS grant PICS-608 and grant RFBR-98-01-22033), as well as by KBN grant 2P03B13012 (J. Lukierski) and INTAS-99-1705 (S. Khoroshkin).
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# 1 Introduction
## 1 Introduction
In , a proposal was made to extend the wave-particle duality where a particle can be described by the Lagrangian $`_1`$ or terms of a Klein-Gordon field $`_2`$
$$_1=\sqrt{\left(\frac{X^\mu }{\tau }\right)^2},_2=\frac{1}{2}\left(\frac{\varphi }{x_\mu }\right)^2$$
(1)
to extended objects so that strings and branes with Lagrangian $`_3`$ could also be described by a theory with Lagrangian $`_4`$.
$$_3=\sqrt{\mathrm{d}et\left|\frac{X^\mu }{\sigma ^i}\frac{X^\mu }{\sigma ^j}\right|},_4=\mathrm{d}et\left|\frac{\varphi ^i}{x_\mu }\frac{\varphi ^j}{x_\mu }\right|$$
(2)
The number of fields $`\varphi ^i`$ is the same as the number of world-sheet co-ordinates $`\sigma ^i`$. The new Lagrangian will be referred to as the companion Lagrangian. Although $`_4`$ is the natural extension of the Klein-Gordon Lagrangian, it would be preferable to consider $`\sqrt{}_4`$ since this possesses general covariance. The following result allows us to link these two Lagrangians.
In it was noted that the equations of motion for the companion Lagrangian without a square root, when subjected to some constraints, reduce to the equations of motion for the companion Lagrangian with a square root in one dimension less. Computer calculations had verified this in several cases and an analytic proof was given for the case of two fields. Here we give an analytic proof of this observation for any number of fields $`n`$ in any number of dimensions $`d`$ where $`d>n`$. This result can most easily be checked for the case of one field.
## 2 Conventions and Notation
Partial derivatives are denoted by
$$\frac{\varphi ^i}{x_\mu }=\varphi _\mu ^i,\frac{^2\varphi ^i}{x_\mu x_\nu }=\varphi _{\mu \nu }^i$$
(3)
Totally antisymmetric tensors $`ϵ_{\nu _1\nu _2\mathrm{}\nu _d}`$ are used throughout the proof with the convention $`ϵ_{12\mathrm{}d}=+1`$. When indices have an arrow above them then they represent several indices. They can be thought of as vectors with several components.
$`\stackrel{}{\mu },\stackrel{}{\nu },\stackrel{}{\rho },\stackrel{}{\sigma }`$ each have $`(n1)`$ components. For example $`\stackrel{}{\mu }`$ denotes $`\{\mu _2,\mu _3,\mathrm{},\mu _n\}`$.
$`\stackrel{}{\tau },\stackrel{}{\kappa }`$ each have $`(dn)`$ components. For example $`\stackrel{}{\kappa }`$ denotes $`\{\kappa _1,\kappa _2,\mathrm{},\kappa _{dn}\}`$.
$`\stackrel{}{\kappa }^{}`$ denotes $`\{\kappa _2,\kappa _3,\mathrm{},\kappa _{dn}\}`$ and $`\stackrel{}{\kappa }^{\prime \prime }`$ denotes $`\{\kappa _3,\mathrm{},\kappa _{dn}\}`$.
For the product of $`(n1)`$ fields we use the notation $`\mathrm{\Phi }_\stackrel{}{\nu }=\varphi _{\nu _2}^2\varphi _{\nu _3}^3\mathrm{}\varphi _{\nu _n}^n`$.
A useful identity which will be used later on is
$`ϵ_{\mu \nu _2\nu _3\mathrm{}\nu _d}ϵ_{\rho _1\rho _2\mathrm{}\rho _d}=ϵ_{\rho _1\nu _2\nu _3\mathrm{}\nu _d}ϵ_{\mu \rho _2\mathrm{}\rho _d}+ϵ_{\rho _2\nu _2\nu _3\mathrm{}\nu _d}ϵ_{\rho _1\mu \rho _3\mathrm{}\rho _d}+\mathrm{}`$
$`\mathrm{}+ϵ_{\rho _d\nu _2\nu _3\mathrm{}\nu _d}ϵ_{\rho _1\rho _2\mathrm{}\rho _{d1}\mu }`$ (4)
It amounts to swapping the index $`\mu `$ from the first epsilon with each index from the second epsilon.
## 3 Lagrangians and Equations of Motion
Consider the Lagrangian for $`n`$ fields $`\varphi ^i`$ in $`d`$ space-time dimensions $`x^\mu `$ which does not involve a square root.
$$=\mathrm{d}et\left|\frac{\varphi ^i}{x_\mu }\frac{\varphi ^j}{x_\mu }\right|$$
(5)
The equations of motion for this Lagrangian are
$$\frac{^2}{\varphi _\mu ^i\varphi _\nu ^j}\varphi _{\mu \nu }^j=0$$
(6)
These determinantal Lagrangians can be written as the sum of squares of Jacobians. The Jacobians will be denoted as
$$J_\stackrel{}{\kappa }=J_{\kappa _1\kappa _2\mathrm{}\kappa _{dn}}=ϵ_{\kappa _1\kappa _2\mathrm{}\kappa _{dn}\nu _1\nu _2\mathrm{}\nu _n}\varphi _{\nu _1}^1\varphi _{\nu _2}^2\mathrm{}\varphi _{\nu _n}^n$$
(7)
For the square root case the Lagrangian is
$$=\sqrt{\mathrm{d}et\left|\frac{\varphi ^i}{x_\mu }\frac{\varphi ^j}{x_\mu }\right|}=\sqrt{\frac{1}{(dn)!}J_\stackrel{}{\kappa }J_\stackrel{}{\kappa }}$$
(8)
The equations of motion for this can be written as
$$J_{\mu \stackrel{}{\kappa }^{}}J_{\nu \stackrel{}{\kappa }^{}}\varphi _{\mu \nu }^i=0$$
(9)
## 4 The Constraints
The equations of motion for the non-square root case will be subject to the following constraints.
$$\frac{}{\varphi _\mu ^i}\varphi _{\mu \nu }^i=0.$$
(10)
There is no summation over the index $`i`$. The Lagrangian must also vanish.
The idea is to reduce the number of dimensions from $`d`$ to $`d1`$. The constraints (10) can be used to eliminate all second derivatives of the fields which involve a partial derivative with respect to $`x_d`$, the $`d`$th dimension. i.e From the constraints
$`\varphi _{d\beta }^i={\displaystyle \frac{\frac{}{\varphi _\alpha ^i}}{\frac{}{\varphi _d^i}}}\varphi _{\alpha \beta }^i,\varphi _{dd}^i={\displaystyle \frac{\frac{}{\varphi _\alpha ^i}\frac{}{\varphi _\beta ^i}}{\left(\frac{}{\varphi _d^i}\right)^2}}\varphi _{\alpha \beta }^i`$ (11)
Putting these into the equations of motion (6) we have:
$`{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \frac{1}{\left(\frac{}{\varphi _d^j}\right)^2}}[\left({\displaystyle \frac{}{\varphi _d^j}}\right)^2{\displaystyle \frac{^2}{\varphi _\alpha ^i\varphi _\beta ^j}}{\displaystyle \frac{}{\varphi _d^j}}{\displaystyle \frac{}{\varphi _\beta ^j}}{\displaystyle \frac{^2}{\varphi _\alpha ^i\varphi _d^j}}{\displaystyle \frac{}{\varphi _d^j}}{\displaystyle \frac{}{\varphi _\alpha ^j}}{\displaystyle \frac{^2}{\varphi _\beta ^i\varphi _d^j}}`$
$`+{\displaystyle \frac{}{\varphi _\alpha ^j}}{\displaystyle \frac{}{\varphi _\beta ^j}}{\displaystyle \frac{^2}{\varphi _d^i\varphi _d^j}}]\varphi ^j_{\alpha \beta }=0`$ (12)
The indices $`\alpha ,\beta =1,2,\mathrm{},(d1)`$ throughout the paper.
## 5 The Proof
For the moment we shall consider the equation of motion with respect to field $`\varphi =\varphi ^1`$ and are only looking at the component which involves the terms $`\varphi _{\alpha \beta }`$. The other components will work in the same way. Using the definition of the Lagrangian which involves the Jacobians then we can write
$$=\frac{1}{(dn)!}\varphi _\nu \varphi _\rho B_{\nu \rho }\mathrm{w}hereB_{\nu \rho }=ϵ_{\nu \stackrel{}{\kappa }\stackrel{}{\nu }}ϵ_{\rho \stackrel{}{\kappa }\stackrel{}{\rho }}\mathrm{\Phi }_\stackrel{}{\nu }\mathrm{\Phi }_\stackrel{}{\rho }$$
(13)
so the numerator of the coefficient of $`\varphi _{\alpha \beta }`$ in (12) becomes
$$[B_{\mu d}(B_{\nu d}B_{\alpha \beta }B_{\nu \beta }B_{\alpha d})+B_{\nu \alpha }(B_{\mu \beta }B_{dd}B_{\mu d}B_{\beta d})]\varphi _\mu \varphi _\nu $$
(14)
Now,
$`B_{\nu d}B_{\alpha \beta }B_{\nu \beta }B_{\alpha d}`$ $`=`$ $`[ϵ_{\nu \stackrel{}{\kappa }\stackrel{}{\mu }}ϵ_{d\stackrel{}{\kappa }\stackrel{}{\nu }}ϵ_{\alpha \stackrel{}{\tau }\stackrel{}{\rho }}ϵ_{\beta \stackrel{}{\tau }\stackrel{}{\sigma }}ϵ_{\nu \stackrel{}{\kappa }\stackrel{}{\mu }}ϵ_{\beta \stackrel{}{\kappa }\stackrel{}{\nu }}ϵ_{\alpha \stackrel{}{\tau }\stackrel{}{\rho }}ϵ_{d\stackrel{}{\tau }\stackrel{}{\sigma }}]\mathrm{\Phi }_\stackrel{}{\mu }\mathrm{\Phi }_\stackrel{}{\nu }\mathrm{\Phi }_\stackrel{}{\rho }\mathrm{\Phi }_\stackrel{}{\sigma }`$ (15)
$`=`$ $`ϵ_{\nu \stackrel{}{\kappa }\stackrel{}{\mu }}ϵ_{\alpha \stackrel{}{\tau }\stackrel{}{\rho }}[ϵ_{d\stackrel{}{\kappa }\stackrel{}{\nu }}ϵ_{\beta \stackrel{}{\tau }\stackrel{}{\sigma }}ϵ_{\beta \stackrel{}{\kappa }\stackrel{}{\nu }}ϵ_{d\stackrel{}{\tau }\stackrel{}{\sigma }}]\mathrm{\Phi }_\stackrel{}{\mu }\mathrm{\Phi }_\stackrel{}{\nu }\mathrm{\Phi }_\stackrel{}{\rho }\mathrm{\Phi }_\stackrel{}{\sigma }`$
Using the epsilon identity (4) to move the index $`\beta `$ around
$`ϵ_{d\stackrel{}{\kappa }\stackrel{}{\nu }}ϵ_{\beta \stackrel{}{\tau }\stackrel{}{\sigma }}=ϵ_{\beta \stackrel{}{\kappa }\stackrel{}{\nu }}ϵ_{d\stackrel{}{\tau }\stackrel{}{\sigma }}+ϵ_{d\beta \kappa _2\mathrm{}\kappa _r\stackrel{}{\nu }}ϵ_{\kappa _1\stackrel{}{\tau }\stackrel{}{\sigma }}+\mathrm{}+ϵ_{d\kappa _1\mathrm{}\kappa _{r1}\beta \stackrel{}{\nu }}ϵ_{\kappa _r\stackrel{}{\tau }\stackrel{}{\sigma }}+`$
$`ϵ_{d\stackrel{}{\kappa }\beta \nu _3\mathrm{}\nu _n}ϵ_{\nu _2\stackrel{}{\tau }\stackrel{}{\sigma }}+\mathrm{}+ϵ_{d\stackrel{}{\kappa }\nu _2\mathrm{}\nu _{n1}\beta }ϵ_{\nu _n\stackrel{}{\tau }\stackrel{}{\sigma }}`$ (16)
The first term on the right hand side is just the other term in expression (15). The last $`(n1)`$ terms will all vanish due to symmetry conditions. This only leaves the middle terms. But by relabelling
$`ϵ_{\nu \kappa _1\mathrm{}\kappa _r\stackrel{}{\mu }}ϵ_{d\kappa _1\mathrm{}\beta \mathrm{}\kappa _r\stackrel{}{\nu }}ϵ_{\kappa _i\stackrel{}{\tau }\stackrel{}{\sigma }}`$ $`=`$ $`ϵ_{\nu \kappa _i\mathrm{}\kappa _1\mathrm{}\kappa _r\stackrel{}{\mu }}ϵ_{d\kappa _i\mathrm{}\beta \mathrm{}\kappa _r\stackrel{}{\nu }}ϵ_{\kappa _1\stackrel{}{\tau }\stackrel{}{\sigma }}`$ (17)
$`=`$ $`ϵ_{\nu \stackrel{}{\kappa }\stackrel{}{\mu }}ϵ_{d\beta \kappa _2\mathrm{}\kappa _r\stackrel{}{\nu }}ϵ_{\kappa _1\stackrel{}{\tau }\stackrel{}{\sigma }}`$
There are $`r=dn`$ of these terms. Therefore,
$$B_{\nu d}B_{\alpha \beta }B_{\nu \beta }B_{\alpha d}=rϵ_{\nu \stackrel{}{\kappa }\stackrel{}{\mu }}ϵ_{\alpha \stackrel{}{\tau }\stackrel{}{\rho }}ϵ_{d\beta \stackrel{}{\kappa }^{}\stackrel{}{\nu }}ϵ_{\kappa _1\stackrel{}{\tau }\stackrel{}{\sigma }}\mathrm{\Phi }_\stackrel{}{\mu }\mathrm{\Phi }_\stackrel{}{\nu }\mathrm{\Phi }_\stackrel{}{\rho }\mathrm{\Phi }_\stackrel{}{\sigma }$$
(18)
Now as in (16), using identity (4) to swap subscript $`\kappa _1`$
$`ϵ_{\nu \stackrel{}{\kappa }\stackrel{}{\mu }}ϵ_{\alpha \stackrel{}{\tau }\stackrel{}{\rho }}=ϵ_{\nu \alpha \kappa _2\mathrm{}\kappa _r\stackrel{}{\mu }}ϵ_{\kappa _1\stackrel{}{\tau }\stackrel{}{\rho }}+ϵ_{\nu \tau _1\kappa _2\mathrm{}\kappa _r\stackrel{}{\mu }}ϵ_{\alpha \kappa _1\tau _2\mathrm{}\tau _r\stackrel{}{\rho }}+\mathrm{}+ϵ_{\nu \tau _r\kappa _2\mathrm{}\kappa _r\stackrel{}{\mu }}ϵ_{\alpha \tau _1\mathrm{}\tau _{r1}\kappa _1\stackrel{}{\rho }}`$
$`+ϵ_{\nu \rho _2\kappa _2\mathrm{}\kappa _r\stackrel{}{\mu }}ϵ_{\alpha \stackrel{}{\tau }\kappa _1\rho _3\mathrm{}\rho _n}+\mathrm{}+ϵ_{\nu \rho _n\kappa _2\mathrm{}\kappa _r\stackrel{}{\mu }}ϵ_{\alpha \stackrel{}{\tau }\rho _2\mathrm{}\rho _{n1}\kappa _1}`$ (19)
And by relabelling indices and using the antisymmetric property of the epsilons
$$ϵ_{\nu \tau _i\kappa _2\mathrm{}\kappa _n\stackrel{}{\mu }}ϵ_{\alpha \tau _1\mathrm{}\kappa _1\mathrm{}\tau _r\stackrel{}{\rho }}ϵ_{\kappa _1\stackrel{}{\tau }\stackrel{}{\sigma }}=ϵ_{\nu \stackrel{}{\kappa }\stackrel{}{\mu }}ϵ_{\alpha \stackrel{}{\tau }\stackrel{}{\rho }}ϵ_{\kappa _1\stackrel{}{\tau }\stackrel{}{\sigma }}$$
(20)
so,
$$(1+r)ϵ_{\nu \stackrel{}{\kappa }\stackrel{}{\mu }}ϵ_{\alpha \stackrel{}{\tau }\stackrel{}{\rho }}ϵ_{\kappa _1\stackrel{}{\tau }\stackrel{}{\sigma }}=ϵ_{\nu \alpha \stackrel{}{\kappa }^{}\stackrel{}{\mu }}ϵ_{\kappa _1\stackrel{}{\tau }\stackrel{}{\rho }}ϵ_{\kappa _1\stackrel{}{\tau }\stackrel{}{\sigma }}$$
(21)
which gives
$$B_{\nu d}B_{\alpha \beta }B_{\nu \beta }B_{\alpha d}=\frac{r}{r+1}B_{\kappa \kappa }[ϵ_{\nu \alpha \stackrel{}{\kappa }^{}\stackrel{}{\mu }}ϵ_{d\beta \stackrel{}{\kappa }^{}\stackrel{}{\nu }}\mathrm{\Phi }_\stackrel{}{\mu }\mathrm{\Phi }_\stackrel{}{\nu }]$$
(22)
Substituting this into the expression (14) we find
$`B_{\tau \tau }[B_{\mu d}ϵ_{\nu \alpha \stackrel{}{\kappa }^{}\stackrel{}{\mu }}ϵ_{d\beta \stackrel{}{\kappa }^{}\stackrel{}{\nu }}+B_{\nu \alpha }ϵ_{\mu d\stackrel{}{\kappa }^{}\stackrel{}{\mu }}ϵ_{\beta d\stackrel{}{\kappa }^{}\stackrel{}{\nu }}]\mathrm{\Phi }_\stackrel{}{\mu }\mathrm{\Phi }_\stackrel{}{\nu }\varphi _\mu \varphi _\nu `$ (23)
$`=`$ $`B_{\tau \tau }[ϵ_{\mu \stackrel{}{\tau }\stackrel{}{\rho }}ϵ_{d\stackrel{}{\tau }\stackrel{}{\sigma }}ϵ_{\nu \alpha \stackrel{}{\kappa }^{}\stackrel{}{\mu }}ϵ_{d\beta \stackrel{}{\kappa }^{}\stackrel{}{\nu }}ϵ_{\mu \stackrel{}{\tau }\stackrel{}{\rho }}ϵ_{\alpha \stackrel{}{\tau }\stackrel{}{\sigma }}ϵ_{\nu d\stackrel{}{\kappa }^{}\stackrel{}{\mu }}ϵ_{d\beta \stackrel{}{\kappa }^{}\stackrel{}{\nu }}]\mathrm{\Phi }_\stackrel{}{\mu }\mathrm{\Phi }_\stackrel{}{\nu }\mathrm{\Phi }_\stackrel{}{\rho }\mathrm{\Phi }_\stackrel{}{\sigma }\varphi _\mu \varphi _\nu `$
Now, using (4) to move subscript $`d`$,
$`ϵ_{d\stackrel{}{\tau }\stackrel{}{\sigma }}ϵ_{\nu \alpha \stackrel{}{\kappa }^{}\stackrel{}{\mu }}ϵ_{\alpha \stackrel{}{\tau }\stackrel{}{\sigma }}ϵ_{\nu d\stackrel{}{\kappa }^{}\stackrel{}{\mu }}=ϵ_{\nu \stackrel{}{\tau }\stackrel{}{\sigma }}ϵ_{d\alpha \stackrel{}{\kappa }^{}\stackrel{}{\mu }}+ϵ_{\kappa _2\stackrel{}{\tau }\stackrel{}{\sigma }}ϵ_{\nu \alpha d\kappa _3\mathrm{}\kappa _n\stackrel{}{\mu }}+\mathrm{}+ϵ_{\kappa _r\stackrel{}{\tau }\stackrel{}{\sigma }}ϵ_{\nu \alpha \kappa _2\mathrm{}\kappa _{r1}d\stackrel{}{\mu }}`$
$`+ϵ_{\mu _2\stackrel{}{\tau }\stackrel{}{\sigma }}ϵ_{\nu \alpha \stackrel{}{\kappa }^{\prime \prime }d\mu _3\mathrm{}\mu _n}+\mathrm{}+ϵ_{\mu _n\stackrel{}{\tau }\stackrel{}{\sigma }}ϵ_{\nu \alpha \stackrel{}{\kappa }^{\prime \prime }\mu _2\mathrm{}\mu _{n1}d}`$ (24)
As before the last terms will vanish due to symmetry considerations. For the middle terms, by relabelling and using antisymmetry
$$ϵ_{\kappa _i\stackrel{}{\tau }\stackrel{}{\sigma }}ϵ_{\nu \alpha \kappa _2\mathrm{}d\mathrm{}\kappa _r\stackrel{}{\mu }}ϵ_{d\beta \kappa _2\mathrm{}\kappa _i\mathrm{}\kappa _r\stackrel{}{\nu }}=ϵ_{\kappa _2\stackrel{}{\tau }\stackrel{}{\sigma }}ϵ_{\nu \alpha d\stackrel{}{\kappa }^{\prime \prime }\stackrel{}{\mu }}ϵ_{d\beta \stackrel{}{\kappa }^{}\stackrel{}{\nu }}$$
(25)
There are $`(r1)=(dn1)`$ of these terms. We now have
$$\frac{r}{r+1}B_{\tau \tau }[ϵ_{\mu \stackrel{}{\tau }\stackrel{}{\rho }}ϵ_{\nu \stackrel{}{\tau }\stackrel{}{\sigma }}ϵ_{d\alpha \stackrel{}{\kappa }^{}\stackrel{}{\mu }}ϵ_{d\beta \stackrel{}{\kappa }^{}\stackrel{}{\nu }}+(r1)ϵ_{\mu \stackrel{}{\tau }\stackrel{}{\rho }}ϵ_{\kappa _2\stackrel{}{\tau }\stackrel{}{\sigma }}ϵ_{\nu \alpha d\stackrel{}{\kappa }^{\prime \prime }\stackrel{}{\mu }}ϵ_{d\beta \stackrel{}{\kappa }^{}\stackrel{}{\nu }}]\mathrm{\Phi }_\stackrel{}{\mu }\mathrm{\Phi }_\stackrel{}{\nu }\mathrm{\Phi }_\stackrel{}{\rho }\mathrm{\Phi }_\stackrel{}{\sigma }\varphi _\mu \varphi _\nu $$
(26)
Again, rewriting the epsilons, this time moving subscript $`\kappa _2`$,
$`ϵ_{\mu \stackrel{}{\tau }\stackrel{}{\rho }}ϵ_{d\beta \stackrel{}{\kappa }^{}\stackrel{}{\nu }}=ϵ_{\kappa _2\stackrel{}{\tau }\stackrel{}{\rho }}ϵ_{d\beta \mu \stackrel{}{\kappa }^{\prime \prime }\stackrel{}{\nu }}+ϵ_{\mu \kappa _2\tau _2\mathrm{}\tau _r\stackrel{}{\rho }}ϵ_{d\beta \tau _1\stackrel{}{\kappa }^{\prime \prime }\stackrel{}{\nu }}+\mathrm{}+ϵ_{\mu \tau _1\mathrm{}\tau _{r1}\kappa _2\stackrel{}{\rho }}ϵ_{d\beta \tau _r\stackrel{}{\kappa }^{\prime \prime }\stackrel{}{\nu }}+`$
$`\mathrm{}+ϵ_{\mu \stackrel{}{\tau }\kappa _2\rho _3\mathrm{}\rho _n}ϵ_{d\beta \rho _2\stackrel{}{\kappa }^{\prime \prime }\stackrel{}{\nu }}+\mathrm{}+ϵ_{\mu \stackrel{}{\tau }\rho _2\mathrm{}\rho _{n1}\kappa _2}ϵ_{d\beta \rho _n\stackrel{}{\kappa }^{\prime \prime }\stackrel{}{\nu }}`$ (27)
And again, by relabelling
$$ϵ_{\mu \tau _1\mathrm{}\kappa _2\mathrm{}\tau _r\stackrel{}{\rho }}ϵ_{d\beta \tau _i\stackrel{}{\kappa }^{\prime \prime }\stackrel{}{\nu }}ϵ_{\kappa _2\stackrel{}{\tau }\stackrel{}{\sigma }}=ϵ_{\mu \stackrel{}{\tau }\stackrel{}{\rho }}ϵ_{d\beta \stackrel{}{\kappa }^{}\stackrel{}{\nu }}ϵ_{\kappa _2\stackrel{}{\tau }\stackrel{}{\sigma }}$$
(28)
There are $`r=dn`$ of these terms. So our expression is now
$`{\displaystyle \frac{r}{r+1}}B_{\tau \tau }[ϵ_{\mu \stackrel{}{\tau }\stackrel{}{\rho }}ϵ_{\nu \stackrel{}{\tau }\stackrel{}{\sigma }}ϵ_{d\alpha \stackrel{}{\kappa }^{}\stackrel{}{\mu }}ϵ_{d\beta \stackrel{}{\kappa }^{}\stackrel{}{\nu }}+{\displaystyle \frac{r1}{r+1}}ϵ_{\kappa _2\stackrel{}{\tau }\stackrel{}{\rho }}ϵ_{\kappa _2\stackrel{}{\tau }\stackrel{}{\sigma }}ϵ_{\nu \alpha d\stackrel{}{\kappa }^{\prime \prime }\stackrel{}{\mu }}ϵ_{d\beta \stackrel{}{\kappa }^{\prime \prime }\stackrel{}{\nu }}]\mathrm{\Phi }_\stackrel{}{\mu }\mathrm{\Phi }_\stackrel{}{\nu }\mathrm{\Phi }_\stackrel{}{\rho }\mathrm{\Phi }_\stackrel{}{\sigma }\varphi _\mu \varphi _\nu \varphi _{\alpha \beta }`$
$`={\displaystyle \frac{r}{r+1}}B_{\tau \tau }[ϵ_{d\alpha \stackrel{}{\kappa }^{}\stackrel{}{\mu }}ϵ_{d\beta \stackrel{}{\kappa }^{}\stackrel{}{\nu }}\mathrm{\Phi }_\stackrel{}{\mu }\mathrm{\Phi }_\stackrel{}{\nu }r!{\displaystyle \frac{r1}{r+1}}B_{\kappa \kappa }J_{d\alpha \stackrel{}{\kappa }^{\prime \prime }}J_{d\beta \stackrel{}{\kappa }^{\prime \prime }}]\varphi _{\alpha \beta }`$
$`={\displaystyle \frac{rr!}{r+1}}\left({\displaystyle \frac{J_\stackrel{}{\mu }}{\varphi _\tau }}{\displaystyle \frac{J_\stackrel{}{\mu }}{\varphi _\tau }}\right)\left[\left({\displaystyle \frac{J_{d\stackrel{}{\kappa }^{}}}{\varphi _\alpha }}{\displaystyle \frac{J_{d\stackrel{}{\kappa }^{}}}{\varphi _\beta }}\right){\displaystyle \frac{r1}{(r+1)!}}\left({\displaystyle \frac{J_\stackrel{}{\nu }}{\varphi _\kappa }}{\displaystyle \frac{J_\stackrel{}{\nu }}{\varphi _\kappa }}\right)J_{d\alpha \stackrel{}{\kappa }^{\prime \prime }}J_{d\beta \stackrel{}{\kappa }^{\prime \prime }}\right]\varphi _{\alpha \beta }`$
A very similar calculation can be carried out to rewrite the coefficients of $`\varphi _{\alpha \beta }^j`$ $`(j1)`$ from (12). These become, for $`\varphi ^j=\psi `$, say.
$$\frac{rr!}{r+1}\left(\frac{J_\stackrel{}{\mu }}{\varphi _\tau }\frac{J_\stackrel{}{\mu }}{\psi _\tau }\right)\left[\left(\frac{J_{d\stackrel{}{\kappa }^{}}}{\psi _\alpha }\frac{J_{d\stackrel{}{\kappa }^{}}}{\psi _\beta }\right)\frac{r1}{(r+1)!}\left(\frac{J_\stackrel{}{\nu }}{\psi _\kappa }\frac{J_\stackrel{}{\nu }}{\psi _\kappa }\right)J_{d\alpha \stackrel{}{\kappa }^{\prime \prime }}J_{d\beta \stackrel{}{\kappa }^{\prime \prime }}\right]\psi _{\alpha \beta }$$
(30)
When the condition that the Lagrangian vanishes is put into the equations of motion, they can just be written as
$$J_{d\alpha \stackrel{}{\kappa }^{\prime \prime }}J_{d\beta \stackrel{}{\kappa }^{\prime \prime }}\varphi _{\alpha \beta }^i=0$$
(31)
as required. Comparing (31) with (9), these are the equations of motion for the Lagrangian involving a square root (8) in $`(d1)`$ dimensions.
## 6 Conclusion
It has been proved that the equations of motion for the companion Lagrangian without a square root when subject to some constraints are equivalent to the equations of motion for the companion Lagrangian with a square root in one less dimension.
## Acknowledgements
I am grateful to David Fairlie for useful discussions and to EPSRC for a postgraduate research award.
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# Dipolar Interactions in Superconductor-Ferromagnet Heterostructures
## I Introduction
The interplay between superconductivity and ferromagnetism in bulk materials has been the subject of active research since 1957 when Ginzburg published a paper in which he considered the effect of the field created by a bulk distribution of magnetisation on a superconductor, which was described by the London equations . He has concluded that for a ferromagnetic induction field of the sample larger than its superconducting critical field, this field would destroy superconductivity, but he also pointed out that in thin films or wires where the induction field is much smaller (due to demagnetisation effects) and the critical field higher (due to the small diamagnetic energy) than in bulk superconductors, it should be possible to observe the coexistence of the two phenomena. Experiments carried out by Mathias et al. on Lanthanum with several rare-earth paramagnetic impurities dissolved at low concentrations suggested that the interaction responsible for the depletion of the superconducting critical temperature of Lanthanum is the exchange interaction between the paramagnetic impurity spins and the superconducting electrons. This interaction induces an effective ferromagnetic interaction between the (antiparallel) spins in the Cooper pair, which tends to destroy it and hence destroy superconductivity. Anderson and Suhl have shown that the RKKY interaction between the ferromagnetic spins due to the conduction electrons is significantly reduced in the superconducting state, but pointed out that ferromagnetism could coexist with superconductivity if the ferromagnetic atoms formed small domains. The dependence of the superconducting critical temperature on the concentration of magnetic impurities due to exchange scattering of electrons from these impurities was addressed with the microscopic theory of superconductivity by Abrikosov and Gor’kov . de Gennes and Sarma have estimated that tipically, the exchange interaction between localised moments and superconducting electrons would be $`10^210^3`$ larger than the dipolar interaction considered by Ginzburg. The detailed form of the Landau-Ginzburg theory of ferromagnetic superconductors was worked out by Suhl . Despite their conflicting character, superconductivity and ferromagnetism are seen to coexist in bulk systems, e.g. in HoMo<sub>6</sub>S<sub>8</sub> and ErRh<sub>4</sub>B<sub>4</sub> .
Another system in which the coexistence of superconductivity and ferromagnetism has been observed in the bulk is the nuclear magnet AuIn<sub>2</sub> . This compound shows a superconducting phase transition at $`T_c=207mK`$ and an ordering transition to a ferromagnetic state at an even lower temperature $`T_M=35\mu K`$. This particularly low temperature can be explained by the weakness of the interaction between the nuclear spins (which is primarily due to an indirect exchange via the conduction electrons).
On a different perspective, the development of epitaxial growth of crystals has permitted the creation of artificial superlattices composed of superconducting and ferromagnetic materials, e.g. Fe/V, Ni/V, Ni/Mo, EuS/Pb, EuO/Al and Nb/Gd . In these superlattices, one can experimentally study the interaction between superconductivity and ferromagnetism when these two effects occur in neighbouring spatial regions and also study the supression of superconductivity as a function of the relative proportion (i.e. layer tickness) of the two materials. More recently, the cuprates RuSr<sub>2</sub>GdCu<sub>2</sub>O<sub>8-δ</sub> and RuSr<sub>2</sub>Gd<sub>1+x</sub>Ce<sub>1-x</sub>Cu<sub>2</sub>O<sub>10</sub> have been found to show superconductivity and a ferromagnetism below their critical temperatures, $`T_c=15`$-$`40K`$ for RuSr<sub>2</sub>GdCu<sub>2</sub>O<sub>8-δ</sub> and $`T_c=37K`$ (for an optimal $`x=0.2`$) for RuSr<sub>2</sub>Gd<sub>1+x</sub>Ce<sub>1-x</sub>Cu<sub>2</sub>O<sub>10</sub>, the Curie temperatures for magnetic ordering being $`T_M133K`$ for the first compound and $`T_M100K`$ for the second. The experimental analysis shows that these materials, like all cuprates, have a layered structure and that superconductivity and ferromagnetism seem to occur in different layers. However, a detailed analysis has been hindered by difficulties with the growth of single crystals .
Motivated by such experiments, in which magnetism and superconductivity are seen to occur in different spatial regions of the studied materials, we wish to address the problem of a thin ferromagnetic layer, placed between two bulk superconducting layers (see fig. 1), in which the thickness of the superconducting layers is much larger than the London penetration depth of the superconducting material and the thickness of the ferromagnetic layer is very small compared to this quantity, which is a condition that can be easily obtained with the modern techniques of epitaxial growth . In this limit, the results obtained can also be applied to superlattices of the two materials, given that the ferromagnetic layers are decoupled from one another.
We consider in this paper a simple model system composed of a very thin ferromagnetic film, with an arbitrary distribution of magnetisation in the plane of the film, which is placed in a spatial gap of size $`2d`$ between two semi-infinite superconductors described by the London equations. The film is coupled to the superconductors by the electromagnetic interaction, i.e. we neglect the proximity effect and we consider the Josephson current flowing between the two superconductors to be zero (the limitations of these approximations will be discussed in section IV). Having made the approximations indicated above, we are able to solve the problem exactly, by firstly considering the simpler problem of a single dipole in the spatial gap and then superimposing the different solutions, due to the linearity of the London equations. One can then compute the dipolar energy of the distribution of magnetisation. It turns out that for wave-vectors much larger than the inverse London penetration depth the form of the dipolar energy in momentum space is unchanged by the presence of the superconductors. On the other hand, for wave vectors much smaller than the inverse London penetration depth, the dipolar energy in momentum space resembles the energy of a distribution of dipoles in a two dimensional space. This behaviour can be traced to the Meissner effect which confines the magnetic flux lines within the spatial gap.
One can think of several possible ways to detect the effects of this change of behaviour of the dipolar interaction at low wave vectors. If one were able to choose the materials composing a layered geometry of superconductor/ferromagnetic film/superconductor in such a way that the Curie transition temperature of the magnetic film to the ferromagnetic state is lower than the critical temperature $`T_c`$ of the superconductor, one should be able to measure the critical properties of the system at the paramagnetic-ferromagnetic transition, in particular such quantities as the specific heat and the magnetic susceptibility, with the superconductors already displaying the Meissner effect and therefore with a modified form of the dipolar energy.
It was shown by Aharony and Fisher that in a $`d`$-dimensional system, a $`d`$-dimensional dipolar interaction (such as the one occuring in a bulk tridimensional ferromagnet or in the layered geometry superconductor/ferromagnet/superconductor) is a relevant interaction (in the sense of RG) near a paramagnetic-ferromagnetic transition, leading to a crossover between the critical exponents of the short-range ferromagnet and the critical exponents of the dipolar system when one approaches the critical temperature . Pelcovits and Halperin have also shown that in the case of a $`d`$-dimensional systems with a $`d+1`$-dimensional dipolar interaction (such as the one ocurring in a thin magnetic film film in free space) the universality class of the dipolar system is the same as above. This is due to the fact that, at the fixed point, the ‘effective’ (renormalised) dipolar coupling constant is infinite, making the susceptibility of the system independent of the longitudinal degree of freedom of the magnetisation, which is the one sensitive to the nature of the dipolar interaction. However, in real systems, measurements are not taken exactly at the critical point and one always probes the crossover region. In this region, the dipolar coupling constant is finite and one should be able to detect the distinct character of the transition if the ferromagnetic film is included in a layered geometry with superconductors or if the film is grown in a non-superconducting substrate, due to the different character of the dipolar interaction at small wavevectors. The ideal experiment to detect such a distinction would presumably be a measurement of the longitudinal susceptibility using polarised neutrons . Experiments done with films of EuS/SrS grown on a Si substrate have shown that the low Curie temperature of EuS (16.5 K) is further reduced in these geometries. The authors of Ref. have also performed experiments with films of EuS/Pb, probing the transition between the superconducting state and the normal state in the Pb layer as a function of the applied magnetic field. Therefore, EuS stands as a good candidate for a material to be used in the ferromagnetic layer. It has the further advantage of being an insulator (see below).
Another possibility would be the study of spin-spin correlation functions in a magnetic film in the ordered phase and outside the critical region. Indeed, Kashuba has shown that the static correlation functions of an XY model with 2d dipolar interactions would display a behaviour analogous to that of the dynamic correlations functions of the stochastic process described the KPZ equation in 1+1 dimensions , for which the form of these correlation functions is known. An adequate experiment to probe these correlation functions at low momentum compared with the inverse London penetration depth (where such length is typically of the order of a thousand Angstroms) would presumably be low-angle neutron scattering from the magnetic fluctuations in a layered geometry. Other possible experiments which could probe the magnetic properties of the system in the ordered phase would be the use of the magnetooptical Kerr effect or of the Faraday effect on samples with a single magnetic layer to image such a layer.
The structure of this paper is as follows: in section II, we define our model in terms of the geometry of the system and the equations which describe it. We also describe the type of boundary conditions we have to consider. In section III, we present the solution of the equations for a single dipole and construct the solution for a general in-plane distribution of magnetisation by linear superposition. In section IV, we compute the dipolar energy of the distribution of magnetisation and discuss the physical limitations of the model we have considered. Finally, in section V, we present our conclusions.
## II Geometry of the model and relevant equations
The geometry of the model is as follows: an infinite distribution of in-plane magnetisation is placed in the plane $`z=0`$. This distribution is constituted by single magnetic dipoles, placed in an arbitrary fashion with respect to one another (see fig. 2). The in-plane constraint implies that all the dipoles point in a direction within the plane. Above and below the distribution are two bulk superconductors which extend from $`z=d`$ (respectively $`z=d`$) to $`z=\mathrm{}`$ (respectively $`z=\mathrm{}`$). The spatial gap with size $`2d`$ is supposed to be filled with an insulator with magnetic permitivity $`\mu _0`$. The two superconductors are identical and have a magnetic permitivity $`\mu `$ (i.e. they are paramagnetic, with relative permitivity $`\mu _r=\mu /\mu _0`$). These supercondutors are described by the London equations (see below), which imply a linear relation between the current and the magnetic field.
This linear relation allows us to consider instead a simpler problem, the one of a single magnetic dipole, placed at the origin of the coordinate system and oriented along the $`x`$ axis. Once this problem has been solved, one can construct the solution for the general case simply by using translational and rotational invariance in the plane and by adding the different solutions. The linearity of the equations will guaranty that the linear combination is also a solution. Furthermore, a uniqueness theorem proved by London guaranties that this solution is unique.
In the spatial gap, the system is described by the following equations ,
$`\times \text{h}`$ $`=`$ 0 (1)
$`\text{b}`$ $`=`$ $`0`$ (2)
b $`=`$ $`\mu _0(\text{h}+\text{m}),`$ (3)
where the first two equations are the Maxwell equations for the magnetic field h and the magnetic induction b and the third equation is the constitutive relation between the two. For the case of a single dipole oriented along the $`x`$ axis, the magnetisation $`\text{m}(\text{r})=m\widehat{𝒙}\delta ^3(\text{r})`$, where $`m`$ is the magnitude of the magnetic dipole.
The superconductors are described by the equations
$`\times \text{h}`$ $`=`$ $`\mathit{ȷ}`$ (4)
b $`=`$ $`\mathrm{\Lambda }\times \mathit{ȷ}^s`$ (5)
b $`=`$ $`\mu \text{h}`$ (6)
where the first equation is the Maxwell equation which relates the magnetic field with the ‘free’ current, the second relation is the second London equation which relates the supercurrent $`\mathit{ȷ}^s`$ with the magnetic induction and the third equation is the constitutive relation between the magnetic induction and the magnetic field. The constant $`\mathrm{\Lambda }`$ is dependent on the type of the superconductor. In a static situation such as the one we are considering, the electric field $`\text{e}=\text{0}`$ in the superconductor and the total current $`\mathit{ȷ}=\mathit{ȷ}^s`$, i.e. there is no normal component of the current.
Substituting equations (5), (6) in equation (4), one obtains
$`^2\mathit{ȷ}\lambda _L^2\mathit{ȷ}`$ $`=`$ 0 (7)
$`^2\text{b}\lambda _L^2\text{b}`$ $`=`$ 0 (8)
where $`\lambda _L=(\mathrm{\Lambda }/\mu )^{1/2}`$ is the London penetration depth and where equation (8) follows from taking the curl of (7) and using (5) and where we have used the fact that $`\mathit{ȷ}=0`$ (equation of continuity) and $`\text{b}=0`$. These two equations show that the magnetic flux density and the supercurrent penetrate a layer of thickness $`\lambda _L`$ at the surface of the superconductor (Meissner effect).
These equations have to be suplemented by boundary conditions at the surface of the superconductors. These conditions are the continuity of the normal component of b, of the tangential components of h and of the normal component of the current $`\mathit{ȷ}`$ at the boundary surfaces of the two superconductors . If one chooses $`\mathit{ȷ}=\times (g\widehat{𝒛})`$, where $`g(\text{r})`$ is a solution of
$$^2g\lambda _L^2g=0$$
(9)
then one can satisfy the equations (7) and (8) and the boundary condition $`ȷ_z=0`$ at $`z=\pm d`$. Notice that this choice implies that $`ȷ_z=0`$ throughout the material which is physically reasonable, since $`ȷ_z=0`$ at the surfaces $`z=\pm d`$ and also for $`z=\pm \mathrm{}`$.
On the other hand, in the spatial gap, we obtain from equation (1) $`\text{h}=\mathrm{\Phi }_M`$. Substituting this result in equation (3) and using equation (2), we obtain
$$^2\mathrm{\Phi }_M=(m\widehat{𝒙}\delta ^3(\text{r}))$$
(10)
which is Poisson’s equation. Since we know the solution of this equation in free space (i.e. in the absence of the superconductors) to obtain the solution in this case, we can write
$$\mathrm{\Phi }_M(\text{r})=\frac{m\rho \mathrm{cos}\varphi }{4\pi (\rho ^2+z^2)^{3/2}}+\chi (\rho ,\varphi ,z)$$
(11)
where the first term on the rhs is the solution in free space and the function $`\chi (\rho ,\varphi ,z)`$ is a solution of the Laplace equation, $`^2\chi (\rho ,\varphi ,z)=0`$ and where we have used cylindrical polar coordinates for later convenience.
Therefore we need to solve the modified Helmholtz equation (9) in the superconductors and the Laplace equation for $`\chi `$ in the gap and then fit the two solutions using the continuity conditions for b and h at the boundary.
We can further simplify the problem if we notice that the system is invariant under a $`\pi `$ rotation around the $`x`$ axis. This invariance imposes the conditions
$`g_{}(\rho ,\varphi ,z)=g_+(\rho ,\varphi ,z)`$ (12)
$`\chi (\rho ,\varphi ,z)=\chi (\rho ,\varphi ,z)`$ (13)
where $`g_+`$ (respectively $`g_{}`$) is the solution of the Helmholtz equation in the upper (respectively lower) superconductor. Since the magnetic flux in the superconductor is given in terms of $`g`$ by
$`b_z`$ $`=`$ $`\mu \lambda _L^2\left({\displaystyle \frac{1}{\rho }}{\displaystyle \frac{}{\rho }}\left(\rho {\displaystyle \frac{g}{\rho }}\right)+{\displaystyle \frac{1}{\rho ^2}}{\displaystyle \frac{^2g}{\varphi ^2}}\right)`$ (14)
$`b_\rho `$ $`=`$ $`\mu \lambda _L^2{\displaystyle \frac{^2g}{\rho z}}`$ (15)
$`b_\varphi `$ $`=`$ $`{\displaystyle \frac{\mu \lambda _L^2}{\rho }}{\displaystyle \frac{^2g}{\varphi z}}`$ (16)
and we have, in the spatial gap $`\text{b}=\mu _0(\mathrm{\Phi }_M(\text{r})+m\widehat{𝒙}\delta ^3(\text{r}))`$, then the continuity conditions for b and h imply that at $`z=d`$,
$`{\displaystyle \frac{\mathrm{\Phi }_M}{z}}|_{z=d}`$ $`=`$ $`\mu _r\lambda _L^2\left({\displaystyle \frac{^2g_+}{^2z}}\lambda _L^2g_+\right)_{z=d}`$ (17)
$`{\displaystyle \frac{\mathrm{\Phi }_M}{\rho }}|_{z=d}`$ $`=`$ $`\lambda _L^2\left({\displaystyle \frac{^2g_+}{\rho z}}\right)_{z=d}`$ (18)
$`{\displaystyle \frac{\mathrm{\Phi }_M}{\varphi }}|_{z=d}`$ $`=`$ $`\lambda _L^2\left({\displaystyle \frac{^2g_+}{\varphi z}}\right)_{z=d}`$ (19)
where we have used the fact that $`g_+`$ is a solution of the Helmholtz equation and that, in the superconductor, $`\text{b}=\mu \text{h}`$. A similar set of conditions is valid at $`z=d`$ but they are trivially related to these conditions by equations (12) and (13).
The above equations and boundary conditions are sufficient to determine the solution of the problem within the London approximation.
## III The single dipole solution and the general solution for an arbitrary distribution of magnetisation
We concluded in the previous section that in order to find the field and current distributions for the case of a single dipole, one needs to find a joint solution of the Laplace and Helmholtz equations, which satisfies the appropriate boundary conditions (17)-(19). Such a solution can be most easily found using cylindrical polar coordinates and is given in terms of Fourier-Bessel integral transforms by
$`\mathrm{\Phi }_M(\rho ,\varphi ,z)`$ $`=`$ $`{\displaystyle \frac{m}{4\pi }}({\displaystyle \frac{\rho \mathrm{cos}\varphi }{(\rho ^2+z^2)^{3/2}}}+{\displaystyle _0^{\mathrm{}}}dkkJ_1(k\rho )\mathrm{cosh}(kz)\mathrm{cos}\varphi `$ (21)
$`\times {\displaystyle \frac{e^{kd}(\mu _r^1\sqrt{k^2+\lambda _L^2}k)}{k\mathrm{cosh}(kd)+\mu _r^1\sqrt{k^2+\lambda _L^2}\mathrm{sinh}(kd)}})`$
for $`\mathrm{\Phi }_M(\text{r})`$ and
$`g_+(\rho ,\varphi ,z)`$ $`=`$ $`{\displaystyle \frac{m}{4\pi \mu _r\lambda _L^2}}{\displaystyle _0^{\mathrm{}}}𝑑kke^{\sqrt{k^2+\lambda _L^2}(zd)}J_1(k\rho )\mathrm{cos}\varphi `$ (23)
$`\times {\displaystyle \frac{1}{k\mathrm{cosh}(kd)+\mu _r^1\sqrt{k^2+\lambda _L^2}\mathrm{sinh}(kd)}}`$
for $`g_+(\text{r})`$ with $`g_{}(\rho ,\varphi ,z)=g_+(\rho ,\varphi ,z)`$ and where $`J_1(x)`$ is the Bessel function of order one. These integrals can only be calculated explicitly in the particular case $`d0`$, $`\mu _r=1`$. We obtain
$`\mathrm{\Phi }_M(\rho ,\varphi ,0)`$ $`=`$ $`{\displaystyle \frac{m}{4\pi }}\left({\displaystyle \frac{1}{\rho \lambda _L}}+{\displaystyle \frac{e^{\rho /\lambda _L}}{\rho ^2}}\right)\mathrm{cos}\varphi `$ (24)
$`g_+(\rho ,\varphi ,0)`$ $`=`$ $`{\displaystyle \frac{m\mathrm{cos}\varphi }{4\pi \lambda _L^2\rho }}.`$ (25)
The magnetic potential $`\mathrm{\Phi }_M(\rho ,\varphi ,0)`$ corresponds to the magnetic potential of a dipole which produces a 3d field at short distances and that at distances $`\lambda _L`$ produces a 2d field, i.e. the field produced by a dipole in a two dimensional space. This behaviour can be traced to the Meissner effect, which confines the flux lines to the spatial gap and to a region of size $`\lambda _L`$ in each of the superconductors. Although the $`d=0`$ case is somewhat unphysical (the superconductor would have to withstand an infinite field), we shall see that as long as $`d\lambda _L`$, this type of behaviour is essentially unchanged.
Now, in order to generalise this solution to the case of an arbitrary distribution of magnetisation, we represent the magnetisation per unit of area in the form
$$\text{m}(𝝆)=\underset{i}{}\text{m}_i\delta ^2(𝝆𝝆_i)$$
(26)
where $`\text{m}_i`$ is a dipole situated at $`𝝆_i=(x_i,y_i)`$ and $`𝝆=(x,y)`$. In real systems, $`𝝆_i`$ will correspond to the sites of a two dimensional lattice where the dipoles are situated. The solution corresponding to this distribution of magnetisation is given by the linear superposition of the solutions corresponding to each $`\text{m}_i`$, i.e.
$`\mathrm{\Phi }_M(\rho ,\varphi ,z)`$ $`=`$ $`{\displaystyle \underset{i}{}}{\displaystyle \frac{\text{m}_i(𝝆𝝆_i)}{4\pi (𝝆𝝆_i^2+z^2)^{3/2}}}`$ (29)
$`+{\displaystyle \underset{i}{}}{\displaystyle \frac{\text{m}_i(𝝆𝝆_i)}{4\pi 𝝆𝝆_i}}{\displaystyle _0^{\mathrm{}}}𝑑kkJ_1(k𝝆𝝆_i)\mathrm{cosh}(kz)`$
$`\times {\displaystyle \frac{e^{kd}(\mu _r^1\sqrt{k^2+\lambda _L^2}k)}{k\mathrm{cosh}(kd)+\mu _r^1\sqrt{k^2+\lambda _L^2}\mathrm{sinh}(kd)}}`$
for $`\mathrm{\Phi }_M(\text{r})`$ and
$`g_+(\rho ,\varphi ,z)`$ $`=`$ $`{\displaystyle \underset{i}{}}{\displaystyle \frac{\text{m}_i(𝝆𝝆_i)}{4\pi \mu _r\lambda _L^2𝝆𝝆_i}}{\displaystyle _0^{\mathrm{}}}𝑑kkJ_1(k𝝆𝝆_i)e^{\sqrt{k^2+\lambda _L^2}(zd)}`$ (31)
$`\times {\displaystyle \frac{1}{k\mathrm{cosh}(kd)+\mu _r^1\sqrt{k^2+\lambda _L^2}\mathrm{sinh}(kd)}}`$
for $`g_+(\text{r})`$. The function $`g_{}(\text{r})`$ is constructed from the single dipole solution in an analogous manner. Note that one cannot use the equations (12) and (13) because we no longer have the $`\pi `$ rotation symmetry around the $`x`$ axis . Using the magnetic potential $`\mathrm{\Phi }_M(\text{r})`$, we can now compute the dipolar energy of the system.
## IV The magnetic energy of the system
The dipolar energy of the system can be obtained by substituting the dipole distribution by an equivalent distribution of loops of current, i.e. one which will produce the same field distribution. The easiest way to compute the energy necessary for the formation of such a current distribution is to compute it with the currents of the individual loops which form the distribution kept constant. The dipolar energy of the magnetisation distribution is equal to the energy necessary to form the current distribution when the fluxes in each loop are kept constant, which is minus the energy computed with constant currents and is given by
$$E_m=\frac{1}{2}\mu _0d^2\rho \text{m}(𝝆)\text{h}(\rho ,\varphi ,z=0).$$
(32)
Substituting $`\text{h}=\mathrm{\Phi }_M`$, with $`\mathrm{\Phi }_M(\text{r})`$ given by (29) and approximating the discrete sums on the lattice by integrals, we obtain
$`E`$ $`=`$ $`{\displaystyle \frac{\mu _0}{8\pi }}{\displaystyle _{BZ}}{\displaystyle \frac{d^2k}{(2\pi )^2}}{\displaystyle \frac{2\pi }{k}}(\text{k}\text{m}(\text{k}))(\text{k}\text{m}(\text{k}))\times `$ (34)
$`{\displaystyle \frac{(\mu _r^1\sqrt{k^2+\lambda _L^2}+k)e^{kd}+(\mu _r^1\sqrt{k^2+\lambda _L^2}k)e^{kd}}{(\mu _r^1\sqrt{k^2+\lambda _L^2}+k)e^{kd}(\mu _r^1\sqrt{k^2+\lambda _L^2}k)e^{kd}}},`$
where $`\text{m}(\text{k})=_i\text{m}_ie^{i\text{k}𝝆_i}`$ and where the integrals over $`k`$ are over the first Brillouin zone of the reciprocal lattice. The expression (34) was obtained in the limit in which one can disregard the lattice structure of the dipole distribution. If such a structure has to be taken into account, then one has to use Ewald summation methods to handle the discrete sums in equation (29).
The expression (34) has two important limits. The first is when $`\mu _r=1`$ and $`\lambda _L\mathrm{}`$, or when $`d\mathrm{}`$. We obtain
$$E=\frac{\mu _0}{8\pi }_{BZ}\frac{d^2k}{(2\pi )^2}\frac{2\pi }{k}(\text{k}\text{m}(\text{k}))(\text{k}\text{m}(\text{k}))$$
(35)
which is the familiar result for the dipolar energy of a thin film. The second limit is when the largest contribution to the energy comes from modes $`\text{m}(\text{k})`$ for which $`kL^1`$ where $`L`$ is a length such that $`L\lambda _Ld`$ (we take $`\lambda _Ld`$). In this case, we obtain
$`E`$ $``$ $`{\displaystyle \frac{\mu _0}{4\mu _r\lambda _L}}{\displaystyle _{kL^1}}{\displaystyle \frac{d^2k}{(2\pi )^2}}{\displaystyle \frac{(\text{k}\text{m}(\text{k}))(\text{k}\text{m}(\text{k}))}{k^2}}\left(\mathrm{\hspace{0.17em}1}+{\displaystyle \frac{1}{2}}(k\lambda _L)^2+O(k^4)\right)`$ (36)
which shows that for $`k\lambda _L^1`$ the energy of the system has the same form as the energy of a system of dipoles in a two dimensional space. Also, comparing equation (34) (with $`\mu _r=1`$) with (35), one can conclude that the dipolar interaction has effectively been enhanced with respect to the simple film situation, since the fraction in (34) is always larger than 1 when $`\mu _r=1`$. This result can be easily understood from the fact that the energy given in (34) also includes the kinetic energy of the supercurrent .
The modified dipolar kernel $`F(k)=\frac{1}{k}\frac{(\sqrt{k^2+\lambda _L^2}+k)e^{kd}+(\sqrt{k^2+\lambda _L^2}k)e^{kd}}{(\sqrt{k^2+\lambda _L^2}+k)e^{kd}(\sqrt{k^2+\lambda _L^2}k)e^{kd}}`$ given in (34) with $`\mu _r=1`$, is plotted against the kernels $`1/k`$ and $`1/k^2`$, which appear respectively, in (35) and (36), in figure 3.
The question now arises if one can indeed detect such change of behaviour in the dipolar interaction in artificial superlattices of superconductors and ferromagnetic materials or even in naturally layered systems like RuSr<sub>2</sub>GdCu<sub>2</sub>O<sub>8-δ</sub> or RuSr<sub>2</sub>Gd<sub>1+x</sub>Ce<sub>1-x</sub>Cu<sub>2</sub>O<sub>10</sub>. The detection of such a change of behaviour relies on the possibility of finding an experimental system which obeys a series of constraints.
Firstly, the London approximation, which was used and which postulates a local relation between the current and the magnetic induction, is only valid sufficiently close to the critical temperature at which the superconducting phase transition occurs. But we believe that the qualitative features of this solution should be valid even when the London approximation is not (i.e. at low temperatures compared to the critical temperature). Physically, one should expect this type of behaviour as long as the superconductor displays the Meissner effect.
Secondly, we have considered a ferromagnetic layer with zero thickness. The solution obtained can only be valid if one can neglect the proximity effect in the ferromagnetic layer, i.e. the induction of superconductivity in the ferromagnetic material by the superconductor . Otherwise, in a layer with finite thickness, an effective exchange interaction will be induced between the magnetic moments and this interaction also has to be taken into account. However, this effect is negligible in insulators. Therefore, an appropriate material to choose for the ferromagnetic layer would be a ferromagnetic insulator . We have also neglected the supression of the superconducting order parameter at the boundary between the two media, which always occurs in the presence of a film of finite thickness. In the case of a film made of a ferromagnetic insulator, the de Gennes boundary conditions tell us that the order parameter is effectively quenched to zero at the boundary . Since the London penetration depth depends on this parameter, this will mean that the effective London penetration depth will be greater than the penetration depth measured in the absence of the ferromagnetic film. Nevertheless, the London approximation remains valid (provided that we are working in a weak field situation, see below). The opposite limit, in which the coupling between the superconducting and the ferromagnetic layers is primarily due to the proximity effect has been considered, using a microscopic approach by Radović et al. and by Schinz and Schwabl using a phenomenological description. In both cases, the problem treated reduces to the one of decoupled thin superconducting layers embedded in a ferromagnet. Experiments done in superlattices of Fe/V and Nb/Gd have confirmed these results.
Thirdly, one cannot have any Josephson currents flowing through the insulating junction, otherwise the boundary condition $`ȷ_z=0`$ at $`z=\pm d`$ is not valid. This implies that the phase difference between the superconducting order parameters of the two superconductors is zero. This is the case for most superconducting systems in equilibrium but in certain junctions containing ferromagnetic materials a non-zero difference between the phases of the superconducting order parameter in each side of the junction, leading to a current in equilibrium, has been predicted .
Finally, we have not considered the thermodynamical stability of the system we are working with. This implies that we are working in a weak field situation, i.e. the field produced by the distribution of magnetisation is much lower than the critical field at which the superconducting system undergoes a transition to the intermediate state (in the case of type I superconductors) or to the mixed state (in the case of type II superconductors). One way to achieve this condition is to choose either a magnetic material with a low saturation magnetisation or a superconductor with a high critical field (in the case of type I superconductor), e.g. V, or a high field $`H_{c1}`$ (in the case of a type II superconductor), e.g. Nb.
In summary, despite the restrictions pointed above, we think that there is room for believing that one could manufacture systems composed of alternating magnetic and superconducting layers where the above effect could be detected using the methods discussed in the introduction or others.
## V Conclusions
We have computed the magnetic field distribution of an arbitrary distribution of planar magnetisation in a spatial gap between two superconductors. The purpose of this calculation is to provide a simple model for a superlattice of ferromagnetic materials and superconductors where the dipolar interaction between the magnetic moments in the ferromagnetic system is taken into account.
We have also computed the dipolar energy of such a system and we have shown that for low momenta compared to the London inverse penetration depth, the system has a dipolar energy which resembles the energy of a magnetisation distribution in 2 dimensions and hence it leads to an enhancement of the dipolar energy compared to that of a simple ferromagnetic film.
As to possible future directions of research one can point out the possibility to consider a treatment of the same problem in the line of what was done in reference , i.e. consider the Landau-Ginzburg theory of the layered system including the free energy of the superconducting system and of the ferromagnetic layer and the interaction between the two via the electromagnetic interaction and the proximity effect, with the appropriate boundary conditions. However, in this case one has to take into account the non-linearity of the problem and the three dimensional character of a solution involving a non-uniform distribution of magnetisation in the plane.
A slightly different model which can also be treated exactly with the methods developed in this paper is the one of a thin ferromagnetic layer between two bulk type II superconductors above $`H_{c1}`$, i.e. with flux penetration in the form of an Abrikosov vortex lattice oriented along the $`z`$ axis. In this case, the boundary conditions (17-19) would be unchanged, but the equation determining $`g(\text{r})`$, equation (9), would be replaced by an inhomogeneous equation where the role of sources is played by the vortices. Also, the $`\pi `$ rotational invariance of the system along the $`x`$ axis is no longer present, but one can still draw useful conclusions from the reflexion properties of the system in the $`xy`$ plane (see ). What makes the model much more difficult to solve is the fact that one cannot exclude the presence of Neumann functions in the Fourier-Bessel integrals which determine $`g(\text{r})`$ and $`\mathrm{\Phi }_M(\text{r})`$, and also that in this case one cannot use the equation (32) as the expression for the total magnetic energy of the system (compare with ). One should nevertheless mention that a method for computing the magnetic energy in the presence of vortices has been considered by Erdin et al. for the case of a thin magnetic film interacting with a thin superconducting film.
Acknowledgements: We acknowledge many helpful discussions with M. Kulić, D. Rainer, H. Braun, H. Kinder, G. Eilenberger, M. A. Santos, C. Bracher, M. Riza and M. Kleber. J.E.S. acknowledges the support of the European Commission, Contract No. ERB-FMBI-CT 97-2816 and from the Deutsche Forschungsgemeinschaft Schwerpunktprogramm ‘Strukturgradienten in Kristallen’, contract no. Schw 348/12-1 (from 01/03/00). E. F. acknowledges the support of the Deutsche Forschungsgemeinschaft through an Heisenberg fellowship, contract no. FR850/3. E. F. also acknowledges the hospitality of the Institut für Theoretische Physik, LMU München, where part of this work was done. F. S. acknowledges the support of the Deutsche Forschungsgemeinschaft Einzelprojekt Schw. 348/10-1 and of the BMBF Verbundprojekt 03-SC5-TUM0.
Figure Captions
Figure 1. Thin ferromagnetic film between bulk superconductors (schematic).
Figure 2. In-plane distribution of magnetisation between bulk superconductors (schematic).
Figure 3. The modified dipolar kernel $`F(k)`$ (the continous grey plot) is plotted against the kernels $`1/k`$ (the long-dashed plot) and $`1/k^2`$ (the short-dashed plot). We have taken $`\mu _r=1`$ and $`d=0.1`$, $`\lambda _L=1`$ in $`F(k)`$ (in arbitrary units). It is seen that the function $`F(k)`$ interpolates between $`1/k^2`$ and $`1/k`$, the crossover length being the London penetration depth $`\lambda _L=1`$.
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# Enhancement of the inverse-cascade of energy in the two-dimensional averaged Euler equations
## Abstract
For a particular choice of the smoothing kernel, it is shown that the system of partial differential equations governing the vortex-blob method corresponds to the averaged Euler equations. These latter equations have recently been derived by averaging the Euler equations over Lagrangian fluctuations of length scale $`\alpha `$, and the same system is also encountered in the description of inviscid and incompressible flow of second-grade polymeric (non-Newtonian) fluids. While previous studies of this system have noted the suppression of nonlinear interaction between modes smaller than $`\alpha `$, we show that the modification of the nonlinear advection term also acts to enhance the inverse-cascade of energy in two-dimensional turbulence and thereby affects scales of motion larger than $`\alpha `$ as well. This latter effect is reminiscent of the drag-reduction that occurs in a turbulent flow when a dilute polymer is added.
The two-dimensional incompressible, Euler equations are
$$_t\omega +\mathbf{}\left(\text{u}\omega \right)=0,\mathbf{}\text{u}=0,\omega (t=0)=\omega _0,$$
(1)
where $`\omega =\mathbf{}\times \text{u}`$ is the vorticity, u is the spatial velocity vector field, $`t`$ denotes time, and all the dependent variables depend on $`t`$ and $`\text{x}=(x_1,x_2)`$, the Cartesian coordinates in the plane. An inversion of the vorticity-velocity relation yields $`\text{u}=\text{K}(\text{x},\text{y})\omega (\text{y})𝑑\text{y}`$, where $`\text{K}=\mathbf{}^{}G`$, $`G`$ is the solution of $`\mathrm{}G=\delta `$, and $`\mathbf{}^{}=(_{x_1},_{x_2})`$.For fluid motion over the entire plane, $`\text{K}(\text{x},\text{y})=(2\pi )^1\mathbf{}^{}\mathrm{log}|\text{x}\text{y}|`$. Let $`𝜼_t`$ denote the flow of $`\text{u}_t=\text{u}(t,)`$, so that $`d𝜼_t/dt=\text{u}_t(𝜼(t))`$. Because $`\text{u}_t`$ is divergence-free, the flow map $`𝜼_t`$ is an area-preserving transformation for each $`t`$. It follows that
$`{\displaystyle \frac{d𝜼_t}{dt}}`$ $`={\displaystyle \text{K}(𝜼_t(\text{x}),𝜼_t(\text{y}))\omega (𝜼_t(\text{y}))𝑑\text{y}}`$ (2)
$`={\displaystyle \text{K}(𝜼_t(\text{x}),𝜼_t(\text{y}))\omega _0(\text{y})𝑑\text{y}},`$ (3)
where the last equality is a consequence of the pointwise conservation of vorticity along Lagrangian trajectories, $`\omega (𝜼_t(\text{x}))=\omega _0(\text{x})`$. Thus, the initial vorticity field completely determines the fluid motion. Choosing the initial vorticity to be a sum of $`N`$ point vortices $`\delta _i`$ positioned at the points $`\text{x}_i`$ in the plane with circulations $`\mathrm{\Gamma }_i`$, $`\omega _0=_{i=1}^N\mathrm{\Gamma }_i\delta _i`$, equation (2) produces the classical point-vortex approximation to (1). This approximation is known to be highly unstable, as finite-time collapse of vortex centers may occur .
Chorin’s vortex blob method alleviates the instability of the point-vortex scheme by smoothing each delta function $`\delta _i`$ with a vortex blob $`\chi `$, a function that decays at infinity, and whose mass is mostly supported in a disc of diameter $`\alpha `$. Thus, instead of using the integral kernel $`\text{K}(\text{x},\text{y})`$, one uses the smoother kernel $`\text{K}^\alpha =\mathbf{}^{}G^\alpha `$ where $`G^\alpha `$ is the solution of $`\mathrm{}G^\alpha =\chi .`$ The vortex-blob method then evolves the point-vortex initial data, which we shall now call $`q_0`$, by the ordinary differential equation
$$\frac{d𝜼_t^\alpha }{dt}=\text{K}^\alpha (𝜼_t^\alpha (\text{x}),𝜼_t^\alpha (\text{y}))q_0(\text{y})𝑑\text{y}.$$
(4)
Henceforth, to keep the notation concise, we will drop the superscript $`\alpha `$ when there is no ambiguity.
When the vortex-blob $`G^\alpha `$ is the modified Bessel function of the second kind $`K_0`$, it is the fundamental solution of the operator $`(1\alpha ^2\mathrm{})`$ in the plane, and the vorticity $`q`$ is related to the smoothed velocity vector field u by $`q=(1\alpha ^2\mathrm{})\mathbf{}\times \text{u}`$. Thus, Chorin’s vortex method for this choice of smoothing is given by the partial differential equation
$$_tq+\mathbf{}\left(\text{u}q\right)=0,\mathbf{}\text{u}=0,q(t=0)=q_0,$$
(5)
The system of equations (5) is also known as the two-dimensional isotropic averaged Euler equations, and are derived by averaging over Lagrangian fluctuations of order $`\alpha `$ about the macroscopic flow field . When the constant $`\alpha >0`$ is interpreted as a material parameter which measures the elastic response of the fluid due to polymerization instead of as a spatial length scale, then (5) are also exactly the equations that govern the inviscid flow of a second-grade non-Newtonian fluid . According to Noll’s theory of simple materials, (5) are obtained from the unique constitutive law that satisfies material frame-indifference and observer objectivity. Consequently, the vortex method with the Bessel function $`K_0`$ smoothing naturally inherits these characteristics .
We find the connections between averaging Euler equations over Lagrangian fluctuations, a constitutive theory for polymeric fluids, and a classical numerical algorithm to be quite intriguing and suspect that these equations will be important from a modeling standpoint. However, most previous studies of the averaged Euler equations have been of a mathematical nature, and we are aware of only a few cases where this system has been used as a (dynamic) modeling tool: Chen et al. used a viscous version of the three-dimensional averaged Euler equations to simulate isotropic turbulence and found that they could reproduce large-scale features without fully resolving the flow. Nadiga considered the inviscid two-dimensional form of the averaged Euler equations and demonstrated that for suitably chosen values of $`\alpha `$, the large-scale spectral-scalings of the Euler equations could be preserved while achieving a faster spectral decay at the smaller scales. Finally, Nadiga & Margolin used an extension of the two-dimensional averaged Euler equations in a geophysical context to model the effects of mesoscale eddies on mean flow.
Before we go on to consider numerical simulation of the averaged Euler system, we wish to point out that there is also a beautiful geometric structure to (5) which follows the framework developed by Arnold and Ebin and Marsden . While the details of this particular issue are far outside the scope of this article, it is, nevertheless, worthwhile to state the result. Arnold showed that the appropriate configuration space for a perfect incompressible fluid is the group of all area preserving diffeomorphisms of the fluid container, and that solutions of the Euler equations are geodesics on this group with respect to a certain kinetic energy metric, characterized by the inner-product $`\left(\text{u}\text{v}\right)𝑑\text{x}`$ for two divergence-free vector fields u and v. The system (5) also has this geometric property, but now the metric is instead characterized by $`\left(\text{u}\text{v}+2\alpha ^2\text{Def}(\text{u})\text{Def}(\text{v})\right)𝑑\text{x},`$ where $`\text{Def}(\text{u})`$ is the rate of deformation tensor $`(\mathbf{}\text{u}+(\mathbf{}\text{u})^T)/2`$ . Equations (5) thus preserve the Hamiltonian structure of the Euler equations. In particular, vorticity remains pointwise conserved by the smooth Lagrangian flow $`𝜼_t^\alpha `$ so that $`q(𝜼_t^\alpha (\text{x}))=q_0(\text{x}),`$ the vorticity momenta $`q^p𝑑\text{x}`$ are conserved, and so the Kelvin circulation theorem remains intact as well.
Since, in each of the three different scenarios—averaged Euler, vortex-blob method, and inviscid second-grade fluid—the essential modification of the original equations is a change of the advective nonlinearity of fluid dynamics, we will now consider forced-dissipative simulations of the system (5) to demonstrate the effect of such an inviscid modification. If we use a vorticity-stream function formulation, the evolution of both the Euler and averaged Euler sytems can be represented by
$$\frac{\omega }{t}+(1\alpha ^2\mathrm{})^1J[\psi ,(1\alpha ^2\mathrm{})\omega ]=F+D;$$
(6)
where $`\omega =\mathrm{}\psi `$, $`J`$ is the Jacobian, $`F`$ the forcing, and $`D`$ the dissipation. (The Euler system corresponds to $`\alpha =0`$.) Our numerical scheme consists of a fully dealiased pseudospectral spatial discretization and a (nominally) fifth-order, adaptive timestep, embedded Runge-Kutta Cash-Karp temporal discretization of (6) (see Nadiga for details). With such a scheme, among the infinity of inviscid ($`F=D=0`$) conserved quantities for (6), the only two conservation properties that survive are those for the kinetic energy $`E_{H^1}`$ and enstrophy $`Z_{H^2}`$ given respectively by
$`E_{H^1}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \left(|\text{u}|^2+\alpha ^2|\mathbf{}\text{u}|^2\right)𝑑\text{x}}\left(=\text{u}_{H^1}^2\right),`$ (7)
$`Z_{H^2}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \left[\left(1\alpha ^2\mathrm{}\right)\omega \right]^2𝑑\text{x}}\left(=\omega _{H^2}^2\right).`$ (8)
In the forced-dissipative runs to be considered, the forcing $`F`$ is achieved by keeping the amplitudes of modes with wavenumbers in the small wavenumber band $`10k<10.001`$ constant in time. The dissipation, $`D`$, is a combination of a fourth order hyperviscous operator and a large-scale friction term: $`D=\delta \psi \left(\nu \mathrm{}\right)^4\omega `$, as has been used in numerous previous studies of two-dimensional turbulence. The form and value of the forcing and dissipation are held exactly the same for all the runs to be presented, irrespective of the resolution and the value of $`\alpha `$.
On the one hand, it could be argued that since the energy and enstrophy that are conserved (in an unforced-inviscid setting) are $`E_{H^1}`$ and $`Z_{H^2}`$ respectively, it is their dynamics which is of primary importance. On the other, it could be argued that in the context of (6), the interest in small scales is only in so much as it affects the larger scales and to that extent $`\alpha `$ has no primary significance and that it is really the large scale components of energy and enstrophy in (7) that are of primary interest. While both these points of view are reasonable, in this short article, we proceed with the latter and concern ourselves with the dynamics of the usual kinetic energy and usual enstrophy as given by
$$E=\frac{1}{2}|\text{u}|^2𝑑\text{x},Z=\frac{1}{2}\omega ^2𝑑\text{x}.$$
(9)
Fig. 1 shows the evolution of the kinetic energy $`E`$ with time for four different values of $`k_\alpha `$. For these computations, 512 physical grid points were used in each direction, resulting in, after accounting for dealiasing, a maximum, circularly-symmetric wavenumber, $`k_{max}`$, of 170. The four runs correspond to $`k_\alpha `$ of $`\mathrm{}`$ (dissipative Euler) and 42, 21, and 14. This figure shows that for identical forcing and dissipation, the tendency with increasing $`\alpha `$ (equivalently decreasing $`k_\alpha `$) is to achieve an overall balance which makes the flow less viscous.
While the kinetic energy of the runs with different $`\alpha `$ shows a definite trend (increasing with increasing $`\alpha `$), such is not the case with the enstrophy shown for the same four cases in the inset of Fig. 1. Here interestingly, all the runs with nonzero $`\alpha `$ seem to display approximately the same level of enstrophy which is lower than for $`k_\alpha =0`$. This indicates:
* That the small-scale behavior is quite different when $`\alpha =0`$ and when $`\alpha `$ is nonzero (as noted in previous studies ), but that this difference is not sensitively dependent on the value of $`\alpha `$ for the interesting range of values of $`\alpha `$.
* That the more significant change with $`\alpha `$ is the behavior of the large scales.
Therefore, to further examine the nature of this (reduced-viscous) behavior of the large scales, we examine the energy-wavenumber spectra in Fig. 2. Here, the average of the one dimensional energy spectrum $`E(k)`$ between times 5 and 20 is plotted against the scalar wavenumber $`k`$. Figure 2 shows that the reduced-viscous behavior for increasing $`\alpha `$ is achieved by systematically increasing the energy in modes larger in scale than the forcing scale and decreasing the energy in modes smaller in scale compared to the forcing scale.
The larger energy content in the larger scales implies an enhancement of the inverse cascade of energy of two-dimensional turbulence by the nonlinear-dispersive modification of the advective nonlinearity when $`\alpha >0`$ in (6). So also, the decreased energy content in the smaller scales is attributable to the same nonlinear-dispersive modification. In the following, we give a simple dimensional argument to explain the observed behavior. For this, consider the governing equations in the form (5). In close analogy with the classical picture for the inertial ranges of two-dimensional dissipative Euler equations , a Kolmogorov-like cascade picture for (5) shows that the inertial range consists of two subranges, the enstrophy cascade subrange where there is a down-scale cascade of the $`Z_{H^2}`$ enstrophy defined in (7), and the energy cascade subrange where there is an up-scale cascade of the $`E_{H^1}`$ energy defined in (7). $`E_{H^1}`$ and $`Z_{H^2}`$ are the relevant energy and enstrophies since these are the ones which are conserved in an inviscid and unforced case.
If we assume that the wavenumber $`k_\alpha `$ only appears in the Helmholtz operator, as it does in the governing equations, then
* In the enstrophy cascade subrange,
$$E(k)\beta _{H^2}^ak^b,$$
(10)
where $`\beta _{H^2}`$ is the rate of dissipation of $`Z_{H^2}`$ enstrophy, and $`a`$ and $`b`$ are exponents to be determined by dimensional analysis. If $`L`$ and $`T`$ are characteristic length and time scales in the enstrophy cascade subrange, (10) implies
$$L^3T^2=T^{3a}\left(1+\alpha ^2L^2\right)^{2a}L^b,$$
from which $`a=2/3`$. However, even in the enstrophy cascade subrange, the value of $`b`$ depends on the the ratio $`\alpha /L`$. For $`\alpha L`$, of course, $`b=3`$, and the classical $`E(k)k^3`$ is recovered; when $`\alpha L`$, $`E(k)k^{\frac{17}{3}}`$. Finally, when $`\alpha `$ is comparable to $`L`$, it is easy to see that $`E(k)`$ decays faster than for Euler, but slower than $`k^{\frac{17}{3}}`$ (as may be seen in Fig. 2).
* In the energy cascade subrange,
$$E(k)ϵ_{H^1}^ak^b,$$
(11)
where $`ϵ_{H^1}`$ is the rate of dissipation of $`E_{H^1}`$ energy, and $`a`$ and $`b`$ are exponents to be determined by dimensional analysis. If $`L`$ and $`T`$ are characteristic length and time scales, now, in the energy cascade subrange, (11) implies
$$L^3T^2=T^{3a}\left(1+\alpha ^2L^2\right)^aL^{2ab},$$
from which $`a=2/3`$. Again, even in the energy cascade subrange, the value of $`b`$ depends on the the ratio $`\alpha /L`$. For $`\alpha L`$, of course, $`b=\frac{5}{3}`$, and the classical $`E(k)k^{\frac{5}{3}}`$ is recovered; when $`\alpha L`$, $`E(k)k^3`$. When $`\alpha `$ is comparable to $`L`$, it is easy to see that the inverse cascade of energy is enhanced, that is $`E(k)`$ increases with decreasing $`k`$ faster than $`k^{\frac{5}{3}}`$ (Euler) but slower than $`k^3`$.
While the effects of the enhancement of the inverse cascade of energy is clear in Fig. 2, we defer the verification of the asymptotic values of the exponent $`b`$ to later studies when we can afford much larger simulations with a good dynamic range in each of the inertial subranges.
The steeper fall-off of the energy spectrum with $`k`$ in the enstrophy cascade range of wavenumbers when $`\alpha >0`$, compared to Euler may, at first, suggest that a coarser resolution may be sufficient to resolve the flow when $`\alpha >0`$ (for the same forcing and dissipation). However, this is not the case, as should be clear from Fig. 3. In this figure, the spectra of the cases previously discussed is replotted together with the corresponding spectra when the resolution is reduced by a quarter ($`k_{max}=128`$) and a half ($`k_{max}=85`$). (The spectra for the different values of $`\alpha `$ are offset to improve clarity.) The degree of non-resolution of the flow due to the reduced resolution is indicated by the deviation of that spectrum from that for the fully resolved case. With a 25% reduction in resolution, the flows are almost resolved for all values of $`\alpha `$, while with a 50% reduction, the flows are not fully resolved anymore. Importantly, the degree of non-resolution is independent of $`\alpha `$ to the lowest order.
Besides their use in describing mean motion, the averaged Euler equations have arisen independently in at least two other contexts—second grade polymeric fluids and vortex blob methods. In this note, we make two observations that are likely to be of fundamental importance in understanding the relevance of these models in describing more realistic flows: While it has been previously noted that with these equations, nonlinear interactions at scales small compared to $`\alpha `$ are suppressed, we have shown here that the modification of the nonlinear advection term in these equations also leads to an enhancement of the inverse cascade of energy in two dimentions—a characteristic feature of two-dimensional turbulence. This in turn implies (1) an overall reduced-viscous behavior and (2) a significant modification of the dynamics of scales larger than $`\alpha `$, both reminiscent of the phenomenon of drag reduction in a turbulent flow when a dilute polymer is added (e.g., see and references therein). Furthermore, we point out that the limiting of the energy spectrum at small scales due to $`\alpha `$ does not, in itself, allow the flow to be resolved on a coarser grid.
The latter notwithstanding, we remark that the averaged Euler equations are useful in better understanding the limit of inviscid fluid flow, since the averaged Euler equations with viscosity, unlike the Euler equations, converge regularly to the solutions of the inviscid system. That is, for an arbitrary but fixed time interval, we can choose $`\alpha `$ small enough so that the solution of the averaged Euler equations are uniformly within any a priori chosen error of the Euler equations and then consider the zero viscosity limit of the viscous, averaged Euler equations.
Finally, we note that as a theoretical model of fluid turbulence, the averaged Euler equations in two dimensions possess other remarkable features. For example, for any initial condition and fixed time interval, one may choose the number of modes $`k_{max}`$ large enough so as to be arbitrarily close to the exact solution of the averaged Euler equations without the addition of viscosity. For such large $`k_{max}`$, and in simulations of unforced decaying turbulence, the averaged Euler equations exhibit a fundametal feature of two-dimensional turbulence: a sharp decrease in enstrophy $`Z`$ during the first few large eddy turnover times. This is extremely interesting, because, while it is necessary to add viscosity to the Euler equations to obtain similar behavior, the averaged Euler equations can reproduce this behavior while exactly conserving an energy. We shall report further on such inviscid simulations in future publications.
The authors thank Jerry Marsden, and David Montgomery for extensive discussions on a variety of issues related to this article. BTN was supported by the Climate Change Prediction Program of the Department of Energy, and SS was partially supported by NSF-KDI grant ATM-98-73133 and the Alfred P. Sloan Research Fellowship.
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# Medium/high field magnetoconductance in chaotic quantum dots
## I Introduction
Magnetoconductance in chaotic quantum dots have attracted a great deal of attention in recent years . Weak localization effects have been thoroughly investigated searching for differences between chaotic and regular cavities . More recently the selfsimilar character of magnetoconductance fluctuations in chaotic quantum dots has deserved several experimental studies . In Ref. it was reported that in cavities with sufficiently soft walls, rather wide leads and a high zero field conductance of the order of 40 conductance quanta, fluctuations were very weak and the magnetoconductance increased steadily in approximately 20% over 50 flux quanta . Although an increase in the magnetoconductance as a function of the magnetic flux in cavities with wide leads may not be that surprising, no theoretical analysis of this result is yet available. This is so despite of the wealth of experimental and theoretical information on magnetoconductance in related mesoscopic systems .
The purpose of the present work is to investigate the magnetoconductance of chaotic quantum dots over a wide range of magnetic field. Quantum dots are described by means of a tight–binding Hamiltonian on $`L\times L`$ clusters of the square lattice. Non–regular (chaotic) behavior is induced by introducing disorder at surface sites, a procedure that has been shown to reproduce all properties of quantum chaotic cavities . The most outstanding conclusions derived from our study are the following. For sufficiently open systems, large leads width $`W`$ or, alternatively, high zero field conductance, the magnetoconductance increases steadily as a function of the magnetic flux, reaching a maximum at a magnetic flux $`\mathrm{\Phi }_{max}`$ proportional to the the leads width. This effect, which is in agreement with the experimental observations of , does not show up in regular or bulk disordered cavities. The average magnetoconductance versus magnetic flux curve shows four clearly differentiated regions: i) At small fluxes (typically below 1–2 flux quanta) the weak localization peak with the typical Lorentzian shape is observed. ii) This is followed by a flux range over which the magnetoconductance shows a non–universal behavior which depends on leads configuration. iii) The latter lasts until the cyclotron radius becomes of the order of $`L/2`$. Beyond this point the average magnetoconductance increases linearly with the magnetic flux with a slope which increases with the leads width (for very small $`W`$ the slope is nearly zero and the magnetoconductance remains constant in a large flux range) . iv) At large fluxes the magnetoconductance decreases stepwise (each step of one flux quanta), due to the successive crossings of the Fermi energy of the transversal modes that contribute to the current. As discussed below these results are compatible with a theory proposed by Beenakker and van Houten to interpret the experimental results for the magnetoconductance of two point contacts in series .
The rest of the paper is organized as follows. Section II includes a description of our model of chaotic quantum dot and of the method we used to compute the current. The results are presented in Section III, and discussed in terms of the theory of Ref. in Section III. Section IV is devoted to summarize the conclusions of our work.
## II Model and Methods
### A Model
The quantum dot is described by means of a tight-binding Hamiltonian with a single atomic level per lattice site on $`L\times L`$ clusters of the square lattice:
$`\widehat{H}`$ $`=`$ $`{\displaystyle \underset{m,nϵIS}{}}\omega _{m,n}|m,n><m,n|`$ (2)
$`{\displaystyle \underset{<m,n;m^{},n^{}>}{}}t_{m,n;m^{},n^{}}|m,n><m^{},n^{}|,`$
where $`|m,n>`$ represents an atomic orbital on site $`(m,n)`$. Indexes run from 1 to $`L`$, and the symbol $`<>`$ denotes that the sum is restricted to nearest-neighbors. Using Landau’s gauge the hopping integral is given by,
$`t_{m,n;m^{},n^{}}`$ $`=`$ $`\mathrm{exp}\left(2\pi i{\displaystyle \frac{m}{(L1)^2}}\mathrm{\Phi }\right),m=m^{}`$ (4)
$`=1,\mathrm{otherwise}`$
where the magnetic flux $`\mathrm{\Phi }`$ is measured in units of the quantum of magnetic flux $`\mathrm{\Phi }_0=h/e`$. The energy $`\omega _{m,n}`$ of atomic levels at impurity sites ($`IS`$) is randomly chosen between $`\mathrm{\Delta }/2`$ and $`\mathrm{\Delta }/2`$, whereas at other sites $`\omega _{m,n}=0`$. Impurities were taken on all surface sites but those coinciding with the leads entrance sites to avoid excessive (unphysical) scattering. This model has been proposed to simulate cavities with rough boundaries and its properties closely follow those characterizing quantum chaotic systems . Some calculations were also carried out on clusters with $`2L`$ bulk impurities .
### B Conductance
The conductance (measured in units of the quantum of conductance $`G_0=e^2/h`$) was computed by using an efficient implementation of Kubo formula. The method is described in , while applications to mesoscopic systems can be found in . For a current propagating in the $`x`$–direction, the static electrical conductivity is given by:
$$G=2\left(\frac{e^2}{h}\right)\mathrm{Tr}\left[(\mathrm{}\widehat{v}_x)\mathrm{Im}\widehat{}G(E)(\mathrm{}\widehat{v}_x)\mathrm{Im}\widehat{}G(E)\right],$$
(5)
where $`\mathrm{Im}\widehat{}G(E)`$ is obtained from the advanced and retarded Green functions:
$$\mathrm{Im}\widehat{}G(E)=\frac{1}{2i}\left[\widehat{}G^R(E)\widehat{}G^A(E)\right],$$
(6)
and the velocity (current) operator $`\widehat{v}_x`$ is related to the position operator $`\widehat{x}`$ through the equation of motion $`\mathrm{}\widehat{v}_x=[\widehat{H},\widehat{x}]`$, $`\widehat{H}`$ being the Hamiltonian.
Numerical calculations were carried out connecting quantum dots to semiinfinite leads of width W in the range 1–$`L`$. The hopping integral inside the leads and between leads and dot at the contact sites is taken equal to that in the quantum dot (ballistic case). Assuming the validity of both the one-electron approximation and linear response, the exact form of the electric field does not change the value of $`G`$. An abrupt potential drop at one of the two junctions provides the simplest numerical implementation of the Kubo formula since, in this case, the velocity operator has finite matrix elements on only two adjacent layers and Green functions are just needed for this restricted subset of sites. Assuming this potential drop to occur at the left contact ($`lc`$) side, the velocity operator can be explicitly written as,
$$i\mathrm{}v_x=\underset{j=1}{\overset{W}{}}\left(|lc,j><1,j||1,j><lc,j|\right)$$
(7)
where $`(|lc,j>`$ are the atomic orbitals at the left contact sites nearest neighbors to the dot.
Green functions are given by:
$$[E\widehat{I}\widehat{H}\widehat{\mathrm{\Sigma }}_1(E)\widehat{\mathrm{\Sigma }}_2(E)]\widehat{}G(E)=\widehat{I},$$
(8)
where $`\widehat{\mathrm{\Sigma }}_{1,2}(E)`$ are the selfenergies introduced by the two semiinfinite leads . The explicit form of the retarded selfenergy due to the mode of wavevector $`k_y`$ is:
$$\mathrm{\Sigma }(E)=\frac{1}{2}\left(Eϵ(k_y)i\sqrt{4(Eϵ(k_y))^2}\right),$$
(9)
for energies within its band $`|Eϵ(k_y)|<2`$, where $`ϵ(k_y)=2\mathrm{c}\mathrm{o}\mathrm{s}(k_y)`$ is the eigenenergy of the mode $`k_y`$ which is quantized as $`k_y=(n_{k_y}\pi )/(W+1)`$, $`n_{k_y}`$ being an integer from 1 to $`W`$. The transformation from the normal modes to the local tight–binding basis is obtained from the amplitudes of the normal modes, $`<n|k_y>=\sqrt{2/(W+1)}\mathrm{sin}(nk_y)`$. Note that in writing Eq. (9) we assumed that the magnetic field was zero outside the dot . This point will have some relevance in relation to the interpretation of our numerical results in terms of the theory of Ref. .
Calculations were carried out on clusters of linear size $`L=47394`$ (in units of the lattice constant) and at a fixed arbitrarily chosen, Fermi energy $`E=\pi /3`$. In some cases averages over disorder realizations were also done. Although most calculations were performed with input/output leads of width $`W`$ connected from site $`(1,1)`$ to site $`(1,1+W)`$, and from $`(L,L)`$ to $`(L,LW)`$, respectively, other input/output leads configurations were also explored. All calculations on disordered systems were done at a fixed value of the disorder parameter $`\mathrm{\Delta }`$=6.
## III Results
We first discuss results for a $`W/L`$ similar to that used in the experiments of Ref. . In that work conductance measurements were taken on a stadium cavity with a lithographic radius of 1.1 $`\mu `$m and leads 0.7 $`\mu `$m wide, which gives a $`W/L`$ ratio of 0.64. Fig. 1 shows the results for the magnetoconductance in cavities of linear sizes in the range $`L`$=47–394 and leads of width $`W=0.65L`$. The results correspond to a single realization of disorder. The most interesting result is the steady increase of the conductance with the magnetic field. The conductance reaches a maximum at a magnetic flux which, as illustrated in the inset of the Figure, increases linearly with the leads width, or the linear size of the system. Note that, as the results of Fig. 1 correspond to $`WL`$, they cannot allow to identify which of the two parameters ($`W`$ and $`L`$) control the maximum in $`G`$. We have checked, however, that it is in fact $`W`$ the one that matters (see also below). The increase in $`G`$ until the maximum is reached can be as high as 30%. It is interesting to remark that the increase in the conductance occurs with relatively small fluctuations due to the large $`W/L`$ ratio (or degree of opening) of the cavity (see ). Although the experimental data were taken at fields not high enough to observe the maximum shown in Fig. 1, it can be safely assessed that our results are compatible with those of Ref. . In particular the zero field conductance reported in that work was around 38 quanta, increasing up to 46 quanta over approximately 50 flux quanta. This is rather similar to the magnetoconductance curve for $`L=197`$ shown in Fig. 1.
In Fig. 2 we compare the results for the cavity with surface disorder with those for a regular cavity and for a cavity with $`2L`$ bulk impurities. The results indicate that the cavity having surface disorder is the only one that reproduces the experimental results . In regular cavities the conductance does not increase steadily due to the large amplitude oscillation discussed in . This behavior clearly differentiates regular and chaotic cavities. At large fields the result for the cavity with surface disorder coincides with that for the regular cavity. This is a consequence of the fact that for sufficiently high fields the current is dominated by edge–like states which are not affected by surface disorder. Semiclassically one can view carriers motion as short orbits bouncing off the same boundary. The associated quantum states have chirality and are thus commonly refer to as chiral states or edge states. The stepwise decrease of the magnetoconductance observed in regular and chaotic cavities with surface disorder is a consequence of the overall depopulation of Landau levels. It is interesting to note that the cavity with bulk disorder shows a markedly different behavior as this type of disorder can, instead, scatter carriers between opposite sides of the cavity. The results shown in Fig. 2 illustrate the only difference we have found up to now between cavities with surface impurities or with a number of bulk impurities proportional to $`L`$ . Appart from this difference the two behave much alike and in line with what one expects to be the behavior of quantum chaotic cavities .
We have investigated how this steady increase of the magnetoconductance is affected by the leads width. In order to reduce fluctuations, which are particularly important at small $`W`$ , we have averaged the conductance over disorder realizations (600 realizations were included in the calculations). The results for a cavity with rather narrow leads attached at the dots in three different ways are illustrated in Fig. 3. At small fluxes (typically below 1–2 flux quanta) the expected Lorentzian peak characteristic of chaotic cavities is obtained (not clearly visible in the Figure). At higher fluxes, a range over which the conductance behaves in a way that is strongly dependent on leads configuration, is observed. Beyond, the conductance increases linearly with a slope which is very similar in the three cases shown in the Figure. The crossover to the linear behavior occurs at a flux for which the radius of the classical cyclotron orbit $`r_c`$ is roughly $`L/2`$. In our units (flux and conductance quanta $`\mathrm{\Phi }_0=G_0=1`$ and energy $`\mathrm{}^2/(2ma^2)=t=1`$, where $`t`$ is the hopping integral), $`r_c`$ is given by,
$$r_c=\mathrm{}v(E)\frac{(L1)^2}{4\pi }\frac{\mathrm{\Phi }_0}{\mathrm{\Phi }},$$
(10)
where $`\mathrm{}v(E)=2\sqrt{\mathrm{sin}^2k_x+\mathrm{sin}^2k_y}_E`$. At the energy chosen here $`\mathrm{}v(E)2.2`$ . Then the flux at which $`r_c=L/2`$ is, for the size of Fig. 3, $`\mathrm{\Phi }=17\mathrm{\Phi }_0`$, which is very close to the flux at which the mentioned crossover occurs.
Fig. 4 shows the averaged magnetoconductance for cavities of linear size $`L=47`$ and leads widths in the range $`W`$=4–20. It is noted that the linear behavior discussed above appears in all cases at roughly the same magnetic flux indicating that it is only related to the cavity size $`L`$, as suggested by the discussion above. Instead the slope increases with the leads width. At small $`W`$ (systems with a low conductance) the conductance increases very slowly as a function of the magnetic flux. For sufficiently large $`W`$ the conductance reaches a maximum and then decreases stepwise. As remarked above, the latter is a consequence of the overall depopulation of transversal modes (or Landau levels). The slope of the linear part of the magnetoconductance is plotted as a function of leads width in Fig. 5. The results can be accurately fitted by a $`W^2`$ law. The increase of the magnetoconductance can be understood in terms of the increase of the transmission probability of the transversal modes as their edge–like character increases and, consequently, its sensitivity to surface disorder is reduced. The steady increase in the conductance takes place until the mentioned depopulation begins to reduce the number of modes that participate in the current. Although this argument seems plausible, it cannot explain quantitative features of the results such as the linear relation between the conductance and the flux or the increase of the slope as the square of the leads width. This issue is addressed in the following Section.
## IV Discussion
The results discussed in the previous subsection ressemble those predicted by Beenakker and van Houten for the magnetoconductance of two point contact in series and the related experiments of Staring et al . Under the hypothesis that transmission between point contacts occurs with intervening equilibration of the current–carrying edge states, the authors of derived the following expression for the conductance (in the following we take the conductance and flux quanta $`G_0=\mathrm{\Phi }_0=1`$ and do not include spin degeneracy),
$$G(\mathrm{\Phi })=\left[\frac{1}{N_1}+\frac{1}{N_2}\frac{1}{N_L}\right]^1$$
(11)
where $`N_i`$ are the number of occupied subbands in the two contacts or leads ($`i=1,2`$) and in the region between the contacts, or in the present case in the dot ($`i=L`$). Disregarding discretness , $`N_i`$ can be written as,
$$N_i=\frac{n_i}{2B}f(\xi _i)$$
(12)
where the function $`f(\xi _i)`$ is,
$`f(\xi _i)`$ $`=`$ $`{\displaystyle \frac{2}{\pi }}\left[\mathrm{arcsin}\xi _i+\xi _i(1\xi _i^2)^{1/2}\right],\mathrm{if}\xi _i<1`$ (14)
$`=1,\mathrm{if}\xi _i>1`$
with $`\xi _i=l_i/2r_c`$, $`l_i`$ being a characteristic linear dimension in the three regions, in the present case, $`l_i=W_1,W_2,L`$.
We have used Eqs. (11) and (12) to fit the numerical results of Fig. 4. We took as fitting parameter the density or the Fermi velocity, $`\mathrm{}v=2\sqrt{2\pi n}`$, and assumed the same density in the leads and dot. As shown in Fig. 6 a satisfactory fitting is obtained for $`\mathrm{}v=3.65`$, almost twice the actual Fermi velocity in our model (see above). The theory reproduces the three regions that characterize our numerical results: an almost constant $`G`$ for small flux or $`r_c>L/2`$, a steadily increasing $`G`$ up to $`r_cW/2`$ followed by a steep decrease at higher fluxes. It is interesting to note that the theory of Ref. show a better agreement with the numerical results for the case in which the leads are attached at opposite corners of the dot, than with those for the other two lead configurations of Fig. 3. A possible reason for this behavior relies upon the equilibration assumption in Beenhakker and van Houten theory. Equilibration is most likely when the leads are not facing each other, as is the case of leads attached to opposite corners. Instead, when leads are attached at contiguous corners direct transmission is more probable and equilibration requires higher magnetic fields to take place. This is a pictorial illustration of the assumptions under which the theory of Ref. holds.
In order to check whether a linear relationship between the conductance and the flux in a rather wide range of fluxes, as indicated by the fittings of Fig.4, can be understood in terms of this theory, we have expanded the conductance for small $`\xi `$ and $`\xi _L>1`$. The result for leads having the same width and the same Fermi velocity (or density) in the leads and dot, is
$$G(\mathrm{\Phi })\frac{n_1}{2\pi B}\left[2\xi +\frac{4}{\pi }\xi ^2+\left(\frac{8}{\pi ^2}\frac{1}{3}\right)\xi ^3\right]$$
(15)
As the maximum in $`G`$ occurs for $`\xi `$ slightly smaller than unity, checking whether $`G`$ varies linearly with the magnetic field below the maximum requires only to calculate the ratio between the coefficients of the second and third power of $`\xi `$ in Eq.(15). This ratio is 2.82, indicating that the linear term dominates in agreement with our numerical results. This equation also shows that the slope of the straight line is proportional to $`W^2`$. To make a quantitative comparison with the result of Fig. 5 we rewrite Eq. (15) introducing the actual expression for $`\xi `$; the result is,
$$G(\mathrm{\Phi })\frac{\mathrm{}v}{4\pi }W+\frac{W^2}{\pi L^2}\mathrm{\Phi }+\frac{\pi }{2\mathrm{}v}\left(\frac{8}{\pi ^2}\frac{1}{3}\right)\frac{W^3}{L^2}\mathrm{\Phi }^2$$
(16)
The coefficient of the linear term results to be $`1.44\times 10^4W^2`$ not too far from the numerical result of Fig. 5.
Before ending it is worth to comment on several differences between our model calculation and the theory of Ref. . We first note that we have assumed that the magnetic field is zero outside the dot, which may at first sight invalidate the use of Eq. (12). Nevertheless, the Green function of the whole system is calculated through Dyson’s equation (see Eq. (8). This means that the region of the leads close to the dot is distorted by the magnetic field, which seems to be enough to validate the calculation of the number of occupied subbands by means of Eq. (12). On the other hand whereas in the theory of Ref. ) each channel contributes with one quanta to the conductance, in our case this contribution is approximately halved (remember that we work on chaotic cavities, see Ref. ). This seems, however, to be irrelevant as far as the qualitative behavior of the conductance is concerned.
## V Concluding Remarks
Summarizing, we have presented a numerical analysis of the magnetoconductance of quantum chaotic cavities in a wide range of magnetic fields. For sufficiently open cavities the magnetoconductance increases steadily reaching a maximum at a flux proportional to the leads width. This steady increase of $`G`$ agrees with the experimental observations reported in . Neither regular nor bulk disordered cavities behave in this way. Numerical results for the average magnetoconductance indicate that, for magnetic fluxes larger than that for which the cyclotron radius is approximately $`L/2`$ and smaller than the flux at which the mentioned maximum is reached, it increases linearly with the magnetic flux $`\mathrm{\Phi }`$ with a slope proportional to the square of the leads width. At higher fluxes the conductance decreases stepwise. These results admit a satisfactory explanation in terms of the theory proposed by Beenhakker and van Houten to interpret the experimental results for the magnetoconductance of two contacts in series. The fact that our results for small magnetic fluxes ($`r_c>L/2`$) show a better agreement with the theory in the case that the two contacts are attached to opposite corners of the dot (and, thus, are not facing each other), is related to the stronger equilibration of edge–states promoted by this lead configuration with respect to the other two geometries explored in this work.
###### Acknowledgements.
We thank C.W.J. Beenakker and D. Khmelnitskii for very useful correspondence, and L. Brey, C.Tejedor and J. Palacios for some interesting comments and remarks. This work was supported in part by the Spanish CICYT (grants PB96-0085 and 1FD97-1358).
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# Absolute non-Archimedean polyhedral expansions of ultrauniform spaces.
## 1 Introduction.
To expansions of topological and uniform spaces as inverse limits several works were devoted . Irreducible polyhedral expansions over the field $`𝐑`$ were very important in . In the case of locally compact groups expansions into inverse limits of Lie groups were investigated in . Such inverse mapping systems expansions may be used for the proof of isomorphisms of definite locally convex spaces over the fields $`𝐑`$ or $`𝐂`$ . In the case of ultrauniform hereditarily disconnected spaces polyhedral expansions over $`𝐑`$ do not take into account many features of such spaces.
This work is devoted to the investigation of the problem about inverse mapping systems expansions of ultrauniform spaces using polyhedra over non-Archimedean locally compact fields. In §2 preliminary results, definitions and notations are given. In §3 Theorems about expansions of complete ultrametric and ultrauniform spaces are proved. In §4 absolute polyhedral expansions and inverse mapping systems of expansions for non-complete spaces are investigated. There are elucidated cases, when in inverse mapping systems of expansions polyhedra are finite-dimensional over non-Archimedean fields $`𝐋`$. In particular locally compact ultrauniform groups are considered. The main results are Theorems 3.18 and in §4.
The considered here topological spaces $`X`$ are hereditarily disconnected and zero-dimensional, but in general they are not necessarily strongly zero-dimensional, that is this class has representatives with covering dimensions from $`0`$ to $`\mathrm{}`$ (see Chapters 6 and 7 in ). This article also contains results about a relation of $`dim(X)`$ and dimensions of polyhedra over $`𝐋`$. With the help of polyhedral expansions it is shown in §5 that there exists a natural correspondence between ultrauniform $`X`$ and uniform spaces $`Y`$, such that there exists a continuous quotient mapping $`\theta :XY`$. That is for each $`X`$ there exist $`\theta `$ and $`Y`$, vice versa for each $`Y`$ there exist $`X`$ and $`\theta `$ such that $`d(X)=d(Y)`$ and $`w(X)=w(Y)`$. Then $`X`$ is dense in itself if and only if $`Y`$ is such. Moreover, there are natural compactifications $`cX`$ and $`c_\theta Y`$ such that $`\theta `$ has an extension on $`cX`$ and $`\theta (cX)=c_\theta Y`$. On the other hand, it gives a relation between dimensions $`dim_𝐋P`$ of polyhedras $`P`$ over $`𝐋`$ in an expansion and $`dim[\theta (P)]`$, also between $`\mathrm{max}_{P_m}dim_𝐋P_m`$ and $`dim[\theta (X)]`$. Then an important particular case of polyhedral expansions of manifolds modelled on non-Archimedean locally convex spaces is investigated. Few results from §3 and §4 were announced in .
## 2 Notations and preliminary results.
2.1. Let us recall that by an ultrametric space $`(X,\rho )`$ is implied a set $`X`$ with a metric $`\rho `$ such that it satisfies the ultrametric inequality: $`\rho (x,y)\mathrm{max}(\rho (x,z);\rho (z,y))`$ for each $`x,`$ $`y`$ and $`zX`$. A uniform space $`X`$ with ultrauniformity $`U`$ is called an ultrauniform space $`(X,U)`$ such that $`U`$ satisfies the following condition: $`|xz|<V^{}`$, if $`|yz|<V^{}`$ and $`|xy|<V`$, where $`VV^{}U`$, $`x,`$ $`y`$ and $`zX`$ . If in $`X`$ a family of pseudoultrametrics $`𝐏`$ is given and it satisfies conditions $`(UP1,UP2)`$, then it induces an ultrauniformity $`U`$ due to Proposition 8.1.18 .
Let $`𝐋`$ be a non-Archimedean field . We say that $`X`$ is a $`𝐋`$-Tychonoff space, if $`X`$ is a $`T_1`$-space and for each $`F=\overline{F}X`$ with $`xF`$ there exists a continuous function $`f:XB(𝐋,0,1)`$ such that $`f(x)=0`$, $`f(F)=\{1\}`$, where $`B(X,y,r):=\{zX:\rho (y,z)r\}`$ for $`yX`$ and $`r0`$. From $`ind(𝐋)=0`$ it follows that the small inductive dimension is $`ind(X)=0`$ (see §6.2 and ch. 7 in ). Since the norm $`||_𝐋:𝐋\stackrel{~}{\mathrm{\Gamma }}_𝐋`$ is continuous, then $`X`$ is the Tychonoff space, where $`\stackrel{~}{\mathrm{\Gamma }}_𝐋:=\{|x|_𝐋:x𝐋\}[0,\mathrm{})`$. Vice versa if $`X`$ is a Tychonoff space with $`ind(X)=0`$, then it is also $`𝐋`$-Tychonoff, since there exists a clopen (closed and open at the same time) neighbourhood $`Wx`$ with $`WF=\mathrm{}`$ and as the locally constant function $`f`$ may be taken with $`f(x)=0`$ and $`f(XW)=\{1\}`$.
Let us consider spaces $`C(X,𝐋):=\{f:X𝐋|f\text{ is continuous }\}`$ and $`C^{}(X,𝐋):=\{fC(X,𝐋):|f(X)|_𝐋\text{ is bounded in }𝐑\}`$, then for each finite family $`\{f_1,\mathrm{},f_m\}C(X,𝐋)`$ (or $`C^{}(X,𝐋)`$) the following pseudoultrametric is defined: $`\rho _{f_1,\mathrm{},f_m}(x,y):=\mathrm{max}(|f_j(x)f_j(y)|_𝐋:j=1,\mathrm{},m\}`$. Families $`𝐏`$ or $`𝐏^{}`$ of such $`\rho _{f_1,\mathrm{},f_m}`$ induce ultrauniformities $`𝖢`$ or $`𝖢^{}`$ respectively and the initial topology on $`X`$. If a sequence $`\{V_j:j=0,1,\mathrm{}\}U`$ is such that $`V_0=X^2`$, $`pV_{j+1}V_j`$ for $`j=1,2,\mathrm{},`$ where $`p`$ is a prime number, then there exists a pseudoultrametric $`\rho `$ such that $`\rho (x,y):=0`$ for $`(x,y)_{j=0}^{\mathrm{}}V_j`$, $`\rho (x,y)=p^j`$ for $`(x,y)V_jV_{j+1}`$, so $`V_i\{(x,y):\rho (x,y)p^i\}V_{i1}`$. Indeed, from $`(x,y)V_iV_{i+1}`$ and $`(y,z)V_jV_{j+1}`$ for $`ji`$ it follows $`(x,z)V_i`$ and $`\rho (x,z)p^i=\rho (x,y)`$. Therefore, ultrauniform spaces may be equivalently characterised by $`U`$ or $`𝐏`$ (see §8.1.11 and §8.1.14 in ).
Henceforth, locally compact non-discrete non-Archimedean infinite fields $`𝐋`$ are considered. If the characteristic $`char(𝐋)=0`$, then due to for each such $`𝐋`$ there exists a prime number $`p`$ such that $`𝐋`$ is a finite algebraic extension of the field $`𝐐_𝐩`$ of $`p`$-adic numbers. If $`char(𝐋)=p>0`$, then $`𝐋`$ is isomorphic with the field $`𝐅_𝐩(\theta )`$ of the formal power series by the indeterminate $`\theta `$, where $`p`$ is the prime number, each $`z𝐋`$ has the form $`z=_{j=k(z)}^+\mathrm{}a_j\theta ^j`$ with $`a_j𝐅_𝐩`$ for each $`j`$, $`|z|_𝐋=p^{k(z)}`$ up to the equivalence of the non-Archimedean valuation, $`k(z)𝐙`$, $`𝐅_𝐩`$ is the finite field consisting of $`p`$ elements, $`𝐅_𝐩𝐋`$ is the natural embedding.
For an ordinal $`A`$ with its cardinality $`m=card(A)`$ by $`c_0(𝐋,A)`$ it is denoted the following Banach space with vectors $`x=(x_a:aA,x_a𝐋)`$ of a finite norm $`x:=sup_{aA}|x_a|_𝐋`$ and such that for each $`b>0`$ a set $`\{aA:|x_a|_𝐋b\}`$ is finite. It has the orthonormal in the non-Archimedean sense basis $`\{e_j:=(\delta _{j,a}:aA):jA\}`$, where $`\delta _{j,a}=1`$ for $`j=a`$ and $`\delta _{j,a}=0`$ for each $`ja`$ . If $`card(A_1)=card(A_2)`$ then $`c_0(𝐋,A_1)`$ is isomorphic with $`c_0(𝐋,A_2)`$. In particular $`c_0(𝐋,n)=𝐋^𝐧`$ for $`n𝐍`$. Then $`card(A)`$ is called the dimension of $`c_0(𝐋,A)`$ and it is denoted by $`card(A)=dim_𝐋c_0(𝐋,A)`$.
Let $`𝖴`$ be a uniform covering of $`(X,𝐏)`$ (see §8.1 in ), then the collection $`\{St(U,𝖴):U𝖴\}=:\text{ }^{}𝖴`$ is called the star of $`𝖴`$. If $`\text{ }^{}𝖴`$ is a refinement of a uniform covering $`𝖵`$, then $`𝖴`$ is called the star refinement of $`𝖵`$.
2.2. Lemma. Let $`(X,\rho )`$ be an ultrametric space, then there exists an ultrametric $`\rho ^{}`$ equivalent to $`\rho `$ such that $`\rho ^{}(X,X)\stackrel{~}{\mathrm{\Gamma }}_𝐋`$.
Proof. Let $`\rho ^{}(x,y):=sup_{b\stackrel{~}{\mathrm{\Gamma }}_𝐋,b\rho (x,y)}b`$, where $`x`$ and $`yX`$, either $`𝐋𝐐_𝐩`$ or $`𝐋=𝐅_𝐩(\theta )`$. Then $`\rho ^{}`$ is the ultrametric such that $`\rho ^{}(x,y)\rho (x,y)p\times \rho ^{}(x,y)`$ for each $`x`$ and $`yX`$.
2.3. Lemma. Let $`(X,𝐏)`$ be an ultrauniform space, then there exists a family $`𝐏^{}`$ such that $`\rho ^{}(X,X)\stackrel{~}{\mathrm{\Gamma }}_𝐋`$ for each $`\rho ^{}𝐏^{}`$; $`(X,𝐏)`$ and $`(X,𝐏^{})`$ are uniformly isomorphic, the completeness of one of them is equivalent to the comleteness of another.
Proof. In view of Lemma 2.2 for each $`\rho 𝐏`$ there exists an equivalent pseudoultrametric $`\rho ^{}`$. They form a family $`𝐏^{}`$. Evidently, the identity mapping $`id:(X,𝐏)(X,𝐏^{})`$ is the uniform isomorphism. The last statement follows from 8.3.20 .
2.4. Theorem. For each ultrametric space $`(X,\rho )`$ there exist an embedding $`f:XB(c_0(𝐋,A_X),0,1)`$ and an uniformly continuous embedding into $`c_0(𝐋,A_X)`$, where $`card(A_X)=w(X)`$, $`w(X)`$ is the topological weight of $`X`$.
Proof. In view of Theorem 7.3.15 there exists an embedding of $`X`$ into the Baire space $`B(m)`$, where $`m=w(X)\mathrm{}_0`$. In the case $`w(X)<\mathrm{}_0`$ this statement is evident, since $`X`$ is finite. In view of Lemma 2.2 we choose in $`B(m)`$ an ultrametric $`\rho `$ equivalent to the initial one with values in $`\stackrel{~}{\mathrm{\Gamma }}_𝐋`$ such that $`\rho (\{x_i\},\{y_i\})=p^k`$, if $`x_ky_k`$ and $`x_i=y_i`$ for $`i<k`$, $`\rho (\{x_i\},\{y_i\})=0`$, if $`x_i=y_i`$ for all $`i`$, where $`\{x_i\}B(m)`$, $`i𝐍`$, $`x_iD(m)`$, $`D(m)`$ denotes the discrete space of cardinality $`m`$. Let $`A=𝐍\times C`$, $`card(C)=m`$, $`\{e_{i,a}|(i,a)A\}`$ be the orthonormal basis in $`c_0(𝐋,A)`$. For each $`\{x_i\}B(m)`$ we have $`x_iD(m)`$ and we can take $`x_i=e_{i,a}`$ for suitable $`a=a(i)`$, since $`D(m)`$ is isomorphic with $`\{e_{i,a}:aC\}`$. Let $`f(\{x_i\}):=_{i𝐍,a=a(i)}p^ie_{i,a}`$, consequently, $`f(\{x_i\})f(\{y_i\})_{c_0(𝐋,A)}=\rho (\{x_i\},\{y_i\})`$.
The last statement of the Theorem follows from the isometrical embedding of $`(X,\rho )`$ into the corresponding free Banach space, which is isomorphic with $`c_0(𝐋,A)`$ (see Theorem 5 in and Theorems 5.13 and 5.16 in ), since each Banach space over a non-Archimedean spherically complete field $`𝐋`$ is isomorphic with $`c_0(𝐋,A)`$ with the corresponding $`A`$. Instead of the free Banach space generated by $`(X,\rho ^{})`$ (see also Lemma 2.2), it can be used an embedding into the topologically dual space $`\stackrel{~}{C}_b^0(X,𝐋)^{}`$ to a space $`\stackrel{~}{C}_b^0(X,𝐋)`$, which is the space of bounded continuous Lipschitz functions $`f:X𝐋`$ with $`f:=\mathrm{max}\{sup_{xX}|f(x)|,sup_{xyX}(|f(x)f(y)|/\rho (x,y))\}<\mathrm{}.`$ For each $`SX`$ there exists a function $`f\stackrel{~}{C}_b^0(X,𝐋)`$ such that $`\rho ^{}(x,S)/[1+\rho ^{}(x,S)]|f(x)|p\rho ^{}(x,S)/[1+\rho ^{}(x,S)]`$, where $`\rho ^{}(x,S):=inf_{yS}\rho ^{}(x,y)`$. To each $`xX`$ there corresponds a continuous functional $`G_x`$ such that $`G_x(f)=f(x)`$, so $`G_x(f)G_y(f)=f(x)f(y)`$. In view of the $`𝐋`$-regularity of $`X`$ we have $`G_xG_y=\rho (x,y)`$.
2.5. Corollary. Each ultrauniform space $`(X,𝐏)`$ has a topological embedding into $`_{\rho 𝐏}B(c_0(𝐋,A_\rho ),0,1)`$ and an uniform embedding into $`_{\rho 𝐏}c_0(𝐋,A_\rho )`$ with $`card(A_\rho )=w(X_\rho )`$.
Proof. For each ultrametric $`\rho 𝐏`$ of an ultrauniform space $`(X,𝐏)`$ there exists the equivalence relation $`R_\rho `$ such that $`xR_\rho y`$ if and only if $`\rho (x,y)=0`$. Then there exists the quotient mapping $`g_\rho :XX_\rho :=(X/R_\rho )^{}`$, where $`X_\rho `$ is the ultrametric space with the ultrametric also denoted by $`\rho `$, $`\stackrel{~}{X}`$ denotes the completion of $`X`$. Then $`X`$ has the uniform embedding into the limit of the inverse spectra $`lim_\rho \stackrel{~}{X}_\rho =\stackrel{~}{X}`$.
2.6. Corollary. For each $`j𝐍`$ and $`f(X)c_0(𝐋,A_X)`$ from Theorem 2.4 there are coverings $`U_j`$ of the space $`f(X)`$ by disjoint clopen balls $`B_l`$ with diameters not greater than $`p^j`$ and with $`inf_{lk}dist(B_l,B_k)>0`$.
Proof follows from the consideration of the covering of the Banach space $`c_0(𝐋,A)`$ by balls $`B(c_0(𝐋,A),x,r)`$ with $`0<rp^j`$ and $`xc_0`$, since such balls are either disjoint or one of them is contained in another and $`\mathrm{\Gamma }_𝐋`$ is discrete in $`(0,\mathrm{})`$, where $`\mathrm{\Gamma }_𝐋:=\stackrel{~}{\mathrm{\Gamma }}_𝐋\{0\}`$. Then $`_{qJ}B(c_0,x_q,r_q)`$ is the clopen ball in $`c_0`$ with $`rp^j`$, if all balls in the family $`J`$ have non-void pairwise intersections. Taking $`B(c_0,x,r)f(X)`$ we get the statement for $`f(X)`$ using the transfinite sequence of the covering. To satisfy Condition $`inf_{lk}dist(B_l,B_k)>0`$ we take $`r_l>r_0>0`$ for each ball $`B_l`$, where $`r_0`$ is a fixed positive number.
2.7. Note. A simplex $`s`$ in $`𝐑^𝐧`$ may be taken with the help of linear functionals, for example, $`\{e_j:j=0,\mathrm{},n\}`$, where $`e_j=(0,\mathrm{},0,1,0,\mathrm{},0)`$ with $`1`$ in the $`j`$-th place for $`j>0`$ and $`e_0=e_1+\mathrm{}+e_n`$, $`s:=\{x𝐑^𝐧:e_j(x)[0,1]\text{ for }j=0,1,\mathrm{},n\}`$. In the case of $`𝐋^𝐧`$ or $`c_0(𝐋,A)`$, if to take $`B(𝐋,0,1)`$ instead of $`[0,1]`$, then conditions $`x_j:=e_j(x)B(𝐋,0,1)`$ for $`j=1,\mathrm{},n`$ imply $`e_0(x)=x_1+\mathrm{}+x_nB(𝐋,0,1)`$ or $`_{j\lambda }e_j(x)B(𝐋,0,1)`$ for each finite subset $`\lambda `$ in $`A`$ respectively due to the ultrametric inequality (since $`B(𝐋,0,1,)`$ is the additive group ), that is $`s=B(𝐋^𝐧,0,1)`$ or $`s=B(c_0(𝐋,A),0,1)`$ correspondingly. Moreover, its topological border is empty $`Fr(s)=\mathrm{}`$ and $`Ind(Fr(s))=1`$. Let us denote by $`\pi _𝐋`$ an element from $`𝐋`$ such that $`B(𝐋,0,1^{}):=\{x𝐋:|x|_𝐋<1\}=\pi _𝐋B(𝐋,0,1)`$ and $`|\pi _𝐋|_𝐋=sup_{b\mathrm{\Gamma }_𝐋,b<1}b=:b_𝐋`$.
2.8. Definitions. (1). A subset $`P`$ in $`c_0(𝐋,A)`$ is called a polyhedron if it is a disjoint union of simplexes $`s_j`$, $`P=_{jF}s_j`$, where $`F`$ is a set, $`s_j=B(c_0(𝐋,A^{}),x,r)=x+\pi _𝐋^kB(c_0(𝐋,A^{}),0,1)`$ are the clopen balls in $`c_0(𝐋,A^{})`$, $`A^{}A`$, $`r=b_𝐋^k`$, $`k𝐙`$. For each $`𝐋`$ we fix $`\pi _𝐋`$ and such affine transformations. The polyhedron $`P`$ is called uniform if it satisfies conditions $`(i,ii)`$:
$`(i)`$ $`sup_{iF}diam(s_i)<\mathrm{}`$,
$`(ii)`$ $`inf_{ij}dist(s_i,s_j)>0`$, where $`dist(s,q):=inf_{xs,yq}\rho (x,y)`$.
By vertexes of the simplex $`s=B(c_0(𝐋,A),0,1)`$ we call points $`x=(x_j)c_0(𝐋,A)`$ such that $`x_j=0`$ or $`x_j=1`$ for each $`jA`$, $`dim_𝐋(s):=card(A)`$. For each $`EA`$, $`EA`$ and a vertex $`e`$ by verge $`v`$ of the simplex $`s`$ we call a subset $`e+B(c_0(𝐋,E),0,1)s`$. A mapping $`f:CD`$ is called affine, if $`f(ax+by)=af(x)+bf(y)`$ for each $`a,bB(𝐋,0,1)`$, $`x,yC`$, where $`C`$ and $`D`$ are absolutely convex subsets of vector spaces $`X`$ and $`Y`$ over $`𝐋`$. For an arbitrary simplex its verges and vertexes are defined with the help of affine transformation as images of verges and vertexes of the unit simplex $`B(c_0(𝐋,A^{}),0,1)`$.
Then in analogy with the classical case there are naturally defined notions of a simplexial complex $`K`$ and his space $`|K|`$, also a subcomplex and a simplexial mapping. The latter is characterized by the following conditions: its rectrictions on each simplex of polyhedra that are affine mappings over the field $`𝐋`$ and images of vertexes are vertexes. A simplexial complex $`K`$ is a collection of verges $`v`$ of a simplex $`s`$ such that each verge $`v^{}`$ of a simplex $`v`$ is in $`K`$. The space $`|K|`$ of $`K`$ is a subset of $`|s|`$ consisting of those points which belong to simplexes of $`K`$, where $`|s|`$ is a topological space of $`s`$ in the topology inherited from $`c_0(𝐋,A)`$. This defines topology of $`|K|`$ also. A simplexial complex has $`dim_𝐋K:=sup_{vK}dim_𝐋v`$, where $`v`$ are simplexes of $`K`$. A subcomplex $`L`$ of $`K`$ is a subcollection of the simplexes of $`K`$ such that each verge of $`v`$ of a simplex $`s`$ in $`L`$ is also in $`L`$, so $`L`$ is a simplexial complex. If $`e`$ is a vertex of $`K`$, the open star of $`e`$ is the subset $`St_K(e)`$ of $`K`$ which is the union of all simplexes of $`K`$ having $`e`$ as a vertex.
Instead of the barycentric subdivision in the classical case we introduce a $`p^j`$-subdivision of simplexes and polyhedra for $`j𝐍`$ and for the field $`𝐋`$, that is a partition of each simplex $`B(c_0(𝐋,A^{}),x,r)`$ into the disjoint union of simplexes with diameters equal to $`rp^j`$, where $`p=b_𝐋^1`$. Each simplex $`s`$ with $`dim_𝐋s=card(A^{})`$ may be considered also in $`c_0(𝐋,A)`$, where $`A^{}A`$, since there exists the isometrical embedding $`c_0(𝐋,A^{})c_0(𝐋,A)`$ and the projector $`\pi :c_0(𝐋,A)c_0(𝐋,A^{})`$. By a dimension of a polyhedron $`P`$ we call $`dim_𝐋P:=sup_{(sP,s\text{ is a simplex })}dim_𝐋s`$. The polyhedron $`P`$ is called locally finite-dimensional if all simplexes $`sP`$ are finite dimensional over $`𝐋`$, that is, $`dim_𝐋s𝐍`$. For a simplex $`s=B(c_0(𝐋,A^{}),x,r)`$ by a $`𝐋`$-border $`s`$ we call the union of all its verges $`q`$ with the codimension over $`𝐋`$ equal to $`1`$ in $`c_0(𝐋,A^{})=:X`$, that is, $`q=e+B^{}`$, where $`B^{}`$ are balls in $`c_0(𝐋,A\mathrm{"})=:Y`$, $`(A^{}A\mathrm{"})`$ is a singleton, $`A\mathrm{"}A^{}`$, $`XY=𝐋`$, $`YX`$. For the polyhedron $`P=_{jF}s_j`$ by the $`𝐋`$-border we call $`P=_{jF}s_j`$, where $`F`$ is the set.
(2). A continuous mapping $`f`$ of a set $`M`$, $`Mc_0(𝐋,A)`$ into a polyhedron $`P`$ we call essential, if there is not any continuous mapping $`g:MP`$ for which are satisfied the following conditions:
$`(i)`$ $`g(M)`$ does not contain $`P`$;
$`(ii)`$ there exists $`M_0M`$, $`M_0M`$ with $`f(M)P=f(M_0)=g(M_0)P`$ and their restrictions coincide $`f|_{M_0}=g|_{M_0}`$;
$`(iii)`$ if $`f`$ is linear on $`[x,y]:=\{tx+(1t)y|tB(𝐋,0,1)\}M`$, then $`g`$ is also linear on $`[x,y]`$ such that $`g(x)g(y)`$ for each $`f(x)f(y)`$.
(3). The function $`f`$ from §2.8.(2) is inessential, if there exists such $`g`$.
(4). Let $`f:MN`$ be a continuous mapping, $`c_0(𝐋,A)NP`$, $`P`$ be a polyhedron. Then $`P`$ is called essentially (or inessentially) covered by $`N`$ under the mapping $`f`$, if $`f|_{f^1(P)}`$ is essential (or inessential respectively).
(5). Let $`f:MP`$ and $`g:MP`$ are continuous mappings, where $`M`$ is a set, $`P`$ is a polyhedron. Then $`g`$ is called a permissible modification of $`f`$, if three conditions are satisfied:
$`(i)`$ from $`aM`$ and $`f(a)s`$ it follows $`g(a)s`$, where $`s`$ is a simplex from $`P`$;
$`(ii)`$ if $`x`$ and $`yM`$, $`[x,y]M`$ and $`f:[x,y]P`$ is linear, then $`g:[x,y]P`$ is also linear (over $`𝐋`$) and $`g(x)g(y)`$ for each $`f(x)f(y)`$;
$`(iii)`$ $`f(M)=g(M)`$, when $`M`$ contains a polyhedron $`Q`$ with $`dim_𝐋Q>0`$ such that $`Q`$ is clopen in $`M`$.
(6). The mapping $`f:MP`$ is called reducible (or irreducible), when it may (or not respectively) have the permissible modification $`g`$ such that $`f(M)`$ is not the subset of $`g(M)`$.
(7). A mapping $`f:PQ`$ for polyhedra $`P`$ and $`Q`$ is called normal, if
$`(i)`$ $`\rho _Q(f(x),f(y))\rho _P(x,y)`$ for each $`x`$ and $`yP`$, that is $`f`$ is a non-stretching mapping;
$`(ii)`$ there exists a $`p^j`$-subdivision $`Q^{}`$ of polyhedra $`Q`$ such that $`f:PQ^{}`$ is a simplexial mapping, that is, $`f|_s`$ is affine on each simplex $`sP`$ and $`f(s)`$ is a simplex from $`Q^{}`$.
(8). Let $`X=lim_j\{X_j,f_i^j,E\}`$ be an expansion of $`X`$ into the limit of inverse spectra of polyhedra $`X_j`$ over $`𝐋`$. This expansion is called
$`(a)`$ irreducible, if for each open $`VX`$ there exists a cofinal subset $`E_VE`$ such that $`\{X_j,f_i^j,E_V\}`$ is also the irreducible polyhedral representation of the space $`V`$, that is $`f_i^j:X_jX_i`$ are irreducible and surjective for each $`ijE_V`$. The polyhedral system (representation) $`\{X_j,f_i^j,E\}`$ is called
$`(b)`$ $`n`$-dimensional, if $`dim_𝐋X_jn`$ for each $`jE`$, $`sup_{jE}dim_𝐋X_j=n`$, where $`n`$ is a cardinal number. The polyhedral spectrum is called
$`(c)`$ non-degenerate, if projections $`f_i^j`$ are non-degenerate, that is $`f_i^j(s)`$ is not a singleton for each simplex $`s`$ of $`X_j`$.
2.9. Notes. Conditions $`2.8.(2(iii),(5(ii,iii))`$ and restrictions on $`f_i^j`$ in $`2.8.(8(a))`$ are additional in comparison with the classical case. They are imposed in §2.8, since there exists a continuous non-linear retraction $`r:s_js_j`$ for the non-Archimedean field $`𝐋`$, which may be constructed with the help of a $`p`$-subdivision and projections of $`c_0(𝐋,A)`$ are on its subspaces $`c_0(𝐋,A^{})`$ associated with the standard basis $`\{e_j:`$ $`jA\}`$, where $`A^{}A`$. If $`f`$ is simplexial on each polyhedron $`M`$ and $`2.8.(2(iii))`$ is accomplished, then $`dim_𝐋g(M)=dim_𝐋f(M)`$.
In were studied uniform spaces and ANRU. Here we mean by ANRU an ultraunifrom space $`X`$ such that under embedding into an ultrauniform space $`Y`$ there exists the uniformly continuous retract $`r:VX`$ of its uniform neighbourhood $`V`$, $`XVY`$. We denote by $`U(X,Y)`$ for two ultrauniform spaces $`X`$ and $`Y`$ an ultrauniform space of uniformly continuous mappings $`f:XY`$ with the uniformity generated by a base of the form $`W=\{(f,g)|(f(x),g(x))V`$ for each $`xX\}`$, where $`V𝖵`$, $`𝖵`$ is a uniformity on $`Y`$ corresponding to $`𝐏_Y`$.
Here we call an ultrauniform space $`Y`$ injective, if for each ultrauniform space $`X`$ with $`HX`$ and an uniformly continuous mapping $`f:HY`$ there exists an uniformly continuous extension $`f:XY`$.
2.10. Theorem. $`B(c_0(𝐋,A),0,1)`$ is an injective space for $`card(A)<\mathrm{}_0`$.
Proof may be done analogously to Theorem 9 on p. 40 in .
2.11. Theorem. Each uniform polyhedron $`P`$ over $`𝐋`$ is ANRU.
Proof. Since $`a=inf_{ijF}dist(s_i,s_j)>0`$, $`b=sup_{iF}diam(s_i)<\mathrm{}`$, then there exists a minimal uniform covering $`𝖴`$ such that for each $`s_i`$ there exists $`a/p`$-clopen neighbourhood $`V_i𝖴`$, where $`V_i=s_i^ϵ`$, $`ϵ=a/p`$, $`A^ϵ=\{yY:`$ $`d_Y(y,A)=inf_{aA}d_Y(y,a)ϵ\}`$ for an ultrametric space $`(Y,d_Y)`$ is an $`ϵ`$-enlargement of a subset $`A`$ in $`Y`$. In this case $`P`$ is a subset of a Banach space $`c_0(𝐋,A)`$, where $`card(A)=w(P)`$. Each $`s_i`$ is an uniform retract $`V_i`$, $`r_i:V_is_i`$, consequently, there exists uniformly continuous retraction $`r:S=_{V_i𝖴}V_iP`$ such that $`r|_{s_i}=r_i`$ for each $`i`$, since $`sup_{iF}diam(s_i)<\mathrm{}`$ and $`r_ir_i=r_i`$, $`r_i|_{s_i}=id`$, hence $`rr=r`$ and $`r|_P=id`$. The neighbourhood $`S`$ is uniform, since $`St(P,𝖴):=_{(A𝖴,AP\mathrm{})}A=S`$ .
2.12. Note. Further for uniform spaces are considered uniformly continuous mappings and uniform polyhedra and for topological spaces continuous mappings and polyhedra if it is not specially outlined.
2.13. Lemma. Let $`(X,𝐏)`$ be an ultrauniform strongly zero-dimensional space, $`P`$ be a polyhedron over $`𝐋`$, $`A_1`$,…,$`A_q`$ are non-intersecting closed subsets in $`X`$, $`q𝐍`$, $`card(P)q`$. Then there exists an uniformly continuous mapping $`f:XP`$ such that $`f(A_i)f(A_j)=\mathrm{}`$ for each $`ij`$.
Proof. There exists a disjoint clopen covering $`V_j`$ for $`X`$ satisfying $`A_jV_j`$ for each $`j=1,\mathrm{},q`$ (see Theorem 6.2.4 ). Then we can take locally constant or locally affine mapping $`f`$ with $`f(V_j)s_jP`$.
2.14. Lemma. Each non-stretching mapping $`f:EP`$ has non-stretching continuation on $`(\stackrel{~}{X},\rho )`$, where $`P`$ is a uniform polyhedron over $`𝐋`$, $`EX`$, $`\stackrel{~}{X}`$ is a completion of an ultrametric space $`(X,\rho )`$.
Proof. There exists an embedding $`Pc_0(𝐋,A)`$ for $`card(A)=w(P)`$ due to Lemma 4 or Theorem 2.4 above. We choose $`f:Ec_0(𝐋,A)`$ with $`f=(f^i:iA)`$, $`f^i:E𝐋e_i`$, where $`(e_i:iA)`$ is the orthonormal basis in $`c_0(𝐋,A)`$ and $`inf_{(ij,s_i\text{ and }s_jP)}dist(s_i,s_j)>0`$, $`s_i`$ are simplexes from $`P`$. Each mapping $`f^i:E𝐋e_i`$ is non-stretching, that is $`f^i(x)f^i(y)\rho (x,y)`$ for each $`x,yE`$. Evidently, $`f`$ has non-stretching extension on $`\stackrel{~}{E}`$ such that $`\stackrel{~}{E}`$ is closed $`\stackrel{~}{X}`$. On the other hand, due to Theorem 2.4 there exists $`ϵ>0`$ such that $`f`$ has a non-stretching extension $`f:E^ϵP^ϵ`$. By Theorem 2.11 for sufficiently small $`ϵ>0`$ there exists a retraction $`r:P^ϵP`$ such that $`rf:E^ϵP`$ is is non-stretching, $`rf=f`$ on $`\stackrel{~}{E}`$. There exists $`V`$ clopen in $`\stackrel{~}{X}`$ such that $`EVE^ϵ`$ due to Lemma 2.2, hence $`rf|_V`$ has a non-stretching extension of $`f`$ from $`E`$, since $`(\stackrel{~}{X},\rho )`$ is strongly zero-dimensional. This follows also from the fact that a free Banach space $`B(E,\rho ^{},𝐋)`$ over $`𝐋`$ generated by $`E`$ has the isometric embedding into $`B(X,\rho ^{},𝐋)`$, where $`\rho ^{}`$ is from Lemma 2.2 (see also ).
2.15. Definitions. An ultrauniform space $`(X,𝐏)`$ is called $`LE`$-space, if each uniformly continuous $`f:Y𝐋`$ has an uniformly continuous extension on all $`X`$, where $`YX`$.
Let $`X`$ be a Tychonoff space. If $`X`$ has the finest uniformity compatible with its topology, then it is called fine.
2.16. Theorem. An ultrametric space $`X`$ is an $`LE`$-space if and only if $`\stackrel{~}{X}=lim\{X_m,f_n^m,E\}`$, where $`X_m`$ are fine spaces and bonding mappings $`f_n^m:X_mX_n`$ are uniformly continuous for each $`mnE`$, $`\{X_m,f_n^m,E\}`$ is an inverse mapping system.
Proof. Let us consider $`X_A=B(c_0(𝐋,A),0,1)`$ for $`card(A)\mathrm{}_0`$ or $`X_A=𝐋^𝐧`$ which are not fine spaces, since on $`X_A`$ there are continuous $`f:X_A𝐋`$ which are not uniformly continuous. If $`X_A`$ is fine, then each continuous function $`f:X_A𝐋`$ is uniformly continuous. Then for the embedding of the non-fine space $`X_Ac_0(𝐋,A)`$ there is not compact $`H`$ with $`X_AHc_0(𝐋,A)`$. Therefore, in $`X_A`$ there exists a countable closed subset of isolated points $`Y=\{x_i:i𝐍\}`$ and a continuous function $`f:Y𝐋`$ with $`|f(x_j)f(x_i)|/|x_jx_i|>c_i>0`$ for each $`ji`$ and $`lim_i\mathrm{}c_i=\mathrm{}`$ such that $`f`$ has a continuous extension $`g`$ on $`X_A`$. Indeed, there exists a clopen neighbourhood $`W`$, that is a $`b`$-enlargement (with $`b>0`$) relative to the ultrametric $`\rho _A`$ in $`X_A`$, so there is a continuous retraction $`r:WY`$, $`g(x)=const𝐋`$ on $`(X_AW)`$, $`g(x)=f(r(x))`$ on $`W`$, $`g|_Y=f`$. Therefore, $`g`$ is continuous and uniformly continuous on $`X_A`$ (see also Theorem 1.7 ).
## 3 Polyhedral expansions.
3.1. Theorem. Each complete ultrauniform space $`(Y,𝐏)`$ is a limit of an inverse spectra of ANRU $`Y_j`$, where $`Y_j`$ are embedded into complete locally $`𝐋`$-convex spaces.
Proof. In view of Corollary 2.5 there exists a uniform embedding $`Y_{\rho 𝐏}c_0(𝐋,A_\rho )=:H`$. In each $`c_0(𝐋,A_\rho )`$ may be taken the orthonormal basis $`\{e_{j,\rho }:jA_\rho \}`$, $`card(A_\rho )=w(Y_\rho )`$ and define canonical neighbourhoods $`U(f,b;(j_1,\rho _1),\mathrm{},(j_n,\rho _n)):=\{qH:|\pi _{j_i,\rho _i}(q)\pi _{j_i,\rho _i}(f)|<b`$ for $`i=1,\mathrm{},n\}`$, where $`\pi _{j,\rho }:H𝐋e_{j,\rho }`$ are projectors, $`fH`$, $`b>0`$, $`n𝐍`$. Each clopen subset $`Z_b:=HU(f,b;(j_1,\rho _1),\mathrm{},(j_n,\rho _n))`$ is uniformly continuous (non-stretching) retract of $`Z_{b/p}`$, that is $`Z_b`$ is ANRU. Analogously each finite intersections $`_{l=1}^qZ(f_l,b_l,(j_1^l,\rho _1^l),\mathrm{},(j_n^l,\rho _n^l))=_{l=1}^kZ_{l,b_l}`$ are also ANRU, since $`Z_{k,b_k}_{l=1}^kZ_{l,b_l}`$ and a retraction $`g_k:Z_{k,b_k/p}Z_{k,b_k}`$ produces non-stretching relative to the corresponding pseudoultrametric retraction $`g_k:_{l=1}^kZ_{l,b_l/p}_{l=1}^kZ_{l,b_l}`$. Hence $`g_q\mathrm{}g_1=g:_{l=1}^qZ_{l,b_l/p}_{l=1}^qZ_{l,b_l}`$ is the (non-stretching) retraction. Let $`S`$ be the ordered family of such $`_{l=1}^qZ_{l,b_l}HY`$, then $`_{ZS}(HZ)=Y`$. Further as in the proof of Theorem 7.1 .
3.2. Lemma. Let $`(X,\rho _X)`$ and $`(Y,\rho _Y)`$ be ultrametric spaces, $`f:XY`$ be a continuous mapping such that $`f|_H`$ is a $`b`$-mapping (that is $`\rho _Y(f(x),i(x))b`$ for each $`xH`$), where $`H`$ is dense in $`X`$, $`b>0`$, $`i:HY`$ is an embedding. Then $`X`$ and $`Y`$ may be embedded into a Banach space $`W`$ over $`𝐋`$ such that $`X`$ with equivalent ultrametric in $`W`$ and $`f(x)xb`$ for each $`xX`$ (that is $`f`$ is a $`b`$-mapping on $`X`$).
Proof. From Theorem 2.4, Lemma 2.3 and Corollary 2.6 it follows that there exists such embedding of $`X`$ into the corresponding $`W=c_0(𝐋,A)`$ with $`card(A)=w(X)`$ with the disjoint clopen covering $`V=\{B(c_0(𝐋,A),x_j,b_j)=:B_j:x_jH,\mathrm{}>b_j>0,jF\}`$. Let $`Y_j:=f(B_j)Y`$, consequently, $`diam(Y_j)b`$, since $`f`$ is the $`b`$-mapping, that is $`f`$ realizes the covering of $`X`$ consisting from subsets of diameters not greater than $`b`$. Then $`\rho _X(x^{},x\mathrm{"})\mathrm{max}(\rho (x^{},x_j),\rho (x_j,x\mathrm{"}))`$, where $`x\mathrm{"}f^1(y)`$, $`yY_j`$, $`x^{}B(c_0(𝐋,A),x_j,b_j)`$, $`f(x^{})=y`$, $`diam(f^1(y))b`$. Let $`x_ix_j`$, this is equivalent to $`B_iB_j=\mathrm{}`$ and is equivalent to $`\rho _X(x_i,x_j)>b`$, consequently, $`Y_iY_j=\mathrm{}`$ and $`Y=_{jF}Y_j`$. In view of continuity of $`f`$ there exists the embedding of $`Y_j`$ into $`B_j`$, since $`f`$ is the $`b`$-mapping, $`w(Y)w(X)`$, $`0<b_jb`$ and due to Theorem 2.4.
3.3.1. Lemma. Let there exists a non-stretching (uniformly continuous) mapping $`f:RP`$ and a non-stretching (uniformly continuous) permissible modification $`g:MP`$, where $`R`$ is a complete ultrametric space ($`LE`$-space respectively), $`P`$ is a polyhedron, $`M`$ is a subspace in $`R`$; if $`R`$ is a polyhedron let $`M`$ be a subpolyhedron. Then $`g`$ has a non-stretching (uniformly continuous respectively) extension on the entire $`R`$ and this extension is a permissible modification of $`f`$.
Proof. For a complete ultrametric space $`R`$ the mapping $`g`$ has the non-stretching (uniformly continuous) extension on the completion of $`M`$ which coincides with the closure $`\overline{M}`$ of $`M`$ in $`R`$, since $`R`$ and $`P`$ are complete. The space $`P`$ is complete, since it is ANRU by Theorem 2.11 (see also Theorem 1.7 , Theorems 8.3.6 and 8.3.10 ). For the embedding $`Rc_0(𝐋,A)`$ with $`card(A)=w(R)`$ by Theorem 2.4 it may be taken due to Corollary 2.6 the disjoint clopen covering $`V`$ such that each $`WV`$ has the form $`W=RB(c_0(𝐋,A),x,r_W)`$, where $`r_W>0`$, $`xR`$, $`sup_{WV}r_Wp^j`$ (for each $`j𝐍`$ there exists such $`V`$).
In view of uniform continuity of $`f`$ and uniformity of the polyhedron $`P`$ there exists $`V`$ such that for each $`W`$ from $`V`$ there exists a simplex $`TP`$ with $`f(W)T`$. The space $`c_0(𝐋,A)`$ has the orthonormal basis $`\{e_j:jA\}`$. If $`f`$ is linear on no any $`[a,b]R`$, then for the construction of $`g`$ may be used Lemma 2.14 or Theorem 2.16. This may be done with the help of transfinite induction by the cardinality of sets of vertexes of simplexes from $`P`$ (or using the Teichmüller-Tukey Lemma), since $`P`$ is a disjoint union of simplexes. In general let $`g(\overline{M})`$ contains each zero-dimensional over $`𝐋`$ simplex $`T^0P`$. If it is not so, then there exists a point $`b=f^1(T^0)`$ in which $`g`$ is not defined. If $`f`$ is non-linear on each $`[a,b]R`$ with $`a\overline{M}`$, then $`g(b)=f(b)`$. If $`f`$ is linear on such $`[a,b]`$, then from $`\rho (a,b)<\mathrm{}`$ it follows that $`[a,b]`$ is homeomorphic with the compact ball $`B(𝐋,a,\rho (b,a))`$ in $`𝐋`$ and $`f([a,b])`$ is homeomorphic with $`B(𝐋,f(a),f(b)f(a))`$. Then in $`\overline{M}[a,b]=Y`$ there exists $`y`$ such that $`\rho (y,a)=sup_{xY}\rho (a,x)=t`$. The subspace $`Y`$ is covered by the finite number of $`WV`$. For $`ya`$, that is $`t>0`$, $`f([a,b])`$ is compact and is contained in the finite number of simplexes from $`P`$, consequently, there exists the permissible modification of $`g`$ on $`[a,b]`$ also.
Let $`E`$ and $`F`$ be two subsets in $`R`$. We denote by $`sp(E,F,f)`$ the subspace $`cl((_{(aE,bF,f|_{[a,b]}\text{ is }𝐋\text{ linear })}[a,b])EF)`$ in $`R`$, where $`cl(S)`$ denotes the closure of $`S`$ in $`R`$ for $`SR`$. If $`B=\{q:f^1(T^0)=q,T^0P,g`$ $`\text{is not defined in }q\}`$, then by the Teichmüller-Tukey Lemma there exists the extension of $`g`$ on $`sp(\overline{M},B,f)=:M_0`$. Let $`M_j:=sp(_{i<j}M_i,B,f)`$, where $`B=_{(T^j\text{ is not the subset of }g(M))}T^j`$, $`T^j`$ are simplexes from $`P`$ with the cardianlity of sets of vertexes equal to $`j`$, where $`jw(P)`$. From Lemma 3.2 it follows that conditions $`2.8.(5(iiii))`$ may be satisfied on $`M_jR`$. Considering vertexes of $`s`$ from $`Rf^1(T^j)_{i<j}M_i`$ we construct $`g`$ on $`M_jR`$. Since $`Rc_0(𝐋,A)`$, then $`sup_jM_j=R`$, where $`M_j`$ are ordered by inclusion: $`M_jM_i`$ for each $`ij`$.
3.3.2. Note. Henceforth, we assume that for a uniformly continuous mapping $`f:YP`$ are satisfied conditions: $`P`$ is a uniform polyhedron, bonding mappings $`f_n^m:X_mX_n`$ are uniformly continuous (see §§2.8, 2.16), where $`YX`$, $`(X,\rho )`$ is an ultrametric space.
3.4. Lemma. Let $`f:MP`$ is irreducible, $`M`$ and $`P`$ are polyhedrons, $`N`$ is a subpolyhedron in $`M`$, $`Q`$ is a subpolyhedron in $`P`$, $`f(N)Q`$, then a mapping $`f:NQ`$ is irreducible.
Proof. Let $`f:NQ`$ be reducible, that is there exists a permissible modification $`g:NQ`$ with $`g(N)`$ not contained in $`f(N)`$, $`qf(N)g(N)`$. In view of Lemma 3.3 there exists the extension $`g:MP`$. Let $`r>0`$ is sufficiently small and $`U:=N_r=\{yM:`$ there exists $`xN`$ with $`\rho (x,y)r\}`$ be the $`r`$-enlargement of the subspace $`N`$ such that $`qg(N_r)`$. Since $`M`$ and $`P`$ are (uniform) polyhedra and $`f`$ is (uniformly) continuous and $`U`$ is clopen in $`M`$, then there exists a $`p^j`$-subdivision $`M^{}`$ and a clopen polyhedron $`H`$ in $`M^{}`$ with $`N_{r/p}NUM^{}`$ such that $`h|_H=g|_H`$, $`M^{}H`$ is the subpolyhedron in $`M^{}`$, $`h|_{MH}=f|_{MH}`$. Then $`qh(H)`$ and from the irreducubility of $`f`$ it follows that $`qf(MH)`$, consequently, $`h`$ is the permissible modification of $`f`$ and $`h(M)`$ is not the subset of $`f(M)`$. This contradiction lead to the statement of this Lemma.
3.5. Lemma. Let $`f:MP`$, $`M`$ and $`P`$ be polyhedrons over $`𝐋`$. Then the condition of irreducibility of $`f`$ is equvalent to that each subpolyhedron in $`QP`$ is essentially covered.
Proof. If $`f`$ is irreducible, then due to Lemma 3.4 each subpolyhedron $`Q`$ is essentially covered. Let vice versa each $`Q`$ is essentially covered and $`f`$ has the permissible modification $`g`$ with $`P=f(M)`$ not contained in $`g(M)`$. From $`f(M)=g(M)`$ and that $`f`$ is essential the contradiction with the statement $`\{P`$ is not contained in $`g(M)\}`$ follows, consequently, $`f`$ is irreducible.
3.6. Lemma. Let $`P`$ and $`M`$ be polyhedrons. If $`f:MP`$ has only addmissible modifications, then $`f`$ is irreducible.
Proof. From $`f(M)=g(M)`$ and that $`g`$ is the permissible modification of $`f`$ it follows that each subpolyhedron $`Q`$ from $`P`$ is essentially covered due to Lemma 3.3. In view of Lemma 3.5 $`f`$ is irreducible.
3.7. Lemma. Let there is an inverse spectra $`S=\{R_m,f_n^m,E\}`$ of polyhedra $`R_m`$ over $`𝐋`$, $`f_n^m`$ are simplexial mappings, $`g^l:R_lP`$, $`g^n=g^lf_l^n`$ for a fixed $`l`$, $`f_n=lim_mf_n^m`$, $`g=g^lf_l`$, $`R=limS`$, $`f_n:RR_n`$, $`E`$ is linearly ordered. If $`g:RP`$ is reducible for a polyhedron $`P`$, then for almost all $`n`$ (that is, there exists $`kE`$ such that for each $`nk`$) $`g^n:R_nP`$ are reducible.
Proof. In view of Lemma 3.3 there exists a permissible modification $`h`$ of $`g`$, then $`g`$ and $`h`$ define mappings $`g^n`$ and $`h^n`$ such that $`h^n`$ on $`R_n`$ are permissible modifications of $`g^n`$. If $`g`$ is reducible, then by Lemma 3.5 $`g^n`$ are reducible for almost all $`nE`$, since $`h(R)`$ is not a subset of $`g(R)`$ and $`h(R)=g(R)`$.
3.8. Lemma. If $`f:MN`$, $`g:NT`$, $`f(M)=N`$, where $`M`$ and $`N`$ are polyhedra, then from $`g`$ is inessential (reducible) it follows $`fg`$ is inessential (reducible).
Proof. If $`f`$ is $`𝐋`$-linear on $`[a,b]M`$ and $`g`$ is $`𝐋`$-linear on $`[f(a),f(b)]`$, then $`fg`$ is $`𝐋`$-linear on $`[a,b]`$. If $`h`$ is a permissible modification of $`g`$, then $`fh`$ is a permissible modification of $`fg`$, hence $`fg`$ is reducible. If $`g`$ is inessential, then there exists a mapping $`h`$ such that
$`(i)`$ $`h(N)`$ does not contain $`T`$;
$`(ii)`$ there exists $`N_0N`$ such that $`N_0N`$, $`g(N)T=g(N_0)=h(N_0)`$ and $`h|_{N_0}=g|_{N_0}`$;
$`(iii)`$ if $`g`$ is linear on $`[x,y]N`$, then $`h`$ is linear on $`[x,y]`$ such that $`h(x)h(y)`$ for each $`g(x)g(y)`$. Therefore, $`fh:MT`$ satisfies Conditions $`(iiii)`$ for $`M`$ instead of $`N`$, $`fh`$ instead of $`h`$ and $`fg`$ instead of $`g`$, consequently, $`fg`$ is inessential.
3.9. Lemma. Suppose that there is given an irreducible inverse mapping system $`S=\{P_m,f_n^m,E\}`$ of polyhedra $`P_m`$ over $`𝐋`$, $`M`$ is closed in $`P=limS`$ such that $`M_n:=f_n(M)`$ are subpolyhedra in $`P_n`$ and for each $`m>l`$ permissible modifications $`g_l^m`$ for $`f_l^m`$ are given (on the entire $`P_m`$) and for each $`n>m`$ mappings $`g_l^mf_m^n`$ are permissible modifications of $`f_l^m`$. Then the inverse mapping system $`S_M:=\{M_n,g_l^mf_m^n,E_l\}`$ is irreducible, where $`E_l=\{nE:nl\}`$.
Proof. From surjectivity and irreducibility of $`f_l^m`$ it follows surjectivity and irreducibility of $`g_l^mf_m^n`$ due to Lemmas 3.4, 3.5 and 3.8. Since $`E_l`$ is cofinal with $`E`$, then for each $`W`$ open in $`M`$ there exists $`V`$ open in $`P`$ such that $`W=MV`$. Let $`dim_𝐋P_l=d`$. We use a disjoint covering of $`P_l`$ by simplexes clopen in it. There exists $`lm_dE`$ such that each $`d`$-dimensional simplex $`T^d`$ from $`T_l`$ over $`𝐋`$ by means of $`(f_l^n,M_n)`$ is either essentially for each $`nl`$ or inessentially for each $`nm_d`$ covered. In the second case there exists a permissible modification $`k_l^{m_d}`$ such that $`k_l^{m_d}[(f_1^{m_d})^1(T^d)]`$ does not contain any point from $`T^dT^d`$, hence we get a permissible modification $`k_l^{m_d}`$ of $`f_l^{m_d}`$ on the entire $`P_{m_d}`$. If $`T^d`$ is essentially covered by $`(f_l^{m_d},M_{m_d})`$, then $`T^d`$ is essentially covered by $`(f^n,M_n)`$ for each $`nm_d`$.
Let $`q<d`$, where $`q`$ is a cardinal. Then there exists $`m_qm_d`$ such that for each simplexes $`T^q`$ from $`P_l`$ with $`dim_𝐋T^q=q`$ is either essentially for each $`nl`$ or inessentially for each $`nm_q`$ covered by $`(f_l^n,M_n)`$. Since $`P_l`$ has a disjoint covering by such simplexes and $`S`$ is irreducible, then $`S_M`$ is irreducible.
3.10. Lemma. If $`T`$ is a simplex from a $`p^j`$-subdivision $`P^{}`$ of polyhedron $`P`$ then for each clopen neighbourhood $`UT`$ such that $`U`$ is a subpolyhedron in $`P^{}`$ there exists a mapping $`k:PP`$ such that $`k|_U`$ is simplexial and $`k(U)=T`$.
Proof. In view of Theorem 2.11 there exists the retraction $`r:PU`$ (it is uniform if $`P`$ is the uniform polyhedron). In view of Lemmas 3.2 and 3.3 there exists a simplexial mapping $`f:UT`$ such that $`f(U)=T`$, hence $`k=fr`$ is the desired mapping, since $`r|_U=id`$ and $`rr=r`$, where $`id(x)=x`$ for each $`xU`$. To construct $`f`$ we consider simplexes $`s`$ of $`P^{}`$ forming its disjoint clopen covering. For each such $`s`$ either $`sT=\mathrm{}`$ or $`s=T`$. If $`s=B(c_0(𝐋,A),x,r)`$ and $`T=B(c_0(𝐋,A^{}),y,r^{})`$ with $`A^{}A`$, then $`f|_s`$ is evidently defined such that $`f(s)=T`$. For $`A^{}A`$ and $`A^{}A`$ we can map $`s`$ on the corresponding verge (or vertex) of $`T`$. For $`s=T`$ we can take $`f|_T=id`$.
3.11. Lemma. Let $`P`$ be a uniform polyhedron, $`f:MP`$ and $`g:MP`$ are two $`b`$-close mappings (that is, $`\rho (f(a),g(a))b`$ for each $`aM`$), where $`b<inf_{sP}diam(s)`$, $`s`$ are clopen simplexes in $`P`$, then there exists $`k`$ from Lemma 3.10 such that $`kg`$ and $`f`$ are $`b`$-close and $`kg`$ is a permissible modification of $`f`$, where $`g|_{[x,y]}`$ is $`𝐋`$-linear or $`g(x)g(y)`$ if and only if $`f|_{[x,y]}`$ is $`𝐋`$-linear or $`f(x)f(y)`$ respectively for each $`[x,y]M`$.
Proof. Let $`P^{}`$ be a subdivision of $`P`$ constructed with the help of $`p^j`$-subdivisions of $`P`$ such that $`j𝐍`$. In view of $`b<inf_{sP}diam(s)`$ we have that $`g(a)\tau `$ if and only if $`f(a)\tau `$ for each simplex $`\tau `$ in $`P`$ and $`aM`$. If $`M_1`$ is a clopen polyhedron in $`M`$ and $`s`$ is a simplex (or its verge) in $`M_1`$, then $`f|_s`$ is affine if and only if $`g|_s`$ is the affine mapping, moreover $`f(s)`$ and $`g(s)`$ have the same dimension over $`𝐋`$ and $`dim_𝐋g(q)=dim_𝐋f(q)`$ for each verge $`q`$ in $`s`$. Therefore, for a polyhedron $`M_1`$ we can choose $`k:PP`$ such that $`kg(M)=f(M)`$, since $`P`$ is a uniform polyhedron and there exists $`j𝐍`$ such that $`p^jsup_{sP}diam(s)<b`$.
3.12. Lemma. Let $`f:PQ`$ be a non-stretching mapping of a uniform polyhedron $`P`$ onto a uniform polyhedron $`Q`$ over $`𝐋`$. Then there exists a $`p^j`$-subdivision $`P^{}`$ of $`P`$ and a normal mapping $`g:PQ`$ such that $`g`$ is a permissible modification of $`f`$.
Proof. For each $`b>0`$ there exist a $`p^j`$-subdivision $`P^{}`$ of $`P`$ and a $`p^i`$-subdivision of $`Q`$ and a simplexial mapping $`h:P^{}Q`$, which is $`b`$-close to $`f`$, since $`f`$ is non-stretching. Indeed, simplexes $`T^lP^{}`$ are disjoint and clopen for them due to Lemmas 3 and 4 $`h`$ can be chosen such that:
$`(i)`$ $`|h(e)f(e)|<b`$ for linearly independent vertexes $`e_{l,j}T^l`$, that is, $`\{(e_{l,j}e_{i,0}):jA_l\}`$ are linearly independent, $`e_{l,0}`$ is a marked vertex from $`T^l`$, $`card(A_l)=dim_𝐋T^l`$ and
$`(ii)`$ $`h(T^l)`$ are simplexes in $`Q`$ with $`diam(h(T^l))diam(T^l)`$ for suitable $`b>0`$ and the $`p^i`$-subdivision $`Q^{}`$ of $`Q`$, where $`h|_{T^l}`$ are affine mappings for each $`l`$. This is possible due to uniformity of polyhedra $`P`$ and $`Q`$. Taking $`g=kh`$, where $`k`$ is the suitable mapping from Lemma 3.11 we get the desired $`g`$.
3.13. Lemma. Let $`g`$ be a permissible modification of $`f:RP`$ and $`h:PQ`$ be a normal mapping, where $`P`$ and $`Q`$ are uniform polyhedra. Then $`hg`$ is the permissible modification of $`hf`$.
Proof. If $`f(R)=g(R)`$, then $`hf(R)=hg(R)`$. If $`f|_{[x,y]}`$ is $`𝐋`$-linear, then $`g|_{[x,y]}`$ is $`𝐋`$-linear, where $`[x,y]R`$, hence if $`hf|_{[u,v]}`$ is $`𝐋`$-linear, then $`hg|_{[u,v]}`$ is $`𝐋`$-linear and $`hg(u)hg(v)`$ for each $`hf(u)hf(v)`$, since $`h`$ is affine for some $`p^j`$-subdivision $`Q^{}`$ of $`Q`$. If $`aR`$ and $`f(a)s`$, where $`s`$ is a simplex from $`P`$, then $`g(a)s`$, hence if $`hf(a)\tau `$, then $`hg(a)\tau `$, where $`\tau `$ is a simplex in $`Q`$, since $`\rho _Q(hf(a),hg(a))\rho _P(f(a),g(a))`$ and $`h:PQ^{}`$ is simplexial.
3.14. Lemma. Let $`\{P_n,f_m^n,E\}=S`$ be an irreducible inverse system of uniform polyhedra over $`𝐋`$, $`M`$ be closed in $`P=limS`$, $`M_n=f_n(M)`$ be subpolyhedra in $`P_n`$, $`\{Q_k,g_l^k,F\}`$ be an inverse spectra of polyhedra $`Q_k`$, which appear by $`p^{j(k)}`$-subdivisions of $`P_{n(k)}`$, $`g_l^k`$ be normal and permissible modifications of $`f_{n(l)}^{n(k)}`$, where $`F`$ is cofinal with $`E`$, $`N_k=g_k(M)`$, $`g_l^k|_{N_k}`$ and $`g_l|_M`$ are irreducible. Then $`N_k`$ are polyhedra.
Proof. This follows from Lemmas 3.7-3.12, since an image of a polyhedron under simplexial mapping is a polyhedron (see also the proof of Lemma IV.31.XII and starting the consideration from some fixed $`n_1=qE`$).
3.15. Lemma. Suppose that for a complete ultrametric space $`R`$ there is a non-stretching mapping $`f:RP`$, $`f(R)=P`$, $`P`$ is a uniform polyhedron over $`𝐋`$. Then for each $`b>0`$ there exists a $`b`$-mapping $`g:RQ`$ and $`Q`$ is a uniform polyhedron over $`𝐋`$ such that for sufficiently fine $`p^j`$-subdivision $`Q^{}`$ of a polyhedron $`Q`$ there exists a normal mapping $`k:QP`$ and $`kg`$ is a non-stretching permissible modification of $`f`$.
Proof. For $`R`$ there exists the embedding $`Rc_0(𝐋,A)`$ with $`card(A)=w(R)`$ by Theorem 2.4 and a clopen neighbourhood $`S`$ with $`RSR(r)`$ that is a uniform polyhedron due to Corollary 2.6, where $`R(r)`$ denotes the $`r=b/p`$-enlargement of $`R`$. By Lemma 3.2 there exists the subpolyhedron $`Q`$ with $`RQS`$ and the $`b^{}`$-mapping $`g:RQ`$. If $`[a,z]S`$ and $`f|_{[a,z]}`$ is $`𝐋`$-linear, then we can choose $`g`$ such that it is linear on $`[a,z]`$ and $`g(a)g(z)`$ when $`f(a)f(z)`$. From the completeness of $`R`$ and Lemma 2.14 it follows that $`f`$ has the non-stretching extension $`f:SP`$. Then there are a $`b`$-mapping $`g`$, a uniform polyhedron $`Q`$ and a non-stretching mapping $`h:QP`$ for sufficiently small $`r>0`$ and $`b^{}>0`$. For $`h`$ due to Lemma 3.12 there exists the permissible modification and normal $`k:Q^{}P`$ such that $`kg(R)=P`$, $`f(R)=kg(R)`$ for $`R0`$, that is, when $`R`$ contains a clopen polyhedron $`R_1`$ with simplexes $`T`$ having $`dim_𝐋T>0`$. Such $`k`$, $`h`$ and $`g`$ may be constructed on each simplex $`T`$ from $`Q^{}`$ and then on the entire space.
3.16. Lemma. Suppose that $`R`$ is a complete ultrametric space, $`f_n:RP_n`$ are non-stretching $`b_n`$-mappings, $`P_n`$ are uniform polyhedra over $`𝐋`$, $`nE`$, $`E`$ is an ordered set such that for each $`b>0`$ there exists $`lE`$ with $`0<b_n<b`$ for $`n>l`$. Then there exists a normal irreducible inverse mapping system $`S=\{Q_m,g_m^n,F\}`$, $`F`$ is cofinal with $`E`$, $`limS=R`$, $`Q_m`$ are subpolyhedra of $`p^{j(m)}`$-subdivisions of $`P_{n(m)}`$ and $`g_m`$ are permissible modifications of $`f_{n(m)}`$, where $`g_m=lim_ng_m^n:RQ_m`$.
Proof. In view of Lemmas 3.5, 3.6, 3.13 and 3.14 it is sufficient at first to construct $`S`$ with non-stretching normal normal and surjective mappings $`g_n^m`$. This can be done due Lemma 3.15 with $`g_l^mf_{n(m)}`$ being non-stretching permissible modifications of $`f_{n(l)}`$ and $`lim_mg_l^mf_{n(m)}`$ are non-stretching permissible modifications of $`f_{n(l)}`$.
3.17. Lemma. If the ultrametric space $`(X,\rho )`$ is isomorphic with
$$lim\{(X_m,\rho _m);f_n^m;F\}$$
and the following conditions are satisfied:
$`(1)`$ for each $`m`$ there are embeddings $`q_m:X_m(E,\rho )`$ into a complete space $`(E,\rho )`$;
$`(2)`$ $`f_m:XX_m`$ are projections;
$`(3)`$ $`(X_m,\rho _m)`$ are ultrametric spaces;
$`(4)`$ there is given a family $`\{b_m>0:mF\}`$, $`b_m\rho (X,X)`$ and for each $`b>0`$ the set $`\{m:b_m>b\}`$ is finite, where $`t_m:=inf_{\rho (x,y)>b_m}\rho (q_m(x),q_m(y))`$ and $`lim_mt_m=0`$, for each $`m>n`$ and $`x`$ and the inequality $`\rho (q_m(x),q_nf_n^m(x))<t_n`$ is accomplished. Then the mappings $`q_mf_m`$ converge uniformly to the embedding $`XE`$.
Proof. In view of Lemma 2.2 we may suppose that $`\rho (X,X)`$ and $`\rho _m(X_m,X_m)\stackrel{~}{\mathrm{\Gamma }}(𝐋)`$. If $`x`$ and $`yX`$ and $`\rho (x,y)>b_n`$, then $`\rho (q_nf_n(x),q_nf_n(y))t_n`$. From conditions $`(14)`$ it follows the existence of $`k`$ such that $`\rho (q_mf_m(x),q_mf_m(y))t_n`$ for each $`m>k`$ and for the ultrametric $`\rho `$, consequently, $`q=lim_mq_m`$ is the embedding, since $`\rho (q(x),q(y))t_n`$.
3.18. Theorem. Let $`(X,𝐏)`$ be a complete ultrauniform space and $`𝐋`$ be a locally compact field. Then there exists an irreducible normal expansion of $`(X,𝐏)`$ into the limit of the inverse mapping system
$$S=\{P_n,f_m^n,E\}$$
of uniform polyhedra $`P_n`$ over $`𝐋`$, moreover, $`limS`$ is isomorphic with $`(X,𝐏)`$, in particular for an ultrametric space $`(X,\rho )`$ the system $`S`$ is the inverse sequence.
Proof. From Corollary 2.5 and Theorem 3.1 it follows the existence of the expansion of $`(X,𝐏)`$ into the uniformly isomorphic limit of the inverse mapping system $`R=\{Y_j,f_i^j,F\}`$ of ANRU $`Y_j`$ with non-stretching $`f_i^j`$, where $`Y_j`$ are the complete ultrametric spaces. Each $`Y_j`$ is closed in the finite products of the spaces $`c_0(𝐋,A_k)`$. From Lemmas 3.2-3.17 it follows the existence of the irreducible normal uniform polyhedral expansion for each $`Y_j`$, moreover, using the permissible modifications of $`g_i^j`$ of $`f_i^j`$ and the same Lemmas we can construct the system of the entire space $`(X,𝐏)`$ with the same properties.
We consider further uniform coverings $`V`$ corresponding to the uniform polyhedra $`P=_{WV}W`$, which due to Theorem 2.11 are ANRU. Let $`\rho `$ be the pseudoultrametric in $`X`$ and $`\rho (X,X)\stackrel{~}{\mathrm{\Gamma }}(𝐋)`$. At first, if $`V`$ is given with the help of the chosen pseudoultrametric $`\rho `$, then we can associate with $`V`$ a $`p^k`$-nerve with $`k𝐙`$, that is, an abstract simplexial complex $`N_k`$ vertexes of which are elements from $`V`$. Its simplexes are the spans of (pulled on) vertexes $`W_j`$ satisfying $`\rho (W_j,W_i)p^kb`$, where $`b=sup_{WV}diam(W)<\mathrm{}`$, each rib $`[W_j,W_i]`$ from $`s`$ has the lenghts not less than $`t=inf_{W_jW_iV}\rho (W_j,W_i)>0`$. Then from $`N_kc_0(𝐋,A_k)`$ with $`card(A_k)=w(N_k)`$ it follows that each $`s`$ is uniformly isomorphic with some ball $`B(c_0(𝐋,A),0,1)`$, where $`card(A)=mw(N_k)`$. With each $`V`$ is associated the equivalence relation: $`xRy`$ if and only if there exists $`WV`$ such that $`x`$ and $`yW`$, then the quotient mapping $`f:XX/R`$ is defined (see Propositions 2.4.3 and 2.4.9 ). In particular with each locally finite functionally open covering $`V`$ is associated the partition of the unity $`\{f_W:WV\}`$, $`f_W:XB(𝐋,0,1)`$, $`f_W(x)=1`$ for $`xW`$ and $`f_W(x)=0`$ for $`xW`$, $`\{f_W:WV\}`$ is subordinated to $`V`$ (see aslo §5.1 ). There are the canonical non-stretching mappings $`F_k:XN_k`$. If $`X`$ is compact, then $`X/R`$ is the finite discrete space and $`dim_𝐋N_k=n𝐍`$. Into each $`V`$ may be refined a disjoint clopen uniform covering $`K`$ with $`sup_{WK}diam(W)bp^j`$, where $`j𝐍`$. That is, $`V`$ has the uniform strict shrinking.
In general there exist abstract simplexial complexes $`N_V`$ with $`dim_𝐋N_Vdim(X)0`$, if $`V`$ is an arbitrary functionally open covering of $`X`$ of order not less than $`dim(X)`$. For $`dim(X)=\mathrm{}`$ we interpret orders of coverings as cardinals equal to $`k\mathrm{}_0`$. This is due to the following second procedure. If $`\{W_j:`$ $`j\mathrm{\Lambda }_W\}`$ are elements of a subfamily $`𝖶`$ of a covering $`V`$ such that $`W_j`$ are pairwise distinct and $`_{j\mathrm{\Lambda }_W}W_j\mathrm{}`$, then to $`𝖶`$ corresponds an abstract simplex $`s`$ in $`c_0(𝐋,\mathrm{\Lambda }_W)`$ such that $`W_j`$ correspond to vertexes of $`s`$. Evidently, if the simplex $`s`$ is clopen in the Banach space $`c_0(𝐋,A)`$, then it is characterized by two points $`x,ys`$ such that $`s=B(c_0(𝐋,A),x,|xy|)`$. If $`s`$ with $`Int(s)=\mathrm{}`$, then vertexes of $`s`$ characterize a minimal Banach subspace of $`c_0(𝐋,A)`$ in which $`s`$ is contained.
Using the first procedure we can consider the sequence of such shrinkings: $`V^{m+1}V^m`$ with $`b_m=bp^m`$, where $`m𝐍`$. With each $`V^m`$ is associated $`p^{k(m)}`$-nerve $`N_{k(m)}`$. Let $`k(m)m`$, $`k(m+1)k(m)`$ for each $`m𝐍`$ and $`lim_m\mathrm{}k(m)=\mathrm{}`$. If $`x`$ is an isolated point in $`X`$, then there exists $`n𝐍`$ with $`\mathrm{max}(b_n,p^{k(n)}b_n)<inf_{yX\{x\}}\rho (x,y)`$. Then the simplex $`sN_{k(m)}`$ with $`xs`$ and $`mn`$ is zero-dimensional over $`𝐋`$, that is, $`s=\{x\}`$.
By the construction of $`N_k`$ for each simplex $`s_{m+1}N_{k(m+1)}`$ there exists $`s_mN_{k(m)}`$ with $`f_m^{m+1}(s_{m+1})s_m`$, where $`f_i^j`$ are the bonding mappings of the inverse sequence $`S=\{N_{k(m)},f_i^m,𝐍)`$. Each $`f_m^{m+1}`$ is non-stretching, since decreases the distance at least into $`p`$ times and $`b_m/b_{m+1}p`$. If $`xy`$, then there exists $`n`$ with $`\mathrm{max}(b_n,b_np^{k(n)})<\rho (x,y)`$, consequently, for each $`m>n`$ there are disjoint simplexes $`s`$ and $`s^{}N_{k(m)}`$ with $`xs`$ and $`ys^{}`$. Therefore, there exists the uniformly continuous mapping $`g:XlimS`$, where $`g(x)=lim_m\{s_m,f_i^m\}`$ and $`s_mx`$ for each $`m𝐍`$. Therefore, the uniformly continuous projectors $`f_m:XN_{k(m)}`$ are defined, since for each $`b>0`$ there exists $`r𝐍`$ such that $`b_mp^{k(m+r)k(m)}<b`$ and $`f_m(W)=f_m^{m+r}f_{m+r}(W)`$ and $`diam(f_m(W))<b`$, where $`WV^{m+r}`$, $`f_{m+r}(W)`$ belongs to clopen star of the corresponding vertex $`vN_{k(m+r)}`$. Further (see also Lemma IV.33 in ) we can verify that $`g(X)=limS`$ and $`g`$ is the uniform isomorphism, since $`_{m=l}^{\mathrm{}}g_m^1(s_m)=\{x\}`$ for the family $`\{s_m:m\}`$ corresponding to $`x`$ (see above).
Evidently, $`dim_𝐋N_{k(m)}`$ may be from $`0`$ for $`k(m)=m`$ up to $`card(A)`$ with $`card(A)=w(X)`$. For $`k(m)>m`$ in the inverse sequence $`S`$ the mappings $`f_m^{m+1}`$ may map simplexes $`s`$ from $`N_{k(m+1)}`$ into simplexes $`q`$ from $`N_{k(m)}`$ of lower dimension over $`𝐋`$, for example, when $`W_{m+1}W_m`$, $`W_mV^m`$, $`W_m=B(c_0(𝐋,A_j),x,r)`$, $`W_{m+1}=B(c_0(𝐋,A_n),x^{},r/p)`$, $`card(A_n)>card(X/R_m)dim_𝐋N_{k(m)}card(A_j)`$, since $`dim_𝐋N_{k(m+1)}card(A_n)`$ for $`k(m+1)>m1`$. This fact is based on the homeomorphism of $`D^\mathrm{}_0`$, $`B(𝐋,0,1)`$ and $`B(𝐋,0,1)^\mathrm{}_0`$, where $`D=\{0,1\}`$ is a discrete two-point set (see also ).
For the complete ultrauniform space $`(X,𝐏)`$ we can consider the base of uniform coverings $`\{V_\rho ^n:n𝐍,\rho 𝐏\}`$, where each $`V_\rho ^n`$ is given relative to the considered $`\rho `$, $`S=\{N_{\rho ,k(m)},f_{\rho ^{},k(m^{})}^{\rho ,k(m)},𝐏\times 𝐍\}`$. To each $`V_\rho ^m`$ there corresponds $`N_{\rho ,k(m)}`$; $`\rho ^{}\rho `$ if and only if $`\rho ^{}(x,y)\rho (x,y)`$ for each $`x`$ and $`yX`$; $`(\rho ^{},m^{})(\rho ,m)`$ if and only if $`\rho ^{}\rho `$ and $`m^{}m`$. If to associate $`N_V`$ with $`dim(X)`$ and orders of arbitrary functionally open coverings $`\stackrel{~}{𝖵}^𝗇`$ of the space $`X`$, then there exist an inverse spectrum of dimension equal to $`dim(X)`$, where $`\text{ }^{}\stackrel{~}{𝖵}^𝗆<\stackrel{~}{𝖵}^𝗇`$ for each $`m<nE`$, that is $`\stackrel{~}{𝖵}^𝗆`$ is a star refinement of $`\stackrel{~}{𝖵}^𝗇`$. We take $`\{\stackrel{~}{𝖵}^𝗆:mE\}`$ such that for each $`xX`$ there exists a cofinal subset $`E(x)E`$ and a family $`\{\stackrel{~}{V}_\beta ^m:\beta E(x)\}`$ with $`_{\beta E(x)}\stackrel{~}{V}_\beta ^m=\{x\}`$. Let $`P_m=N_{\stackrel{~}{𝖵}_𝗆}`$ for each $`mE`$ and $`f_n^m:P_mP_n`$ be uniformly continuous bonding mappings, $`g(x)=lim_m\{s_m,f_i^m\}`$, where $`s_m`$ are such that $`s_mx`$ for each $`mE`$. Polyhedra $`P_m`$ can be chosen such that $`sup_{(sP_m)}diam(s)<p^{k(m)}b`$ and $`lim_{mF}k(m)=\mathrm{}`$ for each linearly ordered subset $`F`$ of $`E`$ such that $`F`$ is cofinal with the first countable ordinal $`\omega _0`$, $`s`$ are simplexes in $`P_m`$, $`k(m)k(n)`$ for each $`mnF`$, where $`k(m)𝐙`$, $`b=const>0`$. Then $`g:XlimS`$ is a uniform isomorphism. Using permissible modifications and Lemmas 3.2-3.17 we get the statement of this theorem.
## 4 Absolute polyhedral expansions and their applications.
Many of the results from about absolute polyhedral expansions in the classical case (that is, for polyhedra over $`𝐑`$) may be transformed in the non-Archimedean case. Further the main differences are given. Part of lemmas and definitions from may be naturally reformulated and proved, for example, with the substitution of barycentric subdivision into $`p^j`$-subdivision. We use Theorems 3.1, 3.18 and apply absolute polyhedral expansions at first to each ultrametric space $`Y_j`$, then to the entire ultrauniform space $`(X,𝐏)`$. The polyhedra $`P`$ considered in §4 as well as in §§2, 3 may have in general infinite dimension over $`𝐋`$.
4.1.1. Definitions and Notes. A proper $`p^j`$-subdivision of a complex $`K^a`$ over $`𝐋`$ is called normal. Certainly $`dim_𝐋K^a`$ may be infinite. With the help of affine mappings for each simplexial complex $`K^a`$ the space $`|K^{}|`$ of its $`p^j`$-subdivision is the same as $`|K^a|.`$ This space is denoted by $`K`$ and is supplied with the ultrametric induced from $`|K^a|`$. A support of a subset $`YK`$ relative to complex $`K^a`$ is the least subcomplex containing $`Y`$. A support of a point is a simplex $`s`$ in $`K^a`$ of minimal $`dim_𝐋s`$ containing it. If $`K_1^a`$ is a subcomplex (clopen) of a complex $`K^a`$, then the set $`KK_1`$ is called an open (clopen) subcomplex of a complex $`K^a`$. If $`A`$ is an open (clopen) subcomplex of a complex $`K^a`$, then a set $`KA`$ is a space of some subcomplex (clopen) of $`K^a`$, which is denoted by $`K^aA`$. The closure of an open subcomplex $`A`$ of a complex $`K^a`$ is a subspace of all subcomplexes of $`K^a`$ consisting of all simplexes interiors of which are contained in $`A`$ and in addition of all its verges in the non-Archimedean sense (see §2). The least clopen subcomplex of the complex $`K^a`$ containing $`X`$, $`XK^a`$, is called the polyhedral neighbourhood or the clopen star of a set $`X`$ relative to $`K^a`$ and is denoted $`K^aX`$. They form a pointwise finite covering of $`K`$. Moreover, a point $`xK`$ is contained in a clopen simplex $`\tau `$ if and only if its vertexes $`\{e_j:jJ\}`$ are all vertexes $`\{e_j:jJ\}`$ stars of which contain this point $`x`$. A clopen star of a subcomplex is by definition a clopen star of its space. A clopen star $`K^aK_1`$ of a subcomplex $`K_1^a`$ in a complex $`K^a`$ is by definition a union of all clopen simplexes in $`K^a`$ which have nonvoid verge in the subcomplex $`K_1^a`$. This means that $`K^aK_1`$ is also a union of clopen stars of vertexes of the subcomplex $`K_1^a`$.
Subcomplex of a complex $`K^a`$ a space of which is a clopen star of a set $`X`$ (subcomplex $`X^a`$) is called a star of the set $`X`$ (subcomplex $`X^a`$) and is denoted by $`(K^aX)^a`$ or by $`\overline{K}^aX^a`$, if $`X^a`$ is the subcomplex of the complex $`K^a`$.
4.1.2. Lemma. Let $`K^{}`$ be a normal, that is multiple $`p^j`$-subdivision, of a complex $`K^a`$ and $`t`$ be a simplex in $`K^{}`$. Then there exists a vertex $`eK^a`$ such that $`K^atK^ae`$. Moreover, if $`K_1^a`$ is the subcomplex of the complex $`K^a`$ such that $`|t|K_1\mathrm{}`$, then $`e`$ is the vertex in the subcomplex $`K_1^a`$.
Proof. By the definition for each vertex $`et`$ we have $`K^aeK_1^a`$, if $`|t|K_1^a\mathrm{}`$. Since simplexes $`s`$ in $`K^{}`$ form a disjoint covering of $`|K^{}|`$, then $`K^{}tK^{}e`$.
4.1.3. Definition. A subcomplex $`K_1^a`$ of a complex $`K^a`$ is called complete, if $`sK_1^a`$ for each simpex $`s`$ all vertexes of which are in $`K_1^a`$.
4.1.4. Lemmas. Let $`K_1^a`$ be a complete subcomplex of a complex $`K^a`$.
(i). If $`K^{}`$ is a $`p^j`$-subdivision of a complex $`K^a`$, then a $`p^j`$-subdivision $`K_1^{}`$ of the subcomplex $`K_1^a`$ induced by $`K^{}`$ is the complete subcomplex of the complex $`K^{}`$.
(ii). If $`f:N^aK^a`$ is a simplexial mapping, then $`f^1K_1^a`$ is a complete subcomplex of the complex $`N^a`$.
Proof. $`(i)`$. If $`\tau `$ is a simplex in $`K^a`$, then $`\tau K_1^a`$ is a verge (may be void) of the simplex $`\tau `$. On the other hand, there exists embedding of $`K^a`$ into $`c_0(𝐋,w(|K^a|))`$. Therefore, $`|K^a|`$ has a disjoint paving by clopen balls $`B_l`$ in $`c_0(𝐋,w(|K^a|))`$. For each simplex $`s`$ in $`K^{}`$ there exists a ball $`B`$ in $`c_0(𝐋,w(|K^a|))`$ such that $`B|K^{}|=s`$. Each $`p^j`$-subdivision of the simplex $`\tau `$ in $`K^a`$ consists of a disjoint union of simplexes $`s_i`$ in $`K^{}`$.
$`(ii)`$. If $`\tau `$ is a simplex in $`N^a`$, then $`f|_\tau `$ is an affine mapping, hence $`f(\tau )`$ is a disjoint union of simplexes from $`K_1^a`$. Therefore, $`f^1(f(\tau ))`$ is covered by disjoint family of simplexes from $`N^a`$. The same is true for verges of simplexes from $`N^a`$.
4.1.5. Note. In view of Lemmas 4.1.4 above Propositions I.1.3 and I.1.5 from are true also in the non-Archimedean case.
4.1.6. Definition. Subcomplexes $`K_1^a`$ and $`K_2^a`$ of a complex $`K^a`$ are called adjacent, if $`K_1=KK^aK_2`$ and $`K_2=KK^aK_1`$.
4.1.7. Lemma. Adjacent subcomplexes $`K_1^a`$ and $`K_2^a`$ of a complex $`K^a`$ are complete subcomplexes.
Proof. If $`\tau `$ is a simplex or a verge of a simplex in $`K^a`$, then either $`\tau K_1^a`$ or $`\tau K_2^a`$ and $`K_1^aK_2^a=\mathrm{}`$.
4.1.8. Definition. For complexes $`K_1^a`$ and $`K_2^a`$ their join $`K_1^aK_2^a`$ is defined in the corresponding Banach subspace $`Z`$ of $`c_0(𝐋,A_1)c_0(𝐋,A_2)𝐋=c_0(𝐋,A_1+A_2+1)`$ with $`card(A_j)=w(K_j^a)`$ analogously to the classical case with substitution of $`\{x:0x1\}=[0,1]_𝐑𝐑`$ into $`B(𝐋,0,1)=[0,1]𝐋`$, where $`j=1`$ or $`j=2`$. Each simplexes $`\tau _1`$ in $`K_1^a`$ and $`\tau _2`$ in $`K_2^a`$ are characterized completely by their vertexes. Then each simplex $`\tau `$ in $`K_1^aK_2^a`$ has a set of vertexes $`\mathrm{\Lambda }_\tau `$, which is a union of sets $`\mathrm{\Lambda }_{\tau _j}`$ of vertexes of the corresponding simplexes $`\tau _j`$ in $`K_j^a`$, where $`j=1`$ or $`j=2`$. When $`K_1^aK_2^a=\mathrm{}`$, then $`Z=c_0(𝐋,A_1+A_2+1)`$. In general $`Z=cl[sp_𝐋(K_1K_2)]`$.
4.1.9. Lemma. $`K_1^a`$ and $`K_2^a`$ are adjacent subcomplexes of a complex $`K^a`$ if and only if each simplex $`\tau `$ of a complex $`K^a`$ has the form
$`(i)`$ $`\tau =\tau _1\tau _2`$, where $`\tau _j`$ are simplexes in $`K_j^a`$, $`\tau _1\tau _2=\mathrm{}`$, $`j=1`$ or $`j=2`$.
Proof. If subcomplexes are adjacent, then $`\tau _1\tau _2`$ is obtained by junction of each pair of points $`v_1\tau _1`$ and $`v_2\tau _2`$ by a segment $`[v_1,v_2]`$ in the sense of §2.8.(2).(iii). On the other hand, equation $`(i)`$ is equivalent to $`\mathrm{\Lambda }_{\tau _1}\mathrm{\Lambda }_{\tau _2}=\mathrm{}`$. The last equation for each pair of $`\tau _1`$ in $`K_1^a`$ and $`\tau _2`$ in $`K_2^a`$ implies that these subcomplexes are adjacent.
4.1.10. Note. In view of Lemma 4.1.9 above Propositions I.1.10 and I.1.11 from are true also in the non-Archimedean case.
4.1.11. Definition and Note. Let $`K_1^a`$ and $`K_2^a`$ be adjacent subcomplexes of a complex $`K^a`$ and $`K_1^{}`$ be a $`p^j`$-subdivision of a complex $`K_1^a`$. For each simplex $`\tau _1`$ of a complex $`K_1^a`$ let $`\tau _1^{}`$ denotes be its $`p^j`$-subdivision induced by $`K_1^{}`$. For each simplex $`\tau `$ in $`K^a`$ let $`\tau ^{}=\tau _1^{}\tau _2`$, where $`\tau _1=\tau K_1^a`$, $`\tau _2=\tau K_2^a`$. Then the union $`Sd_{K_1^{}}K^a`$ of complexes $`\tau ^{}`$ for all simplexes $`\tau `$ of $`K^a`$ is a subdivision of a complex $`K^a`$. This gives possibility to consider more general subdivisions than $`p^j`$-subdivisions.
Evidently, analogs of Propositions I.1.12, I.1.13 and I.1.15 from are true in the non-Archimedean case. The subdivision $`Sd_{K_1^{}}K^a`$ of the complex $`K^a`$ is called the minimal continuation of a $`p^j`$-subdivision $`K_1^{}`$ of a complete subcomplex $`K_1^a`$ of the complex $`K^a`$.
4.1.12. Lemma. $`Sd_{K_1^{}}K^a=K_2^aSd_{K_1^{}}(K^aK_1^a)`$, where $`K^aK_1^a`$ is a clopen star of the subcomplex $`K_1^a`$.
Proof. $`K(K_1K_2)`$ is an intersection of clopen stars of complexes $`K_1^a`$ and $`K_2^a`$.
4.1.13. Definitions and Notes. A mapping $`h:KN`$ is called linear relative to complexes $`K^a`$ and $`N^a`$, if it is simplexial for the corresponding polyhedra. Since $`K^aS`$ is the least clopen subcomplex of a complex $`K^a`$ containing a subset $`SK`$, then $`K^aSh^1(N^ah(S))`$, that is equivalent to $`h(K^aS)N^ah(S)`$. For each mapping $`f:XK^a`$ by $`\omega _fK^a`$ is denoted a covering of a space $`X`$ by $`f^1K^ae`$, where $`K^ae`$ are called main stars for vertexes $`e`$ in $`K^a`$. Let $`<N^a,K^a>_f`$ denotes a subcomplex of a join $`N^aK^a`$ of the complexes $`N^a`$ and $`K^a`$ such that its simplexes are simplexes of $`N^a`$ and $`K^a`$ and also $`f(\tau )\tau `$ and all its verges, where $`\tau `$ is a simplex of the complex $`K^a`$. Then $`<N^a,K^a>_f`$ is called a cylinder of the simplexial mapping $`f`$. There exists a retraction $`\stackrel{~}{f}:<N^a,K^a>_fN^a`$ such that $`\stackrel{~}{f}(x)=x`$ for each $`xN`$, $`\stackrel{~}{f}(x)=f(x)`$ for each $`xK`$. Evidently, Propositions I.2.1, I.2.4-10 from are true in the non-Archimedean case also.
4.1.14. Lemma. $`dim_𝐋<N^a,K^a>_f=\mathrm{max}(dim_𝐋N^a,dim_𝐋f(K^a)+dim_𝐋K^a+1)`$.
Proof. $`dim_𝐋f(\tau )\tau dim_𝐋f(\tau )+dim_𝐋\tau +1`$ for each simplex $`\tau `$ in $`K^a.`$
4.2.1. Definitions and Notes. A polyhedral spectrum $`\{X_i;f_j^i;𝐍\}`$ is called standard if it satisfies conditions $`S.1,S.3,S.4`$ from and the following condition:
$`S.2`$ $`f_i^{i+1}`$ is a simplexial mapping of a complex $`X_{i+1}^a`$ onto a normal subdivision (that is a $`p^j`$-subdivision with the corresponding $`j𝐍`$) $`X_i^{}`$ of a complex $`X_i^a`$ for each $`i𝐍`$ such that
$`(i)`$ $`lim_j\mathrm{}sup_{\tau X_j^a}diam(f_i^j(\tau ))=0`$, where $`\tau `$ are simplexes in $`X_j^a`$. The complex $`K^a`$ is called uniform if the corresponding to it polyhedra over $`𝐋`$ is uniform.
An approximation of all subsets $`Y`$ of an arbitrary ultrauniform space $`X`$ simultaneously by one spectrum $`\{X_\alpha ;f_\beta ^\alpha ;𝖴\}`$ is called absolute, that is for each $`YX`$ there exists $`𝖴_Y𝖴`$ such that $`Y=lim\{X_\alpha ;f_\beta ^\alpha ;𝖴\}`$ and for each $`Y^{}Y`$ the following inclusion $`𝖴_Y^{}𝖴_Y`$ is satisfied. This decomposition is called
$`(a)`$ irreducible if for each open $`VX`$ there exists cofinal subset $`𝖴_V𝖴`$ such that $`\{X_\alpha ;f_\beta ^\alpha ;𝖴_V\}`$ is irreducible. It is called
$`(b)`$ uniform if each subspectrum $`\{X_\alpha ;f_\beta ^\alpha ;𝖴_V\}`$ is uniform.
An absolute polyhedral representation is called uniform from above, if polyhedra $`P`$ satisfy Condition $`2.8.(1.(i))`$. A polyhedral spectrum is called
$`(c)`$ non-degenerate if $`f_\beta ^\alpha `$ are non-degenerate, that is for each simplex $`s`$ in $`X_\alpha `$ restrictions $`f_\beta ^\alpha |_s`$ are simplexial and $`dim_𝐋f_\beta ^\alpha (s)=dim_𝐋s`$.
Notations and propositions from §§II.1-II.10 are transferable onto the non-Archimedean case due to above results and definitions with the substitution of barycentric subdivisions onto $`p^j`$-subdivisions (see for example, $`\tau _{s+1}^{}`$ in §II.1.11 ).
Let $`K_i^c`$ be defined as in §II.3.1 . For two families $`Y^c`$ and $`Y_1^c`$ of simplexes from $`K_1^c`$ a mapping $`h:YY_1`$ is called $`c`$-linear over the field $`𝐋`$, if for each simplex $`\tau `$ from $`Y^c`$ (that is from some subcomplex $`Y^cK_{i,i+1}^a`$, where $`i𝐍`$) there exists $`q𝐍`$ such that $`h(\tau )K_{q,q+1}`$ and the mapping $`h|_\tau `$ is simplexial. In view of II.4.7 $`\{K_{1,i};h_{j,i};𝐍\}`$ is an inverse spectrum, its limit space is denoted $`\text{ }^𝐍K`$ and it is called the cone of the spectrum $`\{K_{1,i};h_{j,i};𝐍\}`$. Then $`\alpha (i)`$ is a space of a complete subcomplex $`\alpha ^a(i)X_i^a`$, $`\alpha ^q(i)=f_{i+q,i}^1(\alpha (i))`$, $`\stackrel{~}{\alpha }^1(i):=X_{i+1}^a\alpha ^1(i)`$ is a polyhedral neighbourhood (clopen star) of a set $`\alpha ^1(i)X_{i+1}`$ relative to a complex $`X_{i+1}^a`$, since the space $`X_{i+1}=|X_{i+1}^a|`$ of a complex $`X_{i+1}^a`$ is the disjoint union of the corresponding simplexes $`s`$ forming its covering such that spaces $`|s|`$ of these simplexes are clopen in $`X_{i+1}`$. Propositions analogous to II.1.9, II.7.18,19, II.9.22 follow from Lemma 4.1.9 above.
4.2.2. Lemma. $`dim_𝐋(X_\alpha K_{i1,i})dim_𝐋\alpha (i)`$.
Proof. In view of §II.7.13 we have $`X_\alpha K_{i1,i}=X_{\alpha ,i}K_{i1,i}`$. Let $`M_{\alpha ,i}^a`$ be a subcomplex of a complex $`K_{i1,i}^a`$ such that $`M_{\alpha ,i}^a`$ is a union of simplexes $`\sigma f_{i,i1}(\sigma )`$, where $`\sigma `$ is a simplex of the complex $`X_i^a`$ such that $`f_{i,i1}(\sigma )X_{i1}\stackrel{~}{\alpha }^1(i2)`$. Then $`h_{i,\alpha }:M_{i,\alpha }^aM_{i,\alpha }^a`$ is the simplexial mapping and $`h_{i,\alpha }(M_{i,\alpha }\alpha (i))=X_{\alpha ,i}K_{i1,i}`$ and $`h_{i,\alpha }(M_{\alpha ,i}K_{\alpha ,i})=X_{\alpha ,i}K_{i1,i}`$ (see also §§II.9.14-19 ). For each simplex $`s`$ we have $`dim_𝐋s=card(A_s)`$, where $`cl(sp_𝐋s)=c_0(𝐋,A_s)`$, the closure is taken in the corresponding Banach space $`c_0(𝐋,A)`$ in which the complex is contained, $`sp_𝐋s`$ denotes the $`𝐋`$-linear span of $`s`$.
4.2.3. Lemma. For each $`i𝐍`$ and $`\alpha 𝖴`$ there are mappings $`h_{i,\alpha }`$ satisfying conditions
$`(1)`$ $`h_{i,\alpha }`$ is $`c`$-linear over $`𝐋`$;
$`(2)`$ $`h_{i,\alpha }(x)=h_{i1,\alpha }(x)`$ for each $`xK_{1,i1}`$;
$`(3)`$ $`h_{i,\alpha }(x)=h_{i1,\alpha }h_{i,i1}(x)`$ for each $`xh_{i,i1}^1(\alpha (i1))`$;
$`(4)`$ $`h_{i,\alpha }(x)=x`$ for each $`xX_i\stackrel{~}{\alpha }^1(i1)`$.
Proof. For $`i=1`$ the mapping $`h_{1,\alpha }`$ is identical for each $`\alpha 𝖴`$ in accordance with §II.9.7 , hence it satisfies Conditions $`(14)`$. Let $`h_{i,\alpha }`$ satisfies Conditions $`(14)`$ for each $`in`$. Then $`h_{n+1,\alpha }`$ is defined by Conditions $`(24)`$. Moreover, $`\alpha (n)`$ and $`\stackrel{~}{\alpha }^1(n)`$ may be taken clopen in $`X_n`$ and $`D_1:=h_{n+1,n}^1(\alpha (n))`$ are clopen in $`X_{n+1}`$, then $`D_1`$ and $`D_2:=K_{1,n}`$ and $`D_3:=M_{\alpha ,n+1}`$ are closed in $`K_{1,n+1}`$. For each simplex $`\tau K_{1,n+1}`$ each point $`x\tau `$ has the form $`x=a_1x_1+a_2x_2+a_3x_3`$, where $`a_jB(𝐋,0,1)`$, $`x_jD_j`$. Let $`g:c_0c_0c_0c_0`$ be the mapping of evaluation, that is, $`g(x_1,x_2,x_3)=x_1+x_2+x_3`$, $`h^j:D_jK_{1,n+1}`$ are continuous mappings, so $`h_{n+1,\alpha }(x)=(g(a_1h^1(x_1),a_2h^2(x_2),a_3h^3(x_3))`$ is continuous. If to consider uniform polyhedra, then the simplexial mappings $`f:K^aN^a`$ are considered as uniformly continuous. Then to the cylinder $`<N^a,K^a>_f`$ it corresponds the uniform polyhedron, since $`f(K)=N`$. Then $`K_{i,j}^a`$, $`X_j^a`$, $`X_i^{}`$, $`K_{i,j}^{}`$ and $`K_{i,j}^c`$ from §II.1 in correspond the uniform polyhedra over $`𝐋`$, since $`K_{i,j}^{}=Sd_{X_j}K_{i,j}^a`$ and $`inf_{(VW=\mathrm{},V\text{ and }W\text{ are simplexes from }K_{i,j}^{})}dist(V,W)>0`$. For $`T=K_i=_{j𝐍}K_{j,j+1}`$, $`K_i^c`$, $`K_\alpha ^c`$ and $`X_\alpha ^c`$ the latter condition may be not satisfied, but $`sup_{(V\text{ is simplex from }T)}diam(V)<\mathrm{}`$. At the same time $`T`$ is not necessarily complete, since it is not necessarily ANRU. If $`f_{j,i}=f_i^j:X_jX_i`$, $`f_j:\text{ }^𝐍XX_i`$ are non-stretching (uniformly continuous), then $`h_{j,i}`$ from §II.4 and $`h_i`$ from §II.5 in are also non-stretching (uniformly continuous respectively). Then from $`h_{1,\alpha }=id`$ and $`(14)`$ with induction by $`j`$ for each $`\alpha `$ it follows that $`h_{j,\alpha }`$ are non-stretching (uniformly continuous correspondingly). From this it follows that the mappings $`h_\alpha `$ from §II.10 and hence also $`f_{\alpha ,\beta }`$ and $`f_\alpha `$ from §II.11 and §II.9.27 are non-stretching (uniformly continuous respectively).
4.2.4. Lemma. For each open covering $`V`$ of $`W`$ which is a clopen subset in a limit $`lim_j\{X_j,f_i^j,𝐍\}=:\text{ }^𝐍X`$ of an inverse sequence of polyhedra there are $`\alpha 𝖴`$ such that $`W_\alpha =W`$ and a clopen covering $`w_{f_\alpha }X_\alpha ^{}\{f_\alpha ^1(X_\alpha ^{}q):q\}`$ of $`W_\alpha `$ by $`f_\alpha ^1(X_\alpha ^{}q)`$ of the main stars $`X_\alpha ^{}q`$ of the complex $`X_\alpha ^{}`$ is refined into $`V`$.
Proof. Let $`\alpha ^a(i)`$ be a subcomplex of a complex $`X_i^a`$ such that a simplex $`\sigma `$ of a complex $`X_i^a`$ is a simplex of a complex $`\alpha ^a(i)`$ if and only if $`f_i^1(X_i^ae)`$ is contained in the corresponding element $`U(e)`$ of a covering $`V`$ for each its vertex $`e`$. Then $`\alpha ^a(i)`$ is a complete subcomplex of a complex $`X_i^a`$. That for a sequence $`\alpha (i)`$ to verify condition
$`(A)`$ $`\text{ }\stackrel{~}{\alpha }^1(i)\alpha (i+1)`$ for each $`i𝐍`$ as in §II.11.20 it is sufficient to mention that $`X_i^{}`$ is the $`p^j`$-subdivision of the complex $`X_i^a`$ with $`j1`$, so $`X_i^a\tau X_i^aq`$ for some vertex $`qX_i^a`$. In this case the main stars are clopen, hence $`w_{f_\alpha }X_\alpha ^{}`$ is the clopen covering. If $`\{q_j:`$ $`j=0,1,2,\mathrm{}\}`$ are vertexes from $`X_i^a`$ stars of which contain $`f_i(x)`$, then there is a clopen simplex $`\tau X_i^a`$ with such vertexes and $`f_i(x)\tau `$. At the same time the star $`St(q_j,w_{f_i}X_i^a)`$ with each $`q_j`$ is contained in $`V`$, consequently, $`\tau \alpha ^a(i)`$ and $`f_i^1(\alpha (i))`$ is clopen in $`W_\alpha `$. Each $`p^j`$-subdivision $`X_{i1}^{}`$ of a complex $`X_{i1}^a`$ is proper for each $`j𝐍`$, consequently, $`X_{i1}^{}\sigma X_{i1}^aq`$ for some suitable vertex $`q`$ in the complex $`X_{i1}^a`$, where $`\sigma `$ is the simplex in $`X_{i1}^{}`$. Therefore, $`h_{i1}(X_\alpha ^{}e)X_{i1}^{}q`$, where $`q`$ is a vertex of the complex $`\alpha ^a(i1)`$. Since $`f_\alpha ^1(X_\alpha ^{}q)`$ contains the clopen subcovering, then $`W_\alpha `$ is clopen and $`W_\alpha =W`$. In the case of separable $`f_\alpha ^1(X_\alpha ^{}q)`$ this subcovering can be chosen disjoint due to ultrametrizability of $`f_\alpha ^1(X_\alpha ^{}q)`$ (see Theorem 7.3.3 in ).
Vice versa, let $`xW`$ and $`xUV`$. Coverings $`\omega _{f_i}X_\alpha ^{}`$ satisfy Condition 4.2.1.(i) such that $`St(x,\omega _{f_i}X_i^a)U`$. Let $`\{e_j:`$ $`j\mathrm{\Lambda }_{i,x}\}`$ be vertexes of the complex $`X_i^a`$, stars of which contain the point $`f_i(x)`$. Then $`f_i(x)\tau `$, where $`\tau `$ is a clopen simplex of a complex $`X_i^a`$ with the described above vertexes. Since $`f_i^1(X_i^ae_j)U`$ for each $`j\mathrm{\Lambda }_{i,x}`$, then $`\tau `$ is the simplex in the complex $`\alpha ^a(i)`$. Therefore, $`xf_i^1(\tau )f_i^1(\alpha (i))V_\alpha ,`$ hence $`W=V_\alpha `$.
4.2.5. Lemma. If $`Y\text{ }^𝐍X`$, $`f_i(Y)=X`$, $`f_i^j:X_jX_i`$ are irreducible mappings for each $`ji𝐍`$, then $`f_\alpha ^\beta :X_\beta ^{}X_\alpha ^{}`$ are irreducible and $`f_\alpha (YV_\alpha )=X_\alpha ^{}`$ for each $`\alpha 𝖴`$ and $`Y`$ is dense in $`\text{ }^𝐍X`$. If $`f_i`$ are irreducible for each $`i`$, then $`f_\alpha `$ are irreducible for each $`\alpha 𝖴`$.
Proof. From irredicibility of $`f_i^j`$ and Lemma 3.5 it follows that for each (clopen) simplex $`\tau X_i^a`$ the mapping $`f_i^j|_{(f_i^j)^1(\tau )X_j}`$ is essential. Let $`\beta 𝖴`$ and $`t`$ be a simplex in the complex $`X_\beta ^{}`$, then there exists some $`p^j`$-subdivision $`\sigma ^{}`$ of the simplex $`\sigma `$ from $`K_{i1,i}^aX_{\beta ,i}^c`$ and there exists a simplex $`\tau `$ from $`M_{\alpha ,i}\alpha (i)`$ such that $`\sigma =h_{i,\alpha }(\tau )`$ and $`t`$ is a simplex from $`\sigma ^{}`$. From §II.3.2,10.2 ($`c`$-linearity of $`h_\alpha :K_1K_1`$), §II.7.1.4,7.15 and §II.4 ($`f_{\beta ,\alpha }=h_\alpha |_{X_\beta }`$, $`X_\alpha =_{i𝐍}X_{\alpha ,i}`$, $`X_{\alpha ,i}^c=X_\alpha ^cK_{1,i}^c`$) it follows that $`f_\alpha ^\beta =f_{\beta ,\alpha }`$ is irreducible due to Lemmas 3.4 and 3.5, since $`h_\alpha `$ is essential on $`tX_\beta `$. Indeed, due to II.4.9 we have $`h_{j,i}=f_{j,i}`$ for $`xX_j`$ and $`ji`$. In view of II.10.3 we have $`h_\alpha (x)=h_{i,\alpha }(x)`$ for $`xK_{1,\alpha }`$. In view of II.9.1 we have $`h_\alpha (x):K_{1,i}K_{1,i}`$ is $`c`$-linear and $`h_{1,\alpha }=id`$, in addition, by II.9.3 $`h_{i,\alpha }(x)=h_{i1,\alpha }h_{i,i1}(x)`$ for $`xh_{i,i1}^1(\alpha (i1))`$, by II.9.4 $`h_{i,\alpha }(x)=x`$ for $`xX_i\stackrel{~}{\alpha }^1(i1)`$, also $`\stackrel{~}{\alpha }^1(i)=X_{i+1}^a\alpha ^1(i)`$ is the clopen star relative to the complex $`X_{i+1}^a`$, where $`\alpha ^1(i)`$ is the clopen subcomplex (and the subpolyhedron as the subspace) in $`X_{i+1}^a`$, consequently, we can suppose that $`\stackrel{~}{\alpha }^1(i)=\alpha ^1(i)`$ in this non-Archimedean case. The rest of the statements of this Lemma follows from §II.11.22 in .
4.2.6. Definition. A mapping $`f:XK^a`$ into a simplexial complex $`K^a`$ is called locally finite, if each $`xX`$ has a neighbourhood $`U`$ such that $`f(U)`$ is contained in a finite subcomplex (that is the subcomplex consisting of finite number of simplexes and their verges) of a complex $`K^a`$.
4.2.7. Lemma. If all mappings $`f_i:\text{ }^𝐍XX_i`$ are locally finite and $`X_i`$ are locally finite-dimensional over $`𝐋`$, then all mappings $`f_\alpha :V_\alpha X_\alpha ^{}`$ are locally finite.
Proof. Since $`X_i`$ are locally finite-dimensional over $`𝐋`$, then each its $`p^j`$-subdivision is also locally finite-dimensional over $`𝐋`$, hence $`X_i^{}`$ and $`X_\alpha ^c`$ are locally finite-dimensional over $`𝐋`$, since $`𝐋`$ is locally compact. For $`xV_\alpha `$ and $`\alpha 𝖴`$ we have $`xf^1(\stackrel{~}{\alpha }^1(i))`$, since $`V_\alpha =_{i=1}^{\mathrm{}}f_{i+1}^1(\stackrel{~}{\alpha }^1(i))`$. If $`W`$ is a neighbourhood of $`x`$ such that $`f_{i+1}(W)`$ is a finite subcomplex $`P^a`$ of a complex $`X_{i+1}^a`$. Taking $`Wf_{i+1}^1(\alpha (i+1))`$, then $`P\alpha (i+1)`$. Since $`h_\alpha `$ is $`c`$-linear over $`𝐋`$, then $`h_\alpha (P)`$ is contained in a finite number of simplexes from $`X_\alpha ^c`$, hence in a finite number of simplexes of its subdivision $`X_\alpha ^{}`$, hence $`f_\alpha `$ is locally finite, since $`f_\alpha |_{f_{i+1}^1(\alpha (i+1))}=h_\alpha f_{i+1}|_{f_{i+1}^1(\alpha (i+1))}`$ due to Propositions II.10.4 and II.11.3 .
4.2.8. Lemma. If all $`X_i`$ are finite-dimensional $`dim(X_i)<\mathrm{}`$, then the complexes $`X_\alpha `$ can be chosen locally finite and such that $`dim_𝐋(X_\alpha )sup_{i𝐍}dim_𝐋X_i`$.
Proof. This follows from Theorems 3.1, 3.18, since complexes $`K_1`$ for each $`Y_j`$ can be chosen locally finite and $`X_\alpha K_1`$. The second statement follows from Lemma 4.2.2.
4.2.9. Theorem. Each ultrauniform space $`(X,𝐏)`$ has an irreducible absolute normal and uniform from above polyhedral representation
$$T:=\{X_\beta ,f_\beta ^\alpha ,𝖴\}$$
over $`𝐋`$, that is, for each $`Y(\stackrel{~}{X},𝐏)`$ there exists an ordered subset $`𝖴_Y𝖴`$ such that there exists an uniform isomorphism
$$\text{ }^Yf:Ylim\{X_\beta ,f_\alpha ^\beta ,𝖴_Y\}$$
(that is $`\text{ }^Yf`$ is a bijective surjective mapping, $`\text{ }^Yf`$ and $`\text{ }^Yf^1`$ are uniformly continuous). Moreover, as an absolute polyhedral expansion $`T`$ (not necessarily uniform such that $`X`$ is homeomorphic with $`limT`$) may be chosen locally finite-dimensional $`X_\beta `$ over the field $`𝐋`$. If $`dim(X)=k`$, then there exists $`k`$-dimensional $`T`$ over $`𝐋`$ and
$$\{X_\beta ,f_\beta ^\alpha ,𝖴_X\}$$
is non-degenerate (and irreducible for complete $`X`$) such that there exists a homeomorphism
$$g:Xlim\{X_\beta ,f_\beta ^\alpha ,𝖴_X\}$$
with irreducible surjective mappings $`g_\alpha :=f_\alpha g`$.
Proof. Let $`S=\{X_j,f_i^j\}`$ be an inverse polyhedral sequence over $`𝐋`$. Then sections $`X_\alpha `$ and bonding mappings $`f_\alpha ^\beta `$ for $`\beta \alpha 𝖴`$ form the inverse mapping system
$$R=\{X_\beta ,f_\alpha ^\beta ,𝖴\}$$
that is an expansion of $`S`$ and the absolute polyhedral expansion for $`X=limS`$. If $`X_i`$ are uniform polyhedra and $`f_i^j`$ are non-stretching (or uniformly continuous) irreducible simplexial mappings for each $`ji𝐍`$, then the same is true for $`f_\alpha ^\beta `$ for each $`\beta \alpha 𝖴`$. From Theorem 3.18, Lemmas and Notes 4.1.1-4.2.8 above and §II.12 in it follows the existence of the homeomorphism
$$\text{ }^Yf:lim\{X_\beta ,f_\alpha ^\beta ,𝖴_Y\}Y$$
such that $`\text{ }^Yf`$ is uniformly continuous, since $`f_\alpha ^\beta `$ and $`f_\alpha `$ are non-stretching. From Lemma 3.17 it follows that $`\text{ }^Yf`$ is the uniform embedding, since
$$\underset{i\mathrm{}}{lim}\underset{(s\text{ are simplexes from }X_i^a)}{sup}diam(s)=0$$
and due to the consideration of the sequence $`\{\alpha (j):ji\}=\alpha `$. Indeed, the corresponding $`X_\alpha `$ has the cofinal subset $`E𝖴_Y`$ such that for each $`b>0`$ the set $`\{\alpha :sup_{(s\text{ are simplexes from }X_\alpha )}diam(s)>b`$ is finite and $`V_\alpha =h_\alpha ^1(X_\alpha )\text{ }^𝐍X`$ and $`Y=_{\alpha 𝖴_Y}V_\alpha `$, so that $`\text{ }^Yf^1`$ is uniformly continuous.
Each ultrametric space $`(X,\rho )`$ has the topological embedding into $`B(𝐋,0,1)^{w(X)}`$ and into $`D^{w(X)}`$ which are compact (see §6.2.16 in ). In particular, the separable space $`X`$ has the embedding into $`B(𝐋,0,1)^q`$ for each $`q𝐍`$ and for $`q=\mathrm{}_0`$, but not uniformly in general, for example, $`c_0(𝐋,𝐍)`$ is complete, but it is not compact. From this the last statement of the Theorem follows. That is, for the polyhedral expansions reproducing the uniformity up to an isomorphism with the initial ultrauniform space in general infinite-dimensioanl over $`𝐋`$ polyhedra may be necessary and $`sup_\beta dim_𝐋X_\beta w(X)`$. For the polyhedral expansion reproducing $`X`$ up to a (topological) homeomorphism locally finite-dimensional $`X_\beta `$ are sufficient, that is for each $`xX_\beta `$ there exists a clopen simplex $`s`$ in $`X_\beta `$ such that $`xs`$ and $`dim_𝐋s<\mathrm{}_0`$, since $`𝐋^{w(X)}`$ is a limit of an inverse system of finite-dimensional polyhedra over $`𝐋`$. In view of Lemma 2.5 it is sufficient to take $`𝖴`$ with $`card(𝖴)\mathrm{}_0card(𝐏^{})2^{w(X)}`$, where $`𝐏^{}`$ is a subfamily of $`𝐏`$ of minimal cardinality, which generate the same uniformity in $`X`$, since sections $`X_\alpha `$ of the spectrum for each ultrametric space $`Y_j`$ are spaces of $`X_\alpha ^c`$ for $`Y_j`$, where $`X_\alpha ^c`$ is a collection of simplexes $`s`$ from $`K_1^c=_{q=1}^{\mathrm{}}K_{q,q+1}^a`$ such that all vertexes of $`s`$ are in $`_{q=1}^{\mathrm{}}(\alpha (q)\stackrel{~}{\alpha }^1(q1))`$. If $`card(X)\mathrm{}_0`$, then $`\mathrm{}_0card(𝐏^{})2^{w(X)}=2^{w(X)}`$, since $`2^{w(X)}>\mathrm{}_0card(𝐏^{})`$.
If $`dim(X)=k`$, then by Theorem 3.18 there exists $`k`$-dimensional over $`𝐋`$ inverse mapping system $`S`$ associated with coverings of order $`k`$. For $`dim(X)=\mathrm{}`$ we can use orders of coverings as cardinals equal to $`k\mathrm{}_0`$ and as follows from results above bonding mappings $`f_m^n`$ and polyhedra $`P_m`$ can be chosen such that there exists $`m𝖴_X`$ with $`dim_𝐋P_m=k`$. The final part of the proof is analogous to that of Theorem 12.21 .
4.2.10. Corollary. If in a complete ultrauniform space $`(X,𝐏)`$ a subspace $`R`$ is compact and each closed subset $`G`$, $`GXR`$ is such that $`G`$ may be uniformly embedded into $`B(𝐋^𝐧,0,1)`$ with $`n=n(G)𝐍`$, then $`X`$ is uniformly isomorphic with $`limS`$, where $`S=\{P_m,f_n^m,E\}`$ is an irreducible normal inverse mapping system and $`P_m`$ are uniform finite-dimensional over $`𝐋`$ polyhedra.
Proof. For $`X_\rho R_\rho `$ there is a sequence of closed subsets $`G_\rho (k)`$ with $`_{k𝐍}G_\rho (k)=X_\rho R_\rho `$ and uniform clopen neighbourhoods $`G_\rho (k)U_\rho (k)B(𝐋^𝐧,0,r_n)`$, $`n=n(G_\rho (k))`$ such that there exists a uniformly continuous retraction $`r_{k,\rho }:U_\rho (k)G_\rho (k)`$ and a family
$$b_{\rho ,k}:=\underset{yU_\rho (k)}{sup}\underset{xG_\rho (k)}{inf}\rho (x,y)$$
satisfying conditions of Lemma 3.17. From $`R_\rho B(𝐋^𝐧,0,1)`$ and the existence of the family $`\{U_\rho (k)\}`$ such that $`lim\{U_\rho (k);f_{\rho ^{},k^{}}^{\rho ,k},F\}`$ $`=XInt(R)`$ for some directed $`F𝐏\times 𝐍`$ it follows the existence of the inverse spectra $`S`$. Indeed, $`f_{\rho ,k}`$ may be extended on $`RInt(R)`$ and compacts $`f_{\rho ,k}(RInt(R))`$ may be uniformly embedded into $`U_\rho (k)`$, and spectra for $`R`$ and $`XInt(R)`$ may be done consistent on $`RInt(R)`$ that to get $`S`$ for $`X`$ (see also Lemmas 3.2, 3.3.1).
4.2.11. Note. In view of Theorem 1.7 in an incomplete uniform space $`(X,𝐏)`$ is not ANRU, hence in $`S`$ in general polyhedra $`X_\alpha `$ are incomplete, for example, $`X_\alpha =𝐐\times B(𝐋,0,1)`$ in $`𝐋^\mathrm{𝟐}`$.
4.3. Theorem. Let $`G`$ be a locally compact group with a left ultrauniformity $`𝐏`$ and $`X`$ is closed in $`G`$. Then $`X`$ is uniformly isomorphic with a limit of an inverse mapping system $`limS`$ and $`P_m`$ are finite-dimensional uniform polyhedra over $`𝐋`$ of an irreducible normal inverse mapping system $`S=\{P_m,f_n^m,E\}`$.
Proof. Let $`V`$ be a compact clopen neighbourhood of the unit element $`eG`$, $`𝐏`$ be the family of the left -invariant ultrapseudometrics in $`G`$ ($`\rho (yx,yz)=\rho (x,z)`$ for each $`x,`$ $`y,`$ $`zG`$ and $`\rho 𝐏`$, see ). That is $`V=\{xG:\rho (x,e)b\}`$ for some $`\rho 𝐏`$ and $`b\mathrm{\Gamma }_𝐋`$. The closed subgroup $`H_\rho :=\{xG:\rho (x,e)=0\}`$ defines the equivalence relation $`T_\rho `$ in $`G`$: $`gT_\rho h`$, if $`h^1gH_\rho `$, since the left classes $`hH_\rho `$ are either disjoint or coinside. Therefore, there exists the uniformly continuous quotient mapping $`\pi _\rho :GG_\rho :=G/T_\rho `$, in addition $`VH_\rho `$ and $`VH_\rho =V`$. From the compactness of $`V/T_\rho =:V_\rho `$ it follows that $`G_\rho `$ is a locally compact ultrametric space. From the completeness of $`X`$ it follows that $`X_\rho `$ is closed in $`G_\rho c_0(𝐋,A_\rho )`$ with $`card(A_\rho )=w(G_\rho )`$. Therefore, $`XV`$ has due to Theorem 3.18 the finite-dimensional over $`𝐋`$ polyhedral expansion. In $`G`$ there is a family $`M`$ of elements such that $`gVhV=\mathrm{}`$ for each $`ghM`$ and $`G=_{gM}gV`$. If $`ghM`$, then $`\rho (g,h)>b`$, since from $`\rho (g,h)b`$ it would follow that $`\rho (gV,hV)b`$ and $`gVhV\mathrm{}`$, that contradicts the definition of $`M`$. Then $`Y(g)Y(h)=\mathrm{}`$ for $`ghM`$ with $`Y(g):=\pi _\rho (gV)`$ and $`G_\rho =_{gM}Y(g)`$ and $`X_\rho =_{gM}X_\rho Y(g)`$. From this the statement of the Theorem follows.
4.4. Theorem. An ultrauniform space $`(X,𝐏)`$ is homeomorphic with $`limS`$ with Čech-complete finite-dimensional over $`𝐋`$ polyhedra $`P_m`$ and an irreducible inverse mapping system $`S=\{P_m,f_n^m,F\}`$ if and only if $`X`$ is Čech-complete.
Proof. At first we verify that for $`(X,𝐏)`$ the condition of Čech-completeness is equivalent with the $`𝐋`$-completeness. The latter means that $`X`$ has the embedding as the closed subspace into $`𝐋^\tau `$, where $`card(\tau )w(X)`$. If $`X`$ is $`𝐋`$-complete, then $`X`$ has the embedding as the closed subspace into $`𝐑^\tau `$, since $`𝐋`$ is Lindelöf and $`X`$ is Čech-complete due to §3.11.3, 3.11.6 and 3.11.12 in .
To prove the reverse statement we verify that the $`𝐋`$-completeness is equivalent to the satisfaction of the following conditions: there is not any Tychonoff space $`Y`$ for which
$`(LC1)`$ there is the homeomorphic embedding $`r:XY`$ with $`r(X)cl(r(X))=Y`$ and
$`(LC2)`$ for each continuous function $`f:X𝐋`$ there is a continuous function $`g:Y𝐋`$ such that $`gr=f`$. Evidently, that it is sufficient to consider the case of ultrametrizable $`X`$. For such $`X`$ the family of continuous $`f:X𝐋`$ separates points $`xX`$ and closed subsets $`C`$ for which $`xC`$, for example, $`f(x)=\rho (x,C)=inf_{yC}\rho (x,y)`$. Then substituting $`𝐑`$ on $`𝐋`$ in the proof of Theorem 3.11.3 we get the desired statement.
There is a continuous function $`h:𝐋𝐑`$ with $`h(𝐋)=𝐑`$ and for each continuous $`q:X𝐑`$ there is a continuous function $`t:X𝐋`$ such that $`ht=q`$. The latter is sufficient to verify for $`(X,\rho )`$, since for $`(X,𝐏)`$ this statement may be obtained with the limits of the inverse system. For $`(X,\rho )`$ and each $`b_n>0`$ there exists a locally constant $`t_n:X𝐋`$ such that $`ht_n`$ and $`q`$ are $`b_n`$-close, since $`(X,\rho )`$ is the disjoint union of clopen subsets $`Y(j)`$ with $`sup_jdiam(f(Y(j))b`$, $`h^1(Y(j))𝐋`$, $`h^1(Y(j))h^1(Y(i))=\mathrm{}`$ for $`ij`$. Therefore, there exists a sequence $`\{t_n:n𝐍\}`$ giving the desired $`t=lim_n\mathrm{}t_n`$, where $`lim_n\mathrm{}b_n=0`$. Let $`(X,𝐏)`$ is Čhech-complete and not $`𝐋`$-complete, then there would be a homeomorphic embedding $`r:XY`$ with $`r(X)cl(r(X))=Y`$ and for each continuous function $`f:X𝐋`$ there is a continuous function $`g:Y𝐋`$ with $`gr=f`$. Consequently, each continuous function $`q:X𝐑`$ has continuous functions $`t:X𝐋`$ and $`g:Y𝐋`$ with $`gr=t`$ and $`hgr=ht=q`$, that contradicts Čhech-completeness of $`X`$.
If $`X`$ is homeomorphic with $`limS`$, then $`X`$ is Čhech-complete, since each $`P_m`$ is Lindelöf and Čhech-complete.
If $`X`$ is Čhech-complete, then $`X`$ is homeomorphic with $`lim\{X_\alpha ,g_\beta ^\alpha ,E\}`$, where each $`X_\alpha `$ is locally compact and has the embedding into $`𝐋^{𝐧(\alpha )}`$ as the clopen subspace, where $`n(\alpha )𝐍`$ (see also Corollary 5 from Theorem 9 in and Proposition 3.2.17 in ). From Theorems 3.18 and 4.9 it follows that $`X`$ is homeomorphic with $`limS`$, where $`P_m`$ are finite-dimensional over $`𝐋`$ polyhedra.
## 5 Polyhedral expansions and relations between ultrauniform and uniform spaces.
5.1. Theorem. For each ultrauniform space $`(X,𝖴)`$ there exists a uniform locally connected space $`(Y,𝖵)`$ and a continuous quotient mapping $`\theta :XY`$ such that
$`(i)`$ $`\theta (X)=Y`$, $`d(X)=d(Y)`$, $`w(X)=w(Y)`$;
$`(ii)`$ if $`X`$ is dense in itself, then $`Y`$ is dense in itself;
$`(iii)`$ if $`X=_jX_j`$, then $`Y=_jY_j`$;
$`(iv)`$ there exist natural compactifications $`cX`$ and $`c_\theta Y`$ such that $`\theta `$ has a continuous extension on $`cX`$ and $`\theta (cX)=c_\theta Y`$;
$`(v)`$ if $`Z`$ is a subspace of $`X`$, then there exists a subspace $`W`$ of $`Y`$ such that Condition $`(i)`$ is satisfied for $`Z`$ and $`W`$ instead of $`X`$ and $`Y`$;
$`(vi)`$ if $`X`$ has a $`n`$-dimensional polyhedral expansion over $`𝐋`$, then $`Y`$ has a $`n`$-dimensional polyhedral expansion over $`𝐑`$.
Proof. In view of Theorem 2.3.24 the Cantor cube $`D^𝗆`$ has the weight $`𝗆`$, where $`D`$ is the two-point discrete set, $`D^𝗆=_{sS}D_s`$, $`card(S)=𝗆`$, $`D_s=D`$ for each $`sS`$. In view of Corollary 6.2 for each non-Archimedean infinite locally compact field $`𝐋`$ with a non-trivial valuation a ball $`B(𝐋,x,r):=\{y𝐋:|xy|r\}`$ is homeomorphic to $`D^\mathrm{}_0`$. On the other hand, there exists a well-known quotient mapping of $`D^\mathrm{}_0`$ onto $`[0,1]`$. It can be also realized for $`B(𝐋,0,1)`$ and in particular for $`𝐙_𝐩:=B(𝐐_𝐩,0,1)`$, where $`p`$ is a prime. Let $`s`$ be a unit ball in $`c_0(𝐋,\alpha )`$ and $`\theta :B(𝐋,0,1)[0,1]`$ be a quotient mapping. Then this mapping generates a mapping $`\theta :s\theta (s)c_0(𝐑,\alpha )`$, where $`\theta (s)=\{y\theta (x):`$ $`x=(x^j:j\alpha ),|x^j|1,xc_0(𝐋,\alpha )\}`$, hence $`(y^j:j\alpha )=y`$, $`y^j=\theta (x^j)`$ and for each $`ϵ>0`$ we have $`card\{j:|x^j|>ϵ\}<\mathrm{}_0`$, hence $`card\{j:|y^j|>ϵ\}<\mathrm{}_0`$ such that $`y^j[0,1]`$. If $`\{e_j:j\alpha \}`$ is the standard orthonormal basis in $`c_0(𝐑,\alpha )`$, then under the embedding $`c_0(𝐑,\alpha )c_0(𝐑,\alpha )^{}`$ associated with this basis to each $`e_j`$ there corresponds a continuous linear functional $`e^jc_0(𝐑,\alpha )^{}`$, where $`c_0(𝐑,\alpha )^{}`$ is a topological dual space, such that $`e^j(e_i)=\delta _i^j`$. Therefore, $`\theta (s)=\{yc_0(𝐑,\alpha ):`$ $`0e^j(y)1\}`$, which is by definition a simplex in $`c_0(𝐑,\alpha )`$ such that if $`y\theta (s)`$ and $`0<y<1`$, then $`B(c_0(𝐑,\alpha ),y,\mathrm{min}(y,1y))\theta (s).`$
On the other hand, if $`b`$ is a simplex in $`c_0(𝐑,\alpha )`$ and $`bB(c_0(𝐑,\alpha ),0,1)`$, then $`\theta ^1(Int(b))`$ is open in $`c_0(𝐑,\alpha )`$, $`\theta ^1(Int(b))\mathrm{}`$ if $`Int(b)\mathrm{}`$, since $`\theta `$ is continuous. Choose $`\alpha `$ such that $`b`$ is a canonical closed subset in $`c_0(𝐑,\alpha )`$, then $`\theta ^1(b)`$ is closed in $`c_0(𝐋,\alpha )`$. Therefore, $`cl(\theta ^1(Int(b)))=\theta ^1(b)`$, since $`\theta ^1(Int(b))\theta ^1(b)`$ and if $`c`$ is a canonical closed subset in $`[0,1]`$, then $`\theta ^1(c)`$ is a canonical closed subset in $`D^\mathrm{}_0`$.
If $`X=lim\{X_\alpha ,f_\beta ^\alpha ,E\}`$ is a spectral decomposition of an ultrauniform space and $`X_\alpha `$ are polyhedra in $`c_0(𝐋,\gamma _\alpha )`$, then $`\theta `$ generates $`\stackrel{~}{\theta }_\alpha :X_\alpha \stackrel{~}{\theta }_\alpha (X_\alpha ),`$ where $`\stackrel{~}{\theta }_\alpha (X_\alpha )`$ are polyhedra in $`c_0(𝐑,\gamma _\alpha )`$, since each $`X_\alpha `$ can be presented as a union of simplexes of diameter not greater, than $`1`$. Then $`\stackrel{~}{f}_\beta ^\alpha :=\stackrel{~}{\theta }_\beta f_\beta ^\alpha \stackrel{~}{\theta }_\alpha ^1:`$ $`\stackrel{~}{\theta }_\alpha (X_\alpha )\stackrel{~}{\theta }_\beta (X_\beta )=:\stackrel{~}{X}_\beta `$ are continuous bonding mappings. In view of Lemma 2.5.9 the topological space $`\stackrel{~}{\theta }(X)=lim\{\stackrel{~}{X}_\alpha ,\stackrel{~}{f}_\beta ^\alpha ,E\}`$ is locally connected and $`\stackrel{~}{\theta }(X)`$ is the quotient space of $`X`$.
Now we proof that if $`X`$ is given, then $`Y`$ satisfies Conditions $`(ivi)`$. In view of Proposition 2.5.6 are satisfied $`(iii,v)`$. On the other hand, $`dim_𝐋X_\alpha =dim_𝐑\stackrel{~}{X}_\alpha `$ for each $`\alpha E`$, consequently, $`(vi)`$ is fulfilled. Since $`d(X_\alpha )=d(\stackrel{~}{X}_\alpha )`$ and $`w(X_\alpha )=w(\stackrel{~}{X}_\alpha )`$ for each $`\alpha `$, then $`(i)`$ is fulfilled. If $`x`$ is an isolated point in $`X`$, then there exists $`\alpha E`$ such that $`x_\alpha `$ is an isolated point of $`X_\alpha `$, consequently, $`\stackrel{~}{x}_\alpha `$ is an isolated point of $`\stackrel{~}{X}_\alpha `$.
For prooving $`(iv)`$ we consider $`cX`$ equal to some specific compactification $`\stackrel{~}{c}\stackrel{~}{X}`$ of $`\stackrel{~}{X}`$, where $`\stackrel{~}{X}`$ is a completion of the uniform space $`X`$. Therefore, the consideration reduces to the case of a complete $`X`$. In such case $`X`$ has the uniform polyhedral expansion such that conditions $`2.8(i,ii)`$ are satisfied. Then we take $`cX`$ as $`lim\{c_\alpha X_\alpha ,\overline{f}_\beta ^\alpha ,E\}`$, where $`c_\alpha X_\alpha `$ are compact spaces. Since each $`X_\alpha `$ is a uniform polyhedron we take $`c_\alpha X_\alpha `$ as a topological sum of compactifications of disjoint polyhedra $`s`$ in $`X_\alpha `$. Practically we take the Banaschewski compactification of each $`s`$ . If $`s=B(c_0(𝐋,\gamma _\alpha ),x,r)`$, then its Banaschewski compactification $`s^\zeta `$ is isomorphic with $`B(𝐋,0,1)^a`$, where $`a=2^{\gamma _\alpha \mathrm{}_0}`$. Therefore, $`\stackrel{~}{\theta }_\alpha `$ has a continuous extension onto $`c_\alpha X_\alpha =_{sX_\alpha }c_\alpha s`$. Evidently, $`\stackrel{~}{\theta }_\alpha (s)=[0,1]^a`$ is the Stone-Čech compactification of $`\stackrel{~}{\theta }(s)`$.
It remains to prove statements about compactifications of simplexes. Let $`Bco(s)`$ denotes an algebra of clopen subsets of $`s`$. Instead of a ring homomorphism $`\eta :Bco(s)𝐅_\mathrm{𝟐}`$ in the Banaschewski construcion let us consider a ring homomorphism $`\eta :Bco(s)𝐙_𝐩`$, where $`𝐅_\mathrm{𝟐}`$ is a finite field consisting of two elements. Topologically $`𝐙_𝐩`$ is homeomorphic to $`𝐅_{\mathrm{𝟐}}^{}{}_{}{}^{\mathrm{}_0}`$, where $`a=𝖼^{w(s)}=2^{\gamma _\alpha \mathrm{}_0}`$. The family $`\mathsf{\Omega }`$ of all such distinct $`\eta `$, which has $`\sigma `$-additive extensions onto $`Bf(s)`$ has the cardinality $`|\mathsf{\Omega }|=𝖼^{w(s)}`$, where $`w(s)=\gamma _\alpha \mathrm{}_0`$ and $`Bf(s)`$ denotes the Borel $`\sigma `$-field of $`s`$, since $`|𝐙_𝐩|=𝖼`$ and each such $`\eta `$ is characterised by a countable family of atoms $`a_j`$ such that $`lim_j\mathrm{}\eta (a_j)=0`$ (see Ch. 7 ). Therefore, $`s`$ has an embedding into $`𝐙_{𝐩}^{}{}_{}{}^{a}=𝐅_{\mathrm{𝟐}}^{}{}_{}{}^{a}`$. The family $`\mathsf{\Omega }`$ is a subfamily of a family $`\mathsf{{\rm Y}}`$ of all ring homomorphisms $`\eta `$. Another way to realize a totally disconnected compactification $`\gamma s`$ of $`s`$ is in the consideration of the space $`H:=C_b^0(s,𝐐_𝐩)`$ of bounded continuous functions $`f:s𝐐_𝐩`$. To each $`xs`$ there corresponds a continuous linear functional $`x^{}H^{}`$ such that $`x^{}(f)=f(x)`$ and $`x^{}=1`$ due to Hahn-Banach Theorem, since $`𝐐_𝐩`$ is spherically complete, where $`H^{}`$ denotes the Banach space of continuous linear functionals on $`H`$. In view of Alaoglu-Bourbaki Theorem (see Exer. 9.202 ) the unit ball $`U:=B(H^{},0,1)`$ is $`\sigma (H^{},H)`$-compact. Therefore, the closure $`\gamma s`$ of $`\gamma (s)`$ in $`U`$ is compact, where $`\gamma `$ is an embedding of $`s`$ into $`U`$. From this construction we have $`\gamma s=𝐙_{𝐩}^{}{}_{}{}^{a}`$. From Example 9.3.3 it follows that the Stone-Čhech compactification of $`\theta (s)`$ is equal to $`[0,1]^a`$. Then $`s^\zeta =\gamma s`$, since $`d(s^\zeta )=d(s)`$ and due to Hewitt-Marczewski-Pondiczeri Theorem $`𝐙_𝐩^𝐚s^\zeta 𝐙_𝐩^𝐚`$.
5.2. Theorem. For each locally connected uniform space $`(Y,𝖵)`$ there exists an ultrauniform space $`(X,𝖴)`$ and a continuous quotient mapping $`\theta :XY`$ such that
$`(i)`$ $`\theta (X)=Y`$, $`d(X)=d(Y)`$, $`w(X)=w(Y)`$;
$`(ii)`$ if $`Y`$ is dense in itself, then $`X`$ is dense in itself;
$`(iii)`$ if $`Y=_jY_j`$, then $`X=_jX_j`$;
$`(iv)`$ there exist natural compactifications $`cX`$ and $`c_\theta Y`$ such that $`\theta `$ has a continuous extension on $`cX`$ and $`\theta (cX)=c_\theta Y`$;
$`(v)`$ if $`Y`$ has a $`n`$-dimensional polyhedral expansion over $`𝐑`$, then $`X`$ has a $`n`$-dimensional polyhedral expansion over $`𝐋`$.
Proof. From the proof of Theorems 2.4 and 5.1 it follows that there exists a topological embedding of $`Y`$ into a subspace $`Z`$ of $`H^{}`$, where $`H:=\stackrel{~}{C}_b^0(Y,𝐑)`$, such that $`Z`$ is isomorphic with $`c_0(𝐑,ord(Y))`$, where $`ord(Y)`$ is an ordinal associated with $`Y`$, $`Z=cl(sp_𝐑Y)`$, the closure of the linear span of $`Y`$ over $`𝐑`$ is taken in $`H^{}`$ (see also ). Therefore, polyhedra of a polyhedral expansion of $`Y`$ can be realized in $`Z`$. To each simplex $`q`$ in $`c_0(𝐑,ord(Y))`$ there corresponds a simplex $`s`$ in $`c_0(𝐐_𝐩,ord(Y))`$ such that $`q`$ is a quotient space of $`s`$ and $`dim_𝐑q=dim_{𝐐_𝐩}s`$. Let $`\stackrel{~}{Y}`$ be a completion of $`Y`$ and $`Q_\alpha `$ be polyhedra of the polyhedral expansion of $`\stackrel{~}{Y}`$. Then there are uniform polyhedra $`P_\alpha `$ such that $`\stackrel{~}{X}=lim\{P_\alpha ,f_\beta ^\alpha ,E\}`$ and $`\stackrel{~}{\theta }(\stackrel{~}{X})=lim\{Q_\alpha ,\stackrel{~}{f}_\beta ^\alpha ,E\}=\stackrel{~}{Y},`$ where $`\stackrel{~}{f}_\beta ^\alpha :=\stackrel{~}{\theta }_\beta f_\beta ^\alpha \stackrel{~}{\theta }_\alpha ^1:`$ and then analogously to the proof of Theorem 5.1 we get statements of Theorem 5.2, where $`X=\stackrel{~}{\theta }^1(Y)`$.
5.3. Corollary. Let spaces $`X`$ and $`Y`$ be infinite as in Theorems 5.1, 5.2, then for absolute polyhedral expansions $`T=\{P_\alpha ,f_\beta ^\alpha ,𝖴\}`$ of $`X`$ and $`S=\{Q_\alpha ,g_\beta ^\alpha ,𝖵\}`$ of $`Y`$ there is equality $`card(𝖴)=card(𝖵)`$.
Proof. In view of Theorems 5.1 and 5.2 to the inverse mapping system $`T_X=\{P_\alpha ,f_\beta ^\alpha ,𝖴_X\}`$ there corresponds the inverse mapping system $`T_Y=\{Q_\alpha ,\stackrel{~}{f}_\beta ^\alpha ,𝖴_X\}`$ such that $`X=limT_X`$ and $`Y=limT_Y`$, where $`P_\alpha `$ and $`Q_\alpha `$ are polyhedra over $`𝐋`$ and $`𝐑`$ respectively. The bonding mappings $`f_\beta ^\alpha `$ are simplexial, but $`\stackrel{~}{f}_\beta ^\alpha `$ are not simplexial. Using permissible modifications (see §2.8, Lemmas 3.2 and 3.8) we get a polyhedral expansion $`W=\{Z_\alpha ,h_\beta ^\alpha ,F\}`$ of $`Y`$ such that $`h_\beta ^\alpha `$ are simplexial and $`Z_\alpha `$ are polyhedra. Moreover, $`card(F)=card(𝖴_X)\mathrm{}_0=card(𝖴_X)`$, since polyhedra $`Q_\alpha `$ are metrizable. Using sections of the polyhedral expansion for $`Y`$ we get $`S`$ with $`card(𝖴)=card(𝖵)`$. Vice versa, using $`S`$ we get $`S_X`$, then with the help of permissible modifications we construct by $`S_X`$ polyhedral expansion of $`X`$ and using sections of spectrum we get an absolute polyhedral expansion for $`X`$.
5.4. Note. Certainly there is not in general a mapping of $`T`$ into $`S`$ or vice versa for $`T`$ and $`S`$ from §5.3.
Each simplex $`s`$ is clopen in its polyhedra $`P`$ over $`𝐋`$, but in general $`P`$ may be not open in the corresponding Banach space, because may be cases of $`s`$ in $`P`$ such that $`dim_𝐋s<dim_𝐋P=sup_{sP}dim_𝐋s`$. There is the following particular case of spaces $`X`$, when in their polyhedral expansions $`S=\{P_n,f_m^n,E\}`$ polyhedra $`P_n`$ are clopen in the corresponding Banach spaces.
5.5. Let $`H`$ be a locally $`𝐊`$-convex space, where $`𝐊`$ is a non-Archimedean field. Let $`M`$ be a manifold modelled on $`H`$ and $`At(M)=\{(U_j,f_j):`$ $`jA\}`$ be an atlas of $`M`$ such that $`card(A)w(H)`$, where $`f_j:U_jV_j`$ are homeomorphisms, $`U_j`$ are open in $`M`$, $`V_j`$ are open in $`H`$, $`_{jA}U_j=M`$, $`f_if_j^1`$ are continuous on $`f_j(U_iU_j)`$ for each $`U_iU_j\mathrm{}`$. Let $`\stackrel{~}{𝐊}`$, $`\stackrel{~}{H}`$ and $`\stackrel{~}{M}`$ denote completions of $`𝐊`$, $`H`$ and $`M`$ relative to their uniformities.
Theorem. If $`H`$ is infinite-dimensional over $`𝐊`$ or $`\stackrel{~}{𝐊}`$ is not locally compact, then $`M`$ is homeomorphic to the clopen subset of $`H`$.
Proof. Since $`\stackrel{~}{H}`$ is the complete locally $`\stackrel{~}{𝐊}`$-convex space, then $`\stackrel{~}{H}=prlim\{\stackrel{~}{H}_q,\pi _v^q,F\}`$ is a projective limit of Banach spaces $`\stackrel{~}{H}_q`$ over $`\stackrel{~}{𝐊}`$, where $`qF`$, $`F`$ is an ordered set, $`\pi _v^q:\stackrel{~}{H}_q\stackrel{~}{H}_v`$ are linear continuous epimorphisms. Therefore, each clopen subset $`W`$ in $`\stackrel{~}{H}`$ has a decomposition $`W=lim\{W_q,\pi _v^q,F\}`$, where $`W_q=\pi _v^q(W)`$ are clopen in $`\stackrel{~}{H}_q`$. The base of topology of $`\stackrel{~}{M}`$ consists of clopen subsets. If $`WV_j`$, then $`f_j^1(W)`$ has an analogous decomposition. From this and Proposition 2.5.6 it follows, that $`\stackrel{~}{M}=lim\{\stackrel{~}{M}_q,\stackrel{~}{\pi }_v^q,F\}`$, where $`\stackrel{~}{M}_q`$ are manifolds on $`\stackrel{~}{H}_q`$ with continuous bonding mappings between charts of their atlases. If $`H`$ is infinite-dimensional over $`𝐊`$, then each $`\stackrel{~}{H}_q`$ is infinite-dimensional over $`\stackrel{~}{𝐊}`$ . From $`card(A)w(H)`$ it follows, that each $`\stackrel{~}{M}_q`$ has an atlas $`At^{}(\stackrel{~}{M}_q)=\{U_{}^{}{}_{j,q}{}^{};f_{j,q};A_{}^{}{}_{q}{}^{}\}`$ equivalent to $`At(\stackrel{~}{M}_q)`$ such that $`card(A_{}^{}{}_{q}{}^{})w(H_q)=w(\stackrel{~}{H}_q)`$, since $`w(\stackrel{~}{H})=w(H)`$, where $`At(\stackrel{~}{M}_q)`$ is induced by $`At(\stackrel{~}{M})`$ by the quotient mapping $`\stackrel{~}{\pi }_q:\stackrel{~}{M}\stackrel{~}{M}_q`$. In view of Theorem 2 each $`\stackrel{~}{M}_q`$ is homeomorphic to a clopen subset $`\stackrel{~}{S}_q`$ of $`\stackrel{~}{H}_q`$, where $`h_q:\stackrel{~}{M}_q\stackrel{~}{S}_q`$ are homeomorphisms. To each clopen ball $`\stackrel{~}{B}`$ in $`\stackrel{~}{H}_q`$ there corresponds a clopen ball $`B=\stackrel{~}{B}H_q`$ in $`H_q`$, hence $`S_q=\stackrel{~}{S}_qH_q`$ is clopen in $`H_q`$ and $`h_q:M_qS_q`$ is a homeomorphism. Therefore, $`M`$ is homeomorphic to a closed subset $`V`$ of $`H`$, where $`h:MV`$ is a homeomorphism, $`VH`$, $`h=lim\{id,h_q,F\}`$, $`id:FF`$ is the identity mapping. Since each $`h_q`$ is surjective, then $`h`$ is surjective by Lemma 2.5.9 . If $`xM`$, then $`\stackrel{~}{\pi }_q(x)=x_qM_q`$, where $`\pi :HH_q`$ are linear quotient mappings and $`\stackrel{~}{\pi }_q:MM_q`$ are induced quotient mappings. Therefore, each $`xM`$ has a neighbourhood $`\stackrel{~}{\pi }_q^1(Y_q)`$, where $`Y_q`$ is an open neighbourhood of $`x_q`$ in $`M_q`$. Therefore, $`h(M)=V`$ is open in $`H`$.
5.6. Theorem. Let $`M`$ be a manifold satisfying conditions of Theorem 5.5, where $`𝐊`$ is with non-trivial valuation, $`M`$ and $`𝐊`$ are complete relative to their own uniformities. Then $`M`$ has a polyhedral expansion $`S=\{P_n,f_m^n,E\}`$ consisting of polyhedra $`P_n`$ clopen in the corresponding Banach spaces $`H_q`$ over $`𝐊`$.
Proof. We treat the case of non-complete $`𝐊`$ also. Let $`\stackrel{~}{𝐊}`$ be a completion of $`𝐊`$ relative to its uniformity. If $`\stackrel{~}{𝐊}`$ is not locally compact then $`\stackrel{~}{𝐊}`$ contains a proper locally compact subfield $`\stackrel{~}{𝐋}`$ . Since $`M`$ is homeomorphic to a clopen subset of $`H`$, then a completion $`\stackrel{~}{M}`$ of $`M`$ is homeomorphic to a clopen subset in a completion $`\stackrel{~}{H}`$ of $`H`$, where $`\stackrel{~}{H}`$ is over $`\stackrel{~}{𝐊}`$. If $`\stackrel{~}{𝐋}`$ is a locally compact subfield of $`\stackrel{~}{𝐊}`$, then $`\stackrel{~}{𝐋}`$ is spherically complete and $`\stackrel{~}{H}`$ can be considered as the topological vector space over $`\stackrel{~}{𝐋}`$ . By Lemmas and Theorems of §3 we can choose $`P_n`$ clopen in the corresponding Banach spaces over $`\stackrel{~}{𝐋}`$. We also can choose polyhedra over $`\stackrel{~}{𝐊}`$, repeating proofs for $`\stackrel{~}{M}`$ over $`\stackrel{~}{𝐊}`$. In §3 polyhedra over locally compact field were chosen that to encompass cases of compact spaces $`X`$, but polyhedra can be considered also over $`\stackrel{~}{𝐊}`$. For compact $`X`$ and non locally compact $`\stackrel{~}{𝐊}`$ such polyhedral expansions have little sense, since discrete $`P_n`$ with singletons $`s`$ as simplexes may appear such that $`dim_{\stackrel{~}{𝐊}}s=0`$. But for the manifold $`\stackrel{~}{M}`$ due to embedding into $`\stackrel{~}{H}`$ as the clopen subset we get the polyhedral expansion such that $`P_n`$ are non-trivial over $`\stackrel{~}{𝐊}`$, moreover there exists a polyhedral expansion such that $`P_n`$ are clopen in the corresponding Banach spaces $`\stackrel{~}{H}_q`$ for each $`n`$. Using absolute polyhedral expansions constructed with the help of sections of the initial polyhedral expansion of $`\stackrel{~}{M}`$ from §4, we get polyhedral expansions of $`M`$ over $`\stackrel{~}{𝐋}`$ and also over $`\stackrel{~}{𝐊}`$. In a particular case of complete $`𝐋`$ and $`𝐊`$ relative to their uniformities the manifold $`M`$ is complete together with $`H`$, that is $`\stackrel{~}{M}=M`$ and $`\stackrel{~}{H}=H`$, consequently, $`M`$ as the topological space is homeomorphic to $`limS`$.
The author expresses his gratefulness to Professor G. Tironi for his interest in this job and hospitality at Mathematical Department of Trieste University since 24 May until 24 October 1999, also to Professor V.V. Fedorchuk at the Mathematical Department of the Moscow State University for careful reading of the manuscript.
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# Active Nucleus in a Poststarburst Galaxy : KUG 1259+280
## 1 Introduction
Recent optical (Ho, Filippenko, & Sargent 1997a,b) and X-ray surveys (Grupe et al. 1998) are finding that nuclei of many nearby galaxies are active. Such activities manifest themselves through continuum and emission-line properties. Characterizing these properties helps us to measure the luminosity functions and evolution of different sub-types, and understanding the relationship between different types of activities, for example, between nuclear starbursts and active galactic nuclei (AGNs), and thus assist in creating unified models for active galaxies. Weedman (1983) has argued that starburst events in the nuclei of galaxies would evolve rapidly into compact configurations as dynamically distinct entities. Perhaps the compact configuration further collapses to form a single massive black hole generally favored as the central engine for AGNs (Miller 1985). On the other hand, Goncalves, Veron-Cetty & Veron (1999) suggest that perhaps it is nuclear activity that triggers circumnuclear star-formation. However, evidence for either case is fragmented at present. In this connection, it is very important to study AGNs in starburst and poststarburst galaxies. In this paper, we report new X-ray and optical observations of a poststarburst galaxy KUG 1259+280.
The galaxy, known as KUG 1259+280 (Takase & Miyauchi-Isobe 1985), is a poststarburst galaxy that appears to be in a small group of 4 galaxies within an arcmin, and is a member of the Coma cluster, where it is known variously as D-61 (Dressler 1980) or GMP 1681 (Godwin, Metcalfe, & Peach 1983). KUG 1259+280 is the brightest among these 4 galaxies (Takase & Miyauchi-Isobe 1985) with $`B=15.88\mathrm{mag}`$ and has a heliocentric radial velocity of $`7102\pm 50\mathrm{kms}^1`$ (Caldwell & Rose 1997). From ground-based imaging, KUG 1259+280 has been classified as an SO galaxy. However, the high resolution images of KUG 1259+280 taken with the Hubble Space Telescope have revealed this galaxy to be an edge-on disk galaxy with a very bright unresolved nucleus (Caldwell, Rose, & Dendy 1999, hereafter CRD). Two radial dust lanes in KUG 1259+280, to the east of the centre and about $`1\mathrm{}`$ away, have also been reported by CRD. An unusual bulge that is boxy or peanut-shaped and, additionally, has an X morphology has been found by CRD. They also report that the central ($`0.16\mathrm{}`$ or 74 pc for $`H_0=75\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ ) luminosity of KUG 1259+280 is $`M_B=16.3`$ (CRD) which is higher than any of the star clusters found in the merging galaxies studied by Whitmore et al. (1997) or by Carlson et al. (1998) and is also higher than the luminosity of nuclear star clusters found in normal early-type galaxies studied by Lauer et al. (1995). The nucleus is about 0.2 mag bluer than that of the surrounding bulge (CRD). In their spectroscopic study of galaxies in the Coma cluster, CRD have found that KUG 1259+280 is a very strong poststarburst galaxy with poststarburst age of $`0.5\mathrm{Gyr}`$ which is shorter than the typical value of 1 Gyr for similar galaxies in Coma, thus indicating that the starburst in this galaxy had happened more recently as compared to other such galaxies in Coma. Based on the ratio of flux of \[N II\]$`\lambda \lambda 6548,6583`$ and H$`\alpha `$, CRD have concluded that the emission line spectrum of KUG 1259+280 to be arising due to LINER-like nuclear activity.
The galaxy, KUG 1259+280, was first detected as a serendipitous $`EXOSAT`$ X-ray source, EXO 125938+2803.1, in the region of Coma cluster by Branduardi-Raymont et al. (1985). The galaxy was also identified as an X-ray source in the $`EXOSAT`$ high Galactic latitude survey by Giommi et al. (1991). X-ray emission from KUG 1259+280 was again detected during the $`ROSAT`$ pointings and was reported by Singh et al. (1995) to be an ultra-soft X-ray source – WGA J1301.9+2746 (RX J1301.9+2746) (White, Giommi, & Angelini 1994). Based on this association, Singh et al. (1985) obtained soft X-ray luminosity ($`L_X`$) of $``$3 $`\times `$ 10<sup>42</sup> erg s<sup>-1</sup> for KUG 1259+280 which is similar to that of a low-luminosity AGN. The basic parameters of KUG 1259+280 are summarized in Table 1.
In this paper, we present a detailed analysis of X-ray data on the ultra-soft X-ray source – WGA J1301.9+2746 (RX J1301.9+2746), taken from the $`ROSAT`$ archives. The data taken on 5 different days of observations during 1991 to 1995 have not been presented before. We confirm the identification of the X-ray source with KUG 1259+280 based on high resolution X-ray imaging data. We present X-ray light curves and spectra of KUG 1259+280 taken on 3 occasions. We also present our new optical spectroscopic observations of the galaxy, and discuss its peculiar nature.
Throughout the paper, luminosities are calculated assuming a Hubble constant of $`H_0=75\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ and a deceleration parameter of $`q_0=0`$.
## 2 Observations
### 2.1 $`ROSAT`$ X-ray Observations
The region of the sky containing this source was observed three times with the $`ROSAT`$ (Truemper et al. 1983) Position Sensitive Proportional Counter (PSPC) and twice with the High Resolution Imager (HRI) (Pfeffermann et al. 1987). The PSPC observations were carried out during 1991 June 16-19. Coma cluster was the target source in the PSPC observations, whereas KUG 1259+280 itself was targeted in the HRI observations of 1995 June 5. The galaxy was again in the field of view of the HRI observations during 1994 June 25, when NGC 4921 and NGC 4923 were targeted. The details of $`ROSAT`$ observations are given in Table 2. The off-set of KUG 1259+280 from the field centre is also listed in Table 2. ROSAT X-ray data corresponding to the above observations were obtained from the public archives maintained at the High Energy Astrophysics Science Archive Research Center (HEASARC) in USA, and analysed by us. An X-ray source centered on the position of KUG 1259+280 was detected in all the observations. It is offset from the centre of the field of view of the PSPC by $`0.5^{}`$, however, it was was not obstructed by the PSPC window support structure in the observations identified by the sequence numbers RP800006N00 (1991 June 16) and RP800013N00 (1991 June 18). In the 1991 June 17 observation (RP80005N00), the source is somewhat closer to the PSPC window structure, however, it does not appear to be obstructed in any significant way by the window structure even though the point spread function of the PSPC and the telescope system is very broad at this position.
### 2.2 Optical Observations
Optical V band image of KUG 1259+280 was obtained at the Vainu Bappu Telescope (VBT), Kavalur, India on the night of 1996 March 13. The observation was carried out using a $`1024\times 1024`$ CCD chip placed at the prime focus of the 2.3 m reflector. The pixel size of $`24\mu \mathrm{m}\mathrm{pixel}^1`$ of the CCD gives a scale of $`0.61\mathrm{}`$ and a total field of view $`10.4\mathrm{}\times 10.4\mathrm{}`$.
Spectroscopic observations of KUG 1259+280 were carried out under photometric conditions on the night of 2000 Feb 02 at the Cassegrain focus of 2.12-m telescope of the Guillermo Haro Observatory (GHO), Cananea, Mexico. The Boller and Chivens Spectrograph with TEX $`1024\times 1024`$ CCD and $`150\mathrm{l}\mathrm{mm}^1`$ grating was used. The slit position angle was $`48^{}`$ which passes roughly along the major axis. The slit width was $`2\mathrm{}`$. The airmass of the observations was close to 1.0. The instrumental resolution (FWHM) was $`10\mathrm{\AA }`$. Three spectra of the central region of KUG 1259+280 were acquired in order to increase the signal-to-noise ratio and to reduce unwanted cosmic ray hits. The exposure times were $`1800`$ s for each of the three exposures. To correct for the bias level and nonuniform response of the pixels in the CCD, bias frames and dome flat frames were also acquired on the same night. For the purpose of wavelength calibration, comparison (He-AR) spectra were obtained. The spectrophotometric standard star HR 4963 was observed (2 exposures) on the same night to flux calibrate the object spectra.
Spectroscopic observations ($`2`$ exposures of $`2400\mathrm{s}`$ each) of KUG 1259+280 were also carried out on the night of 1999 Feb 26 at the Cassegrain focus of the VBT using OMR spectrograph (Prabhu, Anupama, & Surendirath 1998) with a slit of width $`2.8\mathrm{}`$, $`600\mathrm{l}\mathrm{mm}^1`$ grating and $`1024\times 1024`$ CCD camera. The instrumental resolution (FWHM) was $`7\mathrm{\AA }`$, and the slit position angle was $`0^{}`$. The sky conditions were close to photometric. However, the spectra obtained were of poor signal-to-noise ratio due to poor reflectivity of the primary. Because of the wider wavelength range covered and better signal-to-noise ratio obtained at GHO, we present only the results obtained from the GHO.
## 3 Analysis & Results
### 3.1 Image Analysis
The analysis of X-ray images has been carried out using the PROS<sup>6</sup><sup>6</sup>6The PROS software package provided by the $`ROSAT`$ Science Data Center at Smithsonian Astrophysical Observatory. software package. X-ray images were extracted from all the observations listed in Table 2. The high resolution (central full width half maximum of $`4\mathrm{}`$) images, obtained with HRI, were smoothed by convolving with a Gaussian of $`\sigma `$=$`5\mathrm{}`$. The PSPC images were smoothed with a Gaussian of $`\sigma =10\mathrm{}`$. X-ray contours of smoothed HRI and PSPC images have been overlaid on the optical $`V`$ band image and are shown in Figures 1 and 2, respectively. It is clear from Fig. 1 that the X-rays are centered on the galaxy, KUG 1259+280. No other X-ray source, within the angular spread comparable to the point spread function of $`ROSAT`$ PSPC, was seen in either of the two HRI observations. Therefore, X-ray emission from KUG 1259+280 is not contaminated by emission from any other source.
We have estimated count rate of KUG 1259+280 from the PSPC observations. Due to the large off-axis angle ($`30\mathrm{}`$) of KUG 1259+280 from the field centre in each of the PSPC observations, it is necessary to correct the PSPC images for exposure and vignetting. For this reason, we have created exposure map for each of the PSPC observation using appropriate devignetted detector map obtained from HEASARC. Exposure corrections to the PSPC images were carried out using the XIMAGE software. Source count rates were estimated from square boxes centered at the source position and after subtracting the background estimated from square annuli centered at the source position but outside the source region. The source box size was optimally selected to maximize the signal-to-noise ratio. The count rates thus estimated are corrected for exposure and vignetting and are given in Table 3. The X-ray intensity of KUG 1259+280 remained almost steady during the 3 pointed PSPC observations. X-rays from KUG 1259+280 were also detected in 1990 during the $`ROSAT`$ All Sky Survey (RASS) (Voges et al. 1999). X-ray intensity of KUG 1259+280 obtained from the $`ROSAT`$ All-Sky Survey Bright Source Catalogue (Voges et al. 1999) is given in Table 3, and is about $`25\%`$ less than that during the pointed mode observations in 1991.
The source counts from the HRI images, extracted from the data observed on 1994 June 25 and 1995 June 5, were estimated from circles of radii $`30\mathrm{}`$, and $`25\mathrm{}`$, respectively, after subtracting the background counts estimated from annular regions centered at the peak positions with inner and outer radii of $`60\mathrm{}`$ and $`120\mathrm{}`$. The optimal extraction radii were chosen from radial profiles for the KUG 1259+280 to include the maximum signal and minimum background noise. The background estimates were checked by changing the width of the annuli and no significant change in total counts was noticed. The HRI count rates are given in Table 3. The count rates in the two HRI observations are unchanged.
In order to find the spatial extent of KUG 1259+280 in X-rays, we have generated a radial intensity profile of the X-ray source in KUG 1259+280 from raw $`ROSAT`$ HRI images extracted from the 1995 June 5 data, and compared it with the point response function of the telescope and detector. We used this HRI observation because the source is centered on the axis, even though the exposure time is smaller than in the other HRI observation. This was done because the HRI point response function broadens at off-axis positions. The radial profile of KUG 1259+280 was generated by averaging azimuthally in radial bins of $`2\mathrm{}`$ from the raw HRI images. Background intensity was estimated from an annular region of width $`60\mathrm{}`$ centered at the galaxy, $`2.5\mathrm{}`$ away from the galaxy centre. The radial profile thus generated was normalized by the intensity at a radius of $`1\mathrm{}`$ and is shown in Figure 3. The point response function of the $`ROSAT`$ telescope and HRI was generated using the analytic form of the on-axis point response function (David et al. 1993). As can be seen in Fig. 3, KUG 1259+280 has not been spatially resolved by $`ROSAT`$ HRI. It should be noted that the slight excess over the theoretical PRF around radii of $`5\mathrm{}`$ to $`10\mathrm{}`$ does not imply that the galaxy is spatially extended. The work of Lehmann et al. (1998) has proved that the more realistic re-calibrated HRI PRF differs slightly but significantly from the theoretical HRI PRF in the range $`10\mathrm{}30\mathrm{}`$. Thus we conclude that there is no evidence for extended X-ray emission from KUG 1259+280 based on the $`ROSAT`$ HRI observations.
### 3.2 X-ray Timing
In order to investigate the time variability of the soft X-ray emission from KUG 1259+280, we have extracted light curves from the three $`ROSAT`$ PSPC observations. The light curves were extracted using the ‘xselect’ package with time bins of 400 s and in the full energy band of $`ROSAT`$ PSPC containing all the X-ray photons. For each PSPC observation, source$`+`$background light curve was extracted from a circle of radii $`3.25\mathrm{}`$ centered on the peak position, while background light curve was extracted from three nearby circular regions away from the Coma cluster with their centres about $`6.5\mathrm{}`$ away from the source. The large radius for source extraction circle was necessitated due to the broadening of the point spread function (psf) at large off-axis angle from the centre of the field of view of the telescope. The background subtractions were carried out after appropriately scaling the background light curves to have the same area as the source extraction area. The light curves of KUG 1259+280 thus generated are shown Figures 4(a), 4(b), and 4(c) for the PSPC observations carried out on 1991 June 16, 17, and 18 respectively. As can be seen in the Figures 4(a) and 4(b), X-ray emission from KUG 1259+280 is steady during the observations on 1991 June 16 and 17. However, the galaxy shows a remarkable X-ray variability on 1991 June 18 when an event with peak X-ray intensity of 2.5 times the average value that subsequently declined to the average value within about 1300 s, was recorded. It should be noted that the variability amplitude and the time scale could be larger than seen, since the observation does not cover the full time span of the variable event due to Earth occultation. After the decline, X-ray intensity from KUG 1259+280 remains almost steady. To investigate any change in the hardness ratio during the flare or outburst, we generated two light curves of KUG 1259+280 in the energy bands 0.1–0.3 keV and 0.3–1.5 keV. From these light curves, we calculated hardness ratio as a function of time and found no changes in the hardness ratio during the outburst or flare. Since most of the X-ray photons are concentrated below 0.5 keV, we confirmed the above result by calculating hardness ratios using light curves extracted in the energy bands 0.1–0.35 keV, 0.35–1.5 keV and 0.1–0.4, 0.4–1.5 keV.
In order to find out any possibility of the presence of long-term variation of soft X-ray emission from KUG 1259+280, we have searched the literature for observations of the galaxy with other X-ray observatories. The galaxy was not observed with the Einstein observatory (Branduardi-Raymont et al. 1985). The $`EXOSAT`$ LE count rate of KUG 1259+280 observed on 1985 June 17 is $`7.5\times 10^3\mathrm{count}\mathrm{s}^1`$ (Giommi et al. 1991). The $`EXOSAT`$ LE count rate of KUG 1259+280 corresponds to an observed flux of $`4\times 10^{13}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$ in the energy band of 0.05-2.0 keV using the spectral parameters measured with the $`ROSAT`$ PSPC ($`\mathrm{N}_\mathrm{H}\text{ }=1.4\times 10^{20}\mathrm{cm}^2`$, $`\mathrm{\Gamma }_X=4.25`$) (see $`\mathrm{\S }3.3`$). Using W3PIMMS, this flux is found to be equivalent $`ROSAT`$ PSPC count rate of $`0.12`$, which is similar to that obtained from the RASS in 1990 but smaller by $`25\%`$ than during the pointed observation in 1991. We have also converted the $`ROSAT`$ HRI count rates of KUG 1259+280 to the PSPC count rates using the best-fit spectral parameters, for the purpose of comparison. The equivalent PSPC count rates are $`(15.2\pm 1.5)\times 10^2\mathrm{count}\mathrm{s}^1`$ and $`(15.1\pm 0.8)\times 10^2\mathrm{count}\mathrm{s}^1`$, for the HRI observations carried out on 1995 June 6 and 1994 June 25, respectively. Thus X-ray emission from KUG 1259+280 is almost steady during the pointed PSPC and HRI observations. In summary, X-ray intensity of KUG 1259+280 was same in 1990 observations as in the 1985 observations, and then increased by $`25\%`$ in 1991.
### 3.3 X-ray Spectra
Photon energy spectra were accumulated from the PSPC events, separately for the three observations shown in Table 2, and using the same source and background regions as stated above. The $`ROSAT`$ PSPC pulse height data obtained in 256 pulse height channels were re-grouped by adding counts in 8 consecutive channels to improve the statistics. The X-ray spectra from the three observations thus obtained are shown together in the upper panel of Figure 5. Bad channels have been ignored.
We used the XSPEC (Version 10.0) spectral analysis package to fit the data with spectral models. This requires a knowledge of the response of the telescope and the detector. Using the source spectral files and information about the source off-axis angle contained within them, we generated an auxiliary response file of the effective area of the telescope, by using the program $`pcarf`$ in the FTOOLS (version 4.2) software package. This program uses the available off-axis calibration of the telescope, and the amount of scattering of X-ray photons and their dependence on energy are taken into account. Appropriate response matrix, provided by the ROSAT GOF at HEASARC, was used to define the energy response of the PSPC. The $`ROSAT`$ PSPC spectra, shown in Fig. 5, were used for fitting spectral models. The spectra from the three observations were first fitted separately with a redshifted power-law model with photon index, $`\mathrm{\Gamma }_X`$, and absorption due to an intervening medium with absorption cross-sections as given by Balucinska-Church and McCammon (1992) and using the method of $`\chi ^2`$-minimization. The results of this fitting and the best-fit spectral model parameters are shown in Table 4. This simple model is a good fit to the data of 3 days as evidenced by the minimum reduced $`\chi ^2`$ values ($`\chi _\nu ^2`$). The spectra corresponding to the first and the second observations are consistent with each other in the sense that both the best-fit photon indices and the column densities of the absorbing material are similar within the errors. The absorber column densities are similar (within errors) to the 21-cm value ($`\mathrm{N}_\mathrm{H}`$ =9.5$`\times 10^{19}`$$`\mathrm{cm}^2`$ ) measured from radio observations in this direction (Dickey & Lockman 1990), indicating that all the X-ray absorption is due to matter in our own Galaxy. Therefore, we have also fitted the power-law models to these spectra after fixing the neutral hydrogen column density to the Galactic value. The best-fit spectral parameters, in this case, are also given in Table 4. The best-fit minimum $`\chi _\nu ^2`$ and the power-law index do not change significantly from those obtained while varying the $`N_H`$, as can be seen in Table 4. The photon indices thus obtained are, however, better constrained and are $`3.80_{0.24}^{+0.14}`$ and $`4.0_{0.3}^{+0.2}`$ for the spectra observed on 1991 June 16 and 1991 June 17, respectively. The errors quoted, here and below, were calculated at the 90% confidence level based on $`\chi _{\mathrm{min}}^2`$+2.71. Based on the best fit model parameters, the observed X-ray flux from the source is estimated to be $`6.3\times `$10<sup>-13</sup> erg cm<sup>-2</sup> s<sup>-1</sup> for the observation on 1991 June 16, and $`5.6\times `$10<sup>-13</sup> erg cm<sup>-2</sup> s<sup>-1</sup> for the observation on 1991 June 17. We have also estimated the intrinsic flux of KUG 1259+290 by setting $`\mathrm{N}_\mathrm{H}`$ to zero. The intrinsic flux of KUG 1259+280 in the energy band 0.1–2.0 keV is $`2.4\times `$10$`{}_{}{}^{12}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$ on 1991 June 16, and $`2.3\times 10^{12}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$ on 1991 June 17.
Power-law model fits to the spectral data obtained on 1991 June 18, however, show significant excess absorption over and above the Galactic value (see Table 4) when $`\mathrm{N}_\mathrm{H}`$ is kept variable. The best-fit value of $`\mathrm{N}_\mathrm{H}`$ is about twice of that of the Galactic value thus suggesting an excess absorption local to the X-ray source in KUG 1259+280. The best-fit photon index is slightly steeper than those obtained previously (see Table 4). In order to verify the presence of excess absorption, we fixed the neutral hydrogen column density to the Galactic value and fitted the absorbed power-law model. Although the photon index becomes flatter, the fit is not acceptable ($`\chi _\nu ^2=2.427`$ for 23 degrees of freedom) thus implying the need for an excess absorption. To estimate this excess absorption, we introduced an additional component for absorption local to KUG 1259+280 in our model, which makes the fit acceptable ($`\chi _\nu ^2=0.979`$ for 22 degrees of freedom) and the best-fit value of excess (local to the source) $`N_H`$ is $`1.0_{0.3}^{+0.5}\times 10^{20}\mathrm{cm}^2`$. The best-fit photon index in this case is $`4.6_{0.3}^{+0.6}`$. The best-fit power-law is steeper on June 18 than on June 16 & 17 (see Table 4). The absorbed X-ray flux of KUG 1259+280 in the energy band 0.1–2.0 keV is $`6.5\times 10^{13}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$ on 1991 June 18, and the corresponding intrinsic X-ray flux is calculated to be $`7.3\times 10^{12}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$, based on the best-fit parameters.
In order to further investigate whether the differences in photon indices and $`\mathrm{N}_\mathrm{H}`$ obtained for the three data sets are indeed real, we have fitted the absorbed and redshifted power-law model simultaneously to all the three data sets. In this joint fit, the photon indices for the three observations were tied together and varied to obtain the best-fit. Similarly, the $`\mathrm{N}_\mathrm{H}`$ and normalizations were also tied together and varied. The best-fit value of minimum $`\chi _\nu ^2`$ is 1.119 for 86 degrees of freedom which corresponds to a probability of 0.21. The best-fit value of the common photon index is $`4.25_{0.22}^{+0.25}`$ and $`\mathrm{N}_\mathrm{H}`$ =$`1.4_{0.2}^{+0.2}\times 10^{20}\mathrm{cm}^2`$ (see Table 4). These results indicate that the shape of the spectrum of KUG 1259+280 has probably not changed over the three observations. However, $`\mathrm{N}_\mathrm{H}`$ column obtained is in excess of the Galactic value suggesting local absorption in all three observation. To find the excess $`\mathrm{N}_\mathrm{H}`$ , we have carried out joint-fit after introducing additional absorption local to KUG 1259+280 apart from the fixed Galactic absorption. The excess $`\mathrm{N}_\mathrm{H}`$ was found to be $`5_2^{+2}\times 10^{19}\mathrm{cm}^2`$. All other parameters do not change from that obtained in the previous joint-fit. The fitted model is shown in Fig. 5, and the confidence contours of excess $`\mathrm{N}_\mathrm{H}`$ and photon index are shown in Figure 6. A photon index ($`\mathrm{\Gamma }_X`$) between 3.8-4.8 and an excess absorption between $`(0.21.0)\times 10^{20}\mathrm{cm}^2`$ are indicated for all three observations. The observed X-ray flux of KUG 1259+280 in the energy band 0.1–2.0 keV is $`6.1\times 10^{13}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$, and the corresponding rest frame intrinsic X-ray luminosity is $`4.7\times 10^{42}\mathrm{erg}\mathrm{s}^1`$.
We have also fitted absorbed blackbody model to all three data sets jointly. The absorbing column was fixed to the Galactic value for all three data sets. The common blackbody temperature was varied to obtain the best-fit. Although the fit is not acceptable ($`\chi _\nu ^2=1.737`$ for 89 degrees of freedom), the temperature ($`kT=0.055\mathrm{keV}`$) thus inferred reflects the ultra soft nature of KUG 1259+280.
### 3.4 Optical Spectroscopy
Optical spectra were reduced and analysed using the IRAF<sup>7</sup><sup>7</sup>7IRAF is distributed by the National Optical Astronomy Observatories, which is operated by the Association of Universities, Inc. (AURA) under cooperative agreement with the National Science Foundation. The IRAF version $`\mathrm{2.11.3}`$ was used. software package. Corrections for bias level and response variation of pixels were carried out using the standard procedures for each of the three spectra. Three different spectra of KUG 1259+280 were then combined to increase the signal-to-noise ratio. One dimensional spectrum of KUG 1259+280 was then optimally extracted from the two dimensional spectrum using an aperture of size $`4.2\mathrm{}`$ by using the method of Horne (1986). The one dimensional object spectrum was wavelength calibrated using appropriate arc spectrum (He-AR). The wavelength calibration was accurate to $`2\AA `$ in all the spectra. Correction for atmospheric extinction was based on the extinction curve for Kitt Peak National Observatory. However, the correction is not critical as the standard star was observed almost at similar airmass as that of the object. Flux calibration was performed using standard $`IRAF`$ procedures and using the observed spectrophotometric standard star HR 4963. The flux calibrated spectrum of KUG 1259+280 was corrected for Galactic extinction by adopting the extinction law in Cardelli et al. (1989). Color excess $`E_{BV}`$ due to Galactic reddening along the line of sight of the galaxy KUG 1259+280 was calculated from the neutral hydrogen column density $`\mathrm{N}_\mathrm{H}`$ using the relation
$$\mathrm{N}_\mathrm{H}\text{ }=5.8\times 10^{21}\times E_{BV}\mathrm{cm}^2$$
(1)
(Bohlin, Savage, & Drake 1978). The dereddened spectrum was then corrected for Doppler shift (z=0.02369). The final reduced spectrum of KUG 1259+280 is shown in Figure 7. The signal-to-noise ratio of the spectrum is $`52`$ measured at 4500 Å– 6600 Å. The spectrum of KUG 1259+280 shows forbidden emission lines \[N II\]$`\lambda 6548`$, \[N II\]$`\lambda 6583`$, \[O III\]$`\lambda 4959`$, \[O III\]$`\lambda 5007`$. The Balmer lines H$`\beta `$, H$`\gamma `$, and H$`\delta `$ are seen in the absorption. No H$`\alpha `$ emission is observed in the spectrum. The absorption lines of Mg I ($`\lambda 5176`$), Fe ($`\lambda \lambda 5270,5335`$), Na I D ($`\lambda 5890`$), and Ca II+Ba II blend ($`\lambda 6497`$) are also present in the spectrum of KUG 1259+280.
We have fitted Gaussian profiles to the emission lines \[N II\]$`\lambda 6583`$, \[N II\]$`\lambda 6548`$, \[O III\]$`\lambda 4959`$, and \[O III\]$`\lambda 5007`$ using the profile fitting feature in ‘splot’ task within IRAF. The spectral line parameters obtained from the fit are shown in Table 5. The error quoted for each parameter is the absolute deviation containing $`68.3\%`$ of the parameter and were estimated by 100 Monte-Carlo simulations assuming constant Gaussian noise. The FWHM of \[O III\]$`\lambda 4959`$ line is about twice of the FWHM of the \[O III\]$`\lambda 5007`$ line and both have similar flux. This feature is not an artifact and is confirmed by spectrum of the same object taken with the VBT, India but with a poorer signal-to-noise ratio ($`30`$ at $`5000\mathrm{\AA }`$), albeit slightly higher resolution ($`7\mathrm{\AA }`$). This is not expected and can be attributed to the fact that the forbidden line \[O III\]$`\lambda 4959`$ is blended with the Fe II lines at $`4924\mathrm{\AA }`$ and $`4970\mathrm{\AA }`$. The FWHM of the forbidden lines \[N II\]$`\lambda 6583`$, \[N II\]$`\lambda 6548`$, and \[O III\]$`\lambda 5007`$ are similar to the instrumental resolution deduced from the night sky emission lines. Thus the forbidden lines have not been resolved and have FWHMs $`600\mathrm{km}\mathrm{s}^1`$. VBT spectrum show FWHM $`300\mathrm{km}\mathrm{s}^1`$ after removing the effect of instrumental resolution. There is some indication for the presence of Fe II blends between 4435–4700$`\mathrm{\AA }`$ in the blue, and between 5070–5600$`\mathrm{\AA }`$. However, the blends seem to be modified by the galaxy absorption features e. g. by the presence of Fe $`\lambda 5335`$ absorption line usually seen in normal galaxy spectrum (Goudfrooij 1998). To confirm the possible presence of Fe II blends higher resolution and higher signal-to-noise spectrum is required. Subtraction of a template galaxy spectrum representing the stellar population of KUG1250+280 from the observed spectrum will help in confirming the presence of Fe II blends. However, this requires data with better signal-to-noise ratio and higher resolution than the data presented here
## 4 Discussion
Recent optical spectroscopic surveys have been finding a large fraction ($``$43%) of active galactic nuclei among nearby galaxies (Ho, Filippenko, & Sargent 1997a). Nearly 10% of all nearby galaxies surveyed by Ho et al. 1997a are found to have Seyfert type nuclei, the rest being LINERS (low ionization nuclear emission region galaxies) or composite LINER/H II nucleus systems. While the optical emission lines are sensitive to many physical parameters e.g., the shape and strength of ionizing continuum arising from nucleus, temperature, density, composition and velocities of line emitting clouds surrounding the nucleus, and dust in the ambient medium, X-rays probe the innermost regions of AGNs. Therefore, X-ray emission properties are more fundamental to infer the nature of galactic nuclei. X-ray emission from KUG 1259+280 is found to be unresolved with the $`ROSAT`$ HRI observations with a resolution of $`4\mathrm{}`$. This suggests that the X-ray emission from KUG 1259+280 is perhaps mostly of nuclear origin. The nuclear nature of X-ray emission is further strengthened by the observation of strong and rapid variability detected in the X-ray observations which also points to a very active nucleus – an AGN inside KUG 1259+280. This is consistent with the point-like appearance of the central region of KUG 1259+280 in the optical band in the $`HST`$ observations with $`0.1\mathrm{}`$ resolution (CRD). In the following, we discuss the observed variability, X-ray and optical spectral characteristics, and broad-band multiwavelength spectra and their implications on models of AGN.
### 4.1 X-ray Characteristics & Variability
The rest frame intrinsic X-ray luminosity of KUG 1259+280 is in the range of $`3.64.7\times 10^{42}\mathrm{erg}\mathrm{s}^1`$. This is about a factor of 10 higher than the X-ray luminosity of low ionization nuclear emission region galaxies (LINERs) studied by Komossa, B$`\ddot{o}`$hringer, & Huchra (1990) using $`ROSAT`$. All the LINERs in the sample of Komossa et al. (1990) have $`L_X10^{41}\mathrm{erg}\mathrm{s}^1`$ in the energy band 0.1–2.4 keV. On the other hand, Seyfert 1 galaxies studied by Rush et al. (1996) span over 4 orders of magnitude in soft X-ray luminosity, from below $`10^{42}\mathrm{erg}\mathrm{s}^1`$ to above $`10^{46}\mathrm{erg}\mathrm{s}^1`$ in the same energy band. Thus X-ray luminosity of KUG 1259+280 is similar to that of a low luminosity Seyfert 1 galaxy.
X-ray emission from KUG 1259+280 is highly variable. On 1991 June 18 X-ray emission from KUG 1259+280 varied by a factor of $`2.5`$ within $``$ 1300 s (Fig. 4). This variability could be due to an outburst or a flare event in the nucleus of KUG 1259+280. This kind of time variability is seen in several AGNs e.g., the low luminosity Seyfert 1 galaxies such as NGC 4051, and narrow line Seyfert 1 galaxies such as IRAS 13224-3809 vary over periods of $`0.11\mathrm{hr}`$ (Pounds 1979; Lawrence et al. 1985; Boller et al. 1993; Singh 1999). On the other hand, rapid variability on the time scales less than a day is generally absent in the case of LINERs (Ptak et al. 1998, Pellegrini et al. 2000). X-ray emission from KUG 1259+280 is also found to vary by $`25\%`$ within a year between RASS and PSPC pointed observations. The large amplitude variation of the nuclear X-ray source in a short time scale is possible only if the nucleus is very compact. Simple light travel time arguments imply that nuclear X-ray source in KUG 1259+280 is very compact ($`size1.2\times 10^5\mathrm{pc}`$). Assuming that the soft X-ray excess emission originates from a few Schwarzschild radii ($`3R_S`$ or more) away from the central massive object, the mass of the central compact object is expected to be $`4.4\times 10^7M_{\mathrm{}}`$. A lower limit of the mass of the central object can be obtained from the Eddington luminosity. Since KUG 1259+280 is an ultra-soft X-ray source, it will be reasonable to assume that bulk of the luminosity for the nuclear source in KUG 1259+280 is in the soft X-ray band. Assuming that the Eddington luminosity is about a factor of 5 larger than the soft X-ray luminosity in the energy band of $`0.12.0\mathrm{keV}`$, the mass of the central object is inferred to be at least $`2.0\times 10^5M_{\mathrm{}}`$. Thus the mass of the central compact object is in the range of $`M10^510^7M_{\mathrm{}}`$. The X-ray luminosity of $`4\times 10^{42}\mathrm{erg}\mathrm{s}^1`$, small amplitude long term variation, and large amplitude (by a factor of $`2.5`$) short term variation all indicate that the galaxy harbours an active nucleus.
The X-ray spectra of KUG 1259+280 obtained from three $`ROSAT`$ PSPC observations are very steep (Table 4). A joint-fit has revealed that all three PSPC spectra are well represented by a single power-law of $`\mathrm{\Gamma }_X4.25`$, with an excess of $`\mathrm{N}_\mathrm{H}\text{ }=5\times 10^{19}`$$`\mathrm{cm}^2`$ over the Galactic value. This suggests the presence of an excess absorption above the Galactic value during all the three observations. Such an absorber is seen convincingly in the observation on 1991 June 18. It is possible that a change in the absorbing column density local to KUG 1259+280 could have occurred within $`1\mathrm{d}`$ as indicated by the individual power-law fits to the PSPC spectra, indicating the presence of moving clouds close to the central nuclear source. The physical state (warm or cold) of the absorbing medium is unclear from the present data. However, in the spectrum of 1991 June 18, there is an indication of an absorption edge at about 0.8 keV (see residuals in Fig. 5), which is very close to O VII and O VIII edges at 0.74 and 0.87 keV commonly seen in $`ASCA`$ and $`ROSAT`$ spectra of AGNs (Reynolds 1997). But because of the poor signal-to-noise ratio of the data and due to the poor energy resolution of the $`ROSAT`$ PSPC, it is not possible to fit an absorption edge. Higher signal-to-noise data with better energy resolution will be required to confirm the presence of an absorption edge. The short time scale of $`1\mathrm{d}`$ for the change in local $`\mathrm{N}_\mathrm{H}`$ also implies that the nuclear X-ray source is very compact, similar to those found in AGNs.
### 4.2 Optical Characteristics
Optical spectrum of nuclear region of KUG 1259+280 is not typical of either the spectra of normal galaxies or AGNs. In the optical spectrum of KUG 1259+280, Balmer emission lines are absent, instead, H$`\beta `$, H$`\gamma `$, and H$`\delta `$ absorption lines are seen. H$`\alpha `$ absorption appears to be almost balanced by emission as was also pointed out by CRD. Caldwell & Rose (1997) have shown that KUG 1259+280 is a very strong poststarburst galaxy in the Coma cluster, based on the presence of enhanced Balmer absorption lines, H$`\delta `$, H$`\gamma `$, etc. and weaker Ca II K line compared to those in normal galaxies. CRD have pointed out that the weak emission line spectrum of KUG 1259+280 is due to nuclear activity and not due to residual star formation. From the equivalent widths of higher order Balmer absorption lines in the spectrum of KUG 1259, CRD have estimated the underlying H$`\alpha `$ absorption to be $`5\mathrm{\AA }`$ and found the ratio of \[N II\]$`\lambda \lambda 6548,6583`$ to H$`\alpha `$ emission to be slightly greater than 1, which is typical value of LINER emission line spectra. It is, however, well known that LINER, Seyfert galaxy, and H II region spectra cannot be unambiguously distinguished from one another on the basis of any single intensity ratio from any pair of lines. The various types of objects with superficially similar emission line spectra can be distinguished by using a diagnostic diagram based on pairs of emission line intensity ratios (Baldwin, Phillips, & Terlevich 1981). In the case of KUG 1259+280, weak Balmer emission lines arising from the centre are dominated by the enhanced absorption lines arising from young stellar populations. Assuming Case B recombination (e.g. Osterbrock 1989), we estimate the H$`\beta `$ emission flux to be $`2.1\times 10^{15}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$ from the equivalent width of H$`\alpha `$ emission line estimated by CRD. Hence the ratio of \[O III\]$`\lambda 5007`$ to H$`\beta `$ emission is about 2 for KUG 1259+280. Ho, Filippenko, & Sargent 1997b have defined LINER nuclei to be those with $`\frac{[OIII]\lambda 5007}{H\beta }<3.0`$, and $`\frac{[NII]\lambda 6583}{H\alpha }0.6`$. Based on this definition, the estimated pair of line intensity ratios, $`\frac{[OIII]\lambda 5007}{H\beta }2`$ and $`\frac{[NII]\lambda 6583}{H\alpha }0.8`$, suggest that nucleus of KUG 1259+280 is, most likely, of LINER type. This confirms the conclusion by CRD. However, it should be noted that LINERs show the presence of strong \[O I\]$`\lambda 6300`$ forbidden line in their spectra, but this line is not detected in the spectrum of KUG 1259+280. Also, in the spectrum of KUG 1259+280, the forbidden line \[O III\]$`\lambda 4959`$ is blended with Fe II $`\lambda \lambda 4924,4970`$. The Fe II lines are seen in the spectra of Seyfert galaxies and not in the spectra of LINERs. These raise question regarding the LINER nature of KUG 1259+280. The lack of Balmer emission lines in the spectrum of KUG 1259+280 can be attributed to the poststarburst phenomena. To investigate quantitatively the origin of optical emission and absorption lines, higher signal-to-noise and higher resolution optical spectrum as well as UV spectrum is required.
### 4.3 Comparison with Steep Spectrum X-ray AGNs
It is well known that the steep spectrum AGNs generally show optical spectra similar to those of optically identified NLS1 galaxies (Boller, Brandt & Fink 1996). In view of the steep X-ray spectrum of KUG 1259+280, if we assume that the intrinsic nuclear optical spectrum is similar to those of NLS1 galaxies modified by enhanced stellar absorption, then one would expect that the optical spectrum parameters, which are not affected significantly by the stellar absorption, should be similar to those of NLS1 galaxies. The interesting parameters are the ratio of X-ray to \[O III\] luminosity, $`\frac{f_x}{[OIII]}`$, and slope ($`\alpha _{opt}`$) of the optical spectrum ($`f_\nu \nu ^\alpha `$). Grupe et al. (1998) and Grupe et al. (1999) have studied soft X-ray and optical properties of 76 steep X-ray spectrum AGNs in detail. Based on their measurements, we have calculated $`\frac{f_x}{[OIII]}`$ for 71 AGNs and found that the ratio, $`\frac{f_x}{[OIII]}`$, spans a large range (6–3157) with a median value of 399.2. (We have excluded the two objects RXJ0136-35 and a Seyfert 2 galaxy, IC 3599, with $`\frac{f_x}{[OIII]}>15000`$ from our analysis. There are only 8 objects with $`\frac{f_x}{[OIII]}>2000`$, and only 2 with $`\frac{f_x}{[OIII]}>3200`$). For KUG 1259+280, $`\frac{f_x}{[OIII]}=156`$. Thus the ratio, $`\frac{f_x}{[OIII]}`$, for KUG 1259+280 is not different from those of steep spectrum X-ray AGNs. Since the optical continuum emission of the nucleus of KUG 1259+280 is contaminated by the stellar emission as can be seen by the presence of strong absorption lines, we have not compared the slope of the observed optical spectrum to those of NLS1 from Grupe et al. (1998). The \[O III\]$`\lambda 5007`$ luminosity of KUG 1259+280 is $`4.8\times 10^{39}\mathrm{erg}\mathrm{s}^1`$, which is about a factor of about 90 smaller than the mean \[O III\]$`\lambda 5007`$ luminosity of the 3 NLS1s (MS 2340-15, RX J0439-45, and RX J2144-39) of Grupe et al. having similar soft X-ray slopes as for the KUG 1259. This could be due to a dilution of X-ray emission in KUG 1259+280, as the soft X-ray luminosity of KUG 1259+280 is smaller by a factor of $`750`$ than the average soft X-ray luminosity of the 3 AGN in the sample of Grupe et al. (1998). The difference in the ratio of \[O III\]$`\lambda 5007`$ to soft X-ray luminosity between KUG 1259+280 and the 3 AGNs with similar soft X-rays slopes could be either due to the contribution of \[O III\] luminosity by H II regions due to hot stars or due to differences in the physical parameters of the line emitting clouds, or differences in the shape of the primary radiation in the UV-EUV region. Based on the above comparison, it is likely that KUG 1259+280 is a low luminosity version of NLS1 galaxies, where due to the diluted power-law continuum, the optical emission lines are weak. Furthermore, the weak permitted lines are heavily modified by the enhanced stellar absorption in the poststarburst region.
### 4.4 Spectral Energy Distribution of KUG 1259+280
In Figure 8, we have plotted far-infrared, optical and soft X-ray flux of KUG 1259+280 obtained from various observations. No far-infrared measurements have been reported from KUG 1259+280 in the IRAS Point Source Catalogue (Joint IRAS Working Group 1985). The IRAS satellite had, however, scanned the sky containing KUG 1259+280 during the surveys. We have obtained the IRAS flux after co-adding the available scan data. The IRAS data were processed using the facilities of IPAC<sup>8</sup><sup>8</sup>8The IPAC is founded by NASA as part of the IRAS extended mission program under contract to JPL.. The measured IRAS flux densities are $`0.07\pm 0.025\mathrm{Jy}`$, $`0.10\pm 0.047\mathrm{Jy}`$, $`0.08\pm 0.024\mathrm{Jy}`$, and $`0.41\pm 0.105\mathrm{Jy}`$ at $`12`$, $`25`$, $`60`$, and $`100\mu `$, respectively. No radio emission is detected within $`15\mathrm{}`$ of KUG 1259+280 in the NRAO/VLA and FIRST sky surveys. The spectral energy distribution of KUG 1259+280 is contaminated by stellar emission in the optical region, otherwise it is similar in shape with those of narrow line Seyfert 1 galaxies reported by Grupe et al. (1999) and indicates a probable shift in the position of the big blue bump, although not covered in any observation, towards higher energies.
### 4.5 AGN–poststarburst connection
As noted earlier, KUG 1259+280 is a nuclear poststarburst galaxy with a poststarburst age of $`0.5\mathrm{Gyr}`$ (CRD). To build a central massive object of mass $`10^6\mathrm{M}_{\mathrm{}}`$ within KUG 1259+280 after the starburst event, the required average mass accretion rate would be only about $`0.002\mathrm{M}_{\mathrm{}}\mathrm{yr}^1`$. On the other hand, the Eddington accretion rate for efficiency factor $`\eta 0.1`$, and $`M=10^6\mathrm{M}_{\mathrm{}}`$ is $`\dot{M_E}0.022\mathrm{M}_{\mathrm{}}\mathrm{yr}^1`$. Therefore, the required mass accretion rate to build a central massive object of mass $`10^6\mathrm{M}_{\mathrm{}}`$ within KUG 1259+280 is about a factor of 10 less than the Eddington accretion rate. Thus it cannot be ruled out that the birth of the AGN in KUG 1259+280 took place after or during the starburst event. In such a scenario, the mass accretion rate onto the young AGN can be expected to be high as the fuel available would be enormous. It should be noted that some authors (e.g. Brandt & Boller 1998) have argued that a high mass accretion rate onto a low-mass ($`10^610^7`$) black hole can result in AGNs with steep soft X-ray spectrum.
## 5 CONCLUSIONS
(i) X-ray emission from KUG 1259+280 is found to be unresolved by $`ROSAT`$ HRI observations suggesting a point-like nuclear X-ray source in it.
(ii) X-ray luminosity of KUG 1259+280 is estimated to be $`3.64.7\times 10^{42}\mathrm{erg}\mathrm{s}^1`$ which is similar to that of low luminosity Seyfert nuclei.
(iii) Variability by a factor of $`2.5`$ within $`1300\mathrm{s}`$ is detected in X-ray emission from KUG 1259+280. Short time scale of variability strongly suggests an AGN in the galaxy.
(iv) The mass of the central massive object in KUG 1259+280 has been estimated to be in the range $`10^510^7M_{\mathrm{}}`$.
(v) X-ray spectrum of KUG 1259+280 is found to be very steep ($`\mathrm{\Gamma }_X4.25`$) in the $`ROSAT`$ band 0.1–2.4 keV. Individual power-law model fits to 3 PSPC spectra indicate a change in intrinsic absorbing column within about a day, while a joint-fit of three PSPC spectra indicates excess absorption (total $`\mathrm{N}_\mathrm{H}\text{ }=1.4\times 10^{20}`$$`\mathrm{cm}^2`$ ) over the Galactic value ($`\mathrm{N}_\mathrm{H}\text{ }=9.5\times 10^{19}`$$`\mathrm{cm}^2`$ ).
(vi) The observed optical spectrum of the nuclear region of KUG 1259+280 is markedly different from that of any class of AGNs. The observed spectrum could be understood if the NLR/BLR emission line spectrum has been substantially modified by enhanced stellar absorption lines. Estimated absorption-free Balmer line strengths and observed forbidden line strengths point towards LINER nature of KUG 1259+280. However, the nature of X-ray emission points towards NLS1 nature of this galaxy. Presence of Fe II lines supports this picture.
(vii) It is quite possible that the central super-massive object in KUG 1259+280 was formed after or during the starburst event $`0.5\mathrm{Gyr}`$ ago.
## 6 Acknowledgments
We thank our referee D. Grupe for his comments and suggestions. This research has made use of $`ROSAT`$ archival data obtained through the High Energy Astrophysics Science Archive Research Center, HEASARC, Online Service, provided by the NASA-Goddard Space Flight Center. The HEASARC database, the Simbad database, operated at CDS, Strasbourg, France; and the NASA/IPAC Extragalactic Database (NED) which is operated by the Jet Propulsion Laboratory, Caltech, under contract with the NASA were used to identify some of the objects. The PROS software package provided by the ROSAT Science Data Center at Smithsonian Astrophysical Observatory, and the FTOOLS software package provided by the High Energy Astrophysics Science Archive Research Center (HEASARC) of NASA’s Goddard Space Flight Center have been used. We thank Booth Hartley for providing IRAS co-added data.
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# EIKONAL SCATTERING OF MONOPOLES AND DYONS IN DUAL QED
## 1 Introduction
The topic of magnetic charge has received enormous attention since Dirac $`^\mathrm{?}`$ demonstrated its existence was consistent with quantum mechanics provided the quantization condition (in rationalized units) $`eg/4\pi =N/2`$ is satisfied. Here $`e`$ and $`g`$ are the strength of electric and magnetic charges, respectively, and $`N`$ denotes an integer. In the case of dyons, particles containing both magnetic and electric charge, the Schwinger generalization $`^{\mathrm{?},\mathrm{?}}`$
$`{\displaystyle \frac{e_ag_be_bg_a}{4\pi }}=\left\{\begin{array}{c}\frac{N}{2},\text{unsymmetric}\hfill \\ N,\text{symmetric}\hfill \end{array}\right\},`$ (3)
is invoked. (“Symmetric” and “unsymmetric” refer to the presence or absence of dual symmetry in the solutions of Maxwell’s equations.)
With the advent of non-Abelian theories, classical composite monopole solutions were (theoretically) discovered. $`^\mathrm{?}`$ Their mass would be of the order of the relevant gauge-symmetry breaking scale, which for grand unified theories (GUT) is of order $`10^{16}`$ GeV or higher. However, there are models where the electroweak symmetry breaking can give rise to monopoles of mass $`10`$ TeV. $`^\mathrm{?}`$ Yet, even the latter are not accessible to accelerator experiments, so limits on heavy monopoles depend either on cosmological considerations, or detection of cosmologically produced (relic) monopoles impinging upon the earth or moon. $`^\mathrm{?}`$ However, a priori, there is no reason that Dirac/Schwinger monopoles or dyons of arbitrary mass might not exist. In this respect, it is important to set limits below the 1 TeV scale in direct accelerator based experiments. <sup>d</sup><sup>d</sup>dSuch an experiment is currently in progress at the University of Oklahoma, $`^\mathrm{?}`$ where we have set limits on direct monopole production at Fermilab up to several hundred GeV. This is an improvement over previous limits. $`^\mathrm{?}`$ See also Ref. $`^\mathrm{?}`$ for critique of theories underlying indirect searches.
It is envisaged that if monopoles are sufficiently light, they would be produced by a Drell-Yan type of process occurring in $`p\overline{p}`$ collisions at the Tevatron. The difficulty is to derive a reasonable estimate of the elementary process $`q\overline{q}\gamma ^{}M\overline{M}`$, where $`q`$ stands for quark and $`M`$ for magnetic monopole. Attempts to incorporate monopoles consistently into relativistic quantum field theory have met with mixed success. Weinberg, and soon thereafter Rabl, $`^\mathrm{?}`$ demonstrated that the charge-monopole scattering amplitude, calculated in the one-photon-exchange approximation, is a function of the Dirac string singularity. Making matters worse, the value of the vertex coupling implied by $`\alpha _g=g^2/4\pi 34N^2`$, calls into question any approach based on a badly divergent perturbative expansion in $`\alpha _g`$. Although the early efforts using a Feynman-rule perturbation theory resulted in string-dependent cross–sections, subsequently ad hoc assumptions were invoked to render the resulting cross-sections string independent (see Ref. $`^\mathrm{?}`$ for details). In contrast, studying the formal behavior of Green’s functions in the relativistic quantum field theory of electrons and monopoles, both Schwinger $`^{\mathrm{?},\mathrm{?}}`$ and Brandt et al. $`^\mathrm{?}`$ demonstrated Lorentz-string and gauge invariance.
However, with the exception of one instance $`^\mathrm{?}`$ such demonstrations have been conspicuously absent at the phenomenological level.<sup>e</sup><sup>e</sup>eThis is surprising because one expects that the invariant non-relativistic scattering result (see Ref. $`^\mathrm{?}`$ and references therein) corresponds in a certain kinematic regime to a infinite summation of a particular subclass of Feynman diagrams. This deficiency stems from the fact that in most phenomenological treatments of charge-monopole processes the “string independence” of the quantum field theory and the strength of the coupling are treated as separate issues. In fact, these two points are intimately related. The lesson to be learned from the formal and non-relativistic demonstrations of Lorentz and string invariance is this: Because the quantization condition is intimately tied to the demonstration of Lorentz invariance, the latter can only be demonstrated using a method which does not rely on perturbation theory (see $`^\mathrm{?}`$ for further discussion).
In view of the necessity of establishing a reliable estimate for monopole production in accelerators in order to be able to set bounds on monopole masses, it is important to put the theory of dual quantum electrodynamics (dual QED) on a firmer foundation. With that in mind we present our results. $`^\mathrm{?}`$
## 2 Dual Electrodynamics
For a spin $`\frac{1}{2}`$ monopole, a minimal generalization of the QED action $`^{\mathrm{?},\mathrm{?}}`$ for charge-monopole interactions expressed in terms of the vector potential $`A_\mu `$ and field strength tensor $`F_{\mu \nu }`$ (i.e. , in a first-order formalism) is
$`W`$ $`=`$ $`{\displaystyle }(dx)\{{\displaystyle \frac{1}{2}}F^{\mu \nu }(x)(_\mu A_\nu \left(x\right)_\nu A_\mu \left(x\right))+{\displaystyle \frac{1}{4}}F_{\mu \nu }(x)F^{\mu \nu }(x)`$ (4)
$`+\overline{\psi }(x)(i\gamma +e\gamma A(x)m_\psi )\psi (x)+\overline{\chi }(x)(i\gamma +g\gamma B(x)m_\chi )\chi (x)\},`$
where it is assumed that the electrically and magnetically charged particles are spin $`1/2`$. The resulting Maxwell’s equations, which imply the dual conservation of electric and magnetic currents, $`j_\mu `$ and $`{}_{}{}^{}j_{\mu }^{}`$, necessitates the introduction of the Dirac string singularity. The Dirac string function satisfies the differential equation $`_\mu f^\mu (x)=\delta (x)`$, which has the formal solution $`f^\mu (x)=n^\mu \left(n\right)^1\delta (x)`$, where $`n^\mu `$ is an arbitrary vector.<sup>f</sup><sup>f</sup>fHere we have chosen the string to satisfy the oddness condition (this is the “symmetric” solution) $`f^\mu (x)=f^\mu (x)`$, corresponding to Schwinger’s integer quantization condition. $`^\mathrm{?}`$ The field equations resulting from $`\delta W=0`$ are $`_\nu F^{\mu \nu }=j^\mu `$ and
$$F_{\mu \nu }=_\mu A_\nu _\nu A_\mu +^{}G_{\mu \nu },$$
(5)
where
$$G_{\mu \nu }(x)=(dy)\left(f_\mu (xy)^{}j_\nu (y)f_\nu (xy)^{}j_\mu (y)\right).$$
(6)
The auxiliary dual field $`B_\mu `$ is defined as a functional of field-strength and depends on the string function,
$$B_\mu (x)=(dy)f^\nu \left(xy\right)^{}F_{\mu \nu }(y).$$
(7)
Of course, the monopole field satisfies the Dirac equation
$$\left(i\gamma +g\gamma B(x)m_g\right)\chi (x)=0.$$
(8)
From Eq. (7) we find that $`B_\mu `$ satisfies $`(dx^{})f^\mu (xx^{})B_\mu (x^{})=0`$, which is a special case of a gauge-fixed vector field defined in terms of the field strength through an inversion formula, Eq. (7). Similarly, we are at liberty to restrict the vector potential, $`A_\mu `$ to a hypersurface in field space embodied in the inversion formula
$$A_\mu (x)=(dy)f^\nu \left(xy\right)F_{\mu \nu }(y),$$
(9)
which we denote as string-gauge, $`(dx^{})f^\mu (xx^{})A_\mu (x^{})=0`$. The photon kernel derived from the corresponding gauge fixed action now possesses an inverse
$$D_{\mu \nu }(x)=\left[g_{\mu \nu }\frac{n_\mu _\nu +n_\nu _\mu }{(n)}+\left(1\frac{1}{\kappa }\frac{(n)^2^2}{n^2}\right)\frac{n^2_\mu _\nu }{(n)^2}\right]D_+(x),$$
(10)
where $`D_+(x)`$ is the massless scalar propagator,
$$D_+(x)=\frac{1}{^2iϵ}\delta (x),$$
(11)
which enables us to write an integral equation, expressing the vector potential, $`A_\mu `$ ($`B_\mu `$) in terms of the electric (magnetic) and magnetic (electric) currents. We generalize these classical integral equations to one point Green functions in obtaining the generating function for Green’s functions in dual QED.
## 3 Quantization of Dual QED: Schwinger-Dyson Equations
Using a path integral formulation to quantize the string-dependent action is by no means straightforward. In order to unambiguously develop the generating functional for physical processes we make use of Schwinger’s quantum action principle, $`^\mathrm{?}`$ where we write the vacuum persistence amplitude for Green functions in the presence of external sources, $`Z(𝒥)=0_+|0_{}^𝒥`$ for the charge-monopole system. That is, under an arbitrary variation,
$$\delta 0_+|0_{}^𝒥=i0_+\left|\delta W(𝒥)\right|0_{}^𝒥,$$
(12)
where $`W(𝒥)`$ is the action given in Eq. (4) externally driven by the sources, $`𝒥`$, which for the present case are given by the set of terms
$$W(𝒥)=W+(dx)\left\{J^\mu A_\mu +{}_{}{}^{}J_{}^{\mu }B_\mu +\overline{\eta }\psi +\overline{\psi }\eta +\overline{\xi }\chi +\overline{\chi }\xi \right\}.$$
(13)
Given the one-point functions ($`𝒪_\mu `$ is the field conjugate to the source $`𝒥^\mu `$)
$$\frac{\delta \mathrm{log}Z(𝒥)}{i\delta 𝒥^\mu (x)}=\frac{0_+|𝒪_\mu (x)|0_{}^𝒥}{0_+|0_{}^𝒥},$$
(14)
we solve the corresponding coupled Schwinger-Dyson equations for the vacuum amplitude $`^\mathrm{?}`$, subject to the gauge conditions
$$(dx^{})f^\nu (xx^{})\frac{\delta 0_+|0_{}_0^𝒥}{\delta J^\nu (x^{})}=0,(dx^{})f^\mu (xx^{})\frac{\delta 0_+|0_{}^𝒥}{\delta {}_{}{}^{}J_{\mu }^{}(x^{})}=0.$$
(15)
Since any expansion in $`\alpha _g`$ or $`eg`$ is not practically useful we recast the solution into a functional form better suited for a nonperturbative calculation of the four-point Green’s function. $`^\mathrm{?}`$ For dyons, the different species of which are labeled by the index $`a`$ this is
$`Z(𝒥)`$ $`=`$ $`\mathrm{exp}\left\{{\displaystyle \frac{i}{2}}{\displaystyle (dx)(dx^{})𝒦^\mu (x)𝒟_{\mu \nu }\left(xx^{}\right)𝒦^\nu \left(x^{}\right)}\right\}`$ (16)
$`\times \mathrm{exp}\left\{{\displaystyle \frac{i}{2}}{\displaystyle (dx)(dx^{})\frac{\delta }{\delta \overline{𝒜}_\mu (x)}𝒟_{\mu \nu }(xx^{})\frac{\delta }{\delta \overline{𝒜}_\nu (x^{})}}\right\}`$
$`\times \mathrm{exp}\left\{i{\displaystyle \underset{a}{}}{\displaystyle (dx)(dx^{})\overline{\zeta }_a\left(x\right)G_a(x,x^{}|\overline{𝒜_a})\zeta _a\left(x^{}\right)}\right\}`$
$`\times \mathrm{exp}\left\{{\displaystyle \underset{a}{}}{\displaystyle _0^1}𝑑q\mathrm{Tr}\gamma \overline{𝒜_a}G_a(x,x|q\overline{𝒜_a})\right\},`$
where $`𝒜_a=e_aA+g_aB`$, $`\zeta _a`$ is the source for the dyon of species $`a`$, and a matrix notation is adopted,
$`𝒦^\mu (x)`$ $`=`$ $`\left(\begin{array}{c}J(x)\\ ^{}J(x)\end{array}\right),{\displaystyle \frac{\delta }{\delta \overline{𝒜}_\mu (x)}}=\left(\begin{array}{c}\delta /\delta \overline{A}_\mu (x)\\ \delta /\delta \overline{B}_\mu (x)\end{array}\right),`$ (21)
$`𝒟_{\mu \nu }\left(xx^{}\right)`$ $`=`$ $`\left(\begin{array}{cc}D_{\mu \nu }\left(xx^{}\right)& \stackrel{~}{D}_{\mu \nu }\left(xx^{}\right)\\ \stackrel{~}{D}_{\mu \nu }\left(xx^{}\right)& D_{\mu \nu }\left(xx^{}\right)\end{array}\right).`$ (24)
We use the shorthand notation for the “dual propagator” that couples magnetic to electric charge
$$\stackrel{~}{D}_{\mu \nu }\left(xx^{}\right)=ϵ_{\mu \nu \sigma \tau }(dx^{\prime \prime })D_+\left(xx^{\prime \prime }\right)^{\prime \prime \sigma }f^\tau \left(x^{\prime \prime }x^{}\right).$$
(25)
The two-point fermion Green’s functions in the background of the stationary photon field $`\overline{A},\overline{B}`$ are given by
$$G(x,x^{}|\overline{𝒜})=x|(\gamma p+m\overline{𝒜})^1|x^{}.$$
(26)
## 4 Eikonal Approximation for Dyon-Dyon and Charge-Monopole Scattering
To calculate the dyon-dyon scattering cross section we obtain the four–point Green’s function for this process from Eq. (16),
$$G(x_1,y_1;x_2,y_2)=\frac{\delta }{i\delta \overline{\zeta }_1(x_1)}\frac{\delta }{i\delta \zeta _1(y_1)}\frac{\delta }{i\delta \overline{\zeta }_2(x_2)}\frac{\delta }{i\delta \zeta _2(y_2)}Z(𝒥)|_{𝒥=0}.$$
(27)
The subscripts on the sources refer to the two different dyons. As a first step in analyzing the string dependence of the scattering amplitudes, we study high-energy forward scattering processes where soft photon exchanges dominate. Diagrammatically, in this kinematic regime we restrict attention to that subclass in which there are no closed fermion loops and the photons are exchanged between fermions. $`^\mathrm{?}`$ This amounts to quenched-ladder approximation where the linkage operators, $`Ł`$, connect two fermion propagators via photon exchange. We can read this off from Eq. (16):
$$\mathrm{e}^{Ł_{12}}=\mathrm{exp}\left\{i(dx)(dx^{})\frac{\delta }{\delta \overline{𝒜}_1^\mu (x)}𝒟^{\mu \nu }\left(xx^{}\right)\frac{\delta }{\delta \overline{𝒜}_2^\nu (x^{})}\right\},$$
(28)
so Eq. (27) takes the form
$`G(x_1,y_1;x_2,y_2)`$ $`=`$ $`\mathrm{e}^{Ł_{12}}G_1(x_1,y_1|\overline{𝒜}_1)G_2(x_2,y_2|\overline{𝒜}_2)|_{\overline{A}=\overline{B}=0},`$ (29)
where we express the two-point function using the proper-time parameter representation of an ordered exponential. The soft, nonperturbative effects of the interaction between electric and magnetic charges dominate in the region where the momentum exchanged by the photons is small compared to the center of mass energy $`s=(p_1+p_2)^2`$, i.e. $`t/s\mathrm{\hspace{0.17em}1}`$. This amounts to the Bloch-Nordsieck $`^\mathrm{?}`$ or eikonal approximation$`^{\mathrm{?},\mathrm{?}}`$ This approximation substantially simplifies evaluating the path-ordered exponentials in Eq. (29). They are now exponentials of linear functionals of the gauge field.
### 4.1 High Energy Scattering Cross Section
Using the identity,
$$\mathrm{e}^Ł=1+_0^1𝑑a\mathrm{e}^{aŁ}Ł$$
(30)
one obtains to the following form of the four-point Green function, Eqs. (29),
$`G(x_1,y_1;x_2,y_2)`$ $`=`$ $`{\displaystyle _0^1}𝑑a{\displaystyle (dz_1)(dz_2)\overline{𝒟}_{\mu \nu }(z_1z_2)\mathrm{e}^{aŁ_{12}}}`$ (31)
$`\times G_1(x_1,z_1|\overline{𝒜}_1)\gamma _\mu G_1(z_1,y_1|\overline{𝒜}_1)G_2(x_2,z_2|\overline{𝒜}_2)\gamma ^\nu G_2(z_2,y_2|\overline{𝒜}_2)|_{\overline{A}=\overline{B}=0},`$
where, $`\overline{𝒟}_{\mu \nu }(x)`$ represents the combination of propagators,
$$\overline{𝒟}_{\mu \nu }(z_1z_2)=𝗊_1𝗊_2D_{\mu \nu }(z_1z_2)𝗊_1\times 𝗊_2\stackrel{~}{D}_{\mu \nu }(z_1z_2),$$
(32)
and the charge combinations invariant under duality transformations are
$$𝗊_1𝗊_2=e_1e_2+g_1g_2\text{and}𝗊_1\times 𝗊_2=e_1g_2g_1e_2.$$
(33)
Choosing a space like string, $`n^\mu =(0,𝐧)`$ and the incoming momenta to be in the $`z`$ direction, in the center of momentum frame, the Møller amplitude, $`M(s,t)`$ is given by
$`M(s,t)`$ $`=`$ $`{\displaystyle \frac{i}{2\pi }}{\displaystyle _0^1}𝑑a{\displaystyle d^2𝐱\mathrm{e}^{i𝐪𝐱}\overline{u}(p_1^{})\gamma ^\mu u(p_1)\overline{u}(p_2^{})\gamma ^\nu u(p_2)\mathrm{e}^{ia\mathrm{\Phi }_n(s,t;x)}}`$ (34)
$`\left\{g_{\mu \nu }𝗊_1𝗊_2K_0\left(\mu \left|𝐱\right|\right)ϵ_{\mu \nu \sigma \tau }𝗊_1\times 𝗊_2n^\tau {\displaystyle \frac{}{n_\sigma }}{\displaystyle \frac{1}{2}}{\displaystyle \frac{d\text{t}}{\text{t}}K_0\left(\mu \left|\left(𝐱+\text{t}𝐧\right)\right|\right)}\right\},`$
where in this kinematic regime, the eikonal phase is ($`\mu `$ is the photon mass; $`\stackrel{~}{\mu }=\mathrm{e}^\gamma \mu /2`$ and $`\gamma `$ is Euler’s constant)
$$\mathrm{\Phi }_n(s,t;x)=\frac{1}{2\pi }\left\{𝗊_1𝗊_2\mathrm{ln}\left(\stackrel{~}{\mu }\left|𝐱\right|\right)𝗊_1\times 𝗊_2\mathrm{arctan}\left[\frac{\widehat{𝐧}𝐱}{\widehat{𝐳}\left(\widehat{𝐧}\times 𝐱\right)}\right]\right\}.$$
(35)
First, we calculate the helicity amplitudes in the high-energy limit, $`p^0m`$. In performing the integral over the impact parameter care must be taken since the arctangent function is discontinuous when the $`xy`$ component of $`\widehat{𝐧}`$ and $`𝐱`$ lie in the same direction. However, requiring that the eikonal phase factor $`e^{i\mathrm{\Phi }_n}`$ be continuous, necessarily leads to the Schwinger quantization condition (3): $`𝗊_1\times 𝗊_2=4N\pi `$. Now using the integral form for the Bessel function of order $`\nu `$
$$i^\nu J_\nu (t)=_0^{2\pi }\frac{d\varphi }{2\pi }\mathrm{e}^{i\left(t\mathrm{cos}\varphi \nu \varphi \right)},$$
(36)
we find the dyon-dyon scattering amplitude to be
$$M(s,t)=\frac{s}{M_1M_2}\frac{2\pi }{q^2}(Ni\stackrel{~}{\alpha })\mathrm{e}^{i2N\psi }\left(\frac{4\stackrel{~}{\mu }^2}{q^2}\right)^{i\stackrel{~}{\alpha }}\frac{\mathrm{\Gamma }\left(1+N+i\stackrel{~}{\alpha }\right)}{\mathrm{\Gamma }\left(1+Ni\stackrel{~}{\alpha }\right)},$$
(37)
where $`\stackrel{~}{\alpha }=𝗊_1𝗊_2/4\pi `$, and $`\psi `$ is the angle between $`𝐪`$ and $`𝐧`$.
This result is almost identical in structure to the non-relativistic form of the scattering amplitude for the Coulomb potential, which result is recovered by setting $`N=0`$ (see, for example, Ref. $`^\mathrm{?}`$.) Following the standard convention we calculate the spin-averaged cross section for dyon-dyon scattering in the high energy limit,
$$\frac{d\sigma }{dt}=\frac{\left(𝗊_1𝗊_2\right)^2+\left(𝗊_1\times 𝗊_2\right)^2}{4\pi t^2}.$$
(38)
For the case of charge-monopole scattering $`e_1=g_2=0`$, this result, coincides with that found by Urrutia <sup>g</sup><sup>g</sup>gUtilizing Schwinger’s functional source theory $`^\mathrm{?}`$ in the context of a zeroth order eikonal approximation Urrutia $`^\mathrm{?}`$ demonstrated string independence of the charge-monopole scattering cross section, although in his treatment the currents are approximated by those of classical point particles. which is also string independent as a consequence of the quantization condition.
## 5 Corrections
To go beyond the regime of soft or infrared photon exchange requires a detailed analysis of factorization of soft and hard contributions and their correlations. Such effects have been widely studied in the context of phenomenologically based hadronic interacting field theories $`^\mathrm{?}`$ and more recently in the context of diffractive scattering $`^\mathrm{?}`$ in QCD and in the world-line formalism. $`^{\mathrm{?},\mathrm{?}}`$
As a reasonable first step we impose corrections on the eikonal amplitude by relaxing the “high energy” approximation on the spinors in Eq. (34) while assuming that the soft contributions are dominated by the eikonal phase. $`^\mathrm{?}`$ This result, which we expect to be a better and better approximation in the high energy limit (for a given $`t`$), i.e. $`t/s0`$, obeys the expected scaling behavior, for electron monopole scattering
$`s^2{\displaystyle \frac{d\sigma }{dt}}={\displaystyle \frac{(eg)^2}{4\pi }}{\displaystyle \frac{1}{t^2}}\left(s+t\right)^2f\left({\displaystyle \frac{t}{s}}\right).`$ (39)
Assuming that we have extended the kinematic range of this result beyond the low $`t^2`$ limit, we can consider using the analytic properties of the scattering amplitude to calculate the Drell-Yan production amplitude in the $`t`$-channel. Detail will be presented in future publication.
## 6 Conclusions
We have given a complete formulation, in modern quantum field theoretic language, of an interacting electron-monopole and dyon-dyon systems. The challenge is to apply these equations to the calculation of monopole and dyon processes. Perturbation theory is useless, not only because of the strength of the coupling, but also because the graphs are fatally string- (or gauge-) dependent. The most obvious nonperturbative technique for transcending these limitations in scattering processes lies in the high energy regime where the eikonal approximation is applicable; in that limit, our formalism generalizes an early lowest-order result of Urrutia and charts the way to include systematic corrections. More problematic is the treatment of monopole production processes. We will apply the techniques and results found here to the Drell-Yan production of monopole-antimonopole processes, and obtain phenomenologically relevant estimates for the accelerator production of monopole-antimonopole pairs.
In addition we have also detailed how the Dirac string dependence disappears from physical quantities. It is by no means a result of string averaging or a result of dropping string-dependent terms (see Ref. $`^\mathrm{?}`$ for details); but rather, a result of summing the soft contributions to the dyon-dyon or electron-monopole process. At the level of the eikonal approximation and its corrections one might suspect the occurrence of a factorization of hard string-independent and soft string-dependent contributions in a manner similar to that argued recently in strong-coupling QCD. There is reason to believe that inclusion of soft-hard correlations in the scattering will not spoil this consistency.
## Acknowledgments
L.G. would like to thank Herb Fried and the organizers of “Fifth Workshop on QCD” for the invitation to present this work. This work is supported in part by the Department of Energy.
## References
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# 1 Introduction
## 1 Introduction
This letter presents results on W-pair production in $`\mathrm{e}^+\mathrm{e}^{}`$collisions using data collected with the ALEPH detector at a centre-of-mass (CM) energy around 189 $`\mathrm{GeV}`$, during the 1998 data taking period. The $`\mathrm{WW}`$ events are identified in all possible W decay channels, thus allowing the determination of the W branching fractions and indirectly, the coupling of the W to cs pairs.
The experimental conditions and data analysis follow those used in the cross section measurements at lower LEP2 energies. As they are already described in detail in , attention is focused here on changes to selection procedures other than a simple rescaling of cuts with the increased collision energy.
A detailed description of the ALEPH detector can be found in Ref. and of its performance in Ref. . The luminosity is measured from small-angle Bhabha events, using lead-proportional wire sampling calorimeters , with an accepted Bhabha cross section of approximately 4.25 nb . An integrated luminosity of 174.20 $`\pm `$ 0.20 $`(\mathrm{stat}.)\pm `$ 0.73 $`(\mathrm{syst}.)\mathrm{pb}^1`$ was recorded at a mean CM energy of 188.63 $`\pm `$ 0.04 $`\mathrm{GeV}`$ .
In this letter, the quoted signal cross sections are the CC03 cross sections , defined as the production of four-fermion final states through two resonating W bosons. Two processes contribute, $`\nu _\mathrm{e}`$ exchange in the $`t`$-channel and $`\mathrm{Z}/\gamma `$ exchange in the $`s`$-channel. The measured cross sections are corrected for the difference, denoted the “4f-CC03 correction” , between the accepted cross sections for CC03 processes and all Standard Model four-fermion final states consistent with W-pair decays.
The CC03 Standard Model cross section ($`\sigma _{\mathrm{WW}}`$), calculated at $`\sqrt{s}=188.63`$ GeV with the program GENTLE is 16.65 pb ($`\pm 2\%`$). The KORALW version 1.21 Monte Carlo event generator is used to simulate the signal events with a normalised cross section in agreement with the GENTLE value. The JETSET package is used for the hadronisation. Comparison samples, generated with EXCALIBUR and grc4f for both CC03 and all four-fermion diagrams, are used for systematic error evaluation. Samples of events are also generated with different W masses, both for CC03 diagrams and for all WW-like four-fermion diagrams, with KORALW.
The KORALZ Monte Carlo program is used to generate $`\mathrm{e}^+\mathrm{e}^{}\mathrm{q}\overline{\mathrm{q}}`$ background events. Other backgrounds are generated with PYTHIA 5.7 for $`\mathrm{ZZ}`$, $`\mathrm{Z}`$$`\mathrm{ee}`$ and $`\mathrm{W}`$$`\mathrm{e}\nu `$ processes, PHOT02 for $`\gamma \gamma `$ interactions, KORALZ for $`\mu `$ and $`\tau `$ pair production and BHWIDE and UNIBAB for Bhabha events.
## 2 Selection of W-pair candidates
### 2.1 $`\mathrm{𝐖𝐖}\mathbf{}\mathbf{}𝝂\mathbf{}𝝂`$ events
The selection of fully leptonic W-pair decays follows exactly the two selections used for the cross section and branching fraction measurements at 183 GeV . The two selections have similar overall efficiencies and background levels but differ in their sensitivities to the individual dilepton channels. The first is based on topological information and is sensitive to all channels. In the second, lepton identification is used to optimise the cuts according to which final state is being considered. Events are accepted as $`\mathrm{𝐖𝐖}`$ candidates if they pass either of the two selections and are then classified into six di-lepton channels making use of electron and muon identification criteria. A jet or a single charged particle track is classified as a tau if no lepton is identified or the identified lepton has an energy less than 25 GeV.
Beam related background, not simulated in Monte Carlo events, affects the efficiency of the cut which removes events depositing energy within $`\mathrm{𝟏𝟐}^{\mathbf{}}`$ of the beam. Random trigger events were used to model these local energy deposits close to the beam. The inefficiency introduced is found to be $`\mathbf{8.2}\mathbf{\pm }\mathbf{0.5}\mathbf{\%}`$. The CC03 efficiencies in the individual $`\mathbf{}𝝂\mathbf{}𝝂`$ channels are given in Table 1 after correction for this beam related effect. The inclusive combination of the two selections has an overall efficiency of $`\mathbf{64.2}\mathbf{\pm }\mathbf{0.4}\mathbf{\%}`$ for the fully leptonic channels when combined assuming lepton universality. The total background amounts to $`\mathrm{𝟏𝟑𝟏}\mathbf{\pm }\mathrm{𝟕}\mathbf{(}\mathrm{𝐬𝐭𝐚𝐭}\mathbf{.}\mathbf{)}\mathbf{\pm }\mathbf{8.5}\mathbf{(}\mathrm{𝐬𝐲𝐬𝐭}\mathbf{.}\mathbf{)}`$ fb and is dominated by $`𝜸𝜸\mathbf{}\mathbf{}\mathbf{}`$ and non-WW-like $`\mathrm{𝐙𝐙}\mathbf{}\mathbf{}\mathbf{}𝝂𝝂`$ events. In the data, the inclusive combination selects 220 events.
All sources of systematic uncertainties are listed in Table 2. They are dominated by the uncertainties on the background cross sections, the uncertainty from the cut on energy detected close to the beam and by Monte Carlo statistics.
A maximum likelihood fit is applied to determine the cross section for each fully leptonic decay channel using the efficiency matrices for signal and backgrounds given in Table 1. For all channels together the 4f-CC03 correction is $`\mathbf{}\mathrm{𝟏𝟎}\mathbf{\pm }\mathrm{𝟏𝟎}`$ fb, where the uncertainty comes from Monte Carlo statistics.
The results of the fit are
$`𝝈\mathbf{(}\mathrm{𝐖𝐖}\mathbf{}𝐞𝝂𝐞𝝂\mathbf{)}`$ $`\mathbf{=}`$ $`\mathbf{0.19}\mathbf{\pm }\mathbf{0.05}\mathbf{(}\mathrm{𝐬𝐭𝐚𝐭}\mathbf{.}\mathbf{)}\mathbf{\pm }\mathbf{0.01}\mathbf{(}\mathrm{𝐬𝐲𝐬𝐭}\mathbf{.}\mathbf{)}\mathrm{𝐩𝐛}\mathbf{,}`$
$`𝝈\mathbf{(}\mathrm{𝐖𝐖}\mathbf{}𝝁𝝂𝝁𝝂\mathbf{)}`$ $`\mathbf{=}`$ $`\mathbf{0.20}\mathbf{\pm }\mathbf{0.05}\mathbf{(}\mathrm{𝐬𝐭𝐚𝐭}\mathbf{.}\mathbf{)}\mathbf{\pm }\mathbf{0.01}\mathbf{(}\mathrm{𝐬𝐲𝐬𝐭}\mathbf{.}\mathbf{)}\mathrm{𝐩𝐛}\mathbf{,}`$
$`𝝈\mathbf{(}\mathrm{𝐖𝐖}\mathbf{}𝝉𝝂𝝉𝝂\mathbf{)}`$ $`\mathbf{=}`$ $`\mathbf{0.22}\mathbf{\pm }\mathbf{0.08}\mathbf{(}\mathrm{𝐬𝐭𝐚𝐭}\mathbf{.}\mathbf{)}\mathbf{\pm }\mathbf{0.02}\mathbf{(}\mathrm{𝐬𝐲𝐬𝐭}\mathbf{.}\mathbf{)}\mathrm{𝐩𝐛}\mathbf{,}`$
$`𝝈\mathbf{(}\mathrm{𝐖𝐖}\mathbf{}𝐞𝝂𝝁𝝂\mathbf{)}`$ $`\mathbf{=}`$ $`\mathbf{0.43}\mathbf{\pm }\mathbf{0.07}\mathbf{(}\mathrm{𝐬𝐭𝐚𝐭}\mathbf{.}\mathbf{)}\mathbf{\pm }\mathbf{0.01}\mathbf{(}\mathrm{𝐬𝐲𝐬𝐭}\mathbf{.}\mathbf{)}\mathrm{𝐩𝐛}\mathbf{,}`$
$`𝝈\mathbf{(}\mathrm{𝐖𝐖}\mathbf{}𝐞𝝂𝝉𝝂\mathbf{)}`$ $`\mathbf{=}`$ $`\mathbf{0.36}\mathbf{\pm }\mathbf{0.08}\mathbf{(}\mathrm{𝐬𝐭𝐚𝐭}\mathbf{.}\mathbf{)}\mathbf{\pm }\mathbf{0.02}\mathbf{(}\mathrm{𝐬𝐲𝐬𝐭}\mathbf{.}\mathbf{)}\mathrm{𝐩𝐛}\mathbf{,}`$
$`𝝈\mathbf{(}\mathrm{𝐖𝐖}\mathbf{}𝝁𝝂𝝉𝝂\mathbf{)}`$ $`\mathbf{=}`$ $`\mathbf{0.38}\mathbf{\pm }\mathbf{0.08}\mathbf{(}\mathrm{𝐬𝐭𝐚𝐭}\mathbf{.}\mathbf{)}\mathbf{\pm }\mathbf{0.01}\mathbf{(}\mathrm{𝐬𝐲𝐬𝐭}\mathbf{.}\mathbf{)}\mathrm{𝐩𝐛}\mathbf{.}`$
The systematic uncertainties are obtained by varying all input parameters in the fit according to their uncertainties.
The total fully leptonic cross section is obtained with the same fit, assuming lepton universality:
$$𝝈\mathbf{(}\mathrm{𝐖𝐖}\mathbf{}\mathbf{}𝝂\mathbf{}𝝂\mathbf{)}\mathbf{=}\mathbf{1.78}\mathbf{\pm }\mathbf{0.13}\mathbf{(}\mathrm{𝐬𝐭𝐚𝐭}\mathbf{.}\mathbf{)}\mathbf{\pm }\mathbf{0.02}\mathbf{(}\mathrm{𝐬𝐲𝐬𝐭}\mathbf{.}\mathbf{)}\mathrm{𝐩𝐛}\mathbf{,}$$
consistent with the sum of the individual channels.
### 2.2 $`\mathrm{𝐖𝐖}\mathbf{}\mathbf{}𝝂𝐪\overline{𝐪}`$ events
As for the lower energy measurements, three $`\mathbf{}𝝂𝐪\overline{𝐪}`$ selection procedures are applied. One selection requires an identified electron or muon. The other two are designed to select $`𝝉𝝂𝐪\overline{𝐪}`$ events, based on global variables or topological properties of the events.
The selection of $`𝐞𝝂𝐪\overline{𝐪}`$ and $`𝝁𝝂𝐪\overline{𝐪}`$ events has been modified with respect to the previous analysis to take into account the greater initial boost of the W’s. The preselection remains similar. It is based on the total charged particle energy and multiplicity and a cut on the longitudinal momentum and visible energy to reject Z$`𝜸`$ events with an undetected photon. The selection of the lepton track relies on the fact that it is in general more energetic and isolated than the charged particles from the hadronic system. Thus the candidate lepton is chosen as the charged track that maximises $`𝒑_{\mathbf{}}^\mathrm{𝟐}\mathbf{(}\mathrm{𝟏}\mathbf{}\mathrm{𝐜𝐨𝐬}𝜽_\mathbf{}𝒋\mathbf{)}`$, where $`𝒑_{\mathbf{}}`$ is the track momentum and $`𝜽_\mathbf{}𝒋`$ the angle of the track to the closest of the jets clustered using the remaining reconstructed charged particles. For the jet clustering the DURHAM-P algorithm with a $`𝒚_{\mathrm{𝐜𝐮𝐭}}`$ of 0.0003 is used.
The same electron or muon identification criteria as for the fully leptonic channels are required for this lepton candidate track. However, no cut is applied on the lepton energy, so that $`𝝉𝝂`$$`𝐪\overline{𝐪}`$ events where the $`𝝉`$ decays to a softer lepton (as $`𝝉\mathbf{}𝐞𝝂𝝂`$ or $`𝝉\mathbf{}𝝁𝝂𝝂`$) are also selected by this analysis.
For electron candidates the lepton energy is corrected for possible bremsstrahlung photons detected in the electromagnetic calorimeter. The isolation of the lepton is defined as $`\mathrm{𝐥𝐨𝐠}\mathbf{(}\mathrm{𝐭𝐚𝐧}𝜽_𝑪\mathbf{/}\mathrm{𝟐}\mathbf{)}\mathbf{+}\mathrm{𝐥𝐨𝐠}\mathbf{(}\mathrm{𝐭𝐚𝐧}𝜽_𝑭\mathbf{/}\mathrm{𝟐}\mathbf{)}`$ where $`𝜽_𝑪`$ and $`𝜽_𝑭`$ are, respectively, the angle of the lepton to the closest charged track, and the opening angle of the largest cone centred on the lepton direction which contains a total energy smaller than 5 $`\mathrm{𝐆𝐞𝐕}`$.
For each event, probabilities that it comes from each of the three signal processes, $`𝐞𝝂`$$`𝐪\overline{𝐪}`$, $`𝝁𝝂`$$`𝐪\overline{𝐪}`$ or $`𝝉𝝂`$$`𝐪\overline{𝐪}`$, are determined (Fig. 1) using Monte Carlo reference samples of signal and backgrounds. These are evaluated from the identity, energy and isolation of the lepton plus the event total transverse momentum. An event is classified as $`𝐞𝝂`$$`𝐪\overline{𝐪}`$ or $`𝝁𝝂`$$`𝐪\overline{𝐪}`$ if its corresponding probability is greater than 0.40; it is then not considered in the tau search.
The selection of $`𝝉𝝂𝐪\overline{𝐪}`$ events is based both on global event variables and a topological selection which attempts to identify the $`𝝉`$ jet. As the selections were described in detail in previous papers , only changes other than a rescaling of the values of the energy-based cuts are described in the following:
* in the global analysis the visible mass is required to be greater than 85 $`\mathrm{𝐆𝐞𝐕}\mathbf{/}𝒄^\mathrm{𝟐}`$ and less than 155 $`\mathrm{𝐆𝐞𝐕}\mathbf{/}𝒄^\mathrm{𝟐}`$ to take account of the boost of the W boson. The estimated energy of the “primary” neutrino must be smaller than 70 $`\mathrm{𝐆𝐞𝐕}`$;
* in the topological analysis, the energy of the most energetic quark jet must be less than 75 GeV and the mass of the hadronic system, i.e., the mass excluding the tau jet, is required to be less than 100 $`\mathrm{𝐆𝐞𝐕}\mathbf{/}𝒄^\mathrm{𝟐}`$.
In addition, if an event with a well defined e or $`𝝁`$ fails the e/$`𝝁`$ probability cut it is considered as a $`𝝉`$ candidate and kept if the $`𝝉𝝂𝐪\overline{𝐪}`$ probability (Fig. 1c) is greater than 0.4.
#### 2.2.1 Results
Table 1 gives the efficiencies for each selection. The inclusive efficiencies for the three semileptonic decay channels are $`\mathbf{87.8}\mathbf{\pm }\mathbf{0.4}`$% for the electron channel, $`\mathbf{91.6}\mathbf{\pm }\mathbf{0.4}`$% for the muon channel and $`\mathbf{66.5}\mathbf{\pm }\mathbf{0.5}`$% for the tau channel, giving an $`\mathrm{𝟖𝟐}\mathbf{\pm }\mathbf{0.3}`$% average efficiency for $`\mathrm{𝐖𝐖}\mathbf{}\mathbf{}𝝂𝐪\overline{𝐪}`$, with a background of $`\mathrm{𝟑𝟏𝟒}\mathbf{\pm }\mathrm{𝟏𝟐}\mathbf{(}\mathrm{𝐬𝐭𝐚𝐭}\mathbf{.}\mathbf{)}\mathbf{\pm }\mathrm{𝟐𝟓}\mathbf{(}\mathrm{𝐬𝐲𝐬𝐭}\mathbf{.}\mathbf{)}`$ fb. The overall 4f-CC03 correction is +8 $`\mathbf{\pm }`$ 40 fb. A total of 1066 events are selected in the data.
The systematic uncertainties on the combined semileptonic cross section are summarised in Table 2. The largest contribution arises from the Monte Carlo statistical error on the 4f-CC03 correction. Uncertainties arising from the choice of the lepton isolation criterium and the probability cut are estimated from the change of efficiency following a bin-by-bin reweighting of the respective Monte Carlo one-dimensional distributions to the data. Background normalisation mainly affects the tau channel although there is also a contribution from residual Bhabha background in the electron channel. Decays of the $`𝐙`$ to electrons and muons are used to evaluate lepton identification uncertainties whilst the contribution from beam related backgrounds is estimated by superimposing the energy deposits from random triggers on to the simulated events.
To evaluate the individual cross sections a similar fit as for the fully leptonic events is used with the corresponding matrix of efficiencies and backgrounds. This yields
$`𝝈\mathbf{(}\mathrm{𝐖𝐖}\mathbf{}𝐞𝝂𝐪\overline{𝐪}\mathbf{)}`$ $`\mathbf{=}`$ $`\mathbf{2.41}\mathbf{\pm }\mathbf{0.14}\mathbf{(}\mathrm{𝐬𝐭𝐚𝐭}\mathbf{.}\mathbf{)}\mathbf{\pm }\mathbf{0.07}\mathbf{(}\mathrm{𝐬𝐲𝐬𝐭}\mathbf{.}\mathbf{)}\mathrm{𝐩𝐛}\mathbf{,}`$
$`𝝈\mathbf{(}\mathrm{𝐖𝐖}\mathbf{}𝝁𝝂𝐪\overline{𝐪}\mathbf{)}`$ $`\mathbf{=}`$ $`\mathbf{2.39}\mathbf{\pm }\mathbf{0.13}\mathbf{(}\mathrm{𝐬𝐭𝐚𝐭}\mathbf{.}\mathbf{)}\mathbf{\pm }\mathbf{0.06}\mathbf{(}\mathrm{𝐬𝐲𝐬𝐭}\mathbf{.}\mathbf{)}\mathrm{𝐩𝐛}\mathbf{,}`$
$`𝝈\mathbf{(}\mathrm{𝐖𝐖}\mathbf{}𝝉𝝂𝐪\overline{𝐪}\mathbf{)}`$ $`\mathbf{=}`$ $`\mathbf{2.23}\mathbf{\pm }\mathbf{0.17}\mathbf{(}\mathrm{𝐬𝐭𝐚𝐭}\mathbf{.}\mathbf{)}\mathbf{\pm }\mathbf{0.08}\mathbf{(}\mathrm{𝐬𝐲𝐬𝐭}\mathbf{.}\mathbf{)}\mathrm{𝐩𝐛}\mathbf{,}`$
where the systematic uncertainties are obtained by varying all input parameters of the fit according to their uncertainties.
The total $`\mathbf{}𝝂`$$`𝐪\overline{𝐪}`$ cross section is extracted by means of the same fit, under the assumption of lepton universality:
$$𝝈\mathbf{(}\mathrm{𝐖𝐖}\mathbf{}\mathbf{}𝝂𝐪\overline{𝐪}\mathbf{)}\mathbf{=}\mathbf{7.07}\mathbf{\pm }\mathbf{0.23}\mathbf{(}\mathrm{𝐬𝐭𝐚𝐭}\mathbf{.}\mathbf{)}\mathbf{\pm }\mathbf{0.12}\mathbf{(}\mathrm{𝐬𝐲𝐬𝐭}\mathbf{.}\mathbf{)}\mathrm{𝐩𝐛}\mathbf{,}$$
again consistent with a simple sum of the three channels.
### 2.3 $`\mathrm{𝐖𝐖}\mathbf{}𝐪\overline{𝐪}𝐪\overline{𝐪}`$ events
The analysis of $`\mathrm{𝐖𝐖}`$ decays to four jets is updated from Ref. and consists of a simple preselection followed by a fit to the distribution of the output of a neural network (NN) with 14 input variables.
In the preselection, the first step is to remove events with a large undetected initial state (ISR) photon from radiative returns to the Z by requiring that the modulus of the total longitudinal momentum of all objects is less than $`\mathbf{1.5}\mathbf{(}𝑴_{𝒗𝒊𝒔}\mathbf{}𝑴_𝒁\mathbf{)}`$ where $`𝑴_{𝒗𝒊𝒔}`$ is the observed visible mass. The particles are then forced into four jets using the DURHAM-PE algorithm and the value of $`𝒚_{\mathrm{𝟑𝟒}}`$, where a four-jet event becomes a three-jet event, is required to be greater than 0.001. To reject $`𝐪\overline{𝐪}`$ events with a visible ISR photon none of the four jets can have more than 95% of electromagnetic energy in a one degree cone around any particle included in the jet. Four-fermion final states where one of the fermions is a charged lepton are rejected by requiring that the maximum energy fraction of a single charged particle in a jet be smaller than 0.9.
At this point, 3438 events are selected in the data while $`\mathrm{𝟑𝟓𝟗𝟑}\mathbf{\pm }\mathrm{𝟕}`$ events are expected from Standard Model processes. This preselection has an efficiency of 98.4% for CC03 events and a purity of 35.9%.
The input variables for the NN are described in the Appendix and are related to the global event properties, the properties of jets, $`\mathrm{𝐖𝐖}`$ kinematics and the b-tag probabilities for the four jets. The NN is simplified with respect to that used for the analysis of the 183 GeV data , with fewer input variables and a slightly improved performance. The output distributions of the NN for the data compared with the signal and backgrounds predicted by the Monte Carlo are given in Fig. 2. In Table 1 results are given for a cut at 0.3 on the NN value.
The cross section is extracted by means of a binned maximum likelihood fit to the full NN output distribution. In the fit only the normalisation of the Monte Carlo signal is allowed to vary; all backgrounds, including $`\mathrm{𝐖𝐖}\mathbf{}\mathbf{}𝝂𝐪\overline{𝐪}`$, are kept fixed both in shape and normalisation. Systematic studies at the Z peak when two jet $`𝐪\overline{𝐪}`$ events are forced into four jets show small discrepancies between data and Monte Carlo. This affects the $`𝐪\overline{𝐪}`$ background shape and is quantified using the 14 NN variables to yield a correction of $`\mathbf{+}\mathrm{𝟓𝟎}\mathbf{\pm }\mathrm{𝟓𝟎}`$ fb to the value from the fit. The corrected fit result is:
$$𝝈\mathbf{(}\mathrm{𝐖𝐖}\mathbf{}𝐪\overline{𝐪}𝐪\overline{𝐪}\mathbf{)}\mathbf{=}\mathbf{6.89}\mathbf{\pm }\mathbf{0.23}\mathbf{(}\mathrm{𝐬𝐭𝐚𝐭}\mathbf{.}\mathbf{)}\mathbf{\pm }\mathbf{0.13}\mathbf{(}\mathrm{𝐬𝐲𝐬𝐭}\mathbf{.}\mathbf{)}\mathrm{𝐩𝐛}\mathbf{.}$$
If only events with a NN output greater than 0.3 are kept a value of $`\mathbf{6.89}\mathbf{\pm }\mathbf{0.23}\mathrm{𝐩𝐛}`$ is found showing no significant bias due to the $`𝐪\overline{𝐪}`$ background.
The sources contributing to the systematic uncertainty are summarised in Table 2. The largest contribution is the uncertainty in the $`𝐪\overline{𝐪}`$ background normalisation where a 5% variation is assumed. The contributions from uncertainties in the values of jet variables used in the event preselection and the NN input are assessed by adjusting their directions and energies according to residual differences observed between Z calibration data and Monte Carlo . Possible miscalibrations of the calorimeters are also taken into account. The beam background is studied in the same way as described in section 2.2.1. The Monte Carlo statistics error is dominated by the 4f-CC03 correction.
Uncertainties in the $`\mathrm{𝐖𝐖}`$ generator are evaluated by comparing samples of fully simulated Monte Carlo events generated with KORALW and EXCALIBUR. To establish a $`\mathrm{𝐖𝐖}`$ fragmentation uncertainty, the HERWIG generator was tuned at the $`𝐙`$ peak both for all flavours and non-b quark flavors. Then, the same samples of signal events generated with KORALW are fragmented with both JETSET and the appropriately tuned HERWIG. For the $`𝐪\overline{𝐪}`$ background uncertainty, a sample of events generated with KORALZ using JETSET is compared with a separate sample of pure HERWIG events to assess the effect of the choice of generator. Colour reconnection effects are estimated using the SK1 model in JETSET with a reconnection probability of 0.3 and the effect of Bose-Einstein correlations is estimated according to the scheme, denoted $`\mathrm{𝐁𝐄}_\mathrm{𝟑}`$, as proposed for the LUBOEI implementation in JETSET. The procedures followed are the same as those used in Ref. .
Several cross checks have been performed on the fit result to search for possible biases arising in the selection and full simulation of the $`\mathrm{𝐖𝐖}`$ events. The previous version of the analysis with a different preselection and NN gives $`\mathbf{6.81}\mathbf{\pm }\mathbf{0.23}\mathrm{𝐩𝐛}`$. Another estimate with a preselection based mostly on charged tracks and a six variable NN using only charged tracks gives $`\mathbf{6.83}\mathbf{\pm }\mathbf{0.25}\mathrm{𝐩𝐛}`$. Also, a selection based purely on calorimeter measurements with six input variables to a linear discriminant gives $`\mathbf{6.91}\mathbf{\pm }\mathbf{0.26}\mathrm{𝐩𝐛}`$. A variety of linear discriminant analyses using from 4 to 14 variables give results which vary from 6.65 to 6.77 pb. All these checks give results which are consistently lower than the GENTLE prediction.
## 3 Total cross section
The total cross section is obtained from a fit to all channels described above assuming the Standard Model branching fractions, the only unknown being the total cross section. The fit uses the matrices of efficiencies and backgrounds for the various analyses and yields
$$𝝈_{\mathrm{𝐖𝐖}}\mathbf{=}\mathbf{15.71}\mathbf{\pm }\mathbf{0.34}\mathbf{(}\mathrm{𝐬𝐭𝐚𝐭}\mathbf{.}\mathbf{)}\mathbf{\pm }\mathbf{0.18}\mathbf{(}\mathrm{𝐬𝐲𝐬𝐭}\mathbf{.}\mathbf{)}\mathrm{𝐩𝐛}\mathbf{.}$$
Assuming no additional unexpected decay mode, the result is not significantly different if the branching fractions of the Standard Model decay modes are unconstrained.
The measurement is 5.5% lower than the GENTLE prediction. However, new calculations including full $`𝓞\mathbf{(}𝜶\mathbf{)}`$ electroweak corrections, calculable in the double pole approximation (DPA) have recently appeared. Two Monte Carlo programs, YFSWW3 and RacoonWW , are being developed. First numerical calculations find cross sections, respectively, 1.9% and 2.4% lower than GENTLE at 189 GeV. The predictions from RacoonWW also include soft-photon exponentiation and leading log corrections for initial state radiation beyond $`𝓞\mathbf{(}𝜶\mathbf{)}`$ in addition to the calculations described in Ref. . The uncertainty in the two new models is expected to be of the order of 0.5% .
Fig. 3 shows the total cross section measured as a function of the CM energy. The predictions of the two more complete YFSWW3 and RacoonWW calculations are also shown and are in better agreement with the experimental results than GENTLE. At 189 $`\mathrm{𝐆𝐞𝐕}`$, the measurement is 3.8% (1.5 standard deviations) below the YFSWW prediction and 3.2% (1.3 standard deviations) below the RacoonWW prediction. Taking into account also the cross section values measured by ALEPH at 172 and 183 $`\mathrm{𝐆𝐞𝐕}`$, the data are $`\mathbf{4.6}\mathbf{\pm }\mathbf{2.0}\mathbf{\%}\mathbf{,}\mathbf{2.8}\mathbf{\pm }\mathbf{2.0}\mathbf{\%}`$, and $`\mathbf{2.3}\mathbf{\pm }\mathbf{2.0}\mathbf{\%}`$ below the GENTLE, YFSWW and RacoonWW predictions respectively. These values use the signal efficiencies determined with KORALW and the quoted systematic uncertainty takes no account of any efficiency difference which may arise from the new calculations.
## 4 Branching fractions and $`𝑽_{\mathrm{𝐜𝐬}}`$
The same fit as for the total cross section is performed combining the data samples collected at 161, 172, 183 and 189 $`\mathrm{𝐆𝐞𝐕}`$ CM energies.
Without assuming lepton coupling universality, the seven unknowns are the three individual leptonic branching fractions and the four total cross sections at 161, 172, 183 and 189 $`\mathrm{𝐆𝐞𝐕}`$. The hadronic branching fraction is set to $`\mathrm{𝟏}\mathbf{}𝐁_𝐞\mathbf{}𝐁_𝝁\mathbf{}𝐁_𝝉`$. The fitted leptonic branching fractions are
$`𝐁\mathbf{(}𝐖\mathbf{}𝐞𝝂\mathbf{)}`$ $`\mathbf{=}`$ $`\mathbf{11.35}\mathbf{\pm }\mathbf{0.46}\mathbf{(}\mathrm{𝐬𝐭𝐚𝐭}\mathbf{.}\mathbf{)}\mathbf{\pm }\mathbf{0.17}\mathbf{(}\mathrm{𝐬𝐲𝐬𝐭}\mathbf{.}\mathbf{)}\mathbf{\%}\mathbf{,}`$
$`𝐁\mathbf{(}𝐖\mathbf{}𝝁𝝂\mathbf{)}`$ $`\mathbf{=}`$ $`\mathbf{11.10}\mathbf{\pm }\mathbf{0.44}\mathbf{(}\mathrm{𝐬𝐭𝐚𝐭}\mathbf{.}\mathbf{)}\mathbf{\pm }\mathbf{0.16}\mathbf{(}\mathrm{𝐬𝐲𝐬𝐭}\mathbf{.}\mathbf{)}\mathbf{\%}\mathbf{,}`$
$`𝐁\mathbf{(}𝐖\mathbf{}𝝉𝝂\mathbf{)}`$ $`\mathbf{=}`$ $`\mathbf{10.51}\mathbf{\pm }\mathbf{0.55}\mathbf{(}\mathrm{𝐬𝐭𝐚𝐭}\mathbf{.}\mathbf{)}\mathbf{\pm }\mathbf{0.22}\mathbf{(}\mathrm{𝐬𝐲𝐬𝐭}\mathbf{.}\mathbf{)}\mathbf{\%}\mathbf{,}`$
and are consistent with lepton universality and the Standard Model expectations. Due to cross-contaminations in the identification of W decays to $`𝝉𝝂`$, $`𝐞𝝂`$ or $`𝝁𝝂`$, the measured $`𝐁\mathbf{(}𝐖\mathbf{}𝝉𝝂\mathbf{)}`$ is 26% anticorrelated with $`𝐁\mathbf{(}𝐖\mathbf{}𝐞𝝂\mathbf{)}`$ and 24% anticorrelated with $`𝐁\mathbf{(}𝐖\mathbf{}𝝁𝝂\mathbf{)}`$. The $`𝐁\mathbf{(}𝐖\mathbf{}𝐞𝝂\mathbf{)}`$ is 4.6% anticorrelated with $`𝐁\mathbf{(}𝐖\mathbf{}𝝁𝝂\mathbf{)}`$.
If lepton universality is assumed a fit for $`𝐁\mathbf{(}𝐖\mathbf{}𝐪\overline{𝐪}\mathbf{)}`$ and the total cross sections at the four energies yields
$$𝐁\mathbf{(}𝐖\mathbf{}𝐪\overline{𝐪}\mathbf{)}\mathbf{=}\mathbf{66.97}\mathbf{\pm }\mathbf{0.65}\mathbf{(}\mathrm{𝐬𝐭𝐚𝐭}\mathbf{.}\mathbf{)}\mathbf{\pm }\mathbf{0.32}\mathbf{(}\mathrm{𝐬𝐲𝐬𝐭}\mathbf{.}\mathbf{)}\mathbf{\%}\mathbf{.}$$
This result can be expressed in terms of the individual couplings of the W to quark-antiquark pairs:
$`{\displaystyle \frac{𝐁\mathbf{(}𝐖\mathbf{}𝐪\overline{𝐪}\mathbf{)}}{\mathrm{𝟏}\mathbf{}𝐁\mathbf{(}𝐖\mathbf{}𝐪\overline{𝐪}\mathbf{)}}}`$ $`\mathbf{=}`$ $`\mathbf{(}\mathbf{|}𝑽_{\mathrm{𝐮𝐝}}\mathbf{|}^\mathrm{𝟐}\mathbf{+}\mathbf{|}𝑽_{\mathrm{𝐜𝐝}}\mathbf{|}^\mathrm{𝟐}\mathbf{+}\mathbf{|}𝑽_{\mathrm{𝐮𝐬}}\mathbf{|}^\mathrm{𝟐}\mathbf{+}\mathbf{|}𝑽_{\mathrm{𝐜𝐬}}\mathbf{|}^\mathrm{𝟐}\mathbf{+}\mathbf{|}𝑽_{\mathrm{𝐮𝐛}}\mathbf{|}^\mathrm{𝟐}\mathbf{+}\mathbf{|}𝑽_{\mathrm{𝐜𝐛}}\mathbf{|}^\mathrm{𝟐}\mathbf{)}\mathbf{(}\mathrm{𝟏}\mathbf{+}𝜶_𝒔\mathbf{(}𝒎_𝐖^\mathrm{𝟐}\mathbf{)}\mathbf{/}𝝅\mathbf{)}\mathbf{.}`$
The least well known of these is $`\mathbf{|}𝑽_{\mathrm{𝐜𝐬}}\mathbf{|}`$. Using the world average value of $`𝜶_𝒔`$($`𝒎_𝐙^\mathrm{𝟐}`$) evolved to $`𝒎_𝐖^\mathrm{𝟐}`$ , $`𝜶_𝒔`$($`𝒎_𝐖^\mathrm{𝟐}`$) = $`\mathbf{0.121}\mathbf{\pm }\mathbf{0.002}`$, and the squared sum of the other measured CKM matrix elements which is $`\mathbf{1.05}\mathbf{\pm }\mathbf{0.01}`$, the measured hadronic branching fraction is
$$\mathbf{|}𝑽_{\mathrm{𝐜𝐬}}\mathbf{|}\mathbf{=}\mathbf{0.951}\mathbf{\pm }\mathbf{0.030}\mathbf{(}\mathrm{𝐬𝐭𝐚𝐭}\mathbf{.}\mathbf{)}\mathbf{\pm }\mathbf{0.015}\mathbf{(}\mathrm{𝐬𝐲𝐬𝐭}\mathbf{.}\mathbf{)}\mathbf{.}$$
## 5 Conclusions
The W-pair production cross section at $`\sqrt{𝒔}`$ = 188.63 $`\mathrm{𝐆𝐞𝐕}`$ has been measured in all decay channels from an integrated luminosity of 174.20 $`\mathrm{𝐩𝐛}^\mathbf{}\mathrm{𝟏}`$. The total cross section is found to be
$$𝝈_{\mathrm{𝐖𝐖}}\mathbf{=}\mathbf{15.71}\mathbf{\pm }\mathbf{0.34}\mathbf{(}\mathrm{𝐬𝐭𝐚𝐭}\mathbf{.}\mathbf{)}\mathbf{\pm }\mathbf{0.18}\mathbf{(}\mathrm{𝐬𝐲𝐬𝐭}\mathbf{.}\mathbf{)}\mathrm{𝐩𝐛}\mathbf{.}$$
This result is 5.5% (1.9 standard deviations) lower than the GENTLE prediction but in better agreement with more recent calculations. It agrees well with the recent measurement by the DELPHI Collaboration at the same CM energy.
After inclusion of the data taken at CM energies of 161, 172 and 183 $`\mathrm{𝐆𝐞𝐕}`$, the hadronic decay branching fraction is found to be $`\mathbf{66.97}\mathbf{\pm }\mathbf{0.65}\mathbf{(}\mathrm{𝐬𝐭𝐚𝐭}\mathbf{.}\mathbf{)}\mathbf{\pm }\mathbf{0.32}\mathbf{(}\mathrm{𝐬𝐲𝐬𝐭}\mathbf{.}\mathbf{)}`$% which is used to determine the CKM matrix element $`\mathbf{|}𝑽_{\mathrm{𝐜𝐬}}\mathbf{|}`$ equal to $`\mathbf{0.951}\mathbf{\pm }\mathbf{0.030}\mathbf{(}\mathrm{𝐬𝐭𝐚𝐭}\mathbf{.}\mathbf{)}\mathbf{\pm }\mathbf{0.015}\mathbf{(}\mathrm{𝐬𝐲𝐬𝐭}\mathbf{.}\mathbf{)}`$.
## Acknowledgements
We would like to thank S. Dittmaier and W. Placzek for helpful discussions on DPA. It is a pleasure to congratulate our colleagues from the CERN accelerator divisions for the successful operation of LEP at 189 GeV. We are indebted to the engineers and technicians in all our institutions for their contributions to the excellent performance of ALEPH. Those of us from non-member countries thank CERN for its hospitality.
## Appendix: Neural network input variables for the $`𝐪\overline{𝐪}`$$`𝐪\overline{𝐪}`$ selection
The neural network hadronic event selection uses 14 variables. These are based on global event properties, heavy quark flavour tagging, jet properties and $`\mathrm{𝐖𝐖}`$ kinematics and are listed below. The four jets are numbered in order of decreasing energy.
Global event properties
* Thrust
* Sphericity
* Missing energy
* Sum of the four smallest interjet angles
Heavy flavour tagging
* Probability of an event being a light quark (uds) event based upon impact parameter significance of charged particles in the event
Jet properties
* Maximum electromagnetic energy fraction of a jet in any one degree cone
* Maximum summed charged particle energy fraction of a jet
* Minimum number of charged particles in a jet
$`\mathrm{𝐖𝐖}`$ kinematics
* Angle between Jet2 and Jet3
The following jet related variables are determined from kinematically fitted jet momenta.
* Energy of Jet1
* Energy of Jet3
* Energy of Jet4
* Smallest jet mass
* Second smallest jet mass
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# Contents
## 1 Introduction and Summary
Matrix models provide a useful framework in which to extract quantitative information about complex statistical systems given only their general symmetry properties. The key feature which enables this description is the existence of universal behaviour in these models. Over the years random matrix models have been utilized in a wide range of studies in condensed matter systems , quantum chromodynamics , and low-dimensional string theory . They are also useful for solving various statistical mechanical models and combinatorial problems . One interpretation of random matrix ensembles is in terms of generators of ’t Hooft diagram expansions which correspond to discretizations of compact Riemann surfaces. The continuum limits of such discretizations are realized as phase transitions in the statistical mechanics of the matrix model and characterized by universal data at the critical points.
Despite the success of matrix model technology over the years, there remain several issues in random surface theory which are not adequately addressed by standard approaches. The simplest, and one of the most fundamental issues regards a self-consistent definition of the vacuum of pure two-dimensional quantum gravity. The conventional matrix model methods for investigating this theory lead to inconsistencies. For instance, bosonic matrix integrals are typically ill-defined in the physically interesting region of parameter space corresponding to the non-perturbative continuum limit. Furthermore, they are unable to describe the coupling of gravity to conformal matter fields of central charge $`c>1`$. It has been suggested that above this $`c=1`$ barrier the two-dimensional geometry degenerates into a tree-like or branched polymer phase.
In this paper we will study a model of discretized random surfaces which is in part motivated by such issues. It is parameterized by a zero-dimensional action involving matrices with Grassmann-valued elements . This model is known as the adjoint fermion one-matrix model and it is defined by the partition function
$$Z_N=\underset{\mathrm{Gr}(N)^c}{}𝑑\psi 𝑑\overline{\psi }\mathrm{e}^{N\mathrm{tr}V(\overline{\psi }\psi )}$$
(1.1)
where
$$V(\overline{\psi }\psi )=\underset{k=1}{\overset{K}{}}\frac{g_k}{k}\left(\overline{\psi }\psi \right)^k$$
(1.2)
is a polynomial potential of order $`K\mathrm{}`$, and $`d\psi d\overline{\psi }=_{i,j}d\psi _{ij}d\overline{\psi }_{ij}`$ is the standard Berezin measure on the complexified Grassmann algebra $`\mathrm{Gr}(N)^c=\mathrm{Gr}(N)_{\text{}}\text{}`$. For any matrix $`X`$ the trace is normalized as $`\mathrm{tr}X=_iX_{ii}`$. The matrix elements $`\psi _{ij}`$ and $`\overline{\psi }_{ij}`$, $`i,j=1,\mathrm{},N`$, are independent, complex-valued Grassmann numbers with the usual rules for complex conjugation, $`\overline{\psi }_{ij}^{}=\psi _{ij}`$, $`\psi _{ij}^{}=\overline{\psi }_{ij}`$, and $`(\overline{\psi }_{ij}\psi _{kl})^{}=\psi _{kl}^{}\overline{\psi }_{ij}^{}`$ (This ensures that the generating function (1.1) is real-valued). The large $`N`$ limit of this model has been studied in and shown to exhibit the range of critical behaviour seen in the usual bosonic one-matrix models. The topological expansion of the matrix integral (1.1) was considered in and some novel critical behaviour was observed in . The fermionic matrix model can also be coupled to an ordinary complex bosonic one to produce a supersymmetric matrix model, which has been used to describe solutions of certain combinatorial problems and also to generate statistical models of branched polymers .
The most attractive feature of the fermionic matrix model (1.1) is that the integration over Grassmann variables is always well-defined and the theory yields finite, computable observables. The dimension $`N`$ of the Grassmann matrices acts as a cutoff on the number of terms in the partition function (1.1) which counts all fat-graphs generated. The fact that the partition function is a polynomial in the coupling constants $`g_k`$ for finite $`N`$, coupled with the fact that inner products in the Grassmann integration measure are effectively calculable, leads to explicit expressions for statistical quantities in this class of models. In the following we will use this property to solve analytically the combinatorial problem of counting the dynamical triangulations generated by the fermionic partition function via an explicit determinant representation of (1.1).
Since the generating function (1.1) is always a well-defined convergent object at finite $`N`$, one might expect that its large $`N`$ limit is also well-defined. In perturbation theory, the Feynman rules for fermionic matrices include a factor of $`1`$ for each fermion loop thereby resulting in cancellations between large numbers of Feynman diagrams. This observation has led to the conjecture that the genus expansion of the matrix integral (1.1) in the large $`N`$ limit is an alternating series which may be Borel summable. The evidence for this striking property of the adjoint fermion one-matrix model is its intimate relationship with generalized Penner matrix models which are defined by the Hermitian matrix integrals
$$Z_{N,\alpha }^\mathrm{H}=\underset{u(N)}{}𝑑\varphi \mathrm{e}^{N\mathrm{tr}\left(V(\varphi )+\alpha \mathrm{log}\varphi \right)}$$
(1.3)
where $`d\varphi =_{i,j}d\varphi _{ij}`$ is the Haar measure on the Lie algebra $`u(N)`$ of the $`N\times N`$ unitary group $`U(N)`$. It has been argued that the matrix model (1.3) is equivalent to (1.1) for $`\alpha =2`$, with the Hermitian matrix $`\varphi `$ identified with $`\overline{\psi }\psi `$. This has been established in the large $`N`$ limit by a saddle-point computation on $`\varphi \overline{\psi }\psi `$ , and by showing that their Schwinger-Dyson (loop) equations are identical order by order in the large $`N`$ expansions of the matrix integrals . However, for $`\alpha <0`$ the integration over Hermitian matrices in (1.3) is not well-defined because of the logarithmic divergence at $`\varphi =0`$. In this region the integral can only be defined by analytical continuation and the matrix model only makes sense at $`N=\mathrm{}`$. This continuation produces complex-valued endpoints for the support of the corresponding distribution of $`\varphi `$ eigenvalues, analogous to what naturally occurs in fermionic matrix models . It is this complexification of the support of the distribution of eigenvalues that is asserted to result in an alternating genus expansion. Moreover, it can be argued that the topological expansion coincides, modulo these sign factors, with the usual Painlevé expansions of Hermitian matrix integrals with polynomial potentials .
The near equivalence of the genus expansions strongly suggests a correspondence between the Feynman graph expansions in the bosonic and fermionic models. One way of understanding these issues is to appeal to a simplified version of the matrix integral (1.1), namely a fermionic vector model defined by anticommuting vector elements $`\psi _i`$ and $`\overline{\psi }_i`$, $`i=1,\mathrm{},N`$ . A fermionic vector model is related to a bosonic vector model, defined with the same polynomial potential $`V`$, by a simple mapping of its Feynman diagrams and an analytical continuation $`N\frac{N}{2}`$ of the vector dimension which results in a bosonic perturbation series with a cutoff on the number of terms generated. This analytical continuation can be understood in terms of a complex-valued saddle-point of the bosonic model in the large $`N`$ limit, leading to an alternating “genus expansion” which produces a well-defined, unique function . As bosonic vector models generate random models of branched polymers, the fermionic vector model thereby produces a statistical theory of branched polymers whose continuum characteristics are equivalent order by order in large $`N`$ to the conventional models, but whose double-scaling expansion is a non-perturbatively well-defined function.
The arguments given in are based on the class of models with odd polynomial potential, i.e. $`V(\varphi )=V(\varphi )`$. In these cases, the endpoints of the support contour of the spectral density are located symmetrically about the origin on the imaginary axis in the complex plane . Generic polynomial potentials are difficult to treat because the corresponding support contours are located asymmetrically in the complex plane and the loop equations become cumbersome to deal with. Moreover, although the equations of motion for the matrix models (1.1) and (1.3) are identical at each order of the large $`N`$ expansion, beyond the leading order they must be solved with different boundary conditions and the non-perturbative solutions are not the same in the two cases . One of the main results of this paper will be the clarification of the role played by generalized Penner models in fermionic matrix integrals. The analytical expressions for the partition function (1.1) at finite $`N`$ will be interpreted in terms of Toeplitz determinants, which will naturally lead to a proof of the equivalence of the fermionic one-matrix model and a unitary version of the Penner one-matrix model,
$$Z_N^\mathrm{U}=\underset{U(N)}{}[dU]\mathrm{e}^{N\mathrm{tr}\left(V(U)\mathrm{log}U\right)}$$
(1.4)
where $`[dU]`$ is the invariant Haar measure on the unitary group $`U(N)`$. The important feature of this equivalence is that it holds for any matrix dimension $`N`$. The finiteness of the fermionic matrix model at finite $`N`$ is captured by the compactness of the field variable in (1.4). As we will see, the logarithmic term in the effective potential of (1.4) gives certain restrictions on the configurations which contribute to the partition function and its observables when the unitary matrix model is treated from a group theoretic perspective. These restrictions naturally reproduce the basic characteristics of the matrix integral (1.1), and moreover show precisely how the analytic, functional forms of the models (1.1) and (1.4) are equivalent for any value of $`N`$. We will also prove directly that the loop equations of the fermionic and unitary matrix models are identical, and that they admit the same, perturbative Gaussian boundary conditions. It is a remarkable fact that a unitary matrix model such as (1.4) admits an interpretation in terms of random surfaces.
The biggest advantage of the mapping of the fermionic matrix model to a unitary one is that we now have an eigenvalue model to analyse. Grassmann-valued matrices are not diagonalizable, and a direct analysis of fermionic matrix models is only possible using the method of loop equations . In this paper we shall exploit this eigenvalue representation to give a systematic and detailed account of the exotic properties possessed by fermionic matrix models. In particular we will develop a generalization of the orthogonal polynomial technique on the circle which is necessary to study the unitary matrix model defined by (1.4), and describe some of their functional analytic properties. These methods will naturally lead, irrespective of the parity of the potential $`V`$, to a precise investigation of the double-scaling limit of the fermionic matrix model which determines a non-perturbative, all-genus continuum theory of two-dimensional quantum gravity coupled to matter. We will show that the fermionic matrix model leads to a Borel summable genus expansion, in contradistinction to Hermitian one-matrix models. Unlike the vector models, the diagrammatic, finite $`N`$ situation is not quite so simple in the case of fermionic matrices, but the final qualitative result will be the same. In particular, we shall obtain a non-perturbative definition of pure two-dimensional quantum gravity which is equivalent to the usual models order by order in the genus expansion, but which is given by a well-defined unique generating function. This non-perturbative model is what we shall call “fermionic quantum gravity” in the following.
Another advantage of the eigenvalue representation is that it will allow us describe the fermionic one-matrix model (1.1) as an integrable system. Generally, matrix models, being examples of exactly solvable systems, are intimately related to certain reductions of well-known integrable hierarchies of differential equations . In the usual bosonic Hermitian one-matrix models, describing two-dimensional quantum gravity, the partition functions are related to the $`\tau `$-functions of the integrable Toda chain hierarchy. In the following we will see that the fermionic gravity theory is related to the $`\tau `$-function of a particular deformation of this hierarchy known as the relativistic Toda chain. We can therefore understand many of the novel aspects of the fermionic theory, such as what sort of matter coupling it involves, in terms of this deformation. Furthermore, we will develop an operator theoretic approach to the equations of motion of the matrix integral (1.1) and gain in this setting a precise picture of how this deformation leads to a well-defined Borel summable partition function of the statistical model. This sets the stage for a description of the integrable differential hierarchies satisfied by the continuum theory which generalizes the usual KdV flow structure of gravity . It will also indicate that the partition function (1.1) in the continuum limit serves as a concrete, non-perturbative definition of two-dimensional quantum gravity.
The organization of the remainder of this paper is as follows. In section 2 we will derive an analytical expression for the partition function (1.1) and use it to describe the solution to the problem of counting fermionic fat-graphs. We shall see that the matrix integral generates a novel class of discretizations of Riemann surfaces whose polygons are two-colourable. We will also briefly describe what sort of matter-coupled gravity theory is represented by the fermionic one-matrix model, and prove its equivalence to the unitary one-matrix model (1.4). In section 3 we will present the formal orthogonal polynomial solution of the fermionic matrix model. In section 4 we will describe the large $`N`$ limit of the model from the perspective of these orthogonal polynomials. We will find that the adjoint fermion one-matrix model actually possesses a sort of “internal” branched polymer critical behaviour which can be understood in terms of the fermionic characteristics of its fat-graph combinatorics. We shall also demonstrate that the usual universality classes of two-dimensional gravity arise, and that they cannot be smoothly connected to the branched polymer phases of the theory. Of course, the polymer effect is washed out at $`N=\mathrm{}`$ because the surface effect is of higher order in $`N`$. We also construct the genus expansion of the theory and establish its Borel summability explicitly. In section 5 we present the operator approach to the fermionic matrix model and use it to interpret the partition function as a $`\tau `$-function. We derive the Virasoro constraints satisfied by this $`\tau `$-function, which at the same time establishes explicitly the equivalence of the loop equations of the matrix models (1.1), (1.3) with $`\alpha =2`$, and (1.4), in the large $`N`$ expansion. We then use these descriptions to give an functional analytic explanation for the nature of the topological expansion in the fermionic one-matrix model. Finally, in appendix A we describe properties of observables of the statistical model (1.1) using the unitary matrix model and its orthogonal polynomials, while in appendix B we show that the orthogonal polynomials of the Gaussian fermionic one-matrix model are given in terms of confluent hypergeometric polynomials.
## 2 Determinant Representations of the Partition Function
In this section we will present an analytic solution to the problem of counting fermionic fat-graphs. We will obtain a closed expression for the generating function of such discretized random surfaces, and at the same time we will be naturally led to the equivalence of the fermionic matrix model and the unitary matrix model which will occupy the bulk of our analysis in subsequent sections. We will also obtain in this way a geometric picture of what sort of theory of random surfaces is represented by the adjoint fermion one-matrix model and also of the precise role of the corresponding Penner interaction.
### 2.1 Hermitian Representation
One intriguing feature of the fermionic one-matrix model (1.1) is the extent to which it can be solved at finite-$`N`$. This is a consequence of the convergence of the integration over Grassmann variables, in contrast to the case of bosonic matrix models. The combinatorics of the fat-graph expansion of the fermionic partition function $`Z_N`$ can be determined by using the observation of that the partition function and observables of the matrix model (1.1) can be evaluated explicitly by introducing two $`N\times N`$ Hermitian matrices $`\varphi `$ and $`\lambda `$ defined by
$$1=\underset{u(N)}{}d\varphi \delta (\varphi \overline{\psi }\psi ),\delta (\varphi )=\underset{u(N)}{}\frac{d\lambda }{(2\pi )^{N^2}}\mathrm{e}^{i\mathrm{tr}\lambda \varphi }$$
(2.1)
By inserting these definitions into (1.1), the fermionic partition function can be written as
$`Z_N`$ $`=`$ $`{\displaystyle \underset{u(N)}{}}𝑑\varphi \mathrm{e}^{N\mathrm{tr}V(\varphi )}{\displaystyle \underset{\mathrm{Gr}(N)^c}{}}𝑑\psi 𝑑\overline{\psi }{\displaystyle \underset{u(N)}{}}{\displaystyle \frac{d\lambda }{(2\pi )^{N^2}}}\mathrm{e}^{i\mathrm{tr}\lambda (\varphi \overline{\psi }\psi )}`$ (2.2)
$`=`$ $`{\displaystyle \frac{1}{(2\pi )^{N^2}}}{\displaystyle \underset{u(N)}{}𝑑\varphi 𝑑\lambda \mathrm{e}^{N\mathrm{tr}V(\varphi )+i\mathrm{tr}\lambda \varphi }\stackrel{N}{det}(i\lambda )}`$
$`=`$ $`{\displaystyle \underset{u(N)}{}}𝑑\varphi \mathrm{e}^{N\mathrm{tr}V(\varphi )}\stackrel{N}{det}\left(\frac{}{\varphi }\right)\delta (\varphi )`$
so that
$$Z_N=\stackrel{N}{det}\left(\frac{}{\varphi }\right)\mathrm{e}^{N\mathrm{tr}V(\varphi )}|_{\varphi =0}$$
(2.3)
A similar expression can also be derived for the correlators of the fermionic one-matrix model . However, it is difficult to write down an explicit expression for (2.3) which is informative and amenable to analysis.
To evaluate the determinant in (2.3) explicitly, we shall instead write the fermionic matrix model as a random theory of the two Hermitian matrices $`\varphi `$ and $`\lambda `$ introduced above, keeping only the second line of the identity (2.2). In this way, the fermionic one-matrix model can be written as the Hermitian two-matrix model
$$Z_N=\frac{N^{N(N+1)}}{(2\pi )^{N^2}}\underset{u(N)}{}𝑑\varphi 𝑑\lambda \mathrm{e}^{N\mathrm{tr}\left(V(\varphi )+\mathrm{log}(i\lambda )+i\lambda \varphi \right)}$$
(2.4)
where we have rescaled $`\lambda _{ij}N\lambda _{ij}`$. The doubling of degrees of freedom is required to compensate for both the matrix $`\psi `$ and its adjoint $`\overline{\psi }`$. The Hermitian two-matrix model (2.4) describes two-dimensional discretized gravity with a single type of matter state at each of the vertices. It does not admit the usual $`\text{}_2`$ Ising symmetry of a classical spin lattice that characterizes the unitary minimal models of rational conformal field theory . Thus the matter states in the worldsheet interpretation of the fermionic matrix model are more complicated degrees of freedom than just simple Ising spins or other conventional types of conformal matter. The logarithmic nature of the effective two-matrix potential in (2.4) is just another reflection of the connection with Penner matrix models. The original Penner model ($`g_k=0`$ for $`k>1`$, $`g_1=1`$ and $`\alpha =1`$ in (1.3)) we used to calculate the virtual Euler characteristics of the moduli spaces of compact Riemann surfaces. For $`\alpha >0`$ the generalized models are intimately related to the coupling of gravity to matter fields of central charge $`c=1`$ . Note that the Penner interaction in (2.4) is well-defined. We may refer to the fermionic matrix model as a model of two-dimensional quantum gravity interacting with “Penner matter”. As we will see later on, it is this worldsheet model which leads to a well-defined nonperturbative theory of quantum gravity in two-dimensions which we will call “fermionic gravity”.
The partition function (2.4) can be written in terms of a double eigenvalue distribution by diagonalizing the Hermitian matrices $`\varphi =U\mathrm{diag}(\varphi _1,\mathrm{},\varphi _N)U^{}`$ and $`\lambda =V\mathrm{diag}(\lambda _1,\mathrm{},\lambda _N)V^{}`$ by unitary transformations, where $`\varphi _i,\lambda _i\text{}`$ are the eigenvalues of $`\varphi ,\lambda `$. Computing the Jacobian for the change of integration measure and using the Itzykson-Zuber formula to integrate out the unitary degrees of freedom, we arrive at the eigenvalue model
$$Z_N=c_N\underset{i=1}{\overset{N}{}}\underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑\varphi _i𝑑\lambda _i\lambda _i^N\mathrm{e}^{NV(\varphi _i)+iN\lambda _i\varphi _i}\mathrm{\Delta }(\varphi _1,\mathrm{},\varphi _N)\mathrm{\Delta }(\lambda _1,\mathrm{},\lambda _N)$$
(2.5)
where $`c_N=(iN^2)^{N(N+1)/2}/(2\pi )^N_{n=1}^Nn!`$ and
$$\mathrm{\Delta }(x_1,\mathrm{},x_N)=\underset{i,j}{det}\left[x_i^{j1}\right]=\underset{i<j}{}(x_ix_j)$$
(2.6)
is the Vandermonde determinant. Using (2.6) we can then write the partition function as
$$Z_N=c_NN!\underset{i,j}{det}\left[(\varphi ^{i1},\lambda ^{j1})\right]$$
(2.7)
where we have introduced the inner product
$$(F(\varphi ),G(\lambda ))\underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑\varphi 𝑑\lambda \lambda ^N\mathrm{e}^{NV(\varphi )+iN\lambda \varphi }F(\varphi )G(\lambda )$$
(2.8)
on the vector space $`\text{}[\varphi ]\text{}[\lambda ]`$ of complex coefficient polynomials in $`(\varphi ,\lambda )`$. A remarkable feature of the fermionic matrix model is the extent to which its inner product (2.8) can be computed. For arbitrary polynomial potential, integrating over $`\lambda `$ first in (2.8) leads to the result
$$(F(\varphi ),\lambda ^k)=\frac{2\pi }{N}\left(\frac{i}{N}\frac{}{\varphi }\right)^{N+k}\left[F(\varphi )\mathrm{e}^{NV(\varphi )}\right]|_{\varphi =0}$$
(2.9)
which is valid for any analytic function $`F(\varphi )`$ and any integer $`k0`$. Using (2.9), the representation (2.7) of the partition function as the determinant of the moment matrix of the inner product (2.8) thereby yields
$$Z_N=(1)^{[N]_2}\underset{i,j}{det}[(\frac{}{\varphi }+NV^{}(\varphi ))^{N+ji}1|_{\varphi =0}]$$
(2.10)
where, for any integer $`K1`$, $`[N]_K`$ denotes the largest integer less than or equal to $`N/K`$. Here and in the following a prime denotes differentiation. The expression (2.10) is the desired explicit expansion of the determinant operator in (2.3).
### 2.2 Combinatorics of Fermionic Ribbon Graphs
In contrast to the determinant form (2.3), the representation (2.10) leads to an explicit expression for the perturbation series of the fermionic one-matrix model. For a polynomial potential (1.2) of order $`K`$, the determinant in (2.10) may be evaluated explicitly by expanding the exponential function in its Taylor series and applying the multinomial theorem to each term in the expansion to get
$`{\displaystyle \frac{^L}{\varphi ^L}}\mathrm{e}^{NV(\varphi )}|_{\varphi =0}`$ $`=`$ $`(Ng_1)^L{\displaystyle \underset{k_2=0}{\overset{[L]_2}{}}}\mathrm{}{\displaystyle \underset{k_K=0}{\overset{[L]_K}{}}}{\displaystyle \frac{L!}{(L2k_2\mathrm{}Kk_K)!k_2!\mathrm{}k_K!}}`$ (2.11)
$`\times {\displaystyle \underset{\mathrm{}=2}{\overset{K}{}}}\left({\displaystyle \frac{g_{\mathrm{}}}{\mathrm{}N^\mathrm{}1(g_1)^{\mathrm{}}}}\right)^k_{\mathrm{}}`$
for any integer $`L0`$. Normalizing by the Gaussian model for which $`V(\varphi )=g_1\varphi `$ and $`Z_N^{\mathrm{Gauss}}=(1)^{[N]_2}(Ng_1)^{N^2}`$, the partition function (2.10) may be written as
$$\frac{Z_N}{Z_N^{\mathrm{Gauss}}}=\underset{i,j}{det}\left[\underset{k_2=0}{\overset{[N+ji]_2}{}}\mathrm{}\underset{k_K=0}{\overset{[N+ji]_K}{}}\frac{(N+ji)!}{\left(N+ji_{\mathrm{}}\mathrm{}k_{\mathrm{}}\right)!k_2!\mathrm{}k_K!}\underset{\mathrm{}=2}{\overset{K}{}}\left(\frac{g_{\mathrm{}}}{\mathrm{}N^\mathrm{}1(g_1)^{\mathrm{}}}\right)^k_{\mathrm{}}\right]$$
(2.12)
We can extract the sums in (2.12) out of the determinant by using the multilinearity of the determinant as a function of its $`N`$ rows to get
$`{\displaystyle \frac{Z_N}{Z_N^{\mathrm{Gauss}}}}`$ $`=`$ $`{\displaystyle \underset{k_2^{(1)},\mathrm{},k_2^{(N)}0}{}}\mathrm{}{\displaystyle \underset{k_K^{(1)},\mathrm{},k_K^{(N)}0}{}}\left[{\displaystyle \underset{m=2}{\overset{K}{}}}{\displaystyle \underset{n=1}{\overset{N}{}}}{\displaystyle \frac{1}{k_m^{(n)}!}}\left({\displaystyle \frac{g_m}{mN^{m1}(g_1)^m}}\right)^{k_m^{(n)}}\right]`$ (2.13)
$`\times \underset{i,j}{det}\left[{\displaystyle \frac{(N+ji)!}{\left(N+ji_{\mathrm{}}\mathrm{}k_{\mathrm{}}^{(i)}\right)!}}{\displaystyle \underset{\mathrm{}=2}{\overset{K}{}}}\mathrm{\Theta }([N+ji]_{\mathrm{}}k_{\mathrm{}}^{(i)})\right]`$
where $`\mathrm{\Theta }`$ denotes the step function with the convention $`\mathrm{\Theta }(0)=1`$.
The expression (2.13) leads to a relatively straightforward expansion of the partition function in powers of the coupling constants $`g_{\mathrm{}}`$, $`\mathrm{}=2,\mathrm{},K`$. It represents the formal solution to the problem of counting the (connected and disconnected) Feynman-’t Hooft diagrams of the adjoint fermion one-matrix model. It is a closed formula, in terms of a sum over a large set of integers, for the generating function of ribbon graphs of the fermionic matrix model. A fermionic ribbon graph has much more structure to it than a conventional fat-graph, and so we shall now formulate the rules for generating them from the matrix integral. The fermionic propagator is $`\overline{\psi }_{ij}\psi _{kl}_{\mathrm{Gauss}}=\delta _{il}\delta _{jk}/Ng_1`$ which we represent in the usual way as a double line, along with an orientation defined by an arrow pointing towards the $`\overline{\psi }`$ matrix. Vertices are likewise given an orientation by assigning an outgoing arrow for a $`\overline{\psi }`$ line and an incoming arrow for a $`\psi `$ line. Fat-graphs are now drawn with the rule that only $`\overline{\psi }`$ and $`\psi `$ lines can contract. Each such graph is dual to some discretization of a Riemann surface in the standard way. The surface may be two-coloured by shading the triangles in a polygon if and only if their outer edge is dual to a propagator with an incoming arrow. In this way black edges are always associated with $`\overline{\psi }`$ matrices and white ones with $`\psi `$. Since the fermionic propagator only connects $`\overline{\psi }`$’s with $`\psi `$’s, it follows that every such discretization can be two-coloured. Thus the fat-graphs generated by the fermionic matrix integral can be obtained by drawing all discretizations associated with fermionic $`k`$-point vertices in terms of the corresponding Hermitian $`2k`$-point vertices, and keeping only those graphs which are two-colourable. The even parity of the Hermitian potential $`V(\varphi ^2)`$ is required because only discretizations with even-sided polygons can be two-coloured.
However, the Grassmann nature of the generating matrices yields important changes to the rules for counting such triangulations. For a given vertex in the Wick expansion of the matrix integral, we fix a chosen $`\overline{\psi }`$ line. Lines are joined into propagators with a left-handed orientation. Every time a $`\psi `$ line is joined into a $`\overline{\psi }`$ line (rather than the other way around), we “twist” the ribbon and thereby reverse its orientation. To each such twisted line, we associate a factor of $`1`$. This standard sign factor for fermionic fields yields significant reductions in the overall number of ribbon graphs which are actually generated by the fermionic matrix model, because of the cancellations which occur, for example, between fermionic diagrams that are topologically the same but which are twisted relative to one another. Thus topologically equivalent diagrams do not necessarily add up, but may have the effect of cancelling each other. As a simple example, consider the quadratic potential, $`K=2`$. There are two graphs which contain only a single four-point vertex, but they contribute with equal magnitude and opposite sign to the generating function. This vanishing contribution is the coefficient of the $`g_2/2N(g_1)^2`$ term in the perturbative expansion of the partition function. Symbolically, it is given by the Wick expansion of the Gaussian correlator
(2.14)
Note the twist in the second propagator of the second graph in (2.14), which induces the relative minus sign. Notice also that the usual toroidal four-point graph does not appear in (2.14), because it cannot be two-coloured. This vanishing contribution is also easily found from the analytical formula (2.13). The leading term is 1 and it comes from the configuration whereby all $`k`$’s are 0. The terms of order $`g_2`$ come from the configuration whereby only one $`k`$ is non-vanishing and equal to unity. In the determinant in (2.13), the only $`\mathrm{\Theta }`$’s which are not equal to unity are those which appear in row $`i`$ with $`k^{(i)}=1`$. The determinant therefore vanishes, reproducing the graphical result (2.14). The remaining fat-graph combinatorics are now readily determined from the graphs of the Hermitian one-matrix model with even potential $`V(\varphi ^2)`$ by keeping only those which are two-colourable and incorporating the appropriate twists. The determinant formula (2.13) is readily seen to reproduce the correct combinatorics, with appropriate minus signs coming from the determinant<sup>1</sup><sup>1</sup>1See for related considerations..
These arguments readily generalize to arbitrary polynomial potentials. Generally, the total number of $`2\mathrm{}`$-valent vertices in a given ribbon graph $`\mathrm{\Gamma }`$ as determined by the formula (2.13) is the power of $`g_{\mathrm{}}`$ which is given by
$$n_2\mathrm{}(\mathrm{\Gamma })=\underset{n=1}{\overset{N}{}}k_{\mathrm{}}^{(n)}$$
(2.15)
and, because of Fermi statistics, this number is bounded as
$$0n_2\mathrm{}(\mathrm{\Gamma })\underset{n=1}{\overset{N}{}}[2Nn]_{\mathrm{}}$$
(2.16)
It follows that there are only a finite number of fermionic fat-graphs generated. The twisting of ribbons and also the finiteness of the perturbation series may be attributed to the inclusion of a logarithmic interaction $`2\mathrm{log}\varphi `$ in the corresponding Hermitian one-matrix model with potential $`V(\varphi ^2)`$. This role of the Penner potential will be demonstrated explicitly in the next subsection.
The determinant expansion (2.13) is therefore an analytic expression of the combinatorial formula
$$\frac{Z_N}{Z_N^{\mathrm{Gauss}}}=\underset{\mathrm{\Gamma }}{}\frac{(1)^{t(\mathrm{\Gamma })}(Ng_1)^{\chi (\mathrm{\Gamma })}}{|\mathrm{Aut}(\mathrm{\Gamma })|}\underset{\mathrm{}=2}{\overset{K}{}}\left(\frac{g_{\mathrm{}}}{g_1}\right)^{n_2\mathrm{}(\mathrm{\Gamma })}$$
(2.17)
where the sum runs through all (connected and disconnected) two-colourable fat-graphs $`\mathrm{\Gamma }`$ with $`n_2\mathrm{}(\mathrm{\Gamma })`$ vertices of order $`2\mathrm{}`$ bounded as in (2.16), and $`\chi (\mathrm{\Gamma })`$ is the Euler characteristic of $`\mathrm{\Gamma }`$. The alternating sign factor in (2.17), with $`t(\mathrm{\Gamma })`$ the total number of twisted ribbons of $`\mathrm{\Gamma }`$, is due to Fermi statistics, while $`\mathrm{Aut}(\mathrm{\Gamma })`$ is the automorphism group of the unmarked graph $`\mathrm{\Gamma }`$. The sum over fat-graphs in (2.17) is finite, since the maximum number of vertices that a given diagram can have is
$$v_{\mathrm{max}}=\underset{\mathrm{}=2}{\overset{K}{}}\underset{n=1}{\overset{N}{}}[2Nn]_{\mathrm{}}$$
(2.18)
Each such ribbon graph is dual to a tessellation $`\mathrm{\Gamma }^{}`$ of a Riemann surface of Euler characteristic $`\chi (\mathrm{\Gamma })`$ by $`n_2\mathrm{}(\mathrm{\Gamma })`$ $`2\mathrm{}`$-valent tiles for $`\mathrm{}=2,\mathrm{},K`$. Because of the twist factors, the overall combinatorial numbers generated by the fermionic matrix integral will be drastically reduced. Moreover, depending on the precise details of the potential $`V`$, the perturbative or topological expansions of the matrix model may be alternating series. The above considerations show that there is clearly no simple mapping between the Hermitian and fermionic matrix models. These features will have remarkable implications later on for the topological expansion of the matrix integral (1.1).
### 2.3 Unitary Representation
We will now derive an alternative “dual” representation of the perturbative expansion of the fermionic partition function. For this, we note that the derivatives in (2.9) can be evaluated in terms of a contour integration about the origin of the complex plane as
$$(F(\varphi ),\lambda ^k)=\left(\frac{i}{N}\right)^{N+k+1}(N+k)!\underset{z=0}{}𝑑z\frac{F(z)\mathrm{e}^{NV(z)}}{z^{N+k+1}}$$
(2.19)
The partition function (2.10) may then be written as
$$Z_N=c_NN!\underset{k=1}{\overset{N}{}}\left(\left(\frac{i}{N}\right)^{N+k}(N+k1)!\right)detA$$
(2.20)
where we have defined the $`N\times N`$ matrix
$$A_{ij}=\underset{z=0}{}𝑑z\frac{\mathrm{e}^{NV(z)}}{z^{N+ji+1}}$$
(2.21)
The crucial observation now is that the matrix elements (2.21) depend only on the differences $`ji`$ of row and column positions, i.e. $`A_{ij}=A_{ji}`$. Matrices with such a symmetry are known as Toeplitz matrices and their determinants can be evaluated in terms of averages over the unitary group . The Toeplitz determinant $`det_{i,j}[A_{ji}]`$ appearing in (2.20) is associated with the Laurent series of a function $`𝒜`$ through the definition
$$𝒜(\omega )=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}A_n\omega ^n$$
(2.22)
The matrix elements $`A_n`$ are then interpreted as the Fourier coefficients of $`𝒜`$ on the unit circle $`|\omega |=1`$,
$$A_n=\frac{d\omega }{2\pi i\omega }\omega ^n𝒜(\omega )$$
(2.23)
Using the fact that a determinant is a linear function of each of its rows, it follows that
$$detA=\underset{k=1}{\overset{N}{}}\frac{d\omega _k}{2\pi i\omega _k}𝒜(\omega _k)\omega _k^k\overline{\mathrm{\Delta }(\omega _1,\mathrm{},\omega _N)}$$
(2.24)
Replacing the right-hand side of (2.24) by its average over all permutations of the integration variables $`\omega _k`$ yields
$$detA=\frac{1}{N!}\underset{k=1}{\overset{N}{}}\frac{d\omega _k}{2\pi i\omega _k}𝒜(\omega _k)\left|\mathrm{\Delta }(\omega _1,\mathrm{},\omega _N)\right|^2$$
(2.25)
which shows that the Toeplitz determinant is given by the unitary matrix integral
$$detA=\frac{_{n=1}^{N1}n!}{(2\pi )^{N(N1)/2}}\underset{U(N)}{}[dU]det𝒜(U)$$
(2.26)
Upon substituting (2.21) into (2.22), we can interchange the sum and contour integration provided we take into account that the latter vanishes for $`N+n<0`$. We then have
$$𝒜(\omega )=\underset{z=0}{}𝑑z\frac{\mathrm{e}^{NV(z)}}{\omega ^N(z\omega )}=2\pi i\mathrm{e}^{N\left(V(\omega )\mathrm{log}\omega \right)}$$
(2.27)
From (2.20), (2.26) and (2.27) it follows that the partition function of the adjoint fermion one-matrix model (1.1) is completely equivalent, for any dimension $`N`$, to the $`N\times N`$ unitary matrix model
$$Z_N=k_N\underset{U(N)}{}[dU]\mathrm{e}^{N\mathrm{tr}\left(V(U)\mathrm{log}U\right)}$$
(2.28)
where $`k_N=N^{N(N+1)/2}/(2\pi )^{N(N1)/2}_{n=1}^N(N+n1)!`$. We see that the fermionic one-matrix model for any $`N`$ is completely equivalent to both a Hermitian two-matrix model, and a unitary one-matrix model, both of which involve Penner-type interactions. In particular, the unitary representation clearly demonstrates that the equivalence with the Hermitian Penner one-matrix model is not true at finite $`N`$, but rather only at $`N=\mathrm{}`$. Note that the matrix integral (2.28) is perfectly well-defined at finite $`N`$, but it involves an integrand which is not a real-valued potential.
With the partition function given as an integral over the unitary group, we can clarify the geometric role of the Penner-type potential that characterizes fermionic matrix models, and also give an alternative method for a perturbative expansion in terms of the coupling constants in the potential $`V`$. For this, we integrate over the $`U(1)`$ factor of the unitary group $`U(N)=U(1)\times SU(N)/\text{}_N`$ and leave an $`SU(N)`$ integral,
$$Z_N=k_N\underset{U(N)}{}𝑑\theta [dU_0]\mathrm{e}^{N\left(\mathrm{tr}V(\mathrm{e}^{i\theta }U_0)iN\theta \right)}=\frac{k_N}{N}\underset{0}{\overset{2\pi }{}}𝑑\theta \mathrm{e}^{iN^2\theta }\underset{SU(N)}{}[dU_0]\mathrm{e}^{N\mathrm{tr}V(\mathrm{e}^{i\theta }U_0)}$$
(2.29)
The role of the $`U(1)`$ integration over $`\theta `$ is to restrict the terms in the $`SU(N)`$ integration to those involving only $`N^2`$ factors of $`U_0`$. Equivalently, we may expand the invariant function which is the integrand of (2.29) in $`U(N)`$ characters $`\chi _\stackrel{}{n}`$ as
$$\mathrm{e}^{N\mathrm{tr}V(U)}=\underset{\stackrel{}{n}}{}c_\stackrel{}{n}\chi _\stackrel{}{n}(U),U=\mathrm{e}^{i\theta }U_0$$
(2.30)
where $`c_\stackrel{}{n}\text{}`$ are functions of the coupling constants of the potential $`V`$, and the sum goes over all irreducible polynomial representations of $`U(N)`$ which may be parameterized by their highest weight components $`n_1\mathrm{}n_N0`$ associated with the lengths of the lines in the corresponding Young tableaux. Since
$$\chi _\stackrel{}{n}(U)=\mathrm{e}^{i\theta C_1(\stackrel{}{n})}\chi _\stackrel{}{n}(U_0)$$
(2.31)
where
$$C_1(\stackrel{}{n})=\underset{k=1}{\overset{N}{}}n_k$$
(2.32)
is the linear Casimir of the $`U(N)`$ representation $`\stackrel{}{n}`$ (the total number of boxes in the Young tableau), it follows that the constraint on the $`SU(N)`$ integration in (2.29) is a restriction to terms with linear $`U(N)`$ Casimir $`C_1`$ equal to $`N^2`$. We may therefore write
$$Z_N=\frac{k_N}{N}\left[\underset{SU(N)}{}[dU_0]\mathrm{e}^{N\mathrm{tr}V(U_0)}\right]_{C_1=N^2}$$
(2.33)
where the bracket means to restrict the character expansion of $`\mathrm{e}^{N\mathrm{tr}V(U_0)}`$ to those $`\stackrel{}{n}`$ with $`C_1(\stackrel{}{n})=N^2`$. The role of the logarithmic interaction in (2.28) may therefore be characterized as restricting the matrix integral to completely filled Young diagrams.
The $`SU(N)`$ integral in (2.33) can be evaluated using standard techniques from lattice gauge theory . For this, we introduce the generating function
$$W(J)=\underset{SU(N)}{}[dU_0]\mathrm{e}^{\mathrm{tr}JU_0}$$
(2.34)
with the property
$$\underset{SU(N)}{}[dU_0]f(U_0)=f\left(\frac{}{J}\right)W(J)|_{J=0}$$
(2.35)
where $`J`$ is an arbitrary $`N\times N`$ matrix and $`f`$ an arbitrary function on $`SU(N)`$. Using the invariance of the Haar measure $`[dU_0]`$ under $`SU(N)`$ rotations, we have $`W(J)=W(UJV)`$, $`U,VSU(N)`$, from which it can be shown that the generating function may be expanded in powers of determinants of $`J`$ as
$$W(J)=\frac{(2\pi )^{\frac{N(N+1)}{2}1}}{N}\underset{k=0}{\overset{\mathrm{}}{}}\frac{(detJ)^k}{_{n=1}^N(n+k1)!}$$
(2.36)
Applying this result to the partition function (2.33), the restriction $`C_1(\stackrel{}{n})=N^2`$ implies that only the $`k=N`$ term of the expansion (2.36) contributes, giving
$$Z_N=\mathrm{e}^{N\mathrm{tr}V\left(\frac{}{J}\right)}\stackrel{N}{det}J|_{J=0}$$
(2.37)
Comparing (2.3) and (2.37), we see that the finite $`N`$ expressions for the fermionic partition function in the Hermitian two-matrix and unitary one-matrix representations are dual to each other. We can make this duality precise by writing (2.37) as
$`Z_N`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^{N^2}}}{\displaystyle \underset{gl(N)}{}𝑑J𝑑\varphi \mathrm{e}^{i\mathrm{tr}\varphi J}\mathrm{e}^{N\mathrm{tr}V\left(\frac{}{J}\right)}\stackrel{N}{det}J}`$ (2.38)
$`=`$ $`{\displaystyle \frac{1}{(2\pi )^{N^2}}}{\displaystyle \underset{gl(N)}{}𝑑J𝑑\varphi \mathrm{e}^{N\mathrm{tr}V(i\varphi )}\stackrel{N}{det}\left(i\frac{}{\varphi }\right)\mathrm{e}^{i\mathrm{tr}\varphi J}}`$
$`=`$ $`{\displaystyle \underset{gl(N)}{}}d\varphi \delta (\varphi )\stackrel{N}{det}\left(i\frac{}{\varphi }\right)\mathrm{e}^{N\mathrm{tr}V(i\varphi )}=\stackrel{N}{det}\left(\frac{}{\varphi }\right)\mathrm{e}^{N\mathrm{tr}V(\varphi )}|_{\varphi =0}`$
Some further aspects of the unitary representation of the fermionic matrix model are discussed in appendix A. In section 5 we will show explicitly that the equations of motion of the unitary and fermionic one-matrix models are identical.
## 3 Orthogonal Polynomial Solution
In this section we shall demonstrate how to generalize the orthogonal polynomial technique to the adjoint fermion one-matrix model (1.1), with the goal of solving the matrix model in the large $`N`$ limit. Although the fermionic matrix model is not naturally an eigenvalue model, since it is not possible to diagonalize a Grassmann valued matrix, we can exploit its equivalence with the unitary one-matrix model (2.28). We will thereby define the orthogonal polynomials of the fermionic matrix model (1.1) to be the orthogonal polynomials associated with the corresponding unitary eigenvalue model. This definition will lead, as we shall see, to a well-defined real-valued solution for the fermionic matrix model.
### 3.1 General Properties
Let us consider the fermionic partition function in the representation (2.20,2.21), which involves the measure
$$d\mu (z)=\frac{\mathrm{e}^{NV(z)}}{z^{N+1}}dz$$
(3.1)
on the unit circle $`z\overline{z}=1`$. The partition function (2.20) may then be written as the determinant of the moment matrix of the measure (3.1),
$$Z_N=d_N\underset{i,j}{det}\left[z^{i1}|\overline{z}^{j1}\right]$$
(3.2)
where $`d_N`$ is an irrelevant numerical constant and we have defined the inner product
$$F(z)|G(\overline{z})=\underset{z=0}{}𝑑\mu (z)F(z)G(\overline{z})$$
(3.3)
on the space $`\text{}[z]\text{}[\overline{z}]`$ of Laurent polynomials on the unit circle. The construction of orthogonal polynomials to evaluate inner products on the circle such as (3.3) with real, positive definite measure $`d\mu `$ was discussed long ago in the mathematics literature and subsequently in the physics literature . Some aspects of generic unitary matrix models are also dealt with in . In the present case, the measure (3.1) is complex-valued, and we must deal accordingly with defining the appropriate system of orthogonal polynomials.
These are the monic polynomials $`(\mathrm{\Phi }_n(z),\mathrm{\Lambda }_m(z))`$ of order $`(n,m)`$ which form a complete set in the space $`\text{}[z]\text{}[\overline{z}]`$ and which are bi-orthogonal in the measure (3.1),
$$\mathrm{\Phi }_n(z)|\mathrm{\Lambda }_m(z)=h_n\delta _{nm}$$
(3.4)
where $`h_n`$ are some constants. They are normalized as
$$\mathrm{\Phi }_n(z)=z^n\underset{k=0}{\overset{n1}{}}p_{n,k}z^k,\mathrm{\Lambda }_m(z)=z^m\underset{k=0}{\overset{m1}{}}l_{m,k}z^k$$
(3.5)
where the coefficients $`p_{n,k}`$ and $`l_{m,k}`$ can be formally obtained from the usual Gram-Schmidt orthogonalization procedure applied to the monomials $`(z^n,z^m)`$ and the inner product (3.3). In fact, by iterating (3.5) we find that the monomials $`(z^n,z^m)`$ can be expanded as
$`z^n`$ $`=`$ $`\mathrm{\Phi }_n(z)+{\displaystyle \underset{j=0}{\overset{n}{}}}\left({\displaystyle \underset{k_1=0}{\overset{n1}{}}}{\displaystyle \underset{k_2=0}{\overset{k_11}{}}}\mathrm{}{\displaystyle \underset{k_j=0}{\overset{k_{j1}1}{}}}p_{n,k_1}p_{k_1,k_2}\mathrm{}p_{k_{j1},k_j}\mathrm{\Phi }_{k_j}(z)\right)`$
$`z^m`$ $`=`$ $`\mathrm{\Lambda }_m(z)+{\displaystyle \underset{j=0}{\overset{m}{}}}\left({\displaystyle \underset{k_1=0}{\overset{m1}{}}}{\displaystyle \underset{k_2=0}{\overset{k_11}{}}}\mathrm{}{\displaystyle \underset{k_j=0}{\overset{k_{j1}1}{}}}l_{m,k_1}l_{k_1,k_2}\mathrm{}l_{k_{j1},k_j}\mathrm{\Lambda }_{k_j}(z)\right)`$ (3.6)
It follows that any polynomial in the space $`\text{}[z]\text{}[\overline{z}]`$ of degree $`(n,m)`$ can be expressed as a linear combination of $`(\mathrm{\Phi }_k(z),\mathrm{\Lambda }_{\mathrm{}}(z))`$ with $`kn`$ and $`\mathrm{}m`$. From (3.5) it follows they define a change of basis $`(z^{k1},z^{\mathrm{}+1})(\mathrm{\Phi }_{k1}(z),\mathrm{\Lambda }_\mathrm{}1(z))`$ of the vector space $`\text{}[z]\text{}[\overline{z}]`$ that leaves the partition function (3.2) invariant and at the same time diagonalizes the inner product (3.3). In particular, from (3.4) we have
$$Z_N=d_N\underset{i,j}{det}\left[\mathrm{\Phi }_{i1}(z)|\mathrm{\Lambda }_{j1}(z)\right]=d_Nh_0^N\underset{n=1}{\overset{N1}{}}R_n^{Nn}$$
(3.7)
where we have introduced the recursion coefficients
$$R_n=\frac{h_n}{h_{n1}}$$
(3.8)
and $`h_0=_{z=0}𝑑\mu (z)`$ is the normalization of the measure (3.1).
The reason why two independent sets of functions corresponding to clockwise and anti-clockwise rotations on the circle are required here is the complexity of the measure (3.1). In the classical case whereby the measure $`d\mu `$ is real and positive , the polynomials $`\mathrm{\Phi }_n(z)`$ and $`\mathrm{\Lambda }_n(z)`$ are related to each other by complex conjugation. The present system of polynomials may be thought of as a deformation of those for positive definite measures on the unit circle. This deformation is essentially encoded in the coefficients $`p_{n+1,0}`$ and $`l_{n+1,0}`$. In order for the Gram-Schmidt process to work for the complex measure (3.1), one needs polynomials in both $`z`$ and $`1/z`$ and a priori these polynomials are unrelated. Alternatively, we can think of this doubling as the usual doubling of degrees of freedom due to the fermionic nature of the original matrix model. As the orthogonal polynomials define a complete bi-orthogonal system of functions, they satisfy the completeness relation
$$\underset{n=0}{\overset{\mathrm{}}{}}\frac{\mathrm{\Phi }_n(z)\mathrm{\Lambda }_n(z)}{h_n}=\delta ^{(\mu )}(z)=z^{N+1}\mathrm{e}^{NV(z)}\delta (z)$$
(3.9)
on the vector space $`\text{}[z]\text{}[\overline{z}]`$, where the Dirac delta-function corresponding to the measure $`d\mu `$ is defined so that
$$\underset{z=0}{}𝑑\mu (z)F(z)\delta ^{(\mu )}(zz^{})=F(z^{})$$
(3.10)
for any function $`F(z)`$ on the unit circle. The relation (3.9) is understood as an analytical continuation (to the distribution space that is the completion of $`\text{}[z]\text{}[\overline{z}]`$) which can be thought of as a representation of the potential $`V(\overline{\psi }\psi )`$ in terms of the orthogonal polynomials.
The problem of evaluating the partition function of the fermionic matrix model therefore boils down to evaluating the coefficients $`R_n`$ appearing in (3.7). For this, we will first derive a few properties of the system of orthogonal polynomials introduced above. Completeness implies the relations
$$z\mathrm{\Phi }_n(z)=\mathrm{\Phi }_{n+1}(z)+\underset{k=0}{\overset{n}{}}𝒫_k^{(n)}\mathrm{\Phi }_k(z),\frac{1}{z}\mathrm{\Lambda }_m(z)=\mathrm{\Lambda }_{m+1}(z)+\underset{k=0}{\overset{m}{}}_k^{(m)}\mathrm{\Lambda }_k(z)$$
(3.11)
The coefficients $`𝒫_k^{(n)}`$ and $`_k^{(m)}`$ completely determine the form of the orthogonal polynomials and hence the partition function. To see this, we use (3.5) and (3.11) to write
$`\mathrm{\Phi }_n(z)`$ $`=`$ $`z\mathrm{\Phi }_{n1}(z){\displaystyle \underset{k=0}{\overset{n1}{}}}𝒫_k^{(n1)}\mathrm{\Phi }_k(z)`$ (3.12)
$`=`$ $`z^n𝒫_0^{(n)}{\displaystyle \underset{k=1}{\overset{n1}{}}}\left(p_{n1,k1}+𝒫_k^{(n1)}\right)z^k{\displaystyle \underset{k=0}{\overset{n2}{}}}𝒫_k^{(n1)}{\displaystyle \underset{j=0}{\overset{k1}{}}}p_{k,j}z^j`$
Comparing the various polynomial coefficients in (3.12) with those of the definition (3.5), we arrive at the iterative relations
$`p_{n,n1}`$ $`=`$ $`p_{n1,n2}+𝒫_{n1}^{(n1)}`$ (3.13)
$`p_{n,k}`$ $`=`$ $`𝒫_k^{(n1)}+p_{n1,k1}+{\displaystyle \underset{j=k+1}{\overset{n1}{}}}p_{j,k}𝒫_j^{(n1)},1kn2`$ (3.14)
$`p_{n,0}`$ $`=`$ $`𝒫_0^{(n)}+{\displaystyle \underset{j=1}{\overset{n1}{}}}p_{j,0}𝒫_j^{(n1)}`$ (3.15)
The relations (3.13) and (3.15) can be iterated straightforwardly and by induction we have
$`p_{n,n1}`$ $`=`$ $`{\displaystyle \underset{j=0}{\overset{n1}{}}}𝒫_j^{(j)}`$ (3.16)
$`p_{n,0}`$ $`=`$ $`𝒫_0^{(n)}+{\displaystyle \underset{k=1}{\overset{n3}{}}}\left({\displaystyle \underset{j_1=1}{\overset{n1}{}}}{\displaystyle \underset{j_2=1}{\overset{j_11}{}}}\mathrm{}{\displaystyle \underset{j_k=1}{\overset{j_{k1}1}{}}}𝒫_{j_1}^{(n1)}𝒫_{j_2}^{(j_11)}\mathrm{}𝒫_{j_k}^{(j_{k1}1)}𝒫_0^{(j_k)}\right)`$ (3.17)
$`+𝒫_0^{(1)}𝒫_1^{(1)}{\displaystyle \underset{j_1=1}{\overset{n1}{}}}{\displaystyle \underset{j_2=1}{\overset{j_11}{}}}\mathrm{}{\displaystyle \underset{j_{n2}=1}{\overset{j_{n3}1}{}}}𝒫_{j_1}^{(n1)}𝒫_{j_2}^{(j_11)}\mathrm{}𝒫_{j_{n2}}^{(j_{n3}1)}`$
which can be substituted into (3.14) to iteratively determine the remaining $`p_{n,k}`$. Analogous relations exist between the coefficients $`l_{m,k}`$ and $`_k^{(m)}`$.
We will also need higher degree versions of the completeness relations (3.11). For this, we define constants $`P_{n,k}^{[\mathrm{}]}`$ by
$$z^{\mathrm{}}\mathrm{\Phi }_n(z)=\underset{k=0}{\overset{n+\mathrm{}}{}}P_{n,k}^{[\mathrm{}]}\mathrm{\Phi }_k(z)\mathrm{with}P_{n,n+\mathrm{}}^{[\mathrm{}]}=1$$
(3.18)
and use (3.11) to write
$$z^{\mathrm{}+1}\mathrm{\Phi }_n(z)=\underset{k=0}{\overset{n+\mathrm{}+1}{}}P_{n,k}^{[\mathrm{}+1]}\mathrm{\Phi }_k(z)=\underset{k=0}{\overset{n+\mathrm{}}{}}P_{n,k}^{[\mathrm{}]}\left(\mathrm{\Phi }_{k+1}(z)+\underset{j=0}{\overset{k}{}}𝒫_j^{(k)}\mathrm{\Phi }_j(z)\right)$$
(3.19)
Using completeness of the orthogonal polynomials we arrive at a recursive relation for the coefficients $`P_{n,k}^{[\mathrm{}]}`$,
$$P_{n,k}^{[\mathrm{}+1]}=P_{n,k1}^{[\mathrm{}]}+\underset{j=k}{\overset{n+\mathrm{}}{}}P_{n,j}^{[\mathrm{}]}𝒫_k^{(j)},0kn+\mathrm{}+1$$
(3.20)
with $`P_{n,1}^{[\mathrm{}]}0`$, $`𝒫_k^{(n)}0`$ for $`k>n`$, and the initial condition $`P_{n,k}^{[0]}=\delta _{nk}`$.
### 3.2 Recursion Relations
We will now determine the coefficients appearing in the completeness relations using the above properties of the orthogonal polynomials. This will lead to a set of recursion equations for the coefficients $`R_n`$ which thereby completely determines the solution of adjoint fermion one-matrix model. Consider the bi-orthogonality relation
$$\mathrm{\Phi }_{n+1}(z)|\mathrm{\Lambda }_m(z)=0,0mn$$
(3.21)
Using the completeness relations (3.11) we can write (3.21) as
$$z\mathrm{\Phi }_n(z)|\mathrm{\Lambda }_m(z)=𝒫_m^{(n)}h_m$$
(3.22)
The inner product in (3.22) can be evaluated by grouping the factor of $`z`$ with the polynomial $`\mathrm{\Lambda }_m(z)`$ with the result
$$z\mathrm{\Phi }_n(z)|\mathrm{\Lambda }_m(z)=\mathrm{\Phi }_n(z)|\mathrm{\Lambda }_{m1}(z)+O\left(z^{m+2}\right)l_{m,0}z=l_{m,0}z\mathrm{\Phi }_n(z)|1$$
(3.23)
where we have used the fact that, aside from the term $`l_{m,0}z`$ which comes from the constant part of the $`\mathrm{\Lambda }`$ polynomials in (3.5), the function $`z\mathrm{\Lambda }_m(z)`$ is a sum of polynomials of degree $`m1`$. The final inner product in (3.23) can be evaluated in a straightforward fashion using the relation
$$z\mathrm{\Phi }_n(z)|\mathrm{\Lambda }_{n+1}(z)=h_{n+1}$$
(3.24)
which follows from (3.11). Again, by grouping the factor of $`z`$ with $`\mathrm{\Lambda }_{n+1}(z)`$ one finds
$$z\mathrm{\Phi }_n(z)|1=\frac{h_nh_{n+1}}{l_{n+1,0}}$$
(3.25)
Substituting (3.25) and (3.23) into (3.22), we find a concise relation between the constant terms of the $`\mathrm{\Phi }`$ polynomials and the completeness coefficients $`𝒫_m^{(n)}`$,
$`𝒫_m^{(n)}`$ $`=`$ $`{\displaystyle \frac{l_{m,0}}{l_{n+1,0}}}{\displaystyle \frac{h_{n+1}h_n}{h_m}}`$ (3.28)
$`=`$ $`\{\begin{array}{c}Q_m\left(R_{n+1}1\right)_{k=m+1}^nQ_kR_k\mathrm{for}0m<n\hfill \\ Q_n(R_{n+1}1)\mathrm{for}m=n\hfill \end{array}`$
where we have defined
$$Q_n=\frac{l_{n,0}}{l_{n+1,0}}$$
(3.29)
The completeness relation (3.11) for the $`\mathrm{\Phi }`$ polynomials can therefore be written as
$$z\mathrm{\Phi }_n(z)=\mathrm{\Phi }_{n+1}(z)+\frac{h_{n+1}h_n}{l_{n+1,0}}\underset{k=0}{\overset{n}{}}\frac{l_{k,0}}{h_k}\mathrm{\Phi }_k(z)$$
(3.30)
Subtracting this relation from itself under the shift of index $`nn1`$ leads to the three-term recursion relation for the $`\mathrm{\Phi }`$ polynomials
$$\mathrm{\Phi }_{n+1}(z)=z\mathrm{\Phi }_n(z)Q_n\frac{R_{n+1}1}{R_n1}\left(\mathrm{\Phi }_n(z)zR_n\mathrm{\Phi }_{n1}(z)\right)$$
(3.31)
Analogous relations for the $`\mathrm{\Lambda }`$ polynomials may also be derived. Equating the constant terms on both sides of (3.31) yields
$$\frac{p_{n+1,0}}{p_{n,0}}=Q_n\frac{R_{n+1}1}{R_n1}$$
(3.32)
so that
$`_m^{(n)}`$ $`=`$ $`{\displaystyle \frac{p_{m,0}}{p_{n+1,0}}}{\displaystyle \frac{h_{n+1}h_n}{h_m}}`$ (3.35)
$`=`$ $`\{\begin{array}{c}(1)^{nm}\frac{1R_m}{Q_m}\frac{R_{n+1}1}{R_{m+1}1}_{k=m+1}^n\frac{R_k}{Q_k}\frac{R_k1}{R_{k+1}1}\mathrm{for}0m<n\hfill \\ \frac{1R_n}{Q_n}\mathrm{for}m=n\hfill \end{array}`$
Substituting (3.35) into the completeness relation (3.11) for the $`\mathrm{\Lambda }`$ polynomials and subtracting the resulting expression from itself under the index shift $`nn1`$, we arrive at the three-term recursion relation
$$\mathrm{\Lambda }_n(z)=\frac{1}{z}R_n\mathrm{\Lambda }_{n1}(z)Q_n\left(\mathrm{\Lambda }_{n+1}(z)\frac{1}{z}\mathrm{\Lambda }_n(z)\right)$$
(3.36)
The recursion relations (3.31) and (3.36) have been derived without any reference to the particular details of the model, i.e. they hold independently of the precise form of the potential $`V`$ of the fermionic matrix model. Consequently, these relations are merely ‘kinematical’ in origin. In order to solve for the dynamics of the matrix model we need more information. From a more pragmatic point of view, in order to solve for the partition function we need to find a solution for the coefficients $`R_n`$, but the recursion relations also involve the quantities $`Q_n`$. More information is needed to link these two objects. As we will see, one only needs two ‘dynamical’ relations to obtain a closed set of recursion relations for any polynomial potential. The first one is given from (3.4) and (3.5) as
$$\mathrm{\Phi }_n(z)|\mathrm{\Lambda }_{n1}^{}(z)=(1n)h_n$$
(3.37)
Integrating by parts on the left-hand side of (3.37) gives
$$(N+1)\mathrm{\Phi }_n(z)|\frac{1}{z}\mathrm{\Lambda }_{n1}(z)NV^{}(z)\mathrm{\Phi }_n(z)|\mathrm{\Lambda }_{n1}(z)\mathrm{\Phi }_n^{}(z)|\mathrm{\Lambda }_{n1}(z)=(1n)h_n$$
(3.38)
Using the completeness relation (3.11) for the $`\mathrm{\Lambda }`$ polynomials and the fact that $`\mathrm{\Phi }_n^{}(z)=n\mathrm{\Phi }_{n1}(z)+O(z^{n2})`$, we arrive at
$$NV^{}(z)\mathrm{\Phi }_n(z)|\mathrm{\Lambda }_{n1}(z)=(N+n)h_nnh_{n1}$$
(3.39)
For a polynomial potential (1.2) of degree $`K`$, we may use (3.18) to write (3.39) as an equation for the recursion coefficients $`R_n`$,
$$R_n=\frac{n}{n+N}+\frac{N}{n+N}\underset{k=1}{\overset{K1}{}}g_{k+1}P_{n,n1}^{[k]}$$
(3.40)
Closely related to the first one, the second ‘dynamical’ relation is given by the bi-orthogonality relation
$$\mathrm{\Phi }_n(z)|\mathrm{\Lambda }_{n2}^{}(z)=0$$
(3.41)
Integrating this by parts gives
$$NV^{}(z)\mathrm{\Phi }_n(z)|\mathrm{\Lambda }_{n2}(z)=\mathrm{\Phi }_n^{}(z)|\mathrm{\Lambda }_{n2}(z)$$
(3.42)
The right-hand side of (3.42) can be evaluated using (3.5), (3.6) and orthogonality, and keeping only those terms in the $`\mathrm{\Phi }`$ polynomials which have a non-vanishing overlap with $`\mathrm{\Lambda }_{n2}(z)`$ in (3.42). This gives
$`\mathrm{\Phi }_n^{}(z)|\mathrm{\Lambda }_{n2}(z)`$ $`=`$ $`nz^{n1}(n1)p_{n,n1}z^{n2}+O\left(z^{n3}\right)|\mathrm{\Lambda }_{n2}(z)`$ (3.43)
$`=`$ $`\left[(n1)p_{n,n1}np_{n1,n2}\right]h_{n2}`$
By defining
$$P_n=\frac{p_{n,n1}}{N}$$
(3.44)
we thus find that the second ‘dynamical’ relation reads
$$nP_{n1}(n1)P_n=\underset{k=1}{\overset{K1}{}}g_{k+1}P_{n,n2}^{[k]}$$
(3.45)
The coefficients $`P_n`$ can be related to the coefficients $`Q_n`$ and $`R_n`$ by using (3.16) and (3.28) to get
$$N(P_{n+1}P_n)=Q_n(R_{n+1}1)$$
(3.46)
To summarize the results of this subsection, we note that the right-hand sides of (3.40) and (3.45) can be evaluated, using the recursion formula (3.20), in terms of the completeness coefficients $`𝒫_m^{(n)}`$. Using (3.28), it follows that, for any polynomial potential, (3.40), (3.45) and (3.46) form a closed set of recursion relations involving only the coefficients $`R_n`$, $`Q_n`$ and $`P_n`$. An explicit system of fermionic orthogonal polynomials is constructed in appendix B.
### 3.3 Free Energy
The final quantity we need for the solution of the fermionic matrix model is the normalization $`h_0`$ of the measure (3.1). This integral can be evaluated explicitly to give
$$h_0=\underset{z=0}{}dz\frac{\mathrm{e}^{NV(z)}}{z^{N+1}}=\frac{2\pi i}{N!}\frac{^N}{z^N}\mathrm{e}^{NV(z)}|_{z=0}$$
(3.47)
We substitute these quantities into (3.7) and normalize by the Gaussian partition function $`Z_N^{\mathrm{Gauss}}=(1)^{[N]_2}(Ng_1)^{N^2}`$ for which $`V(z)=g_1z`$, $`R_n=n/(n+N)`$ and $`h_0=2\pi i(Ng_1)^N/N!`$. For a polynomial potential, we may evaluate (3.47) explicitly using (2.11), and using the multinomial theorem we thus find that the free energy $`F_N=\frac{1}{N^2}\mathrm{log}\frac{Z_N}{Z_N^{\mathrm{Gauss}}}`$ is given by
$`F_N`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(1)^n(n1)!{\displaystyle \underset{0k_0^{(2)},\mathrm{},k_{[N]_2}^{(2)}n}{}}\mathrm{}{\displaystyle \underset{0k_0^{(K)},\mathrm{},k_{[N]_K}^{(K)}n}{}}{\displaystyle \frac{\delta \left(_{\mathrm{},m}k_{\mathrm{}}^{(m)}n\right)}{_{\mathrm{},m}k_{\mathrm{}}^{(m)}!}}`$ (3.48)
$`\times {\displaystyle \underset{j=2}{\overset{K}{}}}{\displaystyle \underset{\mathrm{}_j=1}{\overset{[N]_j}{}}}\left[{\displaystyle \frac{N!}{(N_ii\mathrm{}_i)!\mathrm{}_2!\mathrm{}\mathrm{}_K!}}{\displaystyle \underset{m=2}{\overset{K}{}}}\left({\displaystyle \frac{g_m}{mN^{m1}(g_1)^m}}\right)^\mathrm{}_m\right]^{k_\mathrm{}_j^{(j)}}`$
$`+{\displaystyle \frac{1}{N}}{\displaystyle \underset{n=1}{\overset{N1}{}}}\left(1{\displaystyle \frac{n}{N}}\right)\mathrm{log}{\displaystyle \frac{(n+N)R_n}{n}}`$
where the coefficients $`R_n`$ are determined from (3.40).
Generally, in the large-$`N`$ limit, the variable $`xn/N`$ becomes a continuous parameter in the interval $`[0,1]`$. We assume that in this limit the recursion coefficients $`_n`$ (any of the quantities $`R_n`$, $`Q_n`$ or $`P_n`$), determined as the solutions of the recursion equations derived in the previous subsection in the large-$`N`$ limit, become continuous functions $`(x)`$ of $`x[0,1]`$. This is justified by the form of the recursion relations which suggest the replacement, $`_n(\frac{n}{N})=(x)`$. It follows that shifts in the index $`n`$ are related to sub-leading contributions in the large $`N`$ limit,
$$_{n+a}(x)+\frac{a}{N}\underset{k=1}{\overset{\mathrm{}}{}}\frac{1}{k!}\left(\frac{a}{N}\right)^{k1}\frac{^k(x)}{x^k}$$
(3.49)
and the second, finite sum in (3.48) can be approximated by a simple one-dimensional integral over $`x[0,1]`$ at $`N=\mathrm{}`$. To deal with the first term in (3.48), we need to determine the large $`N`$ limit of the normalization constant $`h_0`$. For this, we note that there is nothing particularly special about the choice of integration contour in (3.47), as it was merely introduced in (2.19) as a means of evaluating the derivatives in (2.9). In particular, since its integrand is an analytic function, we can deform the contour arbitrarily in the complex plane, and in particular to one out at infinity. In the large $`N`$ limit, the one-dimensional integral $`h_0`$ may then be evaluated using the saddle-point approximation. In this way we find that the genus 0 free energy is given by
$$F_0\underset{N\mathrm{}}{lim}F_N=f_0+\underset{0}{\overset{1}{}}𝑑x(1x)\mathrm{log}\frac{(1+x)R(x)}{x}$$
(3.50)
where
$$f_0=V(\zeta )\mathrm{log}g_1\zeta $$
(3.51)
with $`\zeta `$ the solution of the order $`K`$ algebraic saddle-point equation
$$\zeta V^{}(\zeta )=1$$
(3.52)
for the branch which has the perturbative Gaussian limit $`\zeta |_{g_2=\mathrm{}=g_K=0}=1/g_1`$ and for which the function (3.51) is real-valued.
## 4 Topological Expansion
In this section we will solve for the continuum limit of the discretized random surface theory represented by the adjoint fermion one-matrix model. We will consider two limits of the matrix model. The first one is the naive large-$`N`$ limit which captures the leading critical behaviour and gives a non-perturbative solution to the problem of counting fermionic ribbon graphs of spherical topology. The second one is the double-scaling limit in which an appropriate coupling constant is tuned in the limit $`N\mathrm{}`$ in order to obtain contributions from all orders of the $`\frac{1}{N}`$ expansion of the free energy and which yields the complete topological expansion of the fermionic matrix model.
### 4.1 Critical Behaviour
We will start by making some general remarks concerning the large $`N`$ limit of the fermionic matrix model. From (3.50) we see that there are two contributions to the $`N=\mathrm{}`$ free energy. These two quantities represent very different phases of the system. The first one is given by (3.51) which is determined by the saddle-point equation (3.52). This equation is immediately recognized as the critical equation describing the continuum limit of an abstract branched polymer . In fact, the explicit derivative expression for the measure normalization (3.47) coincides exactly with the partition function of the $`N`$ dimensional fermionic vector model with the same polynomial potential $`V`$ . The latter quantity is related to the generating function for branched polymer networks by a simple mapping of its Feynman diagrams and an analytical continuation $`N\frac{N}{2}`$ of the fermionic vector dimension. The analytic reason for the appearance of branched polymer behaviour in the fermionic matrix model is that the normalization of its measure is not of sub-leading order in $`N`$. From a graphical point of view, its appearance is clear from the analysis of section 2.2. The entropy factors associated with the number of Feynman graphs is drastically reduced in the fermionic case because of the cancellations which occur between twisted diagrams. There is a class of diagrams with no twists in the fermionic matrix model, and these have a tree-like growth in the large $`N`$ limit. The reduced entropy of the remaining twisted graphs is then comparable to that of the graphs which produce a polymer-like behaviour. This feature is unique to fermionic matrices, and it is the property which enables the construction of a model of branched polymers using supersymmetric matrix models via the coupling of the fermionic matrix model to an ordinary, complex matrix model which has the effect of cancelling the set of twisted diagrams. In fact, it is precisely the twisting mechanism described in section 2.2 that enables one to isolate graphs with tree-like growth in the fermionic case. The untwisted diagrams may then be mapped onto branched polymers similarly to the case of the cactus diagrams which appear in supersymmetric matrix models .
The twisted graphs are associated with the second contribution in (3.50) and, as we will demonstrate, they lead to the usual surface effects in the theory. Thus the orthogonal polynomial formalism that we have developed naturally splits the generating function into a piece corresponding to the tree-like graphs and a piece corresponding to the twisted diagrams. The continuum limit of the latter ensemble of graphs describes a surface theory and describes an effect which is of order $`N^2`$, while the former ensemble which has a polymer growth produces an effect of order $`N`$. One typical feature of the polymer generating function is that it becomes non-analytic in the large $`N`$ limit at the Gaussian point $`g_k=0`$, $`k2`$. However, this effect is of sub-leading order in $`N`$, so that a large $`N`$ analysis with arbitrary coupling constants will miss the branched polymer phase transition. There is a “barrier” at the Gaussian point which separates the surface-like continuum limit from the branched polymer behaviour. From a geometrical point of view, the fermionic matrix model probes an intermediate lattice phase where the number of Feynman diagrams contributing to a given process is greater than that expected of a vector model but less than that of a matrix model. Thus the matrix integral generates a random surface theory which contains an “internal” branched polymer phase.
This decrease in the number of graphs as compared to the usual matrix models will be responsible for the Borel summability of the all genus expansion of the free energy of the fermionic matrix model that we will prove in this section. As discussed above, the essential critical behaviour of the fermionic matrix model is encoded in the large-$`N`$ limit of the recursion coefficients $`R_n`$. If we introduce the quantities $`\tau _n`$ defined by
$$h_n=\frac{\tau _{n+1}}{\tau _n}$$
(4.1)
with $`\tau _0=d_N`$ as in (3.2), then the recursion coefficients are given by $`R_n=\tau _{n+1}\tau _{n1}/\tau _n^2`$ and the free energy by $`F_N=_N`$, where $`_n=\mathrm{log}(\tau _n/\tau _n^{\mathrm{Gauss}})`$. Replacing sequences by functions of $`x[0,1]`$ as prescribed by (3.49), we then have to leading order in $`N`$,
$$\mathrm{log}\frac{(1+x)R(x)}{x}=\frac{^2(x)}{x^2},F_0=(1)$$
(4.2)
The relations (4.2) show that the function $`R(x)`$ is related to the specific heat $`u_0`$ of the matrix model, and also that any singularity in the free energy will occur at the boundary $`x=1`$ of the domain of $`R(x)`$. These facts will enable us to construct the topological expansion of the fermionic matrix model from the recursion coefficients $`R_n`$ of the orthogonal polynomials. We remark that it is possible to demonstrate from the constrained sum over $`SU(N)`$ group characters of section 2.3 that the partition function is real and has a definite value for any $`N`$. Combining this with our knowledge from large $`N`$ loop equation calculations, wherein the system goes to a definite limit, one can argue that there is a definitive $`\frac{1}{N}`$ expansion of the free energy. The phase transition at large $`N`$ can thereby be understood as a sort of percolation transition in the Young tableaux, whereby the system evolves into a phase in which the number of filled $`SU(N)`$ Young tableaux grows asymptotically.
### 4.2 Main Recursion Relation
We will now begin quantifying the above discussion. For definiteness, in the remainder of this section we will deal primarily with the simplest non-Gaussian model which is given by the quadratic potential
$$V(z)=z+\frac{g}{2}z^2$$
(4.3)
although, as we will discuss, the results we obtain are universal. The potential (4.3) is difficult to treat within the loop equation approach because the endpoints of the support of the corresponding spectral density are located asymmetrically in the complex plane . Indeed, one power of the technique that we develop is that it is insensitive to the parity symmetry (or lack thereof) of the potential. This will account for the universality properties of the solution that we shall find. In this subsection we will derive the main recursion relation corresponding to the potential (4.3) that will be used to construct the topological expansion.
To solve the recursion relations (3.40), (3.45) and (3.46) corresponding to (4.3), we first use (3.20) and (3.28) to determine the coefficients
$`P_{n,n1}^{[1]}`$ $`=`$ $`𝒫_{n1}^{(n)}=Q_{n1}Q_nR_n(R_{n+1}1)`$
$`P_{n,n2}^{[1]}`$ $`=`$ $`𝒫_{n2}^{(n)}=Q_{n2}Q_{n1}Q_nR_{n1}R_n(R_{n+1}1)`$ (4.4)
This leads to the closed set of recursion relations
$`(N+n)R_nn`$ $`=`$ $`gNQ_{n1}Q_nR_n(R_{n+1}1)`$ (4.5)
$`nP_{n1}(n1)P_n`$ $`=`$ $`gQ_{n2}Q_{n1}Q_nR_{n1}R_n(R_{n+1}1)`$ (4.6)
$`N(P_{n+1}P_n)`$ $`=`$ $`Q_n(R_{n+1}1)`$ (4.7)
From the relations (4.6) and (4.7) one can solve for the coefficients $`P_n`$ as
$$NP_n=gNQ_{n2}Q_{n1}Q_nR_{n1}R_n(R_{n+1}1)+nQ_{n1}(R_n1)$$
(4.8)
Substituting (4.8) into (4.7) and using (4.5) to reduce the degree of the $`Q_n`$’s in the resulting equation yields
$`0`$ $`=`$ $`Q_{n1}R_n\left[(N+n+1)R_{n+1}n1\right]Q_{n2}R_{n1}\left[(N+n)R_nn\right]`$ (4.9)
$`+nQ_n(R_{n+1}1)nQ_{n1}(R_n1)`$
We now use (4.5) to define the quantity
$$\mathrm{\Omega }_n\{R\}Q_nQ_{n1}=\frac{1}{g}\frac{(N+n)R_nn}{NR_n(R_{n+1}1)}$$
(4.10)
and multiply the relation (4.9) through by $`Q_n`$. Using (4.6) along with (4.8), we can then solve for $`Q_n`$ in terms of the $`R_n`$’s alone as
$$Q_n\{R\}^2=\mathrm{\Omega }_n\{R\}^2\frac{n(R_n1)R_n\left[(N+n+1)R_{n+1}n1\right]}{n\mathrm{\Omega }_n\{R\}(R_{n+1}1)\mathrm{\Omega }_{n1}\{R\}R_{n1}\left[(N+n)R_nn\right]}$$
(4.11)
We may therefore write down a recursion relation involving only the coefficients $`R_n`$ in the form
$$Q_n\{R\}^2Q_{n1}\{R\}^2=\mathrm{\Omega }_n\{R\}^2$$
(4.12)
where the function $`Q_n\{R\}^2`$ is given by (4.11) and $`\mathrm{\Omega }_n\{R\}`$ by (4.10). Notice that the coupling constant $`g`$ completely cancels out in this final equation which is an involved non-linear recursion relation for the $`R_n`$’s solely in terms of $`n`$ and $`N`$. The dependence of $`R_n`$ on $`g`$ comes about purely as a boundary condition and can only be recovered after solving for the dynamics of $`R_n`$. However, we will see that the entire topological expansion can be constructed without ever knowing this dependence explicitly. All information about the continuum limit will come from the details of the main recursion relation (4.12).
### 4.3 Planar Limit
The spherical continuum limit of the matrix model is found by taking the limit $`N\mathrm{}`$. Let us first describe the polymer contribution $`f_0`$ in (3.50). For the potential (4.3), the solution of the quadratic saddle-point equation (3.52) which is regular at $`g=0`$ is given by
$$\zeta =\frac{1}{2g}\left(\sqrt{1+4g}1\right)$$
(4.13)
so that the corresponding fermionic vector free energy (3.51) is
$$f_0(g)=\frac{1}{2}+\frac{1}{4g}\left(\sqrt{1+4g}1\right)\mathrm{log}\frac{1}{2g}\left(\sqrt{1+4g}1\right)$$
(4.14)
There is a phase transition at $`g=g^{(0)}=\frac{1}{4}`$ below which the free energy becomes complex-valued. Near this critical point, (4.14) behaves as $`f_0(g)(gg^{(0)})^{3/2}`$, which identifies the corresponding string susceptibility exponent as $`\gamma _{\mathrm{str}}^{(0)}=+\frac{1}{2}`$. The phase transition is therefore associated with the evolution of the system into a pure branched polymer phase of two-dimensional quantum gravity. In fact, the function (4.14) is related to the usual bosonic vector model free energy $`_{\mathrm{vec}}`$, which is the generating function for pure, connected branched polymer graphs, by $`f_0=12_{\mathrm{vec}}`$ . The $`N`$ dimensional Euclidean radial coordinate $`r`$ of the bosonic vector model is related to the saddle-point (4.13) by $`r^2=\frac{1}{2}\zeta `$. One can also compute the one-loop fluctuations around the saddle-point (4.13) and show that the resulting free energy is logarithmically divergent at $`g=0`$ . Therefore, in the regime $`g>0`$ there is no polymer behaviour and one may imagine a “barrier” in the large $`N`$ limit of the theory at the Gaussian point $`g=0`$. It is further possible to carry out a double scaling expansion of the fermionic vector model partition function (3.47) in the limit $`N\mathrm{},gg^{(0)}`$ with the parameter $`N(gg^{(0)})^{3/2}`$ held fixed. The resulting expression is an alternating, Borel summable series which can be expressed as a Bessel function .
We now consider the large $`N`$ behaviour associated with the second term in the free energy (3.50). Replacing the discrete $`R_n`$ coefficients with the large-$`N`$ continuous function $`R(x)`$ as prescribed by (3.49), after some algebra we obtain from the recursion relation (4.12), at leading order in $`N`$, the first order non-linear ordinary differential equation
$$\frac{dR}{dx}=\frac{R\left(R1\right)^2\left(x+2xR3(1+x)R^2\right)}{x^22x^2R+4x(1+x)R^22(1+x)(2+3x)R^3+3(1+x)^2R^4}$$
(4.15)
The dependence of $`R`$ on the coupling constant $`g`$ arises only from the constant of integration of (4.15). It is known from the general theory of ordinary differential equations that the solutions $`R(x)`$ to equations such as (4.15) possess algebraic non-analytic behaviour. Since $`dR/dx`$ in (4.15) is given as the quotient of a quintic polynomial in $`R`$ by a quartic one, it follows that the function $`R(x)`$ is finite for all finite values of $`x`$ . The only singularities which can arise from this differential equation occur when the denominator on the right-hand side vanishes and the numerator is non-zero. In light of the discussion of the previous subsection, we will demand that these singular points occur at the boundary $`x=1`$ of the domain of the function $`R(x)`$. At $`x=1`$, the denominator of the right-hand side of (4.15) has four distinct zeros $`R_c`$ determined as the solutions of the quartic equation
$$12R_c+8R_c^220R_c^3+12R_c^4=0$$
(4.16)
Denoting the corresponding critical value of the coupling constant $`g`$ by $`g_c`$, we make an ansatz for the form of the function $`R(x;g)`$ near the critical point,
$$R(x;g)=R_c+a(g_cgx)^{\gamma _{\mathrm{str}}}+\mathrm{}$$
(4.17)
where $`a`$ is some constant and the ellipsis denotes terms which are less singular at the critical point. This ansatz ensures that for $`gg_c`$, the singularities occur at $`x=1`$, or alternatively that at $`x=1`$ the free energy $`(x)`$ becomes non-analytic at $`g=g_c`$. Substituting (4.17) into the differential equation (4.15), we find for each branch $`R_c`$ of the equation (4.16) the leading behaviour
$$\frac{dR}{dx}\frac{1}{R_cR(x)}$$
(4.18)
Comparing with (4.17) fixes the exponent $`\gamma _{\mathrm{str}}=\frac{1}{2}`$, and so the critical point $`g=g_c`$ of the fermionic matrix model describes a continuum limit that lies in the same universality class as pure two-dimensional quantum gravity in the planar limit. Note that this criticality argument requires no knowledge of the precise value of the critical coupling $`g_c`$. The same will be true of the double scaling limit which will be analysed in the next subsection. It is possible to show using loop equations that $`g_c>0`$. Thus the continuum surface behaviour occurs on the opposite side of the Gaussian point relative to the branched polymer phase transition, and the two phases cannot be connected together in a smooth way.
We conclude our analysis of the planar limit of the adjoint fermion one-matrix model by briefly describing how it extends to more general potentials. For example, consider a potential of the generic form $`V(z)=z+\frac{g}{K}z^K`$, $`K2`$. The recursion equations (3.40) and (3.45) in this case require knowledge of the coefficients $`P_{n,n1}^{[K]}`$ and $`P_{n,n2}^{[K]}`$ determined from (3.20) and (3.28). It is straightforward to see from these relations that $`P_{n,n1}^{[K]}`$ is of degree $`K`$ in both the $`Q_n`$’s and the $`R_n`$’s, while $`P_{n,n2}^{[K]}`$ is of degree $`K+1`$. In the large $`N`$ limit, the recursion equations will therefore assume the form
$`(1+x)Rx`$ $`=`$ $`gQ^KW_1[R]`$ (4.19)
$`Px{\displaystyle \frac{P}{x}}`$ $`=`$ $`gQ^{K+1}W_2[R]`$ (4.20)
$`{\displaystyle \frac{P}{x}}`$ $`=`$ $`Q\left(R1\right)`$ (4.21)
where $`W_1[R]`$ and $`W_2[R]`$ are polynomials of degree $`K`$ and $`K+1`$ in the function $`R`$, respectively, which each contain the factor $`R1`$. Using (4.19) we can solve for $`Q^K`$ and write (4.20) as an equation determining the function $`Q`$. Comparing with the solution for $`Q`$ obtained using (4.21), the relation (4.20) then becomes a first order, linear, homogeneous differential equation for the function $`P`$. Substituting the resulting solution for $`P`$ back into (4.19) yields an integral equation for the function $`R`$ which, upon differentiation, can be transformed into a first order non-linear ordinary differential equation of the form
$$\frac{dR}{dx}=\frac{J_2[R;x]}{J_1[R;x]}$$
(4.22)
where $`J_1[R;x]`$ and $`J_2[R;x]`$ are polynomial functions of their arguments of respective degrees $`2K`$ and $`2K+1`$ in $`R`$ which are independent of the coupling constant $`g`$. Arguing as we did above, the differential equation (4.22) will admit critical behaviour for the function $`R`$ of the square root type, i.e. that of pure gravity. Provided that the zeros of the function $`J_1[R;x]`$ are non-degenerate at $`x=1`$, this will be the only type of critical behaviour that the model admits. It is of course natural that pure gravity exists for a generic matrix potential. Multicritical points require more complicated potentials with two or more coupling constants and an appropriate fine-tuning of them. In this case the differential equations derived analogously to those above become highly non-linear and quite involved. But the general picture will be the same, namely the function $`R`$ will be determined by some differential equation whose polynomial coefficients (in $`R`$ and $`x`$) can be tuned to obtain higher order critical points. In this way we can recover, in the planar limit, the standard universality classes of conformal matter coupled to two-dimensional quantum gravity .
### 4.4 Double Scaling Limit
In this subsection we will study the all genus expansion of the specific heat
$$u(g,N)=\underset{h=0}{\overset{\mathrm{}}{}}N^{2h}u_h(g),u_h(g)=F_h^{\prime \prime }(g)$$
(4.23)
where $`F_h(g)`$ is the genus $`h`$ contribution to the free energy $`F_N(g)`$ and the genus $`h`$ susceptibility is a homogeneous function of degree $`2h`$, $`u_h(ag)=a^{2h}u_h(g)`$. From the analysis of the previous subsection, we know that the leading singular behaviour of (4.23) at genus zero is $`u_0(g)\sqrt{gg_c}`$. In order to keep the partition function $`Z_0(g)\mathrm{e}^{N^2F_0(g)}`$ finite in the large $`N`$ limit, we should therefore perform the double scaling limit $`gg_c,N\mathrm{}`$ by keeping fixed the parameter
$$\kappa =N^{4/5}(g_cg)$$
(4.24)
We now define a scaling function $`u(\kappa )`$ which captures the contributions to the string susceptibility from all genera in the double scaling limit $`gg_c,N\mathrm{}`$ with the variable (4.24) fixed. For this, we introduce the large-$`N`$ expansion parameter $`ϵ=\frac{1}{N}`$ and write the total susceptibility (4.23) in the vicinity of the critical point as
$$u(g,N)=(ϵ^2)^{4/5}a_0+ϵ^2u(\kappa )+(ϵ^2)^{6/5}a_1\kappa +O((ϵ^2)^{7/5},\kappa ^2)$$
(4.25)
where $`a_n`$ are constants. By defining the new effective variable
$$\xi =ϵ^{4/5}(g_cgx)$$
(4.26)
in the large $`N`$ limit, this scaling function may be determined from the $`\xi =\kappa `$ limit of a scaling ansatz for the function $`R(x)`$ in terms of an unknown function $`u(\xi )`$,
$$R(x)=R_c+ϵ^{2/5}u(\xi )$$
(4.27)
where $`R_c`$ is a solution of (4.16). By substituting (4.27) into (3.50) and changing variables from $`x`$ to $`\xi `$ in the integral, we find that, up to irrelevant constants (which may be absorbed by suitable rescaling using the homogeneity of the specific heat) and terms which vanish as $`ϵ0`$, the free energy in the double scaling limit is determined as
$$F(\kappa )=f_0+\underset{ϵ^{4/5}g_c}{\overset{\kappa }{}}𝑑\xi (\kappa \xi )u(\xi )$$
(4.28)
It follows that
$$u(\kappa )=F^{\prime \prime }(\kappa )$$
(4.29)
and so the problem of obtaining the topological expansion of the fermionic matrix model is thereby reduced to the task of finding the solution of the main recursion relation (4.12) for the function $`R(x)`$ in this special limit. Note that the polymer free energy (4.14) contributes only an irrelevant constant in the double scaling limit about the positive-valued critical point $`g_c`$.
We now rewrite (4.12) using (3.49) and take the limit of large-$`N`$ while holding the quantity $`\kappa `$ fixed through the relation $`g=g_c\kappa N^{4/5}`$. We then substitute in the scaling ansatz (4.27) and rewrite derivatives according to the change of variables (4.26), i.e. $`_xϵ^{4/5}_\xi `$. After some algebra, we find, as the coefficient of the leading order $`ϵ^{4/5}`$ term, a non-linear differential equation for the specific heat $`u(\xi )`$,
$`0`$ $`=`$ $`6g_c\left(322R_c+80R_c^2192R_c^3+376R_c^4632R_c^5+752R_c^6512R_c^7+144R_c^8\right)`$
$`\times u(\xi )u^{}(\xi )+R_c\left(15R_c+10R_c^210R_c^3+4R_c^4\right)`$
$`\times \left[6R_c\left(R_c1\right)^2\left(6R_c^22R_c1\right)g_c^3\left(12R_c+20R_c^256R_c^3+36R_c^4\right)u^{\prime \prime \prime }(\xi )\right]`$
We now integrate (LABEL:nonlinDE) up once and use an inessential shift of the independent variable $`\xi `$ to eliminate the constant term. After applying (4.16) and a rescaling using the homogeneity of the function $`u(\xi )`$ (which preserves the relationship (4.29)), we arrive at the parameter-free equation
$$u^{\prime \prime }(\xi )+u(\xi )^2=\xi $$
(4.31)
The non-linear differential equation (4.31) governs the behaviour of the partition function of the fermionic random surface model to all orders in the genus expansion, with the parameter $`\xi ^{5/4}`$ identified as the renormalized string coupling constant. We note again the remarkable feature that one never needs to know the precise location of the critical point to arrive at this equation. In the present case the critical coupling constant $`g_c`$ can be consistently rescaled out of the pertinent equations. This parametric independence is indicative of the universality of the double scaling equation (4.31) for the given class of generic polynomial potentials.
The differential equation (4.31) is known as the Painlevé I equation and it is solved by the first Painlevé transcendent. This equation differs from the versions which usually appear in matrix models by a change in sign of the specific heat $`uu`$. The boundary conditions which $`u(\xi )`$ satisfies follow from the analysis of the planar limit of the previous subsection. The genus zero contribution arises in the limit of large, positive $`\xi `$ where the leading behaviour is $`u(\xi )^2=\xi `$. As a consequence, we are led to postulate an asymptotic expansion for $`u(\xi )`$ of the form
$$u(\xi )=\sqrt{\xi }\left(1+\underset{k=1}{\overset{\mathrm{}}{}}u_k\xi ^{5k/2}\right)$$
(4.32)
By substituting (4.32) into (4.31), the first few coefficients are easily calculated to be
$`u_1`$ $`=`$ $`{\displaystyle \frac{1}{8}}`$
$`u_2`$ $`=`$ $`{\displaystyle \frac{49}{128}}`$
$`u_3`$ $`=`$ $`{\displaystyle \frac{1225}{256}}`$
$`u_4`$ $`=`$ $`{\displaystyle \frac{4412401}{32768}}`$
$`u_5`$ $`=`$ $`{\displaystyle \frac{220680075}{32768}}`$
$`u_6`$ $`=`$ $`{\displaystyle \frac{2207064977649}{4194304}}`$
$`\mathrm{}`$
and in general they are determined by the recursive equation
$$u_k=\frac{(5k6)(5k4)}{8}u_{k1}\frac{1}{2}\underset{n=1}{\overset{k1}{}}u_ku_{kn},k2$$
(4.34)
With the normalization of the spherical specific heat taken as $`u_0=1`$, we see from (4.4) that the coefficients of the genus expansion have precisely the same absolute numerical values as those obtained for pure two-dimensional quantum gravity, i.e. from the usual Painlevé expansion . In particular, they have the same high rate asymptotic growth in magnitude. However, the most notable feature of the coefficients (4.4) is the fact that they oscillate in sign (In the Hermitian cases, all $`u_k`$ are negative). This raises the possibility that the asymptotic series solution (4.32) may in fact be Borel summable. If so, then the free energy would give an unambiguous definition of the genus expansion of pure “fermionic” gravity in powers of $`N^2\xi ^{5/2}`$. A numerical solution of the equation (4.31) is possible and the result is shown in fig. 1. This confirms numerically that the free energy exists as a well-defined function and is real-valued on the positive real $`\xi `$-axis. In the next subsection we shall prove that the coefficients $`u_k`$ of the asymptotic series (4.32) alternate in sign to arbitrarily large orders of the genus expansion, and, moreover, that there is a well-defined Borel resummation of $`u(\xi )`$ for $`\xi >0`$.
### 4.5 Borel Summability
In this subsection we will argue analytically, along the lines of , that the asymptotic series (4.32) defines a unique function $`u(\xi )`$. For this, we consider the Borel transform
$$B(s)=\underset{k=1}{\overset{\mathrm{}}{}}\frac{u_k}{(\beta k)!}s^k$$
(4.35)
where the constant $`\beta `$ will be self-consistently determined by the condition that this series has a finite radius of convergence. A solution of the original Painlevé equation (4.31) is then given by
$$u(\xi )=\sqrt{\xi }\left(1\underset{0}{\overset{\mathrm{}}{}}𝑑t\mathrm{e}^tB\left(t^\beta \xi ^{5/2}\right)\right)$$
(4.36)
The crucial issue now is whether or not the integral in (4.36) actually exists. We will show that there is a contour running through the $`\mathrm{Re}t>0`$ region of the complex $`t`$-plane from $`t=0`$ to $`t=\mathrm{}`$ along which the integral transform (4.36) converges. We will thereby argue that the Borel resummation (4.36) of the specific heat is well-defined and real-valued.
The function (4.35) has singularities in the complex $`s`$-plane with branch cuts running between them and infinity. If some of these singularities lie on the positive $`s`$-axis then the integral definition (4.36) is ambiguous. We will now extract the large order behaviour of the coefficients $`u_k`$ and argue that the singularities of the Borel transform $`B(s)`$ all lie on the negative $`s`$-axis. From (4.35) and (4.36) it follows that the coefficients of the asymptotic expansion (4.32) are given by
$$u_k=\frac{1}{2\pi i}\underset{0}{\overset{\mathrm{}}{}}𝑑t\mathrm{e}^t\underset{s=0}{}\frac{ds}{s^{k+1}}B\left(t^\beta s\right)$$
(4.37)
The contour of integration in (4.37) can be extended out to infinity around the cuts of the Borel transform in the complex $`s`$-plane. In the limit $`k\mathrm{}`$, the contour integral in (4.37) is dominated by contributions along the branch cut of $`B(s)`$ in the complex $`s`$-plane which begins at the point $`s=s_0`$ that is closest to the origin. We then have, for large $`k`$,
$$u_k\underset{0}{\overset{\mathrm{}}{}}𝑑t\mathrm{e}^t\underset{s_0/t^\beta }{\overset{\mathrm{}}{}}\frac{ds}{s^{k+1}}\mathrm{Disc}B\left(t^\beta s\right)\underset{0}{\overset{\mathrm{}}{}}\frac{ds}{s^{k+1}}\underset{(s_0/s)^{1/\beta }}{\overset{\mathrm{}}{}}𝑑t\mathrm{e}^t\mathrm{Disc}B\left(t^\beta s\right)$$
(4.38)
where we have defined the discontinuity of the Borel transform
$$\mathrm{Disc}B(s)=B_+(s)B_{}(s),B_\pm (s)=B(s\pm i0)$$
(4.39)
for $`s`$ a point on its cut.
Going back to (4.36), we see from (4.38) that we should study the functions $`u_\pm (\xi )`$ corresponding to the Borel transforms on either side of the dominant cut,
$$u_\pm (\xi )=\underset{(s_0\xi ^{5/2})^{1/\beta }}{\overset{\mathrm{}}{}}𝑑t\mathrm{e}^tB_\pm \left(t^\beta \xi ^{5/2}\right)$$
(4.40)
Since each of the functions (4.40) solves the Painlevé equation (4.31), it is convenient to define their sum and difference
$$u_\mathrm{d}(\xi )=u_+(\xi )u_{}(\xi )=\mathrm{Disc}B(\xi ),u_\mathrm{s}(\xi )=\frac{1}{2}(u_+(\xi )+u_{}(\xi ))$$
(4.41)
which on using (4.31) are seen to obey the coupled system of differential equations
$`u_\mathrm{d}^{\prime \prime }(\xi )+2u_\mathrm{d}(\xi )u_\mathrm{s}(\xi )`$ $`=`$ $`0`$ (4.42)
$`u_\mathrm{d}^{\prime \prime }(\xi )+u_\mathrm{s}(\xi )^2+{\displaystyle \frac{1}{4}}u_\mathrm{d}(\xi )^2`$ $`=`$ $`\xi `$ (4.43)
These equations can be solved in the WKB approximation. In the planar limit $`\xi \mathrm{}`$, we self-consistently assume that $`u_\mathrm{d}(\xi )`$ is of sub-leading order in (4.43) and thereby find, to leading order, the solution $`u_\mathrm{s}(\xi )\sqrt{\xi }`$. By substituting this into (4.42), to leading order we have the first order linear differential equation for the function $`u_\mathrm{d}(\xi )`$,
$$u_\mathrm{d}^{\prime \prime }(\xi )+2\sqrt{\xi }u_\mathrm{d}(\xi )=0$$
(4.44)
which can be solved in terms of Bessel functions as $`u_\mathrm{d}(\xi )=\sqrt{\xi }Z_{\frac{2}{5}}\left(\frac{4\sqrt{2}}{5}\xi ^{5/4}\right)`$. The solution for $`\xi \mathrm{}`$ is therefore given by
$$\frac{u_\mathrm{d}(\xi )}{\sqrt{\xi }}\xi ^{5/8}\mathrm{cos}\left(\frac{4\sqrt{2}}{5}\xi ^{5/4}\frac{5\pi }{4}\frac{\pi }{4}\right)$$
(4.45)
We are finally ready to estimate the large order behaviour of the coefficients $`u_k`$. For this, it is convenient to change variables in order to bring the asymptotic expansion for $`\xi \mathrm{}`$ to $`\omega 0`$ through the definition $`\omega =\xi ^{5/4}`$. Then, from (4.38)–(4.41) and (4.45), we find as $`k\mathrm{}`$,
$$u_k\frac{d\omega }{\omega ^{2k+1}}\omega ^{2/5}u_\mathrm{d}\left(\omega ^{4/5}\right)\frac{d\omega }{\omega ^{2k+1}}\sqrt{\omega }\mathrm{cos}\left(\frac{4\sqrt{2}}{5\omega }\frac{5\pi }{4}\frac{\pi }{4}\right)$$
(4.46)
For each choice of sign the integral (4.46) gives, up to constant factors, the result
$$u_k\left(\frac{32}{25}\right)^k\mathrm{\Gamma }\left(2k\frac{1}{2}\right)\mathrm{for}k\mathrm{}$$
(4.47)
The asymptotic estimate (4.47) contains a large amount of information. Going back to the Borel transform (4.35), the large order behaviour of $`u_k`$ fixes $`\beta =2`$ and gives
$$B(s)\underset{k}{}\left(\frac{32}{25}\right)^ks^k$$
(4.48)
We see therefore that the Borel transform $`B(s)`$ is well-defined in an open region of the complex $`s`$-plane and, moreover, that its first singularity appears on the negative real $`s`$-axis. Assuming that all of its singularities are so restricted, we can expect that the asymptotic series (4.32) can be Borel resummed through the integration in (4.36) to define a unique real function $`u(\xi )`$ on the positive real $`\xi `$-axis. The results of a numerical analysis (fig. 1) support these arguments.
From this analysis it follows that the double-scaled partition function of the adjoint fermion one-matrix model leads to a completely well-defined genus expansion of the corresponding random surface theory. The coefficients (4.47) have precisely the same large order behaviour $`\mathrm{\Gamma }(2k\frac{1}{2})`$ as those of pure two-dimensional quantum gravity. In fact, the topological expansion of the fermionic matrix model is identical term by term to that of the usual Hermitian matrix models. However, the alternating nature of the asymptotic series expansion allows one to define the string susceptibility in an unambiguous way through its Borel resummation (4.36), at least for positive values of the string coupling $`\xi `$. In this way we may think of the fermionic matrix model as lending a well-defined version of the generating functions provided by Hermitian matrix models. It gives an analytic solution to the problem of counting fermionic random triangulations on arbitrary genus Riemann surfaces, whose asymptotic expansion is Borel summable but otherwise coincides with the usual Painlevé expansion of two-dimensional quantum gravity.
The usual movable, double pole singularities of the Painlevé I equation appear as well in the present case, but this time they appear for negative values of the cosmological constant $`\xi `$. The existence of such poles is inconsistent with the loop equations of two-dimensional quantum gravity . Pole-free solutions of Painlevé equations are provided by the triply-truncated Boutroux solutions, but these are complex-valued and do not lead to physically acceptable free energy functions. However, in the present situation there is a unique, well-defined real-valued specific heat on the positive $`\xi `$-axis<sup>2</sup><sup>2</sup>2We thank A.A. Kapaev for pointing out the relevant mathematical literature on this point. which should have a meaningful analytical continuation to $`\xi <0`$. It provides an unambiguous definition of a whole branch of the random surface theory, and, given the relation between the double scaling variable $`\xi `$ and the continuum cosmological constant, the continuum limit of the discretized model is indeed well-defined. This suggests that by restricting our “fermionic gravity” model to $`\xi >0`$, we could define a model of two-dimensional quantum gravity with the canonical characteristics order by order in the genus expansion, but whose full asymptotic series makes perfect sense and provides a well-defined non-perturbative definition of the theory.
In the next and final section we will argue that this is indeed the case. In the usual Hermitian one-matrix models, the problem with the double scaling limit, in which the original matrix integral can only be defined by an analytical continuation, can be traced back to an instability in the corresponding double-scaled eigenvalue model. The non-compactness of the double-scaled eigenvalue space and the form of the effective potential for the eigenvalues demonstrates that the definition of the critical points of the model is unstable to the tunneling of eigenvalues into a different configuration . This leads to instanton solutions in the single-well eigenvalue model, and thereby explains the complexity of the free energy. However, in the present situation the Grassmann matrix integral is completely well-defined at finite $`N`$, in contrast to the Hermitian cases, and there is no reason a priori to expect that the model becomes unstable in any way at large $`N`$. Indeed, with the mapping of the fermionic matrix model onto a unitary matrix model, we see that the matrix integral is determined by an eigenvalue space whose topology is that of a circle. Being a compact space, there is no asymptotic behaviour in the effective eigenvalue action, and thus the compactness could have the effect of eliminating the eigenvalue tunneling problem. We see then that the unitary representation of the fermionic matrix model provides, at least naively, an eigenvalue description of the random surface model which reflects the reasons why the matrix integral is superior to its Hermitian counterparts. We shall tackle this problem using an operator theoretic approach which gives the fermionic analog of the (generalized) KdV hierarchical structure of the Hermitian one-matrix models . The integrable flows that we shall find are similar to those of the bosonic models, but with some very important changes. The most important difference will be the absence of a translational symmetry in eigenvalue space. This feature can be interpreted as allowing one to restrict the double-scaled specific heat to positive values of the cosmological constant $`\xi `$. It may then follow that the partition function of the fermionic matrix model serves as an unambiguous definition of that for two-dimensional quantum gravity.
We stress that the resulting Borel summability of the fermionic case is very different from that of Borel summable bosonic models, such as the conformal field theory of the Yang-Lee edge singularity in which the lattice system couples to an imaginary magnetic field and is therefore non-unitary . The coefficients of the genus expansion of the adjoint fermion one-matrix model are all real-valued<sup>3</sup><sup>3</sup>3In a similar, but non-unitary, alternating genus expansion was used to formulate a model of two-dimensional quantum gravity in terms of integrals over the moduli spaces of punctured spheres.. The stability of the fermionic model is also quite distinct from the stabilizations provided by supersymmetric and stochastic quantization methods , because these latter models violate the KdV flow structure of two-dimensional quantum gravity at a non-perturbative level . In the present case, we will see that the appropriate integrable hierarchical structure is a deformation of that for gravity, and so the KdV flow structure continues to hold in a certain sense. The precise meaning of this deformation comes from the worldsheet interpretation of the fermionic one-matrix model. This is the model of fermionic gravity that we alluded to in section 2.1, wherein “Penner matter states” are located at the vertices of the corresponding surface discretizations.
## 5 Operator Formalism
The discussion at the end of the previous section motivates the development of an operator approach to analyse more carefully the properties of the adjoint fermion one-matrix model within the orthogonal polynomial formalism. The main purpose is two-fold. First of all, it will allow us to relate the partition function of the fermionic matrix model to the $`\tau `$-function of an integrable hierarchy . The usefulness of this association is that it partitions the solutions of the fermionic matrix model into universality classes which are related to well-known integrable hierarchies, and thereby formulates the partition function in an invariant way as the solution to a set of differential equations, in particular in the double scaling limit. Secondly, given this relationship, we will obtain a clear geometric picture of the origin of the alternating, Borel summable genus expansion of the model. This will allow us to clarify the description of the string susceptibility as a well-defined, non-perturbative generating function for the fermionic gravity model, as alluded to at the end of the previous section, and at the same time it will present the appropriate generalization of the KdV flow structure that characterizes the gravity model, lending a complementary characterization to the worldsheet description. In the following we will begin by introducing the formalism and deriving the main equations that will be required. We will then identify the integrable hierarchy of which the fermionic partition function is a $`\tau `$-function, and derive the Virasoro constraints which must be satisfied by the flows.
### 5.1 String Equations
We begin by introducing multiplication operators $`𝐐`$ and $`\overline{𝐐}`$ defined on the space $`\text{}[z]\text{}[\overline{z}]`$ by
$$(𝐐\mathrm{\Phi }_m)(z)z\mathrm{\Phi }_m(z),(\overline{𝐐}\mathrm{\Lambda }_m)(z)\frac{1}{z}\mathrm{\Lambda }_m(z)$$
(5.1)
and differentiation operators $`𝐏`$ and $`\overline{𝐏}`$ by
$$(𝐏\mathrm{\Phi }_m)(z)\mathrm{\Phi }_m^{}(z),(\overline{𝐏}\mathrm{\Lambda }_m)(z)z^2\mathrm{\Lambda }_m^{}(z)$$
(5.2)
The operators $`𝐐`$ and $`\overline{𝐐}`$ are related to each other in a very simple way. From (5.1) it follows that
$$𝐐\mathrm{\Phi }_n(z)|\overline{𝐐}\mathrm{\Lambda }_m(z)=\mathrm{\Phi }_n(z)|\mathrm{\Lambda }_m(z)$$
(5.3)
and so $`𝐐`$ is an invertible operator with inverse
$$𝐐^1=\overline{𝐐}^{}$$
(5.4)
where the adjoint $`𝐎^{}`$ of any operator $`𝐎`$ on $`\text{}[z]\text{}[\overline{z}]`$ is defined by
$$F(z)|(𝐎G)(\overline{z})=(𝐎^{}F)(z)|G(\overline{z})$$
(5.5)
By writing $`[z\mathrm{\Phi }_m(z)]^{}`$ in two different ways as $`[(𝐐\mathrm{\Phi }_m)(z)]^{}`$ and $`\mathrm{\Phi }_m(z)+z\mathrm{\Phi }_m^{}(z)`$, we can infer the canonical commutation relation
$$[𝐏,𝐐]=\mathrm{𝟏}$$
(5.6)
Similarly, by considering $`z^2[\frac{1}{z}\mathrm{\Lambda }_m(z)]^{}`$ we arrive at $`[\overline{𝐐},\overline{𝐏}]=\mathrm{𝟏}`$, which with the unitarity condition (5.4) implies the canonical commutator
$$[\overline{𝐏}^{},𝐐^1]=\mathrm{𝟏}$$
(5.7)
Geometrically, the operator $`𝐏`$ is the canonical conjugate of the operator generating clockwise unit shifts on the circle, while $`\overline{𝐏}`$ serves as the canonical conjugate of that which generates counterclockwise shifts. The operators introduced above are not all independent, but are related to each other through a Schwinger-Dyson equation. This follows from the inner product
$`(𝐏\mathrm{\Phi }_n)(z)|\mathrm{\Lambda }_m(z)`$ $`=`$ $`\mathrm{\Phi }_n^{}(z)|\mathrm{\Lambda }_m(z)`$ (5.8)
$`=`$ $`(N+1)\mathrm{\Phi }_n(z)|\frac{1}{z}\mathrm{\Lambda }_m(z)NV^{}(z)\mathrm{\Phi }_n(z)|\mathrm{\Lambda }_m(z)`$
$`\mathrm{\Phi }_n(z)|\mathrm{\Lambda }_m^{}(z)`$
which implies the operator identity
$$𝐏+\overline{𝐏}^{}𝐐^2=(N+1)𝐐^1NV^{}(𝐐)$$
(5.9)
The equations (5.6), (5.7) and (5.9) define the string equations of the adjoint fermion one-matrix model.
Let us now describe some basic properties of the operators introduced above. For this, we first rewrite their actions on the same basis $`\mathrm{\Phi }_n`$ of the space of polynomials in a single variable on the circle. We have
$`(𝐐^{\pm 1}\mathrm{\Phi }_n)(z)`$ $`=`$ $`{\displaystyle \underset{m0}{}}\left[𝐐^{\pm 1}\right]_{nm}\mathrm{\Phi }_m(z)`$
$`(𝐏\mathrm{\Phi }_n)(z)`$ $`=`$ $`{\displaystyle \underset{m0}{}}\left[𝐏\right]_{nm}\mathrm{\Phi }_m(z)`$
$`(\overline{𝐏}^{}\mathrm{\Phi }_n)(z)`$ $`=`$ $`{\displaystyle \underset{m0}{}}\left[\overline{𝐏}\right]_{nm}\mathrm{\Phi }_m(z)`$ (5.10)
We denote by $`𝐎_+`$ the upper triangular part, including the main diagonal, of the discrete operator $`𝐎`$ when represented in the basis of polynomials in (5.10). Then $`𝐎_{}=𝐎𝐎_+`$ is the lower triangular part of $`𝐎`$.
From the definition (5.1) it follows that the operator $`𝐐`$ defines a Jacobi matrix, because
$$[𝐐]_{nm}=0\mathrm{for}mn>1$$
(5.11)
In particular, its pure upper triangular part is given by
$$𝐐_+\underset{n0}{}[𝐐]_{nn}𝐄_{n,n}=\underset{n0}{}𝐄_{n,n+1}$$
(5.12)
where $`𝐄_{n,m}`$ are the step operators with matrix elements $`[𝐄_{n,m}]_k\mathrm{}=\delta _{kn}\delta _\mathrm{}m`$. The remaining matrix elements of $`𝐐`$ are given by $`[𝐐]_{nm}=𝒫_m^{(n)}`$ in (3.28), along with the recurrence relations derived in section 3. Similarly, the operator $`𝐐^1`$ defines a Jacobi matrix in the basis (5.10) because
$$\left[𝐐^1\right]_{nm}=0\mathrm{for}nm>1$$
(5.13)
Since
$$h_{n1}\left[𝐐^1\right]_{n,n1}=(\overline{𝐐}^{}\mathrm{\Phi }_n)(z)|\mathrm{\Lambda }_{n1}(z)=\mathrm{\Phi }_n(z)|\frac{1}{z}\mathrm{\Lambda }_{n1}(z)=h_n$$
(5.14)
it follows that its lower triangular part is given by
$$𝐐_{}^1=\underset{n1}{}R_n𝐄_{n,n1}$$
(5.15)
The remaining matrix elements of $`𝐐^1`$ are given by $`[𝐐^1]_{nm}=_n^{(m)}`$ in (3.35). Furthermore, from the definition (5.2) we see that $`𝐏=𝐏_{}`$ is purely lower triangular with
$$[𝐏]_{n,n1}=n$$
(5.16)
and since
$$h_{n+1}\left[\overline{𝐏}\right]_{n,n+1}=(\overline{𝐏}^{}\mathrm{\Phi }_n)(z)|\mathrm{\Lambda }_{n+1}(z)=\mathrm{\Phi }_n(z)|z^2\mathrm{\Lambda }_{n+1}^{}(z)=(n+1)h_n$$
(5.17)
it follows that $`\overline{𝐏}^{}=\overline{𝐏}_+^{}`$ is purely upper triangular, $`\left[\overline{𝐏}\right]_{nn}=0`$, with
$$\left[\overline{𝐏}\right]_{n,n+1}=\frac{n+1}{R_{n+1}}$$
(5.18)
Note that, in this basis, the $`(n,n1)`$ matrix element of the Schwinger-Dyson equation (5.9) coincides with (3.40). This follows from the structure of the matrix elements
$$h_{n1}\left[\overline{𝐏}^{}𝐐^2\right]_{n,n1}=\mathrm{\Phi }_n(z)|\mathrm{\Lambda }_{n1}^{}(z)=(n1)h_n$$
(5.19)
which implies that $`\left[\overline{𝐏}^{}𝐐^2\right]_{n,n1}=(n1)R_n`$.
### 5.2 Flow Equations
We will now derive a system of flow equations for the operators above and the partition function (3.7), which will enable us to identify the particular integrable hierarchy at work here. We shall describe the evolutions with respect to the discrete set of “time” variables $`t_kNg_k/k`$ of the generic potential (1.2). Consider the relation
$$\frac{}{t_k}\mathrm{\Phi }_n(z)|\mathrm{\Lambda }_{nm}(z)=0=z^k\mathrm{\Phi }_n(z)|\mathrm{\Lambda }_{nm}(z)+\frac{}{t_k}\mathrm{\Phi }_n(z)|\mathrm{\Lambda }_{nm}(z)$$
(5.20)
where $`1mn`$ and we have used the fact that $`\frac{\mathrm{\Lambda }_{nm}(z)}{t_k}`$ is a polynomial of degree at most $`nm1`$. This relation leads immediately to the evolution equation
$$\frac{\mathrm{\Phi }_n(z)}{t_k}=\underset{m=0}{\overset{n1}{}}\left[𝐐^k\right]_{nm}\mathrm{\Phi }_m(z)=\left(𝐐_{}^k\mathrm{\Phi }_n\right)(z)$$
(5.21)
By taking time derivatives of the equations (5.10), we then arrive at the discrete operator flow equations
$`{\displaystyle \frac{𝐐}{t_k}}`$ $`=`$ $`[𝐐,𝐐_{}^k]`$ (5.22)
$`{\displaystyle \frac{𝐏}{t_k}}`$ $`=`$ $`[𝐏,𝐐_{}^k]`$ (5.23)
$`{\displaystyle \frac{𝐐^1}{t_k}}`$ $`=`$ $`[𝐐^1,𝐐_{}^k]`$ (5.24)
$`{\displaystyle \frac{\overline{𝐏}^{}}{t_k}}`$ $`=`$ $`[\overline{𝐏}^{},𝐐_{}^k]`$ (5.25)
The flow equations (5.22) constitute a discrete KP hierarchy . It is straightforward to show that these flows are mutually commutative, so that this hierarchy is in fact integrable. This follows from taking time derivatives of (5.10) and using (5.21) to get
$$\frac{\left[𝐐^k\right]_{mn}}{t_r}=\underset{\mathrm{}=mk}{\overset{n1}{}}\left[𝐐^k\right]_m\mathrm{}\left[𝐐^r\right]_\mathrm{}n\underset{\mathrm{}=m+1}{\overset{n+k}{}}\left[𝐐^r\right]_m\mathrm{}\left[𝐐^k\right]_\mathrm{}n$$
(5.26)
which leads to the Zakharov-Shabat equations
$$\frac{𝐐_{}^m}{t_n}\frac{𝐐_{}^n}{t_m}+[𝐐_{}^m,𝐐_{}^n]=0$$
(5.27)
The matrix model actually defines a reduced, generalized KP hierarchy, because the flow equations are to be supplemented with the constraints imposed by the string equations.
The knowledge of the solutions $`𝐐`$ to (5.22) completely determines the partition function of the fermionic matrix model. To see this, we consider the $`(n,n1)`$ matrix element of the flow equation (5.24), which using (5.13) and (5.15) leads to the differential equation
$$\frac{\mathrm{log}R_n}{t_k}=\left[𝐐^k\right]_{nn}\left[𝐐^k\right]_{n1,n1}$$
(5.28)
We can use (3.7) and (5.28) to obtain the flow equation for the partition function
$$\frac{\mathrm{log}Z_N}{t_k}=\mathrm{Tr}_{(N)}\left(𝐐^k\right)$$
(5.29)
where we have defined the $`N`$ dimensional trace
$$\mathrm{Tr}_{(N)}(𝐎)=\underset{n=0}{\overset{N1}{}}[𝐎]_{nn}$$
(5.30)
Using the KP equation (5.22) and the Jacobi property (5.11) of the $`𝐐`$ operator, we have
$$\frac{[𝐐]_{nn}}{t_k}=\left[𝐐^k\right]_{n+1,n}\left[𝐐^k\right]_{n,n1}$$
(5.31)
which, on using (5.28) with $`k=1`$, leads to
$$\frac{^2\mathrm{log}Z_N}{t_1t_k}=\left[𝐐^k\right]_{N,N1}$$
(5.32)
This procedure can be iterated to determine any number of time derivatives of $`\mathrm{log}Z_N`$ in terms of the $`𝐐`$ operators. Knowing $`𝐐`$ therefore determines the free energy of the matrix model up to an overall integration constant.
### 5.3 Fermionic $`\tau `$-Function
It is clear from (5.32) that the partition function $`Z_N`$ is a $`\tau `$-function for the discrete KP hierarchy (5.22) . Let us introduce the Baker-Akhiezer functions
$$\mathrm{\Psi }_n[\stackrel{}{t};z]=\frac{\mathrm{e}^{\frac{N}{2}V(z)}}{z^{\frac{N+1}{2}}}\frac{\mathrm{\Phi }_n(z)}{\sqrt{h_n}}$$
(5.33)
and their conjugates
$$\overline{\mathrm{\Psi }}_n[\stackrel{}{t};z]=\frac{\mathrm{e}^{\frac{N}{2}V(z)}}{z^{\frac{N+1}{2}}}\frac{\mathrm{\Lambda }_n(z)}{\sqrt{h_n}}$$
(5.34)
which are bi-orthonormal in the standard line measure of the complex plane,
$$\underset{z=0}{}𝑑z\mathrm{\Psi }_n[\stackrel{}{t};z]\overline{\mathrm{\Psi }}_m[\stackrel{}{t};z]=\delta _{nm}$$
(5.35)
The operator $`𝐐`$ is a discrete Lax type operator which has eigenfunctions $`\mathrm{\Psi }_n[\stackrel{}{t};z]`$ with eigenvalue $`z`$,
$$𝐐\mathrm{\Psi }_n[\stackrel{}{t};z]=z\mathrm{\Psi }_n[\stackrel{}{t};z]$$
(5.36)
The function $`Z_N[\stackrel{}{t}]`$ is then the generating function for the Lax operator $`𝐐`$, in the sense that $`𝐐`$ can be reconstructed from a knowledge of $`Z_N`$. In fact, if we introduce the discrete $`\tau `$-function $`\tau _n[\stackrel{}{t}]`$ for this hierarchy by the formula (4.1) , then from (3.7) we see that the partition function of the fermionic matrix model is given by
$$Z_N[\stackrel{}{t}]=\tau _N[\stackrel{}{t}]$$
(5.37)
Formally, the $`\tau `$-function is a section of the determinant line bundle over a Sato Grassmannian associated with Riemann surfaces of the spectral parameters $`z`$ . It may be characterized as the unique function depending on the times $`t_k`$ and satisfying the given hierarchy of constrained differential equations. It is the generating function for solutions of the generalized KdV equations. Thus once the $`\tau `$-function is known, everything about the matrix model is likewise known. The usefulness of the relationship (5.37) is that there are a large number of identities satisfied by $`\tau _N[\stackrel{}{t}]`$ that can be used to characterize the partition function of the fermionic matrix model.
The flow equations (5.22) or (5.26) constitute the Lax representation of a discrete integrable hierarchy of differential equations of Toda type . It is tempting therefore to characterize the fermionic partition function as the $`\tau `$-function of a certain reduction of the integrable Toda chain. However, as we will now demonstrate, this is not the case. By equating the constant terms on both sides of the flow equation (5.21) for $`k=1`$ and using (3.28) and (3.29), we have
$$\frac{p_{n,0}}{t_1}=\frac{1R_{n+1}}{l_{n+1,0}}\underset{m=0}{\overset{n1}{}}p_{m,0}l_{m,0}\underset{k=m+1}{\overset{n}{}}R_k$$
(5.38)
By iterating the relation (3.32) using (3.29) it follows that
$$p_{n,0}l_{n,0}=(1)^{n+1}(R_n1)$$
(5.39)
where we have defined $`R_00`$. Substituting (5.39) into (5.38) we obtain
$$\frac{p_{n,0}}{t_1}=p_{n+1,0}R_n=p_{n+1,0}\left(1+(1)^{n+1}p_{n,0}l_{n,0}\right)$$
(5.40)
Similarly, one can derive the evolution equation
$$\frac{l_{n,0}}{t_1}=l_{n1,0}\left(1+(1)^{n+1}p_{n,0}l_{n,0}\right)$$
(5.41)
The flow equations (5.40) and (5.41) are the first members of an integrable Hamiltonian system known as the Toeplitz chain hierarchy . The general evolution equations may be similarly derived using the flow equations of the previous subsection and are given by the Hamilton equations of motion
$$\frac{p_{n,0}}{t_k}=(1+(1)^{n+1}p_{n,0}l_{n,0})\frac{H_k}{l_{n,0}},\frac{l_{n,0}}{t_k}=(1+(1)^{n+1}p_{n,0}l_{n,0})\frac{H_k}{p_{n,0}}$$
(5.42)
for the system of Hamiltonians
$$H_k=\frac{1}{k}\mathrm{Tr}_{(N)}\left(𝐐^k\right)=\frac{1}{k}\frac{\mathrm{log}Z_N}{t_k},k1$$
(5.43)
and the symplectic structure
$$\omega =\underset{n1}{}\frac{dp_{n,0}dl_{n,0}}{1+(1)^{n+1}p_{n,0}l_{n,0}}$$
(5.44)
The Toeplitz chain is a particular reduction of the two-dimensional Toda lattice hierarchy, the latter of which characterizes generic matrix integrals . Thus the integrable hierarchy to which the adjoint fermion one-matrix model belongs is not the reduction of a one-dimensional (Toda) hierarchy, but rather of a two-dimensional one. This is of course anticipated from the lattice interpretation of the fermionic matrix model and its associated doubling of degrees of freedom due to the Penner matter states.
However, the Toeplitz chain hierarchy can be viewed as a deformation of the standard Toda chain hierarchy. To see this, we use (5.28) for $`k=1`$ and (3.28) for $`m=n`$ to obtain the flow equation
$$\frac{\mathrm{log}h_n}{t_1}=[𝐐]_{nn}=Q_n\left(R_{n+1}1\right)$$
(5.45)
On the other hand, by using (5.31) for $`k=1`$ and (3.28) for $`m=n1`$ we have
$$\frac{[𝐐]_{nn}}{t_1}=Q_n\left[Q_{n+1}R_{n+1}(R_{n+2}1)Q_{n1}R_n(R_{n+1}1)\right]$$
(5.46)
The two flow equations (5.45) and (5.46) may be combined together to give
$$\frac{^2\mathrm{log}h_n}{t_1^2}=\frac{\mathrm{log}h_n}{t_1}\left(\frac{R_{n+1}}{R_{n+1}1}\frac{\mathrm{log}h_{n+1}}{t_1}\frac{R_n}{R_n1}\frac{\mathrm{log}h_{n1}}{t_1}\right)$$
(5.47)
By introducing time-dependent functions $`q_n[\stackrel{}{t}]`$ through
$$h_n[\stackrel{}{t}]=(1)^n\epsilon ^{2n}\mathrm{e}^{t_1/\epsilon }\mathrm{e}^{q_n[\stackrel{}{t}]}$$
(5.48)
where $`\epsilon `$ is a time-independent parameter, we can write (5.47) as
$`{\displaystyle \frac{^2q_n}{t_1^2}}`$ $`=`$ $`\left(1+\epsilon {\displaystyle \frac{q_n}{t_1}}\right)\left(1+\epsilon {\displaystyle \frac{q_{n+1}}{t_1}}\right){\displaystyle \frac{\mathrm{e}^{q_{n+1}q_n}}{1+\epsilon ^2\mathrm{e}^{q_{n+1}q_n}}}`$ (5.49)
$`\left(1+\epsilon {\displaystyle \frac{q_{n1}}{t_1}}\right)\left(1+\epsilon {\displaystyle \frac{q_n}{t_1}}\right){\displaystyle \frac{\mathrm{e}^{q_nq_{n1}}}{1+\epsilon ^2\mathrm{e}^{q_nq_{n1}}}}`$
The differential equation (5.49) is the first member of an integrable hierarchy known as the relativistic Toda chain . In the “non-relativistic” limit $`\epsilon 0`$, it reduces to the Toda chain hierarchy which characterizes generic Hermitian one-matrix models. Thus, the partition function of the adjoint fermion one-matrix model is a $`\tau `$-function of not the Toda chain hierarchy, but rather of its deformation to the relativistic Toda chain. The latter chain is itself a reduction of the two-dimensional Toda lattice. From this point of view, the fermionic matrix model may be thought of as a “deformation” of the Hermitian one-matrix model, the role of the deformation being played by the Penner interaction.
The remaining differential equations of the relativistic Toda chain hierarchy are given by the Lax equations (5.26), by (5.28), and by exploiting the Jacobi properties (5.11) and (5.12) of the operator $`𝐐`$. It is important to realize that, because the present model contains only a single set $`\stackrel{}{t}`$ of times, the ordinary Toda lattice does not itself appear in the hierarchy satisfied by the partition function. Thus the standard Toda lattice structure which underlies all integrable models (and in particular matrix integrals) only appears in a very subtle way through the reductions obtained above by eliminating one set of its time variables. This reduced structure leads to a certain degeneracy in the $`\tau `$-function as compared to the usual two-dimensional Toda lattice structure. Of course, the relations to integrable hierarchies described in this subsection are only ‘kinematical’, as they merely rely on the very basic recursion properties satisfied by the orthogonal polynomials. To incorporate the ‘dynamical’ aspects, and in particular the effects of the Penner potential, we need to impose the constraints implied by the string equations. Indeed, there are many $`\tau `$-function solutions of the above hierarchies, and the string equations select the one appropriate to describe the nonperturbative dynamics of the adjoint fermion one-matrix model. This is the topic of the next subsection.
### 5.4 Virasoro Constraints
We will now begin examining the constraints imposed on the fermionic $`\tau `$-function $`Z_N[\stackrel{}{t}]`$ as dictated by the string equations. Proceeding as in (5.8) for the operators $`𝐏𝐐^{n+1}`$, $`n1`$, we can express the string equation (5.9) as an infinite system of equations
$$\mathrm{Tr}_{(N)}\left[(𝐏+\overline{𝐏}^{}𝐐^2(N+1)𝐐^1+\underset{k1}{}kt_k𝐐^{k1})𝐐^{n+1}\right]=0,n1$$
(5.50)
The strategy now is to rewrite the equations of motion (5.50) using the flow equations (5.29) and (5.32). In this way, the string equations will be represented as the annihilation of the partition function by a system of differential operators in the coupling constants of the potential $`V`$. In this subsection we will deal with the cases $`n0`$.
Let us start with the $`n=0`$ string equation of (5.50). Using the canonical commutation relations (5.6) and the matrix elements
$$h_n[𝐐𝐏]_{nn}=z\mathrm{\Phi }_n^{}(z)|\mathrm{\Lambda }_n(z)=nh_n$$
(5.51)
we may compute the first trace in (5.50) for $`n=0`$ as
$$\mathrm{Tr}_{(N)}(𝐏𝐐)=N+\underset{n=0}{\overset{N1}{}}n=\frac{N(N+1)}{2}$$
(5.52)
The second trace may be similarly computed by using
$$h_n\left[\overline{𝐏}^{}𝐐^1\right]_{nn}=\mathrm{\Phi }_n(z)|z\mathrm{\Lambda }_n^{}(z)=nh_n$$
(5.53)
to get
$$\mathrm{Tr}_{(N)}\left(\overline{𝐏}^{}𝐐^1\right)=\frac{N(N1)}{2}$$
(5.54)
By substituting (5.52) and (5.54), and using (5.29), we can represent the $`n=0`$ constraint of the system (5.50) as the flow equation
$$\left(\underset{k1}{}kt_k\frac{}{t_k}N^2\right)Z_N[\stackrel{}{t}]=0$$
(5.55)
Next, let us consider the $`n=1`$ constraint of (5.50). For the first trace, we use (3.5) and (3.6) to compute the matrix elements
$`h_n\left[𝐐^2𝐏\right]_{nn}`$ $`=`$ $`z^2\mathrm{\Phi }_n^{}(z)|\mathrm{\Lambda }_n(z)`$ (5.56)
$`=`$ $`nz^{n+1}(n1)p_{n,n1}z^n+O\left(z^{n1}\right)|\mathrm{\Lambda }_n(z)`$
$`=`$ $`\left[np_{n+1,n}(n1)p_{n,n1}\right]h_n`$
Using (3.16) and the canonical commutator (5.6) we thereby arrive after a little algebra at
$$\mathrm{Tr}_{(N)}\left(𝐏𝐐^2\right)=2\mathrm{Tr}_{(N)}(𝐐)+\underset{n=1}{\overset{N1}{}}\left(n[𝐐]_{nn}+\mathrm{Tr}_{(n)}(𝐐)\right)=(N+1)\mathrm{Tr}_{(N)}(𝐐)$$
(5.57)
Since $`\left[\overline{𝐏}^{}\right]_{nn}=0`$, upon substituting (5.57) into (5.50) we find that the terms linear in $`𝐐`$ cancel out and applying the flow equation (5.29) we see that the $`n=1`$ constraint can be written as
$$\underset{k1}{}kt_k\frac{}{t_{k+1}}Z_N[\stackrel{}{t}]=0$$
(5.58)
Now we move on to the $`n=2`$ string equation of (5.50). Again using (3.5) and (3.6) we may compute
$`h_n\left[𝐐^3𝐏\right]_{nn}`$ $`=`$ $`z^3\mathrm{\Phi }_n^{}(z)|\mathrm{\Lambda }_n(z)`$
$`=`$ $`nz^{n+2}(n1)p_{n,n1}z^{n+1}(n2)p_{n,n2}z^n+O\left(z^{n1}\right)|\mathrm{\Lambda }_n(z)`$
$`=`$ $`\left[n(p_{n+2,n}+p_{n+2,n+1}p_{n+1,n})(n2)p_{n,n2}(n1)p_{n+1,n}p_{n,n1}\right]h_n`$
Using (3.16) we can iterate the relation (3.14) for $`k=n2`$ to get
$$p_{n,n2}=\underset{k=1}{\overset{n1}{}}\left([𝐐]_{k,k1}+[𝐐]_{kk}\mathrm{Tr}_{(k)}(𝐐)\right)$$
(5.60)
Substituting (5.60) and (3.16) into (LABEL:Q3Pnn) and using (5.6), we arrive after some algebra at
$`\mathrm{Tr}_{(N)}\left(𝐏𝐐^3\right)`$ $`=`$ $`3\mathrm{Tr}_{(N)}\left(𝐐^2\right)+(N1)[𝐐]_{N,N1}+(N+1)\left(\mathrm{Tr}_{(N)}(𝐐)\right)^2`$ (5.61)
$`+(2N3){\displaystyle \underset{n=1}{\overset{N1}{}}}\left([𝐐]_{n,n1}[𝐐]_{nn}\mathrm{Tr}_{(n)}(𝐐)\right)`$
The sums in (5.61) can be simplified by using the Jacobi property (5.11) of the operator $`𝐐`$ to get
$$\mathrm{Tr}_{(N)}\left(𝐐^2\right)=[𝐐]_{N,N1}+2\underset{n=1}{\overset{N1}{}}[𝐐]_{n,n1}+\underset{n=0}{\overset{N1}{}}\left([𝐐]_{nn}\right)^2$$
(5.62)
and in this way we arrive at an expression for the first trace in the $`n=2`$ equation of (5.50),
$$\mathrm{Tr}_{(N)}\left(𝐏𝐐^3\right)=\left(N+\frac{3}{2}\right)\mathrm{Tr}_{(N)}\left(𝐐^2\right)+\frac{1}{2}[𝐐]_{N,N1}+\frac{1}{2}\left(\mathrm{Tr}_{(N)}(𝐐)\right)^2$$
(5.63)
For the second trace in (5.50), we use (3.5) to compute the matrix elements
$`h_n\left[\overline{𝐏}^{}𝐐\right]_{nn}`$ $`=`$ $`\mathrm{\Phi }_n(z)|z^3\mathrm{\Lambda }_n^{}(z)`$ (5.64)
$`=`$ $`\mathrm{\Phi }_n(z)|nz^{(n2)}+O\left(z^{(n3)}\right)+l_{n,1}z`$
$`=`$ $`l_{n,1}z\mathrm{\Phi }_n(z)|1`$
Using (3.25) we obtain
$$\left[\overline{𝐏}^{}𝐐\right]_{nn}=S_nQ_n(1R_{n+1})$$
(5.65)
where we have defined
$$S_n=\frac{l_{n,1}}{l_{n,0}}$$
(5.66)
By equating the coefficients of the $`z^1`$ terms on both sides of the three-term recursion relation (3.36) for the $`\mathrm{\Lambda }`$ polynomials, we can write down an iterative equation for the coefficients (5.66),
$$S_n=R_nQ_{n1}+S_{n+1}Q_n$$
(5.67)
which has solution
$$S_n=Q_{n1}R_n+\underset{k=0}{\overset{n1}{}}Q_k(1R_{k+1})$$
(5.68)
Substituting (5.68) into (5.65) and using (3.28) for $`m=n1`$ and $`m=n`$, we therefore find
$$\left[\overline{𝐏}^{}𝐐\right]_{nn}=[𝐐]_{nn}\mathrm{Tr}_{(n)}(𝐐)[𝐐]_{n,n1}$$
(5.69)
Using the identity (5.62) we then arrive at an expression for the second trace in the $`n=2`$ equation of (5.50),
$$\mathrm{Tr}_{(N)}\left(\overline{𝐏}^{}𝐐\right)=\frac{1}{2}[𝐐]_{N,N1}+\frac{1}{2}\left(\mathrm{Tr}_{(N)}(𝐐)\right)^2\frac{1}{2}\mathrm{Tr}_{(N)}\left(𝐐^2\right)$$
(5.70)
Upon substitution of (5.63) and (5.70) into (5.50) for $`n=2`$, we see that the terms quadratic in $`𝐐`$ cancel out. The remaining terms can be simplified using the flow equations (5.29) and (5.32), which leads to the $`n=2`$ constraint equation
$$\left(\underset{k1}{}kt_k\frac{}{t_{k+2}}+\frac{^2}{t_1^2}\right)Z_N[\stackrel{}{t}]=0$$
(5.71)
This procedure can be easily generalized to higher orders $`n3`$. However, it is possible to determine the general relation from the first three equations (5.55), (5.58) and (5.71) by using the fact that the commutator bracket of the pertinent differential operators must also annihilate the partition function and thereby using commutators of the above operators to generate the higher order ones. The final result can be expressed in a more familiar form by introducing the zeroth time $`t_0`$ through the flow equation
$$\frac{Z_N}{t_0}=NZ_N$$
(5.72)
for the partition function.
In this way we find that the string equations (5.50) for $`n0`$ can be written as the set of discrete Virasoro constraints
$$L_nZ_N[\stackrel{}{t}]=0,n0$$
(5.73)
where the second order linear differential operators
$$L_n=\underset{k1}{}kt_k\frac{}{t_{k+n}}+\underset{k=0}{\overset{n}{}}\frac{}{t_k}\frac{}{t_{nk}}2N\frac{}{t_n}$$
(5.74)
satisfy the Witt algebra
$$[L_n,L_m]=(nm)L_{n+m}$$
(5.75)
for $`n,m\text{}^+`$. These Virasoro constraints are the usual ones for the fermionic matrix model and they are identical to those of the Hermitian Penner matrix model as defined in (1.3) with $`\alpha =2`$. They represent the full system of equations of motion of the matrix model and can be equivalently derived by demanding the invariance of the matrix integral under arbitrary changes of the matrix variables . The first set of operators in (5.74) comes from the variation of the potential $`V`$ in the action, while the second one comes from the Jacobian of the change of matrix integration measure. The first two terms in (5.74) therefore coincide with the standard Virasoro generators of generic Hermitian one-matrix models . The last operator in (5.74) comes from the variation of the logarithmic Penner interaction.
It is a remarkable fact that the loop equations of the adjoint fermion one-matrix model (1.1), the Hermitian Penner one-matrix model ((1.3) with $`\alpha =2`$), and the unitary Penner matrix model (1.4) are all equivalent. However, the equivalence between the Hermitian model and the other two matrix models only holds order by order in the $`\frac{1}{N}`$ expansions of the matrix integrals. Although the models are equivalent at $`N=\mathrm{}`$, beyond leading order in the large $`N`$ expansion the loop equations should be solved with different boundary conditions and the solutions are different. The fermionic and unitary equations should be solved with boundary conditions appropriate to a perturbative Gaussian limit, in order that the models admit interpretations in terms of random surfaces, while the Hermitian ones should be solved with boundary conditions appropriate to recover the Penner model of the discretized moduli space of Riemann surfaces . The distinction of the fermionic and unitary matrix integrals from the Hermitian one is evident from the structure of the finite $`N`$ correlators (see appendix A). In fact, it is only in the large $`N`$ limit that a full Virasoro symmetry is realized in these models, as well as an infinite hierarchy of generalized KP differential equations. This makes the large $`N`$ limit a very natural ingredient of the fermionic matrix model, as is also apparent from its fat-graph interpretation of section 2.2.
### 5.5 Origin of the Painlevé Expansion
In the previous subsection we have shown that the Schwinger-Dyson equations of the fermionic matrix model are generated by the positive Borel subalgebra of a Virasoro algebra of vanishing central charge. There is a nice algebraic way to characterize the effect of the Penner interaction potential in the Virasoro generators (5.74). Let us write them as $`L_n=L_n^\mathrm{H}+T_n`$, where $`L_n^\mathrm{H}`$ are the standard Virasoro generators of a Hermitian one-matrix model with potential $`V`$, and
$$T_n=2N\frac{}{t_n}$$
(5.76)
are the generators of translations $`t_nt_n2N\epsilon `$ in coupling constant space. Together these operators generate the symmetry algebra
$`[T_n,T_m]`$ $`=`$ $`0`$
$`[L_n^\mathrm{H},L_m^\mathrm{H}]`$ $`=`$ $`(nm)L_{n+m}^\mathrm{H}`$
$`[T_n,L_m^\mathrm{H}]`$ $`=`$ $`2NnT_{n+m}`$ (5.77)
The commutation relations (5.77) characterize a master symmetry of the generalized KdV hierarchy . We see then that the Virasoro constraints of the fermionic matrix model constitute a deformation of those for Hermitian one-matrix models by the master symmetry algebra generators of the KdV flow equations. Note that the particular symmetry generated by the operators (5.76) is restricted to translations in units of a discrete “lattice spacing” $`2N`$, indicating what sort of reduction of this master symmetry is taken. These Virasoro constraints describe the desired reduction of the fermionic $`\tau `$-function, and they lead immediately to the desired geometrical interpretation of the Painlevé expansion of the adjoint fermion one-matrix model. The Painlevé equations may be characterized as following from a reduction of the $`\tau `$-function of the integrable Toeplitz (or relativistic Toda) chain hierarchy, subject to the Virasoro symmetry obtained in the previous subsection. They yield a stability condition $`L_n\tau _N=0`$ on the points of the model Grassmannian, which are associated with the Baker-Akhiezer functions (5.33), that select a particular class of transcendental solutions to the KP equations.
It is a standard fact that a set of operators of the form (5.74) can be embedded into a Virasoro algebra of non-vanishing central charge. However, an important aspect of the present approach to the fermionic matrix model is that there are in fact extra dynamical constraints imposed on the partition function which can be attributed to the invertibility of the Lax operator $`𝐐`$. For instance, the constraint (5.50) makes perfect sense for $`n=1`$, and the calculation of the trace of the operator $`\overline{𝐏}^{}𝐐^2`$ proceeds in an analogous manner to that of (5.56) with the result
$$\mathrm{Tr}_{(N)}\left(\overline{𝐏}^{}𝐐^2\right)=(N1)\mathrm{Tr}_{(N)}\left(𝐐^1\right)$$
(5.78)
Since $`[𝐏]_{nn}=0`$, we see that on substitution of (5.78) into (5.50) for $`n=1`$ we produce an extra non-vanishing term $`2N\mathrm{Tr}_{(N)}(𝐐^1)`$. This trace can be written in an operator form as follows. From (3.35) with $`m=n`$ and (5.45) we have
$$\left[𝐐^1\right]_{nn}=\left(R_n1\right)\left(R_{n+1}1\right)\left(\frac{\mathrm{log}h_n}{t_1}\right)^1$$
(5.79)
This quantity is independent of the first time $`t_1`$, since on using the flow equation (5.24) for $`k=1`$ and the Jacobi properties of the operators $`𝐐`$ and $`\overline{𝐐}`$ we find $`[𝐐^1]_{nn}/t_1=0`$. This suggests defining a negative time $`t_1`$ such that $`\mathrm{Tr}_{(N)}(𝐐^1)`$ is represented as an operator acting on the partition function as
$`{\displaystyle \frac{DZ_N}{Dt_1}}`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{N1}{}}}{\displaystyle \frac{1}{\tau _n\tau _{n+1}}}{\displaystyle \frac{\left(\tau _{n+1}\tau _{n1}\tau _n^2\right)\left(\tau _{n+2}\tau _n\tau _{n+1}^2\right)}{\tau _n\frac{\tau _{n+1}}{t_1}\tau _{n+1}\frac{\tau _n}{t_1}}}`$
$`{\displaystyle \frac{}{t_1}}{\displaystyle \frac{DZ_N}{Dt_1}}`$ $`=`$ $`0`$ (5.80)
where we have used (4.1). Then the constraint (5.50) for $`n=1`$ can be written as
$$\left(\underset{k2}{}kt_k\frac{}{t_{k1}}+Nt_12N\frac{D}{Dt_1}\right)Z_N[\stackrel{}{t}]=0$$
(5.81)
In fact, there are infinitely many negatively moded constraints such as this, associated with higher powers $`𝐐^m`$, $`m1`$. However, the operators such as that in (5.81) which are thereby generated do not form a closed algebra among themselves or with the operators (5.74). Indeed, this set of constraints arises from the previously mentioned reduction from a two-dimensional lattice (containing both negative and positive time parameters) to a (deformed) chain by eliminating all negative couplings. As usual, such reductions break the full Virasoro symmetry of the original integrable model and only the positive Borel subalgebra survives here as a symmetry of the matrix model. The negatively moded constraints are therefore merely an artifact of the lattice reduction or equivalently the chain deformation. In practical terms, all observables in the fermionic model are described by observables conjugate to positive times in the unitary matrix model. Observables conjugate to negative times exist but do not appear to play a role in the equivalence.
We can understand all of these matters more clearly by rewriting the Virasoro constraints (5.73) in terms of the Baker-Akhiezer functions (5.33). For this, we note first of all that (5.21) and (5.28) imply that they obey the flow equations
$$\frac{\mathrm{\Psi }_n[\stackrel{}{t};z]}{t_k}=\frac{1}{2}\left(𝐐_+^k𝐐_{}^k\left[𝐐^k\right]_{nn}\right)\mathrm{\Psi }_n[\stackrel{}{t};z]$$
(5.82)
Now let us compute the action of the scale transformation generator $`z\frac{}{z}`$ in eigenvalue space on the Baker-Akhiezer functions. We have
$$z\frac{\mathrm{\Psi }_n[\stackrel{}{t};z]}{z}=\left(\frac{N+3}{2}+\frac{N}{2}zV^{}(z)\right)\mathrm{\Psi }_n[\stackrel{}{t};z]+\frac{\mathrm{e}^{\frac{N}{2}V(z)}}{z^{\frac{N+1}{2}}}𝐏𝐐\mathrm{\Phi }_n(z)$$
(5.83)
To evaluate the last term in (5.83), we have to take into account the triangularity of the operator $`𝐏𝐐`$. From $`𝐏=𝐏_{}`$ and the Jacobi property (5.11), we have
$$𝐏𝐐=(𝐏𝐐)_{}+\underset{n0}{}[𝐏𝐐]_{nn}𝐄_{n,n}$$
(5.84)
We can therefore evaluate the action of $`𝐏𝐐`$ on $`\mathrm{\Phi }_n(z)`$ by using the string equation (5.9) multiplied on the right by $`𝐐`$, but eliminating the pure upper triangular part of both sides of the equation. Using (5.53) and the fact that
$$\left(\overline{𝐏}^{}𝐐^1\right)_{}=0$$
(5.85)
which can be derived similarly to the other matrix elements of this section, we can write (5.83) as
$`z{\displaystyle \frac{\mathrm{\Psi }_n[\stackrel{}{t};z]}{z}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(2n+N1\right)\mathrm{\Psi }_n[\stackrel{}{t};z]`$ (5.86)
$`+{\displaystyle \frac{N}{2}}\left[\left(𝐐V^{}(𝐐)\right)_+\left(𝐐V^{}(𝐐)\right)_{}\left[𝐐V^{}(𝐐)\right]_{nn}\right]\mathrm{\Psi }_n[\stackrel{}{t};z]`$
The last term in (5.86) can be rewritten using the flow equation (5.82), leading to the eigenvalue equation
$$L_0[\stackrel{}{t};z]\mathrm{\Psi }_n[\stackrel{}{t};z]=n\mathrm{\Psi }_n[\stackrel{}{t};z]$$
(5.87)
where we have defined the Virasoro operator
$$L_0[\stackrel{}{t};z]=\underset{k1}{}kt_k\frac{}{t_k}z\frac{}{z}+\frac{N1}{2}$$
(5.88)
The eigenvalue equation (5.87) shows that the lowest Virasoro constraint is indeed a true symmetry of the system, in that the operators $`𝐐`$ and $`L_0[\stackrel{}{t};z]`$ are simultaneously diagonalizable in the basis of bi-orthonormal Baker-Akhiezer functions. The same is true of all higher Virasoro constraints. This shows precisely what sort of constraints the auxiliary eigenvalue problem (5.36) must possess in order to correctly reconstruct the fermionic $`\tau `$-function. However, an analogous computation using the translation generator $`\frac{}{z}`$ on eigenvalue space leads to the result
$$\left(\underset{k1}{}kt_k\frac{}{t_{k1}}\frac{}{z}\frac{N+1}{2z}\right)\mathrm{\Psi }_n[\stackrel{}{t};z]=(n+N)R_n\mathrm{\Psi }_{n1}[\stackrel{}{t};z]$$
(5.89)
which shows that there is no corresponding $`L_1`$ operator which commutes with the Lax operator $`𝐐`$. This means that there is no translational symmetry in the system, and the leading equations of motion at large $`N`$ are determined instead by the generator $`L_0`$ of scale transformations of the system. These equations can be written in a form comparable to those of sections 3 and 4 by writing a consistency condition between the eigenvalue equations (5.36) and (5.87) in the form
$$\left(L_0[\stackrel{}{t};z]+n\right)\left(𝐐z\right)\mathrm{\Psi }_n[\stackrel{}{t};z]=0$$
(5.90)
Expanding (5.90) using (5.87), (5.88) and completeness of the Baker-Akhiezer functions then leads to the system of equations
$$(1+nm)[𝐐]_{nm}+\underset{k1}{}kt_k\frac{[𝐐]_{nm}}{t_k}=0$$
(5.91)
In particular, setting $`n=m`$ in (5.91) and using the flow equation (5.31), we have
$$[𝐐]_{nn}=\underset{k1}{}kt_k\left(\left[𝐐^k\right]_{n,n1}\left[𝐐^k\right]_{n+1,n}\right)$$
(5.92)
What this means is that it is the operator $`𝐏𝐐`$ (along with its conjugate $`\overline{𝐏}^{}𝐐^1`$) which plays the fundamental role in the dynamics of the adjoint fermion one-matrix model. Thus it is the scale invariance of the system which leads to the fundamental double scaling equations in the large $`N`$ limit. The translational symmetry generated by the operator $`𝐏`$ itself (and its conjugate $`\overline{𝐏}^{}𝐐^2`$) is an extra constraint on the theory which plays no immediate role in the continuum limit of the matrix model. By considering the matrix elements $`\mathrm{\Phi }_n(z)|z\left[z\mathrm{\Lambda }_m(z)\right]^{}`$, it is possible to infer the commutation relation $`[𝐐,\overline{𝐏}^{}𝐐^1]=𝐐`$. Using the canonical commutator (5.6), it follows that the fundamental symmetry operator of the system is given by
$$\mathrm{\Pi }=\frac{1}{2}\left(𝐏𝐐\overline{𝐏}^{}𝐐^1\right)$$
(5.93)
which together with the Lax operator $`𝐐`$ obeys the non-canonical commutation relation
$$[\mathrm{\Pi },𝐐]=𝐐$$
(5.94)
The operator expression (5.94) defines the appropriate string equations of the fermionic matrix model. The upper and lower triangular parts of the scaling operator (5.93) may be computed by using (5.9), (5.51), (5.84) and (5.85) to get
$`\mathrm{\Pi }_+`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{n0}{}}(nN)𝐄_{n,n}+{\displaystyle \frac{1}{2}}{\displaystyle \underset{k1}{}}kt_k𝐐_+^k`$
$`\mathrm{\Pi }_{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{k1}{}}kt_k𝐐_{}^k`$ (5.95)
Substituting (5.95) into the string equation (5.94) and using the KP equation (5.22), we arrive at (5.91).
It is the equation (5.92) which produces the novel Painlevé expansions of the fermionic matrix model. For instance, in the case of the quadratic potential studied in section 4, it is implicit in the manipulations carried out in section 4.2. In fact, the structural form of (5.91) is identical to the so-called “automodel” constraints which arise in the usual Hermitian one-matrix models . On the other hand, in the continuum limit $`𝐐`$ becomes a differential operator of a certain finite degree and the string equations (5.94) can be translated into an equation for pseudo-differential operators which satisfy a generalized KdV hierarchy of flow equations . General solutions of string equations such as (5.94) have been studied before, in the case that the double scaling limit of $`𝐐`$ is described by a Schrödinger operator, with stable, pole-free solutions for the corresponding string susceptibility . We see therefore that the string equations of the fermionic one-matrix model lead to a Painlevé differential equation which has the same structure as that in the usual Hermitian one-matrix models, except that now the automodel form of the equations produce an alternating genus expansion. Moreover, the structural form (5.94) of the string equations allow us to restrict to positive values of the double-scaling variable $`\xi `$ and thereby obtain pole-free solutions for the free energy of the fermionic matrix model.
In the present case, the details of the continuum limit of the operator $`𝐐`$ are somewhat involved. However, the discrete equations (5.91), (5.92) and (5.94) completely characterize the double scaled equations for the partition function in the large $`N`$ limit and the ensuing topological expansion. They serve as the starting point for a characterization of the differential hierarchies of the fermionic one-matrix model. From the generic structure of the flow equations obtained in this section, we can expect that the continuum equations satisfied by the partition function of “fermionic quantum gravity” are closely related to the usual KdV flow structure of two-dimensional quantum gravity in terms of Gelfand-Dikii differential polynomials .
### Acknowledgments
We thank J. Ambjørn, C. Kristjansen, Y. Makeenko, G. Semenoff and J. Wheater for helpful discussions. We also thank A.A. Kapaev, M. Matone and P. Zinn-Justin for comments on the manuscript. The work of L.D.P. was supported in part by the Natural Sciences and Engineering Research Council of Canada and NSF grant PHY98-02484. The work of R.J.S. was supported in part by the Danish Natural Science Research Council.
## Appendix A Spectral Correlation Functions
In this appendix we will briefly describe some properties of correlators of the fermionic one-matrix model within the formalism of this paper. The symmetries of the matrix integral (1.1) restrict its observables to those which are invariant functions of $`\overline{\psi }\psi `$ . Connected correlation functions may be generated by taking derivatives of the free energy with respect to the coupling constants of the potential $`V`$,
$$\underset{j=1}{\overset{L}{}}\mathrm{tr}\left(\overline{\psi }\psi \right)^{p_j}_{\mathrm{F},\mathrm{conn}}=\frac{1}{N^L}\underset{j=1}{\overset{L}{}}p_j\frac{}{g_{p_j}}\mathrm{log}Z_N$$
(A.1)
where the normalized fermionic correlation functions are defined by
$$f\left(\overline{\psi }\psi \right)_\mathrm{F}\frac{1}{Z_N}\underset{\mathrm{Gr}(N)^c}{}𝑑\psi 𝑑\overline{\psi }f\left(\overline{\psi }\psi \right)\mathrm{e}^{N\mathrm{tr}V(\overline{\psi }\psi )}$$
(A.2)
and in (A.1) it is understood that, if necessary, a set of auxiliary coupling constants $`g_k`$ are introduced into the potential $`V`$ and then set to zero after differentiation. Given the equivalence of the generating function $`\mathrm{log}Z_N`$ for the connected correlators of the fermionic and unitary matrix models, we therefore also have complete equivalence of their observables,
$$\underset{j=1}{\overset{L}{}}\mathrm{tr}\left(\overline{\psi }\psi \right)^{p_j}_{\mathrm{F},\mathrm{conn}}=\underset{j=1}{\overset{L}{}}\mathrm{tr}U^{p_j}_{\mathrm{U},\mathrm{conn}}$$
(A.3)
where the normalized unitary correlation functions are defined by
$$f(U)_\mathrm{U}\frac{k_N}{Z_N}\underset{U(N)}{}[dU]f(U)\mathrm{e}^{N\mathrm{tr}\left(V(U)\mathrm{log}U\right)}$$
(A.4)
The correlation functions in (A.3) may be evaluated as in (2.29)–(2.33). For example, one may readily compute
$$\mathrm{tr}\left(\overline{\psi }\psi \right)^p_\mathrm{F}=\mathrm{e}^{ip\theta }\mathrm{tr}U_0^p_\mathrm{U}=\frac{k_N}{NZ_N}\left[\underset{SU(N)}{}[dU_0]\mathrm{tr}U_0^p\mathrm{e}^{N\mathrm{tr}V(U_0)}\right]_{C_1=N^2p}$$
(A.5)
The restriction in (A.5) of the character expansion of $`\mathrm{e}^{N\mathrm{tr}V(U_0)}`$ to Young tableaux with $`N^2p`$ filled boxes implies that single trace correlators are non-vanishing only for $`pN^2`$. The finiteness of the set of non-vanishing correlators is of course natural as a consequence of the anticommuting property of the matrices $`\psi `$ and $`\overline{\psi }`$, but it is a dramatic result in the unitary matrix model. Thus not only does the Penner interaction term in the unitary matrix model (2.28) capture the finite nature of the perturbation expansion of $`Z_N`$, but it also reduces the correlation functions appropriately to reproduce the correct properties of the fermionic correlators at finite $`N`$. Note that, geometrically, the operator $`\mathrm{tr}(\overline{\psi }\psi )^p`$ inserts, on the dual triangulated surface, a hole with $`p`$ boundary lengths. The lattice expansion of such operators thereby generates fermionic ribbon graphs which are dual to tesselations of Riemann surfaces with a given number of boundaries. Again there are only finitely many such fat-graphs, because the maximum number of boundaries that a given discretization can have is $`N^2`$.
The complete set of observables of the fermionic matrix model can be generated by the joint probability distributions
$$\rho _n(z_1,\mathrm{},z_n)=\frac{(Nn)!}{N!}\underset{k=1}{\overset{n}{}}\mathrm{tr}\delta \left(z_k\overline{\psi }\psi \right)_\mathrm{F}$$
(A.6)
where $`z_k`$ are points in the complex plane and $`1nN`$. When the correlation function in (A.6) is mapped to the unitary matrix model as described above, the points $`z_k`$ can be interpreted as eigenvalues of unitary matrices. By diagonalizing the unitary matrices in (A.4), using the identities
$$\mathrm{\Delta }(z_1,\mathrm{},z_N)=\underset{i,j}{det}\left[\mathrm{\Phi }_{j1}(z_i)\right],\overline{\mathrm{\Delta }(z_1,\mathrm{},z_N)}=\underset{i,j}{det}\left[\mathrm{\Lambda }_{j1}(z_i)\right]$$
(A.7)
which follow from (2.6) and (3.5), and by using (3.7), it is straightforward to derive the determinant representation
$$\rho _n(z_1,\mathrm{},z_n)=\frac{(Nn)!}{N!}\underset{k=1}{\overset{n}{}}\mathrm{tr}\delta (z_kU)_\mathrm{U}=\frac{(Nn)!}{N!}\underset{i,j}{det}\left[𝒦(z_i,z_j)\right]$$
(A.8)
where $`𝒦(z,z^{})`$ is the spectral kernel which is defined in terms of the orthogonal polynomials as
$$𝒦(z,z^{})=\frac{\mathrm{e}^{\frac{N}{2}\left(V(z)+V(z^{})\right)}}{(zz^{})^{\frac{N+1}{2}}}\underset{n=0}{\overset{N1}{}}\frac{\mathrm{\Phi }_n(z)\mathrm{\Lambda }_n(z^{})}{h_n}$$
(A.9)
We see therefore that the problem of evaluating correlation functions of the fermionic matrix model reduces to that of determining the spectral kernel (A.9). It is possible to express it in a much simpler form by deriving the appropriate generalization of the Christoffel-Darboux formula . For this, we first need to derive a “mixed” recursion relation between the $`\mathrm{\Phi }`$ and $`\mathrm{\Lambda }`$ polynomials. Consider the bi-orthogonality relations
$`\mathrm{\Phi }_{n+1}(z)z\mathrm{\Phi }_n(z)|z^m`$ $`=`$ $`\mathrm{\Phi }_{n+1}(z)|z^m\mathrm{\Phi }_n(z)|z^{(m1)}=0`$
$`z^n\mathrm{\Lambda }_n(z)|z^m`$ $`=`$ $`z^{nm}|\mathrm{\Lambda }_n(z)=0`$ (A.10)
which are valid for $`1mn`$. From (A.10) it follows that the two polynomials $`\mathrm{\Phi }_{n+1}(z)z\mathrm{\Phi }_n(z)`$ and $`z^n\mathrm{\Lambda }_n(z)`$ of degree $`n`$ are equal up to some constant $`c_n`$. Equating the constant terms of these polynomials gives $`c_n=p_{n+1,0}`$, and we arrive at the mixed three-term recursion relation
$$\mathrm{\Phi }_{n+1}(z)=z\mathrm{\Phi }_n(z)+p_{n+1,0}z^n\mathrm{\Lambda }_n(z)$$
(A.11)
In an analogous way, it is possible to derive the recursion relation
$$\mathrm{\Lambda }_{n+1}(z)=\frac{1}{z}\mathrm{\Lambda }_n(z)+l_{n+1,0}\frac{1}{z^n}\mathrm{\Phi }_n(z)$$
(A.12)
We now multiply (A.11) through by $`\mathrm{\Lambda }_n(z^{})/h_n`$ and (A.12) with the index shift $`nn1`$ and $`zz^{}`$ through by $`z^{}\mathrm{\Phi }_n(z)/h_n`$, subtract the resulting two equations, and then sum over $`n=1,\mathrm{},N1`$. This yields
$`(zz^{}){\displaystyle \underset{n=1}{\overset{N1}{}}}{\displaystyle \frac{\mathrm{\Phi }_n(z)\mathrm{\Lambda }_n(z^{})}{h_n}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Phi }_N(z)\mathrm{\Lambda }_{N1}(z^{})}{h_{N1}}}{\displaystyle \frac{\mathrm{\Phi }_1(z)}{h_1}}`$
$`{\displaystyle \underset{n=1}{\overset{N1}{}}}{\displaystyle \frac{1}{h_k}}\left[p_{n+1,0}z^n\mathrm{\Lambda }_n(z)\mathrm{\Lambda }_n(z^{}){\displaystyle \frac{l_{n,0}}{z^{n1}}}\mathrm{\Phi }_n(z)\mathrm{\Phi }_{n1}(z^{})\right]`$
We can iterate the recursion equations (A.11) and (A.12) to get
$$z^k\mathrm{\Lambda }_k(z)=1+z\underset{j=1}{\overset{k}{}}l_{j,0}\mathrm{\Phi }_{j1}(z),\frac{1}{z^{k1}}\mathrm{\Phi }_{k1}(z^{})=1+\frac{1}{z^{}}\underset{j=1}{\overset{k1}{}}p_{j,0}\mathrm{\Lambda }_{j1}(z^{})$$
(A.14)
Substituting (A.14) into (LABEL:mixeddiff) and comparing the resulting equation with itself under the interchange of arguments $`zz^{}`$, we arrive after some algebra at
$`\left(zz^{}\right){\displaystyle \underset{n=0}{\overset{N1}{}}}{\displaystyle \frac{\mathrm{\Phi }_n(z)\mathrm{\Lambda }_n(z^{})}{h_n}}`$ $`=`$ $`{\displaystyle \frac{z^{N1}\mathrm{\Lambda }_{N1}(z^{})\mathrm{\Phi }_N(z)z^{N1}\mathrm{\Lambda }_{N1}(z)\mathrm{\Phi }_N(z^{})}{h_{N1}z^{N1}}}`$ (A.15)
$`+{\displaystyle \frac{1h_0}{h_0}}\left(zz^{}\right)`$
The expression (A.15) is valid for any pair of complex numbers $`zz^{}`$. We can take the $`zz^{}`$ limit of (A.15) using l’Hospital’s rule to get
$$\underset{n=0}{\overset{N1}{}}\frac{\mathrm{\Phi }_n(z)\mathrm{\Lambda }_n(z)}{h_n}=\frac{z^{N1}\mathrm{\Lambda }_{N1}(z)\mathrm{\Phi }_N^{}(z)\mathrm{\Phi }_N(z)\left(z^{N1}\mathrm{\Lambda }_{N1}(z)\right)^{}}{h_{N1}z^{N1}}+\frac{1h_0}{h_0}$$
(A.16)
The identities (A.15) and (A.16) are the fermionic analogs of the Christoffel-Darboux formula for the usual Hermitian one-matrix model orthogonal polynomials.
The Christoffel-Darboux formula allows us to express the spectral kernel (A.9) in terms of orthogonal polynomials with indices $`n`$ close to the matrix dimension $`N`$. Using the mixed recursion relation (A.11), we obtain
$`𝒦(z,z^{})`$ $`=`$ $`{\displaystyle \frac{\mathrm{e}^{\frac{N}{2}\left(V(z)+V(z^{})\right)}}{(zz^{})^{\frac{N+1}{2}}}}`$
$`\times \left\{{\displaystyle \frac{z^{N+1}}{h_{N1}p_{N,0}(zz^{})}}\left[z\mathrm{\Phi }_{N1}(z)\mathrm{\Phi }_N(z^{})z^{}\mathrm{\Phi }_{N1}(z^{})\mathrm{\Phi }_N(z)\right]+{\displaystyle \frac{1h_0}{h_0}}\right\}`$
It follows that the fermionic matrix model is completely determined by the single set of $`\mathrm{\Phi }`$ polynomials, as expected because of the one-matrix nature of the model and the heuristic identification $`U\overline{\psi }\psi `$. In particular, for the spectral density we find
$`\rho (z)`$ $``$ $`{\displaystyle \frac{1}{N}}𝒦(z,z)`$ (A.18)
$`=`$ $`{\displaystyle \frac{\mathrm{e}^{NV(z)}}{Nz^N}}\{{\displaystyle \frac{1}{h_{N1}p_{N,0}z^N}}[z(\mathrm{\Phi }_N(z)\mathrm{\Phi }_{N1}^{}(z)\mathrm{\Phi }_N^{}(z)\mathrm{\Phi }_{N1}(z))+\mathrm{\Phi }_N(z)\mathrm{\Phi }_{N1}(z)]`$
$`+{\displaystyle \frac{1h_0}{h_0}}\}`$
The formula (A.18) is particularly useful for determining the spectral density, and hence all correlators, of the fermionic matrix model in the large $`N`$ limit. In that case, factorization and symmetry imply that all connected correlation functions vanish and the large $`N`$ limit of the model is completely characterized by the set of correlators
$$\frac{1}{N}\mathrm{tr}\left(\overline{\psi }\psi \right)^p_\mathrm{F}=𝑑z\rho (z)z^p$$
(A.19)
where the integral goes over the support of the function $`\rho (z)`$ in the complex plane. Note that the spectral density is defined a priori in the adjoint fermion matrix model by the formula (A.6) . In the present formalism it has a natural interpretation as the probability density of eigenvalues in the corresponding unitary one-matrix model. As such, it is generally supported on the unit circle. However, as discussed in the previous subsection, the restriction of the integration contour to the unit circle in the large $`N`$ limit is not necessary, and the spectral density can be supported generically in the complex plane. These facts explain the general properties of the spectral densities of fermionic matrix models , for instance how the pole generated by the Penner interaction requires the support contour of $`\rho (z)`$ to be adjusted so as to avoid the origin of the complex plane and why the support endpoints are complex-valued.
## Appendix B Solution of the Gaussian Model
In this appendix we will illustrate how the formalism of section 3 works in a specific example. We shall consider the example of the Gaussian potential
$$V(z)=z$$
(B.1)
for which everything can be obtained explicitly. In this case the recursion relations of section 3.2 are easily solved to give the coefficients
$`R_n`$ $`=`$ $`{\displaystyle \frac{n}{n+N}}`$
$`P_n`$ $`=`$ $`{\displaystyle \frac{n}{N}}`$
$`Q_n`$ $`=`$ $`\left(1+{\displaystyle \frac{n+1}{N}}\right)`$ (B.2)
and the normalization constants are
$$h_n=2\pi iN^N\frac{n!}{(n+N)!}$$
(B.3)
Substituting (B.2) into the three-term recursion relation (3.31), we obtain
$$\mathrm{\Phi }_n(z)=\left(z+\frac{n1+N}{N}\right)\mathrm{\Phi }_{n1}(z)\frac{(n1)z}{N}\mathrm{\Phi }_{n2}(z)$$
(B.4)
To solve the recurrence relation (B.4), we introduce the generating function which is defined as the formal power series
$$\mathrm{\Xi }_N(z;s)=\underset{n=0}{\overset{\mathrm{}}{}}\mathrm{\Phi }_n(z)s^n$$
(B.5)
in a variable $`s`$, with the boundary condition $`\mathrm{\Xi }_N(z;0)=1`$. The recursion relation (B.4) is then equivalent to a first order inhomogeneous linear differential equation for the function $`\mathrm{\Xi }_N`$,
$$\left[\frac{s^2}{N}\left(1sz\right)\frac{}{s}+s\left(z+1\frac{sz}{N}\right)1\right]\mathrm{\Xi }_N(z;s)=1$$
(B.6)
Integrating (B.6) we find that the generating function (B.5) is given by
$$\mathrm{\Xi }_N(z;s)=\frac{1}{s}\underset{k=0}{\overset{\mathrm{}}{}}\frac{(N+k1)!}{(N1)!}\left(\frac{1}{N}\right)^k\left(\frac{s}{1sz}\right)^{k+1}$$
(B.7)
Expanding the function (B.7) as a power series in $`s`$ and equating its coefficient of $`s^n`$ with that of the definition (B.5), we arrive at an explicit form for the polynomials $`\mathrm{\Phi }_n(z)`$,
$$\mathrm{\Phi }_n(z)=\underset{k=0}{\overset{n}{}}\left(\genfrac{}{}{0pt}{}{n}{k}\right)\frac{(N+k1)!}{(N1)!}\left(\frac{1}{N}\right)^kz^{nk}$$
(B.8)
which can be expressed as
$$\mathrm{\Phi }_n(z)=N^Nz^{N+n}U[N,n+1+N;Nz]$$
(B.9)
where
$$U[a,c;x]=\frac{\pi }{\mathrm{sin}\pi c}\left(\frac{{}_{1}{}^{}F_{1}^{}[a,c;x]}{\mathrm{\Gamma }(ac+1)\mathrm{\Gamma }(c)}+\frac{x^{1c}{}_{1}{}^{}F_{1}^{}[ac+1,2c;x]}{\mathrm{\Gamma }(a)\mathrm{\Gamma }(2c)}\right)$$
(B.10)
is a confluent hypergeometric function. The partition function and observables of the Gaussian fermionic one-matrix model are therefore completely determined by a system of confluent hypergeometric polynomials.
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# 1 Introduction
## 1 Introduction
For the reader’s convenience we first remind the definition of quiver varieties. Let $`Q`$ be a quiver, that is a finite oriented graph with vertex set $`I`$. Let $`V=\{V_i\}_{iI}`$ be a collection of finite dimensional $``$-vector spaces. By a representation of $`Q`$ in $`V`$ we mean an assignment of a linear map: $`V_iV_j`$, for any pair $`i,jI`$ and each oriented edge of $`Q`$ with tail $`i`$ and head $`j`$. Let $`(Q,V)`$ denote the set of all representations of $`Q`$ in $`V`$, which is a $``$-vector space. The group $`_{iI}\mathrm{𝙶𝙻}_{_{}}(V_i)`$ acts naturally on $`(Q,V)`$, by conjugation. This action clearly factors through $`G(V):=\left(_i\mathrm{𝙶𝙻}(V_i)\right)/^{},`$ the quotient by the group $`^{}`$ imbedded diagonally, as scalar matrices, into each of groups $`\mathrm{𝙶𝙻}(V_i)`$. Let $`𝔤(V)=\left(𝔤𝔩(V_i)\right)/`$ denote the Lie algebra of the group $`G(V)`$.
Let $`\overline{Q}`$ be the double of $`Q`$, the quiver obtained by adding a reverse arrow $`a^{}`$, for every (oriented) arrow $`aQ`$. For any $`V=\{V_i\}_{iI}`$, the vector space $`(\overline{Q},V)`$ may be identified naturally with $`T^{}(Q,V)=`$ cotangent bundle on $`(Q,V).`$ Hence, $`(\overline{Q},V)`$ has a canonical symplectic structure. Furthermore, the $`G(V)`$-action on $`(\overline{Q},V)`$ is Hamiltonian, and the corresponding moment map $`\mu :(\overline{Q},V)𝔤(V)^{}`$ is given by the following formula:
$$\varrho \mu \left(\varrho \right)=\left\{\left\{\mu \left(\varrho \right)_i\right\}_{_{iI}}𝔤𝔩\left(V_i\right)\right|\mu \left(\varrho \right)__i=\underset{\left\{\genfrac{}{}{0pt}{}{aQ}{\text{head}\left(a\right)=i}\right\}}{}\varrho \left(a\right)\varrho \left(a^{}\right)\underset{\left\{\genfrac{}{}{0pt}{}{aQ}{\text{tail}\left(a\right)=i}\right\}}{}\varrho \left(a^{}\right)\varrho \left(a\right)\}$$
(1.1)
Here and below, we identify $`𝔤(V)^{}`$ with a subspace in $`𝔤𝔩(V_i)`$ by means of the trace pairing: $`x,y_{iI}\mathrm{𝗍𝗋}(x_iy_i).`$ Specifically, we have:
$$𝔤(V)^{}𝔰𝔤(V):=\{x=(x_i)_{iI}𝔤𝔩(V_i)|_{iI}\mathrm{𝗍𝗋}(x_i)=0\}.$$
Example. Let $`Q`$ be the quiver consisting of a single vertex and a single edge-loop at this vertex. Thus $`\overline{Q}`$ is the quiver with two edge-loops at the same vertex. Clearly, giving a representation of $`\overline{Q}`$ in the vector space $`V=^n`$ amounts to giving an arbitrary pair of $`n\times n`$-matrices. Therefore, we have: $`(\overline{Q},^n)=𝔤𝔩_n𝔤𝔩_n`$, and hence: $`G(V)=\mathrm{𝙿𝙶𝙻}_n`$. The moment map (1.1) reduces to the map $`\mu :𝔤𝔩_n𝔤𝔩_n𝔤(V)^{}=𝔰𝔩_n,`$ given by the formula: $`(x,y)[x,y].`$
Next, fix $`𝖮𝔰𝔤(V)`$, a closed $`\mathrm{𝙰𝚍}G(V)`$-orbit, and assume that the group $`G(V)`$ acts freely on the subvariety $`\mu ^1(𝖮)(\overline{Q},V)`$. Then, the orbit space $`_𝖮(\overline{Q},V):=\mu ^1(𝖮)/G(V)`$ is an affine variety, to be called an affine quiver variety. Thus, by definition: $`_𝖮(\overline{Q},V)`$ $`:=\mathrm{𝚂𝚙𝚎𝚌}\left([(\overline{Q},V)]^{G(V)}/^{G(V)}\right),`$ where $`[(\overline{Q},V)]`$ stands for the defining ideal of the subvariety $`\mu ^1(𝖮)`$, and we have used that $`[(\overline{Q},V)]^{G(V)}/^{G(V)}=([(\overline{Q},V)]/)^{G(V)},`$ due to reductivity of $`G(V)`$. If $`\mu ^1(𝖮)`$ is smooth then $`_𝖮(\overline{Q},V)`$ is also smooth, and the symplectic structure on $`(\overline{Q},V)`$ induces, via the symplectic reduction construction, see \[GS\], a canonical symplectic structure on $`_𝖮(\overline{Q},V)`$.
One of the main results of this paper is
###### Theorem 1.2.
In the above setting, the symplectic variety $`_𝖮(\overline{Q},V)`$ can be imbedded as a coadjoint orbit in the dual of $`𝔏(Q),`$ an infinite dimensional Lie algebra canonically attached to the quiver $`Q`$.
It is implicit in the theorem that the symplectic structure on $`_𝖮(\overline{Q},V)`$ goes, under the imbedding, into the canonical Kirillov-Kostant symplectic structure on the coadjoint orbit. Note also that the Lie algebra $`𝔏(Q)`$ does not depend on the representation space $`V`$.
Remark. A choice of Hermitian metric on $`V`$ makes $`(\overline{Q},V)`$ a flat hyper-Kähler space. An equivalence: holomorphic symplectic reduction $``$ hyper-Kähler reduction, see \[Hi\], gives, for many orbits $`𝖮`$, a hyper-Kähler structure on the quiver variety $`_𝖮(\overline{Q},V)`$. Recall further that by a well-known result of Kronheimer \[Kr\], any coadjoint orbit in a complex reductive Lie algebra has a hyper-Kähler structure. Based on this analogy, N. Hitchin asked if the Calogero-Moser space (a special case of quiver variety, see below) is a coadjoint orbit of some infinite dimensional Lie algebra. Hitchin’s question has been motivated by the recent work of Berest-Wilson \[BW\], who constructed a transitive action of $`Aut(A_1)`$, the automorphism group of the Weyl algebra, on the Calogero-Moser space. Theorem 1.2 gives a positive answer to Hitchin’s question and sheds some new light on the Berest-Wilson construction.
Strategy of the proof of Theorem 1.2. The symplectic structure on $`_𝖮(\overline{Q},V)`$ makes the coordinate ring $`[_𝖮(\overline{Q},V)]`$ an infinite dimensional Lie algebra with respect to the Poisson bracket. We will construct a sequence of Lie algebra morphisms:
$$𝔏(Q)\stackrel{\psi }{}[(\overline{Q},V)]^{G(V)}\stackrel{\mathrm{𝚙𝚛}}{}[(\overline{Q},V)]^{G(V)}/^{G(V)}=[_𝖮(\overline{Q},V)],$$
(1.3)
where $`[_𝖮(\overline{Q},V)]`$, the coordinate ring, is viewed as a Lie algebra with respect to the Poisson bracket arising from the symplectic structure on $`(\overline{Q},V)`$, and the map $`\mathrm{𝚙𝚛}`$ stands for the canonical projection.
Now, for any affine symplectic manifold $`X`$ and any point $`xX`$, evaluation at $`x`$ gives a linear function on $`[X]`$, whence induces an evaluation map: $`X\stackrel{\mathrm{𝚎𝚟}}{}[X]^{^{}}`$. Note that the vector space $`[X]^{^{}}`$ is an (infinite dimensional) Poisson manifold with Kirillov-Kostant bracket. It is immediate from the definitions that the map: $`X[X]^{^{}}`$ is a morphism of Poisson varieties, i.e., the induced map on polynomial functions is a morphism of Poisson algebras. Since $`X`$ is smooth and affine, regular functions on $`X`$ separate points of $`X`$ and, moreover, the differentials of regular functions span tangent spaces at each point of $`X`$. This implies that the evaluation map is injective, and that the infinitesimal Hamiltonian action of the Lie algebra $`[X]`$ (with the Poisson bracket) on the image of the evaluation map is infinitesimally transitive. Thus, the evaluation imbedding makes $`X`$ a coadjoint orbit in $`[X]^{^{}}`$.
Applying the considerations above to the symplectic manifold $`X=_𝖮(\overline{Q},V)`$, and dualizing the maps in (1.3), one gets a sequence of Poisson morphisms:
$$_𝖮(\overline{Q},V)\stackrel{\mathrm{𝚎𝚟}}{}[_𝖮(\overline{Q},V)]^{}\stackrel{\mathrm{𝚙𝚛}^{}}{}\left([(\overline{Q},V)]^{G(V)}\right)^{}\stackrel{\psi ^{}}{}𝔏(Q)^{}.$$
It will be shown later that the composite map above is injective, and the image of $`_𝖮(\overline{Q},V)`$ is a coadjoint orbit in $`𝔏(Q)^{}`$. Thus, a key step in proving Theorem 1.2 is the construction of Lie algebra map $`\psi `$ in (1.3).
We now illustrate our construction of $`\psi `$ in a very special case, where $`Q`$ is the quiver consisting of a single vertex and a single edge-loop (see Example above). To define the Lie algebra $`𝔏(Q)`$, it is convenient to introduce an auxiliary 2-dimensional symplectic vector space $`(E,\omega )`$ with basis $`x,y`$ (corresponding to the two loops in $`\overline{Q}`$) such that $`\omega (x,y)=1`$. For any $`p,q0`$, we define a $``$-bilinear map $`\{,\}__\omega :E^p\times E^qE^{(p+q2)}`$ by the formula:
$`\{u_1u_2\mathrm{}u_p,v_1v_2\mathrm{}v_q\}__\omega =`$ (1.4)
$`{\displaystyle \underset{i=1}{\overset{p}{}}}{\displaystyle \underset{j=1}{\overset{q}{}}}\omega (u_i,v_j)u_{i+1}\mathrm{}u_pu_1\mathrm{}u_{i1}v_{j+1}\mathrm{}v_qv_1\mathrm{}v_{j1},`$
where $`u_1,\mathrm{},u_p,v_1,\mathrm{},v_qE`$. Assembled together, these maps give a bilinear pairing $`\{,\}{}_{\omega }{}^{}:TE\times TETE,`$ where $`TE=_{i0}E^i`$ is the tensor algebra of $`E`$. Let $`[TE,TE]TE`$ denote the $``$-linear span of the set $`\{abba\}_{a,bTE}.`$
###### Proposition 1.5.
The pairing $`\{,\}_\omega `$ gives rise to a well-defined Lie algebra structure on the vector space $`𝔏(Q):=TE/[TE,TE]`$.
Remark. One of the goals of the paper is to give an interpretation of the Lie algebra $`(TE/[TE,TE],\{,\}{}_{\omega }{}^{})`$ as a sort of Poisson algebra associated to an appropriate ‘non-commutative’ symplectic variety.
To complete our construction we must define a Lie algebra morphism $`\psi :𝔏(Q)=TE/[TE,TE][(\overline{Q},V)]^{G(V)},`$ see (1.3). As we know, for $`V=^n`$ one has: $`(\overline{Q},V)𝔤𝔩_n()𝔤𝔩_n(),`$ and $`G(V)\mathrm{𝙿𝙶𝙻}_n`$. It is convenient to identify the tensor algebra $`TE`$ with the free associative algebra generated by $`x,y`$. We define a $``$-linear map $`\mathrm{𝗍𝗋}:TE[𝔤𝔩_n𝔤𝔩_n]`$ by assigning to any non-commutative monomial $`f=x^{k_1}y^{l_1}x^{k_2}\mathrm{}TE`$ a polynomial function $`\mathrm{𝗍𝗋}f[𝔤𝔩_n𝔤𝔩_n],`$ given by the formula:
$$\mathrm{𝗍𝗋}f:(X,Y)\mathrm{𝚃𝚛𝚊𝚌𝚎}(X^{k_1}Y^{l_1}X^{k_2}\mathrm{}),X,Y𝔤=𝔤𝔩_n.$$
(1.6)
It is clear that $`\mathrm{𝗍𝗋}f[𝔤𝔩_n𝔤𝔩_n]^{\mathrm{𝙶𝙻}_n}`$, and that $`\mathrm{𝗍𝗋}f=0`$ if $`f[TE,TE]`$, by symmetry of the trace. Thus, the assignment: $`f\mathrm{𝗍𝗋}f`$ gives a well-defined linear map $`\psi :𝔏(Q)=TE/[TE,TE][𝔤𝔩_n𝔤𝔩_n]^{\mathrm{𝙶𝙻}_n}`$. It turns out that this map is a Lie algebra morphism. This completes our construction, and the outline of the proof of Theorem 1.2
Example: Calogero-Moser space. Let $`Q`$ be the quiver consisting of a single vertex and a single edge-loop at this vertex, and assume $`dimV=n`$, as above. Then, $`𝔤(V)=𝔭𝔤𝔩_n`$. We will be concerned with the coadjoint orbit $`𝖮𝔤(V)^{}=𝔰𝔩_n`$, formed by all $`n\times n`$-matrices of the form: $`s\text{-}\mathrm{𝙸𝚍}`$, where $`s`$ is a rank 1 semisimple matrix such that $`\text{Trace}(s)=\text{Trace}(\mathrm{𝙸𝚍})=n`$. Thus, $`𝖮`$ is a closed $`G`$-conjugacy class in $`𝔰𝔩_n`$, and it has been shown in \[W\] that
$$\mu ^1(𝖮)=\{(X,Y)𝔰𝔩_n\times 𝔰𝔩_n|[X,Y]+\mathrm{𝙸𝚍}\text{is a rank one semisimple matrix}\},$$
is a smooth connected algebraic variety, and the $`\mathrm{𝙰𝚍}G`$-diagonal action on $`\mu ^1(𝖮)`$ is free. The reduced space $`𝕄:=\mu ^1(𝖮)/G`$ is, according to \[KKS\] (see also \[W\]), nothing but the phase space of the (rational) Calogero-Moser integrable system. This is a smooth affine algebraic symplectic manifold. Thus, Theorem 1.2 makes $`𝕄`$ a coadjoint orbit in $`(A/[A,A])^{}`$, where $`A=TE=x,y`$. This very special case was the starting point of our analysis.∎
An earlier version of this paper has been greatly motivated by \[BW\], whose question led me to the development of non-commutative geometry in the special case of the Calogero-Moser space. The results presented in §3 below form a natural generalization of the Calogero-Moser case. This generalization has been found simultaneously and independently by L. Le Bruyn \[LB1\] and the author.
Acknowledgements. I am grateful to Yu. Berest and G. Wilson for explaining to me the results of \[BW\] prior to their publication. I have benefited from interesting correspondence with L. Le Bruyn, especially from his letter \[LB1\].
## 2 Non-commutative Symplectic geometry
Throughout this paper we will be working over a ground field $`\mathrm{𝕜}`$ of characteristic zero, and write $`=_\mathrm{𝕜}`$. We fix a commutative unital $`\mathrm{𝕜}`$-algebra $`B`$, and for any $`B`$-bimodule $`M`$, write $`T_B^jM=M___B\mathrm{}___BM`$ ($`j`$ factors $`M`$), which is a $`B`$-bimodule again.
Let $`A`$ be a unital associative $`\mathrm{𝕜}`$-algebra containing the commutative algebra $`B`$ as a subalgebra. Recall that the free differential envelope of $`A`$ over $`B`$ is a graded vector space $`\mathrm{\Omega }__B^{}A=_{j0}\mathrm{\Omega }__B^jA,`$ where $`\mathrm{\Omega }__B^jA=A_BT_B^j(A/B)`$ is the $`B`$-bimodule formed by linear combinations of expressions $`a_0da_1\mathrm{}da_jAT_B^j(A/B)`$. Moreover, it is known, cf. \[L\], that there is a $`B`$-bimodule isomorphism: $`\mathrm{\Omega }__B^{}A_{j0_{}}T_B^j(\mathrm{\Omega }__B^1A),`$ and there is a $`B`$-bimodule super-differential $`d:\mathrm{\Omega }__B^{}A\mathrm{\Omega }__B^{+1}A,`$ making $`\mathrm{\Omega }__B^{}A`$ an associative differential graded algebra.
Given $`\alpha \mathrm{\Omega }__B^iA,\beta \mathrm{\Omega }__B^jA,`$ we put: $`[\alpha ,\beta ]=\alpha \beta (1)^{ij}\beta \alpha ,`$ and write $`[\mathrm{\Omega }__B^{}A,\mathrm{\Omega }__B^{}A]`$ for the $`B`$-linear span of all such super-commutators. Following Karoubi \[Ka\], see also \[L, §2.6\], define the relative non-commutative de Rham complex of the pair $`(A,B)`$ as the differential graded vector space:
$$\mathrm{𝖣𝖱}__B^{}A=\mathrm{\Omega }__B^{}A/[\mathrm{\Omega }__B^{}A,\mathrm{\Omega }__B^{}A],\mathrm{𝖣𝖱}__B^{}A=_{j0}\mathrm{𝖣𝖱}__B^jA,$$
where the differential and the grading are induced from those on $`\mathrm{\Omega }__B^{}A`$. Abusing the notation we will write: $`a_0da_1\mathrm{}da_j\mathrm{𝖣𝖱}__B^jA`$, meaning the corresponding class modulo commutators. We have: $`\mathrm{𝖣𝖱}__B^0A=A/[A,A],`$ and $`H^0(\mathrm{𝖣𝖱}__B^{}A)=\mathrm{ker}(\mathrm{𝖣𝖱}__B^0A\mathrm{𝖣𝖱}__B^1A)=B.`$
Let $`\mathrm{𝙳𝚎𝚛}_{}^{}{}_{_B}{}^{}A`$ denote the Lie algebra of all $`B`$-linear derivations of $`A`$. Given $`\theta \mathrm{𝙳𝚎𝚛}_{}^{}{}_{_B}{}^{}A`$ one introduces, following \[K2\], a Lie operator $`L_\theta :\mathrm{\Omega }__B^{}A\mathrm{\Omega }__B^{}A`$, resp. a contraction operator $`i_\theta :\mathrm{\Omega }__B^{}A\mathrm{\Omega }__B^1A`$, as a derivation, resp. a super-derivation, of the associative algebra $`\mathrm{\Omega }__B^{}A`$ defined on generators by the formulas:
$$L_\theta (a_0)=\theta (a_0),L_\theta (da)=d(\theta (a))\text{and}i_\theta (a_0)=0,i_\theta (da)=\theta (a),a_0,aA.$$
It is straightforward to verify that the induced operators on $`\mathrm{𝖣𝖱}__B^{}A`$, satisfy the following standard commutation relations:
$$L_\theta =i_\theta {}_{^{}}{}^{}d+d{}_{^{}}{}^{}i_{\theta }^{},[L_\theta ,i_\gamma ]=i_{[\theta ,\gamma ]},[L_\theta ,L_\gamma ]=L_{[\theta ,\gamma ]},i_\theta {}_{^{}}{}^{}i_{\gamma }^{}=i_\gamma {}_{^{}}{}^{}i_{\theta }^{},$$
(2.1)
where all the commutation relations but the last one hold already in $`\mathrm{\Omega }__B^{}A`$.
Fix $`\omega \mathrm{𝖣𝖱}__B^2A`$, and set $`\mathrm{𝙳𝚎𝚛}__B(A,\omega )=\{\theta \mathrm{𝙳𝚎𝚛}^{}A|L_\theta \omega =0\}.`$ Clearly, $`\mathrm{𝙳𝚎𝚛}__B(A,\omega )`$ is a Lie subalgebra in $`\mathrm{𝙳𝚎𝚛}_{}^{}{}_{_B}{}^{}A`$. The assignment: $`\theta i_\theta \omega `$ gives a linear map $`i:\mathrm{𝙳𝚎𝚛}_{}^{}{}_{_B}{}^{}A\mathrm{𝖣𝖱}__B^1A`$. The 2-form $`\omega \mathrm{𝖣𝖱}__B^2A`$ is called non-degenerate provided the map $`i`$ is bijective.
###### Lemma 2.2.
Let $`\omega \mathrm{𝖣𝖱}__B^2A`$ be a non-degenerate 2-form such that $`d\omega =0`$ in $`\mathrm{𝖣𝖱}__B^3A`$. Then the map: $`\theta i_\theta \omega `$ induces a bijection $`i:\mathrm{𝙳𝚎𝚛}__B(A,\omega )\stackrel{}{}(\mathrm{𝖣𝖱}__B^1A)_{_{\mathrm{𝖼𝗅𝗈𝗌𝖾𝖽}}},`$ that is: $`\theta \mathrm{𝙳𝚎𝚛}__B(A,\omega )d(i_\theta \omega )=0`$ in $`\mathrm{𝖣𝖱}__B^2A`$.
Proof. Since, $`d\omega =0`$, we have: $`L_\theta \omega =i_\theta d\omega +di_\theta \omega =di_\theta \omega .`$ Hence, $`\theta \mathrm{𝙳𝚎𝚛}__B(A,\omega )\mathrm{\hspace{0.33em}0}=L_\theta \omega =d(i_\theta \omega ).`$
By Lemma 2.2, one may invert the isomorphism $`i`$ to obtain a linear bijection $`i^1:(\mathrm{𝖣𝖱}__B^1A)_{_{\mathrm{𝖼𝗅𝗈𝗌𝖾𝖽}}}\stackrel{}{}`$ $`\mathrm{𝙳𝚎𝚛}__B(A,\omega )`$. Let: $`f\theta _f`$ denote the map given by the composition:
$$A/[A,A]=\mathrm{𝖣𝖱}__B^0A\stackrel{d}{}(\mathrm{𝖣𝖱}__B^1A)_{_{\mathrm{𝖾𝗑𝖺𝖼𝗍}}}(\mathrm{𝖣𝖱}__B^1A)_{_{\mathrm{𝖼𝗅𝗈𝗌𝖾𝖽}}}\stackrel{i^1}{}\mathrm{𝙳𝚎𝚛}__B(A,\omega ).$$
(2.3)
Using the map: $`f\theta _f`$, we define a Poisson bracket on $`A/[A,A]`$ by any of the following equivalent expressions:
$$\{f,g\}__\omega :=i_{_{\theta _f}}(i_{_{\theta _g}}\omega )=i_{_{\theta _f}}(dg)=i_{_{\theta _g}}(df)=L_{_{\theta _f}}g=L_{_{\theta _g}}f.$$
(2.4)
Here, in the first expression for $`\{f,g\}__\omega `$ we have used the composite map: $`i_{_{\theta _f}}{}_{^{}}{}^{}i_{_{\theta _g}}^{}:\mathrm{𝖣𝖱}__B^2A\mathrm{𝖣𝖱}__B^1A\mathrm{𝖣𝖱}__B^0A.`$ Other equalities, e.g.: $`i_{_{\theta _f}}(i_{_{\theta _g}}\omega )=L_{_{\theta _f}}g`$, follow from the equation $`i_{_{\theta _g}}\omega =dg`$ (which is the definition of $`\theta _g`$), the obvious identity: $`i_{_{\theta _f}}(dg)=L_{_{\theta _f}}g`$, and the last equation in (2.1).
###### Theorem 2.5.
The bracket (2.4) makes $`A/[A,A]`$ into a Lie algebra.
Proof. Skew symmetry of the bracket is immediate from (2.4). We begin the proof of the Jacobi identity by observing that, for any $`fA/[A,A]`$ and $`\eta \mathrm{𝙳𝚎𝚛}__BA`$ one has:
$$i_\eta i_{\theta _f}\omega =L_\eta f.$$
(2.6)
Further, for any $`\xi ,\eta __1,\eta __2\mathrm{𝙳𝚎𝚛}__BA`$, commutation relations (2.1) yield:
$$L_\xi i_{\eta __1}i_{\eta __2}\omega i_{\eta __1}i_{\eta __2}L_\xi \omega =i_{[\xi ,\eta __1]}i_{\eta __2}\omega +i_{\eta __1}i_{[\xi ,\eta __2]}\omega .$$
In the special case: $`\xi =\theta _f`$, we have: $`L_\xi \omega =0`$, hence, the above equation reads:
$$L_{\theta _f}i_{\eta __1}i_{\eta __2}\omega =i_{[\theta _f,\eta __1]}i_{\eta __2}\omega +i_{\eta __1}i_{[\theta _f,\eta __2]}\omega .$$
(2.7)
We now choose $`gA/[A,A]`$ and put $`\eta __1=\theta _g`$. Formula (2.6) shows that the LHS of (2.7) equals: $`L_{\theta _f}i_{\eta __2}i_{\theta _g}\omega =L_{\theta _f}L_{\eta __2}g.`$ Similarly, the second summand on the RHS of (2.7) equals:
$$i_{[\theta _f,\eta __2]}i_{\theta _g}\omega =L_{[\theta _f,\eta __2]}g=L_{\eta __2}L_{\theta _f}gL_{\theta _f}L_{\eta __2}g.$$
Thus, writing $`\eta =\eta __2`$, from (2.7) we deduce:
$$i_{[\theta _f,\theta _g]}(i_\eta \omega )=L_\eta L_{\theta _f}g=L_\eta \{f,g\}=i_{\theta _{\{f,g\}}}i_\eta \omega ,$$
where the last equality is due to formula (2.6). We conclude that the derivation $`\delta :=[\theta _f,\theta _g]\theta _{\{f,g\}}\mathrm{𝙳𝚎𝚛}__B(A,\omega )`$ has the property that, for any $`\eta \mathrm{𝙳𝚎𝚛}__BA`$, one has: $`i_\delta i_\eta \omega =0.`$ The bijection: $`\eta \alpha =i_\eta \omega `$ of Lemma 2.2 now implies that, for any 1-form $`\alpha `$, one has: $`i_\delta \alpha =0`$.
To complete the proof, for any $`f,g,h\mathrm{𝖣𝖱}__B^0A,`$ we write the identity:
$$(L_{\theta _f}{}_{^{}}{}^{}L_{\theta _g}^{}L_{\theta _g}{}_{^{}}{}^{}L_{\theta _f}^{})h=L_{[\theta _f,\theta _g]}h=(L_{\theta _{\{f,g\}}+\delta })h=L_{\theta _{\{f,g\}}}h+L_\delta h.$$
The leftmost commutator here equals: $`\{f,\{g,h\}\}\{g,\{f,h\}\},`$ and the term $`L_{\theta _{\{f,g\}}}h`$ on the right equals $`\{\{f,g\},h\}`$, by definition, see (2.4). Finally, we have: $`L_\delta h=i_\delta (dh)=0,`$ because of the property of $`\delta `$ established earlier. Thus, the identity above yields: $`\{f,\{g,h\}\}\{g,\{f,h\}\}=\{\{f,g\},h\},`$ and the Theorem is proved. ∎
Assume that $`B=\mathrm{𝕜}\mathrm{𝕜}\mathrm{}\mathrm{𝕜}`$ (direct sum of $`p`$ copies of the ground field). For each $`i\{1,\mathrm{},p\},`$ let $`\mathrm{𝟏}_iB`$ denote the idempotent corresponding to the $`i`$-th direct summand $`\mathrm{𝕜}`$.
Further, let $`V`$ be a finite dimensional left $`B`$-module. Clearly, giving such a $`V`$ amounts to giving a collection of finite dimensional $`\mathrm{𝕜}`$-vector space $`\{V_i\}_{\mathrm{\hspace{0.17em}1}ip}`$, one for each $`i`$, such that $`V=_iV_i,`$ and such that $`\mathrm{𝟏}_iB`$ acts as the projector onto the $`i`$-th direct summand. We consider the algebra $`\mathrm{𝙴𝚗𝚍}V:=\mathrm{𝙴𝚗𝚍}_\mathrm{𝕜}V`$ of $`\mathrm{𝕜}`$-linear endomorphisms of $`V`$. The action of $`B`$ on $`V`$ makes $`V^{}:=\mathrm{𝙷𝚘𝚖}_\mathrm{𝕜}(V,\mathrm{𝕜})`$ a right $`B`$-module, and gives an algebra imbedding: $`B\mathrm{𝙴𝚗𝚍}V`$. Hence, left and right multiplication by $`B`$ make $`\mathrm{𝙴𝚗𝚍}V`$ a $`B`$-bimodule which is canonically isomorphic to the $`B`$-bimodule $`V__\mathrm{𝕜}V^{}`$. Further, the assignment:
$$f(\mathrm{𝗍𝗋}(\mathrm{𝟏}_1f\mathrm{𝟏}_1),\mathrm{𝗍𝗋}(\mathrm{𝟏}_2f\mathrm{𝟏}_2),\mathrm{},\mathrm{𝗍𝗋}(\mathrm{𝟏}_pf\mathrm{𝟏}_p))\mathrm{𝕜}\mathrm{𝕜}\mathrm{}\mathrm{𝕜}=B$$
gives a canonical $`B`$-bimodule trace map $`\mathrm{𝗍𝗋}:\mathrm{𝙴𝚗𝚍}VB`$.
Representation functor. Given a finitely generated associative $`B`$-algebra $`A`$, let $`\mathrm{𝙷𝚘𝚖}_{_{B\text{-}\mathrm{𝖺𝗅𝗀}}}(A,\mathrm{𝙴𝚗𝚍}V)`$ denote the affine algebraic variety of all associative algebra homomorphisms $`\rho :A\mathrm{𝙴𝚗𝚍}V`$, such that $`\rho |_B=\mathrm{𝙸𝚍}_B`$. Let $`\mathrm{𝖱𝖾𝗉}(A,V):=\mathrm{𝕜}[\mathrm{𝙷𝚘𝚖}_{_{B\text{-}\mathrm{𝖺𝗅𝗀}}}(A,\mathrm{𝙴𝚗𝚍}V)]`$ denote the coordinate ring of $`\mathrm{𝙷𝚘𝚖}_{_{B\text{-}\mathrm{𝖺𝗅𝗀}}}(A,\mathrm{𝙴𝚗𝚍}V)`$. The natural action on $`\mathrm{𝙴𝚗𝚍}V`$ of the group $`G(V)=GL__B(V)`$ (of $`B`$-linear automorphisms of $`V`$) by conjugation induces a $`G(V)`$-action on $`\mathrm{𝙷𝚘𝚖}_{_{B\text{-}\mathrm{𝖺𝗅𝗀}}}(A,\mathrm{𝙴𝚗𝚍}V)`$. This gives a $`G(V)`$-action on $`\mathrm{𝖱𝖾𝗉}(A,V)`$ by algebra automorphisms.
The tautological evaluation map: $`A\times \mathrm{𝙷𝚘𝚖}_{_{B\text{-}\mathrm{𝖺𝗅𝗀}}}(A,\mathrm{𝙴𝚗𝚍}V)\mathrm{𝙴𝚗𝚍}V`$ assigns to any element $`aA`$ an $`\mathrm{𝙴𝚗𝚍}V`$-valued function $`\widehat{a}`$ on $`\mathrm{𝙷𝚘𝚖}_{_{B\text{-}\mathrm{𝖺𝗅𝗀}}}(A,\mathrm{𝙴𝚗𝚍}V)`$. Equivalently, this function may be viewed as an element $`\widehat{a}\left(\mathrm{𝖱𝖾𝗉}(A,V)__B\mathrm{𝙴𝚗𝚍}V\right)^{G(V)}`$. Taking the trace on the second tensor factor, one obtains a $`G(V)`$-invariant $`\mathrm{𝕜}`$-valued function $`\mathrm{𝗍𝗋}(\widehat{a})\left(\mathrm{𝖱𝖾𝗉}(A,V)__BB\right)^{G(V)}=\mathrm{𝖱𝖾𝗉}(A,V)^{G(V)}.`$ The assignment: $`a\mathrm{𝗍𝗋}(\widehat{a})`$ clearly vanishes on $`[A,A]`$ due to the cyclic symmetry of the trace map. Thus, it descends to a well-defined $`B`$-linear map
$$\widehat{\mathrm{𝗍𝗋}}:\mathrm{𝖣𝖱}__B^0A=A/[A,A]\mathrm{𝖱𝖾𝗉}(A,V)^{G(V)},a\mathrm{𝗍𝗋}(\widehat{a}).$$
(2.8)
Remark. More generally, for any $`p0`$, the assignment: $`a_0da_1\mathrm{}da_p\mathrm{𝗍𝗋}(\widehat{a}_0d\widehat{a}_1\mathrm{}d\widehat{a}_p)`$ gives a well-defined map from $`\mathrm{𝖣𝖱}__B^pA`$ to the space of $`G(V)`$-invariant regular $`p`$-forms (in the ordinary sense) on the algebraic variety $`\mathrm{𝙷𝚘𝚖}_{_{B\text{-}\mathrm{𝖺𝗅𝗀}}}(A,\mathrm{𝙴𝚗𝚍}V)`$.
## 3 Lie algebra associated to a Quiver
Fix $`B`$, a commutative $`\mathrm{𝕜}`$-algebra and $`E`$, a finite rank projective $`B`$-bimodule, i.e. a projective $`BB^{op}`$-module. The space $`E^{}:=\mathrm{𝙷𝚘𝚖}_{_{\mathrm{𝗅𝖾𝖿𝗍}B\text{-}\mathrm{𝗆𝗈𝖽}}}(E,B)`$ has a canonical $`B`$-bimodule structure given by: $`(b_1\phi b_2)(e)=\phi (eb_1)b_2,`$ where $`b_1,b_2B,eE,`$ and $`\phi E^{}`$.
A $`B`$-bimodule map $`\omega :E___BEB`$ will be referred to as a $`B`$-bilinear form on $`E`$. For such an $`\omega `$, the assignment: $`e\omega (e)`$ gives a $`B`$-bimodule map $`EE^{}`$. We call $`\omega `$ non-degenerate if the latter map is an isomorphism. If, furthermore, $`\omega `$ is skew-symmetric, i.e. $`\omega (x,y)+\omega (y,x)=0,`$ for any $`x,yE`$, we will say that $`\omega `$ is a symplectic $`B`$-form on $`E`$. For example, for any finite rank projective $`B`$-bimodule $`V`$, the bimodule $`E=VV^{}`$ carries a canonical symplectic $`B`$-form.
Fix a finite dimensional $`B`$-bimodule $`E`$, and let $`A=T__BE:=_{i0}T__B^iE`$ be the tensor algebra, a graded associative algebra such that $`T__B^0E=B`$. For each $`i>0`$, let $`(T__B^iE)_{_{\mathrm{𝖼𝗒𝖼𝗅𝗂𝖼}}}`$ denote the quotient of $`T__B^iE`$ by the $`B`$-sub-bimodule generated by the elements:
$$x_1___Bx_2___B\mathrm{}___Bx_ix_i___Bx_1___B\mathrm{}___Bx_{i1},x_1,\mathrm{},x_iE.$$
The following result was obtained independently by L. Le Bruyn \[LB1\] and the author.
###### Lemma 3.1.
$`(𝗂)`$The de Rham complex of $`A=T__BE`$ is acyclic, i.e., $`H^k(\mathrm{𝖣𝖱}__B^{}A)=0`$, for all $`k1`$. Furthermore, $`H^0(\mathrm{𝖣𝖱}__B^{}A)=B`$.
$`(\mathrm{𝗂𝗂})`$We have: $`\mathrm{𝖣𝖱}__B^0(T__BE)=(T__BE)_{_{\mathrm{𝖼𝗒𝖼𝗅𝗂𝖼}}},`$ and $`\mathrm{𝖣𝖱}__B^1(T__BE)=(T__BE)_BE.`$
Proof. To prove (i), we imitate, following Kontsevich \[K2\], the classical proof of the Poincaré lemma. To this end, introduce a ($`B`$-linear) Euler derivation $`\mathrm{𝖾𝗎}:T__BET__BE`$ by letting it act on generators $`xE=T__B^1E`$ by: $`\mathrm{𝖾𝗎}(x)=x`$. The induced map $`L_{\mathrm{𝖾𝗎}}:\mathrm{𝖣𝖱}__B^{}A\mathrm{𝖣𝖱}__B^{}A`$ is diagonalizable and has non-negative integral eigenvalues. Cartan’s homotopy formula: $`L_{\mathrm{𝖾𝗎}}=d{}_{^{}}{}^{}i_{\mathrm{𝖾𝗎}}^{}+i_{\mathrm{𝖾𝗎}}{}_{^{}}{}^{}d`$ shows that the de Rham complex is quasi-isomorphic to the zero eigen-space of the operator $`L_{\mathrm{𝖾𝗎}}`$, which is the subspace $`B`$ sitting in degree 0. Part (i) follows. Part (ii) is straightforward.∎
From now until the end of the section assume that $`B=\mathrm{𝕜}^I,`$ where $`I`$ is a finite set, and put $`A:=T__BE`$, where $`(E,\omega )`$ is a symplectic $`B`$-bimodule. Using the isomorphism: $`E\stackrel{}{}E^{}=\mathrm{𝙷𝚘𝚖}_{_{\mathrm{𝗅𝖾𝖿𝗍}B\text{-}\mathrm{𝗆𝗈𝖽}}}(E,B),`$ provided by $`\omega `$, one transports the symplectic structure from $`E`$ to $`E^{}`$. Let $`\omega ^{}=_r\varphi _r\psi _rEE`$ be the resulting symplectic $`B`$-form on $`E^{}`$. It is straightforward to see that $`_rd\varphi _rd\psi _r\mathrm{\Omega }__B^2A`$ gives a well-defined closed and non-degenerate class in $`\mathrm{𝖣𝖱}__B^2A`$, to be dented $`\omega _{_{\mathrm{𝖣𝖱}}}`$. Thus, the general construction (2.4) yields a Lie bracket $`\{,\}_{_{\omega _{_{\mathrm{𝖣𝖱}}}}}`$ on $`A/[A,A]`$.
Example. For each $`iI,`$ let $`\mathrm{𝟏}_iB=\mathrm{𝕜}^I`$ denote the idempotent corresponding to the $`i`$-th direct summand. Clearly, giving a finite rank $`B`$-bimodule amounts to giving a finite dimensional $`\mathrm{𝕜}`$-vector space $`E`$ equipped with a direct sum decomposition: $`E=_{i,jI}E_{i,j},`$ where $`E_{i,j}=\mathrm{𝟏}_iE\mathrm{𝟏}_j`$. Thus, one may think of the data $`(B,E)`$ as an oriented graph with vertex set $`I`$ and with $`dimE_{i,j}`$ edges going from the vertex $`i`$ to the vertex $`j`$.
Conversely, let $`Q`$ denote an oriented quiver with vertex set $`I`$. Set $`B=\mathrm{𝕜}^I`$, and let $`E_Q`$ be the $`\mathrm{𝕜}`$-vector space with basis formed by the set of edges $`\{aQ\}.`$ Then $`E_Q`$ has an obvious $`B`$-bimodule structure, and $`T__B(E_Q)`$ is known as the path algebra of $`Q`$. Further, the $`B`$-bimodule $`E_{\overline{Q}}`$ associated with $`\overline{Q}`$, the double of $`Q`$, has a natural symplectic $`B`$-form. The corresponding class in $`\mathrm{𝖣𝖱}__B^2\left(T__B(E_Q)\right)`$ is given by the formula: $`\omega _{_{\mathrm{𝖣𝖱}}}=_{aQ}dada^{}.`$
In the special case $`B=\mathrm{𝕜}`$, the Lie bracket $`\{,\}_{_{\omega _{_{\mathrm{𝖣𝖱}}}}}`$ on $`A/[A,A]`$ has been introduced by Kontsevich \[K2\] in a somewhat different way as follows. Let $`x_1,\mathrm{},x_n,y_1,\mathrm{},y_n`$ be a symplectic basis of the vector space $`E`$, i.e. a $`\mathrm{𝕜}`$-basis such that: $`\omega (x_i,y_j)=\delta _{ij},`$ and $`\omega (x_i,x_j)=\omega (y_i,y_j)=0`$. By Lemma 3.1$`(\mathrm{𝗂𝗂})`$, one has: $`\mathrm{𝖣𝖱}_\mathrm{𝕜}^1AAE`$. Kontsevich exploits this isomorphism to write any 1-form $`\alpha \mathrm{𝖣𝖱}_\mathrm{𝕜}^1A`$ in the form: $`\alpha =_{i=1}^nF_{x_i}(\alpha )x_i+_{j=1}^nF_{y_j}(\alpha )y_j`$, for certain uniquely determined elements $`F_{x_i}(\alpha ),F_{y_j}(\alpha )A`$. He then introduces , for any $`i=1,\mathrm{}n,`$ the following $`\mathrm{𝕜}`$-linear maps:
$$\frac{}{x_i},\frac{}{y_i}:\mathrm{𝖣𝖱}_\mathrm{𝕜}^0AA,\text{given by}\frac{f}{x_i}:=F_{x_i}(df)\text{and}\frac{f}{y_i}:=F_{y_i}(df).$$
Using these maps Kontsevich defines (put another way: gives a coordinate expression for) the Lie bracket $`\{,\}_\omega `$ by the familiar formula:
$$\{f,g\}__\omega :=\underset{i=1}{\overset{n}{}}\left(\frac{f}{x_i}\frac{g}{y_i}\frac{f}{y_i}\frac{g}{x_i}\right)\text{mod }[A,A]A/[A,A]=\mathrm{𝖣𝖱}_\mathrm{𝕜}^0A,$$
(3.2)
where ”dot” stands for the product in $`A`$. We leave to the reader to check that formulas (2.4), and (1.4) give rise to the same bracket on $`\mathrm{𝖣𝖱}_\mathrm{𝕜}^0A=A/[A,A]`$ as formula (3.2)
Remarks.$`(𝗂)`$In the general case of an arbitrary quiver $`Q`$, the analogue of Kontsevich’s formula (3.2) for the Poisson bracket associated with the corresponding algebra $`A=T_B(E_{\overline{Q}})`$, in obvious notation, cf. (1.1), is:
$$\{f,g\}__\omega =_{aQ}\left(\frac{f}{a}\frac{g}{a^{}}\frac{f}{a^{}}\frac{g}{a}\right)\text{mod }[A,A]A/[A,A]=\mathrm{𝖣𝖱}__B^0A.$$
(3.3)
$`(\mathrm{𝗂𝗂})`$The Poisson structure given by (3.3) admits a natural quantization, in which the symplectic manifold $`(\overline{Q},V)=T^{}(Q,V)`$ gets replaced by the algebra $`𝒟((Q,V))`$ of polynomial differential operators on the vector space $`(Q,V)`$. This quantization has been found by M. Holland \[Ho\], even before the non-commutative Poisson structure given by (2.4) and (3.3) has been discovered.∎
The next Proposition gives a non-commutative analogue of the classical Lie algebra exact sequence:
$$0\text{constant functions}\text{regular functions}\text{symplectic vector fields}\mathrm{\hspace{0.33em}0},$$
associated with a connected and simply-connected symplectic manifold.
###### Proposition 3.4.
There is a natural Lie algebra central extension:
$$0BA/[A,A]\mathrm{𝙳𝚎𝚛}__B(A,\omega )\mathrm{\hspace{0.17em}\hspace{0.17em}0}.$$
Proof. It is immediate from formula (2.3) that for the map: $`f\theta _f`$ we have: $`\mathrm{𝙺𝚎𝚛}\{A/[A,A]\mathrm{𝙳𝚎𝚛}__B(A,\omega )\}`$ $`=\mathrm{𝙺𝚎𝚛}d`$. By Lemma 3.1$`(𝗂)`$we get: $`\mathrm{𝙺𝚎𝚛}d=B`$. Further, Lemma 3.1$`(𝗂)`$insures that every closed element in $`\mathrm{𝖣𝖱}__B^1A`$ is exact. This yields surjectivity of the map: $`A/[A,A]\mathrm{𝙳𝚎𝚛}__B(A,\omega ).`$
It remains to show that the map: $`f\theta _f`$ is a Lie algebra homomorphism. To this end, we use the notation: $`\delta =[\theta _f,\theta _g]\theta _{\{f,g\}}\mathrm{𝙳𝚎𝚛}__B(A,\omega )`$ from the proof of Theorem 2.5. The proof of Theorem 2.5 implies that the derivation $`\delta `$ has the property that, for any $`\eta \mathrm{𝙳𝚎𝚛}__BA`$, one has: $`i_\theta i_\eta \omega =0.`$ Equivalently, setting $`\alpha :=i_\delta \omega `$, we get: $`i_\eta \alpha =0,\eta \mathrm{𝙳𝚎𝚛}__BA`$. But the isomorphism: $`\mathrm{𝖣𝖱}__B^1(T__BE)=(T__BE)_BE`$ of Lemma 3.1(ii) clearly forces such a 1-form $`\alpha `$ to vanish. The 2-form $`\omega `$ being non-degenerate, it follows that $`\delta =0`$. Thus, we have shown that $`[\theta _f,\theta _g]=\theta _{\{f,g\}}.`$ This completes the proof. ∎
Representations. We now fix a finite dimensional left $`B`$-module $`V`$, as at the end of §2. Observe that if $`Q`$ is a quiver, and $`E=E_{\overline{Q}}`$ is the symplectic $`B`$-bimodule attached, as has been explained earlier, to the double of $`Q`$, then for $`A=T__B(E_{\overline{Q}})`$, in the notation of the Introduction we have: $`\mathrm{𝙷𝚘𝚖}_{_{B\text{-}\mathrm{𝖺𝗅𝗀}}}(A,\mathrm{𝙴𝚗𝚍}V)=(\overline{Q},V)`$. In general, let $`E`$ be a finite dimensional symplectic $`B`$-bimodule. Then, for $`A=T__BE`$, one has: $`\mathrm{𝙷𝚘𝚖}_{_{B\text{-}\mathrm{𝖺𝗅𝗀}}}(A,\mathrm{𝙴𝚗𝚍}V)=\mathrm{𝙷𝚘𝚖}_{_{B\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}}(E,\mathrm{𝙴𝚗𝚍}V).`$ The latter space can be naturally identified with $`E^{}__B\mathrm{𝙴𝚗𝚍}V.`$ Note that we have the symplectic $`B`$-form $`\omega ^{}`$ on $`E^{}`$, and a non-degenerate symmetric bilinear form $`\mathrm{𝗍𝗋}:\mathrm{𝙴𝚗𝚍}V__B\mathrm{𝙴𝚗𝚍}VB,`$ given by: $`(F_1,F_2)\mathrm{𝗍𝗋}(F_1{}_{^{}}{}^{}F_{2}^{}).`$ By standard Linear Algebra, the tensor product of a skew-symmetric and symmetric non-degenerate forms gives the skew-symmetric non-degenerate bilinear form: $`\omega _{_{\mathrm{𝖱𝖾𝗉}}}:=\omega ^{}\mathrm{𝗍𝗋}.`$ The 2-form $`\omega _{_{\mathrm{𝖱𝖾𝗉}}}`$ makes $`E^{}__B\mathrm{𝙴𝚗𝚍}V`$, hence, $`\mathrm{𝙷𝚘𝚖}_{_{B\text{-}\mathrm{𝖺𝗅𝗀}}}(A,\mathrm{𝙴𝚗𝚍}V)`$, a symplectic $`B`$-bimodule, therefore gives rise to a $`G(V)`$-invariant Poisson bracket $`\{,\}_{\omega _{_{\mathrm{𝖱𝖾𝗉}}}}`$ on the coordinate ring $`\mathrm{𝕜}[\mathrm{𝙷𝚘𝚖}_{_{B\text{-}\mathrm{𝖺𝗅𝗀}}}(A,\mathrm{𝙴𝚗𝚍}V)]=\mathrm{𝖱𝖾𝗉}(A,V)`$. The invariants, $`\mathrm{𝖱𝖾𝗉}(A,V)^{G(V)},`$ clearly form a Poisson subalgebra in $`\mathrm{𝖱𝖾𝗉}(A,V)`$, and we have:
###### Proposition 3.5.
The map $`\widehat{\mathrm{𝗍𝗋}}:A/[A,A]\mathrm{𝖱𝖾𝗉}(A,V)^{G(V)}`$ defined in (2.8) is a Lie algebra homomorphism, that is, for any $`f,gA`$, one has:
$$\{\widehat{\mathrm{𝗍𝗋}}f,\widehat{\mathrm{𝗍𝗋}}g\}_{_{\omega _{_{\mathrm{𝖱𝖾𝗉}}}}}=\widehat{\mathrm{𝗍𝗋}}(\{f,g\}_{\omega _{_{\mathrm{𝖣𝖱}}}}).$$
Proof. Straightforward calculation for $`f,g`$ taken to be non-commutative monomials. $`\mathrm{}`$
We can now complete the proof of Theorem 1.2. As we have mentioned in the Introduction, the $`G(V)`$-action on $`\mathrm{𝙷𝚘𝚖}_{_{B\text{-}\mathrm{𝖺𝗅𝗀}}}(A,\mathrm{𝙴𝚗𝚍}V)=E^{}__B\mathrm{𝙴𝚗𝚍}V`$ turns out to be Hamiltonian, and the corresponding moment map $`\mu :E^{}__B\mathrm{𝙴𝚗𝚍}V𝔤(V)^{}=𝔰𝔤(V)`$, cf. (1.1), is given by the following formula:
$$_i\varphi _iF_i_{j<k}\omega ^{}(\varphi _j,\varphi _k)[F_j,F_k]𝔰𝔤(V),\varphi _iE^{},F_i\mathrm{𝙴𝚗𝚍}V.$$
(3.6)
Fix $`𝖮𝔰𝔤(V)`$, a closed $`\mathrm{𝙰𝚍}G(V)`$-orbit, and assume that the group $`G(V)`$ acts freely on the subvariety $`\mu ^1(𝖮)E^{}__B\mathrm{𝙴𝚗𝚍}V`$. Then, the orbit space $`\mu ^1(𝖮)/G(V)`$ is a smooth affine subvariety in $`\mathrm{𝚂𝚙𝚎𝚌}(\mathrm{𝖱𝖾𝗉}(A,V)^{G(V)})`$.
###### Proposition 3.7.
The composite map:
$$\mu ^1(𝖮)/G(V)\mathrm{𝚂𝚙𝚎𝚌}\left(\mathrm{𝖱𝖾𝗉}(A,V)^{G(V)}\right)\stackrel{_{\mathrm{𝖾𝗏𝖺𝗅𝗎𝖺𝗍𝗂𝗈𝗇}}}{}\left(\mathrm{𝖱𝖾𝗉}(A,V)^{G(V)}\right)^{}\stackrel{\mathrm{𝗍𝗋}^{}}{}(A/[A,A])^{}$$
is injective and makes $`\mu ^1(𝖮)/G(V)`$ a coadjoint orbit in $`(A/[A,A])^{}`$.
Proof. Set $`X=\mu ^1(𝖮)/G(V)`$, a smooth affine variety. As we have argued in §1, proving the proposition amounts to showing that regular functions on $`\mathrm{𝚂𝚙𝚎𝚌}\left(\mathrm{𝖱𝖾𝗉}(A,V)^{G(V)}\right)`$ of the form $`\mathrm{𝗍𝗋}(\widehat{a}),aA/[A,A],`$ separate points and tangents of the variety $`X`$ $`\mathrm{𝚂𝚙𝚎𝚌}\left(\mathrm{𝖱𝖾𝗉}(A,V)^{G(V)}\right)`$. This is clearly true for the whole algebra $`\mathrm{𝖱𝖾𝗉}(A,V)^{G(V)}`$, since it is true for the algebra $`\mathrm{𝕜}[X]`$, and every regular function on $`X`$ is obtained from an element of $`\mathrm{𝖱𝖾𝗉}(A,V)^{G(V)}`$, by restriction.
We now use the result of Le Bruyn- Procesi \[LP\], saying that the algebra $`\mathrm{𝖱𝖾𝗉}(A,V)^{G(V)}`$ is generated by elements of the form:
$$\mathrm{𝗍𝗋}(\widehat{\mathrm{𝟏}}_i\widehat{x}_1\widehat{x}_2\mathrm{}\widehat{x}_k\widehat{\mathrm{𝟏}}_i),i=1,\mathrm{},p,x_jE,k1.$$
The expression above is nothing but $`\mathrm{𝗍𝗋}(\widehat{a}),`$ for $`a=\mathrm{𝟏}_i(x_1\mathrm{}x_k)\mathrm{𝟏}_iT__B^kEA.`$ It follows that, although the map $`\widehat{\mathrm{𝗍𝗋}}:A/[A,A]\mathrm{𝖱𝖾𝗉}(A,V)^{G(V)}`$ is not itself surjective, the algebra $`\mathrm{𝖱𝖾𝗉}(A,V)^{G(V)}`$ is generated by its image. Thus, elements of the image separate points and tangents of the variety $`X`$. ∎
## 4 Stabilization: infinite dimensional limit
We keep the setup of §2; in particular, we let $`B=\mathrm{𝕜}^I`$ and fix $`A`$, a finitely generated associative $`B`$-algebra. Any imbedding: $`VV^{}`$ of finite rank left $`B`$-modules induces a map: $`\mathrm{𝙷𝚘𝚖}_{_{B\text{-}\mathrm{𝖺𝗅𝗀}}}(A,\mathrm{𝙴𝚗𝚍}V)\mathrm{𝙷𝚘𝚖}_{_{B\text{-}\mathrm{𝖺𝗅𝗀}}}(A,\mathrm{𝙴𝚗𝚍}V^{})`$, which is a closed imbedding of affine algebraic varieties. The latter imbedding gives rise to the restriction homomorphism of coordinate rings
$$r_{_{V^{},V}}:\mathrm{𝖱𝖾𝗉}(A,V^{})^{G(V^{})}\mathrm{𝖱𝖾𝗉}(A,V)^{G(V)}.$$
(4.1)
Observe that the collection of all finite rank $`B`$-modules, $`V`$, forms a direct system with respect to $`B`$-module imbeddings, and we set $`V_{_{\mathrm{}}}:=\underset{}{lim}V,`$ and let $`G_{_{\mathrm{}}}:=\underset{}{lim}G(V)`$ be the corresponding ind-group. By definition we put: $`\mathrm{𝖱𝖾𝗉}(A,V_{_{\mathrm{}}})^G_{_{\mathrm{}}}:=\underset{}{lim}\mathrm{𝖱𝖾𝗉}(A,V)^{G(V)}.`$
There is a standard way to introduce a cocomutative coproduct $`\mathrm{\Delta }:\mathrm{𝖱𝖾𝗉}(A,V_{_{\mathrm{}}})^G_{_{\mathrm{}}}\mathrm{𝖱𝖾𝗉}(A,V_{_{\mathrm{}}})^G_{_{\mathrm{}}}__B\mathrm{𝖱𝖾𝗉}(A,V_{_{\mathrm{}}})^G_{_{\mathrm{}}}.`$ To see this, it is convenient to think of $`\mathrm{𝖱𝖾𝗉}(A,V_{_{\mathrm{}}})`$ as some sort of coordinate ring $`\mathrm{𝕜}[\mathrm{𝙷𝚘𝚖}_{_{B\text{-}\mathrm{𝖺𝗅𝗀}}}(A,\mathrm{𝙴𝚗𝚍}V_{_{\mathrm{}}})]`$. Then any choice of $`B`$-module isomorphism: $`V_{_{\mathrm{}}}\stackrel{\varpi }{}V_{_{\mathrm{}}}V_{_{\mathrm{}}}`$ gives a morphism of ind-schemes:
$$\mathrm{𝙷𝚘𝚖}_{_{B\text{-}\mathrm{𝖺𝗅𝗀}}}(A,\mathrm{𝙴𝚗𝚍}V_{_{\mathrm{}}})\times \mathrm{𝙷𝚘𝚖}_{_{B\text{-}\mathrm{𝖺𝗅𝗀}}}(A,\mathrm{𝙴𝚗𝚍}V_{_{\mathrm{}}})\mathrm{𝙷𝚘𝚖}_{_{B\text{-}\mathrm{𝖺𝗅𝗀}}}(A,\mathrm{𝙴𝚗𝚍}(V_{_{\mathrm{}}}V_{_{\mathrm{}}}))\stackrel{\varpi }{\stackrel{_{}}{}}\mathrm{𝙷𝚘𝚖}_{_{B\text{-}\mathrm{𝖺𝗅𝗀}}}(A,\mathrm{𝙴𝚗𝚍}V_{_{\mathrm{}}}).$$
The coproduct $`\mathrm{\Delta }`$ on $`\mathrm{𝖱𝖾𝗉}(A,V_{_{\mathrm{}}})^G_{_{\mathrm{}}}`$ is the one induced by the corresponding algebra map:
$$\mathrm{\Delta }:\mathrm{𝕜}[\mathrm{𝙷𝚘𝚖}_{_{B\text{-}\mathrm{𝖺𝗅𝗀}}}(A,\mathrm{𝙴𝚗𝚍}V_{_{\mathrm{}}})]\mathrm{𝕜}[\mathrm{𝙷𝚘𝚖}_{_{B\text{-}\mathrm{𝖺𝗅𝗀}}}(A,\mathrm{𝙴𝚗𝚍}V_{_{\mathrm{}}})]_B\mathrm{𝕜}[\mathrm{𝙷𝚘𝚖}_{_{B\text{-}\mathrm{𝖺𝗅𝗀}}}(A,\mathrm{𝙴𝚗𝚍}V_{_{\mathrm{}}})].$$
Let $`\mathrm{𝚙𝚛𝚒𝚖}\left(\mathrm{𝖱𝖾𝗉}(A,V_{_{\mathrm{}}})^G_{_{\mathrm{}}}\right)`$ denote the $`B`$-module of primitive elements in $`\mathrm{𝖱𝖾𝗉}(A,V_{_{\mathrm{}}})^G_{_{\mathrm{}}}`$, i.e., the elements $`f\mathrm{𝖱𝖾𝗉}(A,V_{_{\mathrm{}}})^G_{_{\mathrm{}}}`$ such that $`\mathrm{\Delta }(f)=f1+1f.`$
Observe further that the map $`\widehat{\mathrm{𝗍𝗋}}__V:A/[A,A]\mathrm{𝖱𝖾𝗉}(A,V)^{G(V)}`$ given by (2.8) is compatible with restriction morphisms $`r_{_{V^{},V}}`$, see (4.1), that is, for any imbedding $`VV^{}`$, one has a commutative triangle: $`r_{_{V^{},V}}{}_{^{}}{}^{}\widehat{\mathrm{𝗍𝗋}}_{_V^{}}^{}=\widehat{\mathrm{𝗍𝗋}}__V.`$ Therefore, the maps $`\{\widehat{\mathrm{𝗍𝗋}}__V\}`$ give rise to a well-defined limit map $`\widehat{\mathrm{𝗍𝗋}}_{_{\mathrm{}}}:A/[A,A]\mathrm{𝖱𝖾𝗉}(A,V_{_{\mathrm{}}})^G_{_{\mathrm{}}}.`$
We now specialize to the setup of §3 and assume that $`A=T__BE`$, for a certain finite rank projective $`B`$-bimodule $`E`$. The following result is, in a sense, dual to the well-known relationship, see \[LQ\], \[L\], between cyclic homology of an associative algebra $`A`$ and primitive homology of the Lie algebra $`𝔤𝔩_{\mathrm{}}(A)`$.
###### Proposition 4.2.
For $`A=T__BE`$, the map $`\mathrm{𝗍𝗋}_{_{\mathrm{}}}`$ sets up a bijection:
$$\mathrm{𝗍𝗋}_{_{\mathrm{}}}:A/(B+[A,A])\stackrel{}{}\mathrm{𝚙𝚛𝚒𝚖}(\mathrm{𝖱𝖾𝗉}(A,V_{_{\mathrm{}}})^G_{_{\mathrm{}}}).$$
Notice next that, for $`A=T__BE`$, we have: $`\mathrm{𝖱𝖾𝗉}(A,V_{_{\mathrm{}}})=\mathrm{𝕜}[\mathrm{𝙷𝚘𝚖}_{_{B\text{-}\mathrm{𝖻𝗂𝗆𝗈𝖽}}}(E,\mathrm{𝙴𝚗𝚍}V_{_{\mathrm{}}})],`$ is a polynomial algebra with a natural grading, that also induces a grading on $`\mathrm{𝖱𝖾𝗉}(A,V_{_{\mathrm{}}})^G_{_{\mathrm{}}}`$. Furthermore, the coproduct $`\mathrm{\Delta }`$ is compatible with the (graded) algebra structure, hence makes $`\mathrm{𝖱𝖾𝗉}(A,V_{_{\mathrm{}}})^G_{_{\mathrm{}}}`$ a commutative and cocomutative graded Hopf $`B`$-algebra. The structure theorem for commutative and cocommutative graded Hopf algebras implies that $`\mathrm{𝖱𝖾𝗉}(A,V_{_{\mathrm{}}})^G_{_{\mathrm{}}}`$ must be the symmetric algebra (over $`B`$) on the $`B`$-bimodule of its primitive elements. Therefore, Proposition 4.2 yields
###### Corollary 4.3.
For $`A=T__BE`$, the map $`\mathrm{𝗍𝗋}_{_{\mathrm{}}}`$ extends, by multiplicativity, to a graded isomorphism of Poisson algebras: $`\mathrm{𝚂𝚢𝚖}^{^{_{}}}\left(A/(B+[A,A])\right)\stackrel{}{}\mathrm{𝖱𝖾𝗉}(A,V_{_{\mathrm{}}})^G_{_{\mathrm{}}}.`$
Remark. It is interesting to note that, for any finite dimensional $`V`$ such that $`dimV>1`$, the variety $`\mathrm{𝚂𝚙𝚎𝚌}\left(\mathrm{𝖱𝖾𝗉}(A,V)^{G(V)}\right)`$ is quite complicated, e.g., in the Calogero-Moser case. Nonetheless, Corollary 4.3 says that the ‘limiting’ variety $`\mathrm{𝚂𝚙𝚎𝚌}\left(\mathrm{𝖱𝖾𝗉}(A,V_{_{\mathrm{}}})^G_{_{\mathrm{}}}\right)`$ is always a vector space.
Proof of Proposition 4.2. It is clear from definitions, that $`\mathrm{𝗍𝗋}_{_{\mathrm{}}}(f)\mathrm{𝖱𝖾𝗉}(A,V_{_{\mathrm{}}})^G_{_{\mathrm{}}}`$ is a primitive element, for any homogeneous element $`fA`$ such that $`\mathrm{deg}f>0`$. Furthermore, one verifies that any element not contained in the image of the map $`\mathrm{𝗍𝗋}_{_{\mathrm{}}}`$ cannot satisfy the equation $`\mathrm{\Delta }(f)=f1+1f`$, hence, is not primitive. Thus, the map $`\mathrm{𝗍𝗋}_{_{\mathrm{}}}`$ is surjective, and it suffices to prove it is injective.
In order to avoid complicated notation, we restrict ourselves to proving injectivity in the special case of the quiver $`Q`$ consisting of a single vertex and a single edge-loop, that is the Calogero-Moser quiver (the general case goes in a similar fashion with minor modifications). Thus, we assume that $`A=\mathrm{𝕜}x,y`$, and therefore, $`\mathrm{𝖱𝖾𝗉}(A,V_{_{\mathrm{}}})=\mathrm{𝕜}[𝔤𝔩_{\mathrm{}}𝔤𝔩_{\mathrm{}}],`$ where $`𝔤𝔩_{\mathrm{}}:=\underset{}{lim}𝔤𝔩_n(\mathrm{𝕜}).`$ We must show that, given $`fA`$, the equation: $`\mathrm{𝗍𝗋}_{_{\mathrm{}}}(f)=0`$ implies: $`f[A,A]`$. This is proved as follows (the argument below seems to be standard, but we could not find an appropriate reference in the literature).
Let $`𝔸=\mathrm{𝕜}x_1,x_2,\mathrm{},y_1,y_2\mathrm{}`$ be the free associative algebra on countably many variables, and $`[𝔸,𝔸]`$ the $`\mathrm{𝕜}`$-linear subspace of $`𝔸`$ spanned by the commutators. Similarly to formula (1.6), to any element $`F𝔸`$ one assignes a polynomial function $`\mathrm{𝗍𝗋}F`$ in infinitely many matrix variables: $`X_1,X_2,\mathrm{},Y_1,Y_2\mathrm{}𝔤𝔩_{\mathrm{}},`$ by inserting matrices instead of formal variables. We claim that: if $`F`$ is multi-linear in all its variables, and the polynomial function $`\mathrm{𝗍𝗋}F`$ is identically zero on $`𝔤𝔩_{\mathrm{}}`$, then $`F[𝔸,𝔸]`$. To prove this, note that modulo $`[𝔸,𝔸]`$ we can write: $`F(x_1,x_2,\mathrm{},y_1,y_2\mathrm{})=x_1\stackrel{~}{F}(x_2,\mathrm{},y_1,y_2\mathrm{}).`$ Hence, equation: $`0=\mathrm{𝗍𝗋}F(X_1,X_2,\mathrm{},Y_1,Y_2\mathrm{})=\mathrm{𝗍𝗋}\left(X_1\stackrel{~}{F}(X_2,\mathrm{},Y_1,Y_2\mathrm{})\right)`$ implies, since the trace pairing on $`𝔤𝔩_{\mathrm{}}`$ is nondegenerate, that the function $`\stackrel{~}{F}(x_2,\mathrm{},y_1,y_2\mathrm{})`$ is identically zero on $`𝔤𝔩_{\mathrm{}}`$. Furthermore, since $`𝔤𝔩_{\mathrm{}}`$ (viewed as an associative algebra) is known to be an algebra without polynomial identities, we conclude that $`\stackrel{~}{F}=0`$. Thus, $`F[𝔸,𝔸]`$, and our claim is proved.
We can now complete the proof of the Proposition. Fix $`fA`$ such that $`\mathrm{𝗍𝗋}f(X,Y)=0`$ identically on $`𝔤𝔩_{\mathrm{}}𝔤𝔩_{\mathrm{}}`$. Rescaling transformations: $`XtX,YsY,t,s\mathrm{𝕜}^\times ,`$ show that we may reduce to the case where $`f`$ is homogeneous in $`X`$ and $`Y`$ of degrees, say $`p,q`$, respectively. We now use the standard polarisation trick, and formally substitute: $`x=t_1x_1+\mathrm{}t_px_p,y=s_1y_1+\mathrm{}s_qy_q`$ into $`f`$, and then take the term multilinear in $`t_1,\mathrm{},s_q`$. This way we get from $`fA`$ a multilinear element $`F𝔸`$ such that $`\mathrm{𝗍𝗋}F=0`$ identically on $`𝔤𝔩_{\mathrm{}}`$. By the claim of the preceeding paragraph we conclude that $`F[𝔸,𝔸]`$. Observe now that sending all the $`x_i`$’s to $`x`$, and all the $`y_i`$’s to $`y`$ yields an algebra homomorphism $`\pi :𝔸A`$ such that $`\pi (F)=p!q!f`$. Applying this homomorphism to $`F`$ we get: $`f=\frac{1}{p!q!}\pi (F)\pi ([𝔸,𝔸])=[A,A].`$
## 5 The basics of $`𝒫`$-geometry.
Let $`𝒫=\{𝒫(n),n=1,2,\mathrm{}\}`$ be a $`\mathrm{𝕜}`$-linear quadratic operad with $`𝒫(1)=\mathrm{𝕜}`$, see \[GiK\]. Let $`𝕊_n`$ denote the Symmetric group on $`n`$ letters. Given $`\mu 𝒫(n)`$ and a $`𝒫`$-algebra $`A`$, we will write: $`\mu _A(a_1,\mathrm{},a_n)`$ for the image of $`\mu a_1\mathrm{}a_n`$ under the structure map: $`𝒫(n)_{_{𝕊_n}}A^nA.`$ Following \[GiK, §1.6.4\], we introduce an enveloping algebra $`𝒰^𝒫A`$, the associative unital $`\mathrm{𝕜}`$-algebra such that the abelian category of (left) $`A`$-modules is equivalent to the category of left modules over $`𝒰^𝒫A`$, see \[GiK, Thm. 1.6.6\]. The algebra $`𝒰^𝒫A`$ is generated by the symbols: $`u(\mu ,a),\mu 𝒫(2),aA,`$ subject to certain relations, see \[Ba, §1.7\].
An ideal $`I`$ in a $`𝒫`$-algebra $`A`$ will be called $`N`$-nilpotent if, for any $`nN,\mu 𝒫(n)`$, and $`a_1,\mathrm{},a_nA,`$ one has: $`\mu _A(a_1,\mathrm{},a_n)=0,`$ whenever at least $`N`$ among the elements $`a_1,\mathrm{},a_n`$ belong to $`I`$. The following useful reformulation of the notion of a left $`A`$-module is essentially well-known, see e.g., \[Ba, 1.2\]:
###### Lemma 5.1.
Giving a left $`A`$-module structure on a vector space $`M`$ is equivalent to giving a $`𝒫`$-algebra structure on $`A\mathrm{}M:=AM`$ such that the following conditions hold:
$`(𝗂)`$The imbedding: $`aa0`$ makes $`A`$ a $`𝒫`$-subalgebra in $`A\mathrm{}M`$.
$`(\mathrm{𝗂𝗂})`$$`M`$ is a 2-nilpotent ideal in $`A\mathrm{}M`$. ∎
A $`𝒫`$-algebra in the monoidal category of $`/2`$-graded, (resp. $``$-graded) super-vector spaces, see \[GiK, §1.3.17-1.3.18\], will be referred to as a $`𝒫`$-superalgebra, (resp. graded superalgebra). Any $`𝒫`$-algebra may be regarded as a $`𝒫`$-superalgebra concentrated in degree zero. Given a finite dimensional (super-) vector space $`V`$, write $`\overline{V}`$ for the same vector space with reversed parity. Let
$$𝖳__𝒫^{^{_{}}}V:=_{i1}𝒫(i)_{_{𝕊_i}}V^i\text{and}\stackrel{ˇ}{𝖳}__𝒫^{^{_{}}}{}_{}{}^{}V:=_{i1}𝒫(i)_{_{𝕊_i}}\overline{V}^i$$
be the free graded $`𝒫`$-algebra (resp. super-algebra) generated by $`V`$.
Fix a $`𝒫`$-algebra $`A`$, and consider the category of $`A`$-algebras, i.e. of pairs $`(B,p)`$, where $`B`$ is a $`𝒫`$-algebra and $`p:AB`$ is a $`𝒫`$-algebra morphism. Note that such a morphism makes $`B`$ an $`A`$-module. Thus, we get an obvious forgetful functor: $`A`$-algebras $``$ $`A`$-modules. The result below says that this functor has a right adjoint:
###### Lemma 5.2.
$`(𝗂)`$Given a $`𝒫`$-algebra $`A`$, there is a functor: $`MT_A^{^{_{}}}M,`$ (resp. $`M\stackrel{ˇ}{T}_A^{^{_{}}}M`$) assigning to a left $`A`$-module $`M`$ a graded $`𝒫`$-algebra $`T_A^{^{_{}}}M=_{i0}T_A^iM`$ (resp. graded $`𝒫`$-superalgebra $`\stackrel{ˇ}{T}_A^{^{_{}}}M=_{i0}T_A^i\overline{M}`$), such that $`T_A^0M=A`$.
$`(\mathrm{𝗂𝗂})`$For any $`𝒫`$-algebra map: $`AB`$, one has a natural adjunction isomorphism:
$$\mathrm{𝙷𝚘𝚖}_{_{A\text{-}\mathrm{𝗆𝗈𝖽}}}(M,B)\stackrel{}{}\mathrm{𝙷𝚘𝚖}_{_{𝒫\text{-}\mathrm{𝖺𝗅𝗀}}}(T_A^{^{_{}}}M,B).$$
Proof: If $`A`$ is a $`𝒫`$-subalgebra in a $`𝒫`$-algebra $`\stackrel{~}{A}`$, we define a $`𝒫`$-algebra $`T__A\stackrel{~}{A}`$ as the quotient of $`𝖳__𝒫^{}\stackrel{~}{A}`$, a free $`𝒫`$-algebra, modulo two-sided ideal generated by all relations of the form:
$$\mu a\stackrel{~}{a}=\mu _{_{\stackrel{~}{A}}}(a,\stackrel{~}{a}),\mu \stackrel{~}{a}a=\mu _{_{\stackrel{~}{A}}}(\stackrel{~}{a},a),\mu 𝒫(2),aA\stackrel{~}{A},\stackrel{~}{a}\stackrel{~}{A},$$
where $`\mu \stackrel{~}{a}a,\mu a\stackrel{~}{a}𝒫(2)\stackrel{~}{A}^2=𝖳__𝒫^2\stackrel{~}{A}`$, and $`\mu _{_{\stackrel{~}{A}}}(a,\stackrel{~}{a}),\mu _{_{\stackrel{~}{A}}}(\stackrel{~}{a},a)𝒫(1)\stackrel{~}{A}=𝖳__𝒫^1\stackrel{~}{A}`$. We now apply this construction to the algebra $`\stackrel{~}{A}=A\mathrm{}M`$, and put $`T_A^{^{_{}}}M:=T__A\stackrel{~}{A}`$, where the grading on the left accounts for the number of occurrences of elements of $`M`$, which is well-defined since the relations involved in the definition of $`T__A\stackrel{~}{A}`$ are ‘homogeneous in $`M`$’.
A closer look at the construction above shows that
$$T_A^{^{_{}}}M=A(𝖳__𝒫^{}M)/\mu ^{(12)}(a,m_1)m_2m_1\mu (a,m_2),$$
(5.3)
where $`\mathrm{}`$ denotes the two-sided ideal generated by the indicated subset of $`𝒫(2)M^2=𝖳__𝒫^2M`$, for all $`\mu 𝒫(2),aA,m_1,m_2M,`$ and where $`\mu ^{(12)}`$ stands for the action of the transposition $`(12)𝕊_2`$ on $`\mu `$. In particular, we have: $`T_A^0M=A`$ and $`T_A^1M=M`$.∎
Let $`A`$ be a $`𝒫`$-algebra and $`M`$ a left $`A`$-module. By Lemma 5.1, we may (and will) regard $`A\mathrm{}M`$ as a $`𝒫`$-algebra.
###### Definition 5.4.
A $`\mathrm{𝕜}`$-linear map $`\theta :AM`$ is called a derivation if the map: $`ama\theta (a)+m,`$ is an automorphism of the $`𝒫`$-algebra $`A\mathrm{}M`$.
Equivalently, following \[Ba, Definition 3.2.6\], extend $`\theta `$ to a $`\mathrm{𝕜}`$-linear map $`\theta ^{\mathrm{}}:A\mathrm{}MA\mathrm{}M,`$ given by $`\theta ^{\mathrm{}}:am\mathrm{\hspace{0.33em}0}\theta a`$. Then, $`\theta `$ is a derivation if and only if, for any $`\mu 𝒫(n)`$, we have:
$$\theta ^{\mathrm{}}\left(\mu _{_{A\mathrm{}M}}(b_1,\mathrm{},b_n)\right)=\underset{i=1}{\overset{n}{}}\mu _{_{A\mathrm{}M}}(b_1,\mathrm{},b_{i1},\theta ^{\mathrm{}}b_i,b_{i+1},\mathrm{},b_n),b_1,\mathrm{},b_nA\mathrm{}M.$$
Let $`\mathrm{𝙳𝚎𝚛}__𝒫(A,M)`$ denote the $`\mathrm{𝕜}`$-vector space of all derivations from $`A`$ to $`M`$. It is straightforward to see that the ordinary commutator makes $`\mathrm{𝙳𝚎𝚛}__𝒫(A,A)`$ a Lie algebra.
Next we define, following \[Ba, Def. 4.5.2\], an $`A`$-module of Kähler differentials as the left $`𝒰^𝒫A`$-module, $`\mathrm{\Omega }__𝒫^1A`$, generated by the symbols $`da`$, for $`aA`$, subject to the relations:
$`(𝗂)`$$`d(\lambda _1a_1+\lambda _2a_2)=\lambda _1da_1+\lambda _2da_2,\lambda _1,\lambda _2\mathrm{𝕜};`$
$`(\mathrm{𝗂𝗂})`$$`d(\mu (a_1,a_2))=u(\mu ,a_1)da_2+u(\mu ^{(12)},a_2)da_1,\mu 𝒫(2),a_1,a_2A,`$
where $`u(\mu ,a)`$ denote the standard generators of $`𝒰^𝒫A`$, see \[Ba\].
By construction, $`\mathrm{\Omega }__𝒫^1A`$ is a left $`A`$-module, and the assignment $`ada`$ gives a derivation $`d\mathrm{𝙳𝚎𝚛}__𝒫(A,\mathrm{\Omega }__𝒫^1A)`$. Moreover, this derivation is universal in the following sense. Given any left $`A`$-module $`M`$ and a derivation $`\theta :AM`$, there exists an $`A`$-module morphism $`\mathrm{\Omega }^1\theta :\mathrm{\Omega }__𝒫^1AM`$, uniquely determined by the condition that $`(\mathrm{\Omega }^1\theta )(da)=\theta (a).`$ It follows that the $`A`$-module of Kähler differentials represents the functor $`\mathrm{𝙳𝚎𝚛}__𝒫(A,)`$, i.e., we have (see \[Ba, Remark 4.5.4\]):
###### Lemma 5.5.
For any left $`A`$-module $`M`$ there is a natural isomorphism:
$$\mathrm{𝙳𝚎𝚛}__𝒫(A,M)\mathrm{𝙷𝚘𝚖}_{_{A\text{-}\mathrm{𝗆𝗈𝖽}}}(\mathrm{\Omega }__𝒫^1A,M).\mathrm{}$$
In particular, for $`M=A`$, we get an isomorphism: $`\mathrm{𝙳𝚎𝚛}__𝒫(A,A)\stackrel{}{}\mathrm{𝙷𝚘𝚖}_{_{A\text{-}\mathrm{𝗆𝗈𝖽}}}(\mathrm{\Omega }__𝒫^1A,A).`$ We let $`i_\theta \mathrm{𝙷𝚘𝚖}_{_{A\text{-}\mathrm{𝗆𝗈𝖽}}}(\mathrm{\Omega }__𝒫^1A,A)`$ denote the morphism: $`\mathrm{\Omega }__𝒫^1AA`$, corresponding to $`\theta \mathrm{𝙳𝚎𝚛}__𝒫(A,A)`$ under the isomorphism above.
We set $`\mathrm{\Omega }__𝒫^{}A:=\stackrel{ˇ}{T}_A^{^{_{}}}(\mathrm{\Omega }__𝒫^1A),`$ a graded $`𝒫`$-superalgebra generated by the $`A`$-module $`\mathrm{\Omega }__𝒫^1A`$. Recall that the differential envelope of a $`𝒫`$-algebra $`A`$ is a differential graded $`𝒫`$-super-algebra $`D^{}(A)=_{i0}D^i(A),`$ such that $`D^0(A)=A`$, and such that the following universal property holds: For any differential graded $`𝒫`$-superalgebra $`\stackrel{~}{D}^{}=_{i0}\stackrel{~}{D}^i,`$ and a $`𝒫`$-algebra morphism $`\rho :A\stackrel{~}{D}^0`$ there exists a unique DG-superalgebra morphism $`D(\rho ):D^{}(A)\stackrel{~}{D}^{}`$ such that $`D(\rho )|_{_{D^0(A)}}=\rho `$.
###### Proposition 5.6.
$`(𝗂)`$On $`\mathrm{\Omega }__𝒫^{}A`$, there exists a natural super-differential $`d:\mathrm{\Omega }__𝒫^{}A\mathrm{\Omega }__𝒫^{+1}A,d^2=0,`$ such that its restriction: $`A=\mathrm{\Omega }__𝒫^0A\mathrm{\Omega }__𝒫^1A`$ coincides with the canonical $`A`$-module derivation $`d:A\mathrm{\Omega }__𝒫^1A`$.
$`(\mathrm{𝗂𝗂})`$The differential graded $`𝒫`$-superalgebra $`(\mathrm{\Omega }__𝒫^{}A,d)`$ is the differential envelope of $`A`$.
Proof. We first give a direct construction of the differential envelope $`D^{}(A)`$ of a $`𝒫`$-algebra $`A`$, as follows. Let $`\overline{A}`$ denote a second copy of $`A`$ viewed as a $`\mathrm{𝕜}`$-vector space, and write $`\overline{a}`$ for the element of $`\overline{A}`$ corresponding to an element $`aA`$. We form the graded super-vector space $`A\overline{A}`$, where $`A`$ is placed in grade degree zero, and $`\overline{A}`$ is placed in grade degree 1. Let $`\stackrel{ˇ}{𝖳}__𝒫^{^{_{}}}{}_{}{}^{}(A\overline{A}):=_{i1}𝒫(i)_{_{𝕊_i}}(A\overline{A})^i`$ be the free $`𝒫`$-superalgebra generated by $`A\overline{A}`$, viewed as a graded superalgebra with respect to the total grading coming from both the grading on $`A\overline{A}`$ and the grading on the tensor algebra. We put: $`D^{}(A):=\stackrel{ˇ}{𝖳}__𝒫^{^{_{}}}{}_{}{}^{}(A\overline{A})/I,`$ where $`I`$ is the two-sided ideal generated by the following set:
$$\{\mu a_1a_2\mu (a_1,a_2),\mu \overline{a}_1a_2+\mu a_1\overline{a}_2\overline{\mu (a_1,a_2)}\}_{\mu 𝒫(2),a_1,a_2A}.$$
(5.7)
Thus, $`D^{}(A)`$ is a graded $`𝒫`$-superalgebra.
The $`\mathrm{𝕜}`$-linear endomorphism of $`A\overline{A}`$ given by the assignment: $`a\overline{a}_10\overline{a}`$ extends uniquely to a super-derivation $`\stackrel{ˇ}{𝖳}__𝒫^{^{_{}}}{}_{}{}^{}(A\overline{A})\stackrel{ˇ}{𝖳}__𝒫^{^{_{}}}{}_{}{}^{}(A\overline{A}).`$ This derivation descends to a well-defined derivation $`d`$ on $`D^{}(A)`$. Note that, for any $`xA\overline{A}`$ we have: $`d^2(x)=0`$. This implies, since the subspace $`A\overline{A}`$ generates the algebra $`D^{}(A)`$, that $`d^2=0`$ identically on $`D^{}(A)`$. Thus, $`d`$ makes $`D^{}(A)`$ a differential graded $`𝒫`$\- superalgebra.
The zero-degree component, $`D^0(A)`$, of the super-algebra $`D^{}(A)`$ is by construction a $`𝒫`$-subalgebra isomorphic to $`A`$, i.e., there is a canonical superalgebra imbedding $`j:A=D^0(A)D^{}(A)`$. Hence, $`D^{}(A)`$ may be regarded as an $`A`$-module, and the assignment: $`a\overline{a}`$ gives a derivation $`d\mathrm{𝙳𝚎𝚛}__𝒫(A,D^{}(A))`$. This derivation is universal in the sense explained above (for uniqueness property use that the superalgebra $`D^{}(A)`$ is generated by the subspace $`A\overline{A}`$). Hence $`D^{}(A)`$ is the differential envelope of $`A`$.
Observe next that the degree 1 component of $`D^{}(A)`$ is isomorphic, by definition of $`D^{}(A)`$, to the quotient of $`𝒰^𝒫A\overline{A}`$ by the relations $`(𝗂)`$–$`(\mathrm{𝗂𝗂})`$defining the module $`\mathrm{\Omega }__𝒫^1A`$ of Kähler differentials. Therefore, $`D^1(A)`$, the degree 1 component of $`D^{}(A)`$, is isomorphic to $`\mathrm{\Omega }__𝒫^1A`$ and, moreover, the canonical derivation $`d:A\mathrm{\Omega }__𝒫^1A`$ may be identified with the map: $`a\overline{a}D^{}(A)`$.
By the universal property of the tensor algebra, the $`A`$-module imbedding $`\mathrm{\Omega }^1AD^{}(A)`$ can be extended uniquely to a graded super-algebra morphism $`f:\stackrel{ˇ}{T}_A^{^{_{}}}(\mathrm{\Omega }__𝒫^1A)D^{}(A)`$. To show that $`f`$ is an isomorphism we construct its inverse, a map $`g:D^{}(A)\stackrel{ˇ}{T}_A^{^{_{}}}(\mathrm{\Omega }__𝒫^1A)`$, as follows. We have an obvious imbedding of $`\mathrm{𝕜}`$-vector spaces: $`A\overline{A}A\mathrm{\Omega }__𝒫^1A`$, given by: $`a\overline{a}_1ada_1.`$ This imbedding extends, by the universal property of a free $`𝒫`$-algebra, to a $`𝒫`$-superalgebra morphism $`\stackrel{~}{g}:\stackrel{ˇ}{𝖳}__𝒫^{^{_{}}}{}_{}{}^{}(A\overline{A})\stackrel{ˇ}{T}_A^{^{_{}}}(\mathrm{\Omega }__𝒫^1A).`$ The relations defining the ideal $`I`$ in formula (5.7) are designed in such a way that the morphism $`\stackrel{~}{g}`$ descends to a well-defined super-algebra morphism $`g:D^{}(A)\stackrel{ˇ}{T}_A^{^{_{}}}(\mathrm{\Omega }__𝒫^1A).`$ It is straightforward to verify that $`g=f^1`$. ∎
Remark. Our construction agrees with the notion of non-commutative differential forms for an algebra over the associative operad, as defined e.g. in \[L\] and used in §2 above.
From now on we assume, in addition, that $`𝒫`$ is a cyclic Koszul operad, see \[GeK\], with $`𝒫(1)=\mathrm{𝕜}`$. In particular, for each $`n1`$, the space $`𝒫(n)`$ is equipped with an $`𝕊_{n+1}`$-action that extends the $`𝕊_n`$-module structure on $`𝒫(n)`$ arising from the operad structure. Write $`\mathrm{𝚂𝚢𝚖}__\mathrm{𝕜}^2A`$ for the symmetric square of $`A`$. Following an idea of Kontsevich, Getzler and Kapranov introduce a functor $`𝖱:𝒫\text{-algebras}\mathrm{𝕜}\text{-vector spaces},`$
$$𝖱:A𝖱(A):=\frac{\mathrm{𝚂𝚢𝚖}__\mathrm{𝕜}^2A}{a_0\mu (a_1,a_2)\mu (a_0,a_1)a_2}_{|a_0,a_1,a_2A,\mu 𝒫(2)}.$$
Generalizing the Karoubi’s construction \[Ka\] in the associative case, define de Rham complex of $`A`$ as the graded vector space $`\mathrm{𝖣𝖱}^{^{_{}}}A:=𝖱(\mathrm{\Omega }__𝒫^{}A).`$ The differential $`d`$ on $`\mathrm{\Omega }__𝒫^{}A`$ induces a differential on $`\mathrm{𝖣𝖱}^{^{_{}}}A`$.
For any $`\theta \mathrm{𝙳𝚎𝚛}__𝒫A`$, the morphism $`i_\theta :\mathrm{\Omega }__𝒫^1AA`$ introduced after Lemma 5.5 extends to a super-derivation $`i_\theta :\mathrm{\Omega }__𝒫^{}A\mathrm{\Omega }__𝒫^1A`$, called the contraction operator. Further, the derivation $`\theta `$ induces, by a standard argument, a derivation $`L_\theta `$ of the associative algebra $`𝒰^𝒫A`$, and a map $`L_\theta :\mathrm{\Omega }__𝒫^1A\mathrm{\Omega }__𝒫^1A`$. The latter one extends to a derivation $`L_\theta :\mathrm{\Omega }__𝒫^{}A\mathrm{\Omega }__𝒫^{}A`$, called the Lie operator. The maps $`i_\theta `$ and $`L_\theta `$ descend naturally to the corresponding operators on $`\mathrm{𝖣𝖱}^{^{_{}}}A`$. It is straightforward to verify that these latter operators satisfy all the standard commutation relations (2.1).
## 6 Symplectic geometry of a free $`𝒫`$-algebra.
We keep the assumption that $`𝒫`$ is a cyclic Koszul operad. In this section which is a generalization of §3, inspired by works of Drinfeld \[Dr, Proposition 6.1 and above it\] and Kontsevich (private communication, 1994), we consider the case of a free $`𝒫`$-algebra. To avoid unnecessary repetitions and to simplify notation we only consider the ‘absolute’ case, i.e., the case of the ground ring $`B=\mathrm{𝕜}`$.
Fix a finite dimensional $`\mathrm{𝕜}`$-vector space $`E`$, and write $`A=𝖳__𝒫E`$ for the free $`𝒫`$-algebra (note that $`𝒫`$-algebras are algebras without unit, in general). We have:
$$𝖱(A)=\mathrm{𝖣𝖱}^0(A)=_{i1}𝒫_i_{_{𝕊_{i+1}}}E^{(i+1)},\mathrm{𝖣𝖱}^1(A)=AE.$$
(6.1)
Let $`\widehat{A}=_{i0}𝖳__𝒫^iE`$ denote the completion of $`A`$ with respect to the augmentation, and let $`\mathrm{𝙰𝚞𝚝}(\widehat{A})`$ denote the group of continuous algebra automorphisms of $`\widehat{A}`$. Any such automorphism $`\mathrm{\Phi }`$ is determined by its restriction to $`E=𝖳__𝒫^1E`$, a $`\mathrm{𝕜}`$-linear map $`\varphi :E\widehat{A}`$. We have an expansion: $`\varphi (v)=_{i=1}^{\mathrm{}}\varphi _i(v),`$ where $`\varphi _i(v)𝖳__𝒫^iE`$. We write $`d\mathrm{\Phi }:EE`$, for the map: $`v\varphi _1(v)`$; and we let $`\mathrm{𝙰𝚞𝚝}_{}(\widehat{A})`$ be the subgroup of $`\mathrm{𝙰𝚞𝚝}(\widehat{A})`$ formed by all automorphisms $`\mathrm{\Phi }`$ such that $`d\mathrm{\Phi }=\mathrm{𝙸𝚍}_E`$.
Observe further that the obvious grading on the free algebra $`A=𝖳__𝒫E`$ induces a natural grading $`𝖱^{^{_{}}}(A)=_i𝖱(A)_{(i)}`$, and, for each $`p0`$, a similar grading $`\mathrm{𝖣𝖱}^p(A)=_i\mathrm{𝖣𝖱}^p(A)_{(i)}`$. Fix a closed 2-form $`\omega \mathrm{𝖣𝖱}^2A`$, and let $`\omega =\omega _0+\omega _1+\mathrm{},\omega _i\mathrm{𝖣𝖱}^2(A)_{(i)},`$ be its expansion into graded components. We see, in particular, that $`\omega _0`$ may be viewed as an ordinary skew-symmetric $`\mathrm{𝕜}`$-bilinear form: $`E\times E\mathrm{𝕜}`$.
###### Theorem 6.2.
(Darboux theorem)$`(𝗂)`$A closed 2-form $`\omega =\omega _0+\omega _1+\mathrm{}\mathrm{𝖣𝖱}^2A`$ is non-degenerate if and only if so is the associated bilinear form $`\omega _0:E\times E\mathrm{𝕜}`$.
$`(\mathrm{𝗂𝗂})`$If $`\omega `$ is non-degenerate then there exists an automorphism $`\mathrm{\Phi }\mathrm{𝙰𝚞𝚝}_{}(\widehat{A})`$ such that: $`\mathrm{\Phi }^{}\omega =\omega _0.`$
Proof. Part $`(𝗂)`$is clear. Part $`(\mathrm{𝗂𝗂})`$is proved by the standard ‘homotopy argument’. Specifically, we consider a 1-parameter ‘family’:$`\omega _t=\omega _0+t\omega ^{}\mathrm{𝖣𝖱}^2A[[t]],`$ where $`\omega ^{}=\omega \omega _0=\omega _1+\omega _2+\mathrm{}\mathrm{𝖣𝖱}^2A`$. The 2-form $`\omega ^{}`$ being closed, there exists $`\alpha _{p1}\mathrm{𝖣𝖱}^1(A)_{(p)},`$ such that $`\omega ^{}=d\alpha `$. Since $`\omega _0`$ is non-degenerate, there exists a 1-parameter family $`\theta _t\mathrm{𝕜}[[t]]\widehat{}\mathrm{𝙳𝚎𝚛}__𝒫A=\mathrm{𝙳𝚎𝚛}__𝒫A[[t]]`$ determined uniquely from the equation: $`i_{\theta _t}\omega _t=\alpha `$. We define $`\mathrm{\Phi }(t)\mathrm{𝙰𝚞𝚝}(\widehat{A}[[t]]),`$ a formal one-parameter family of automorphisms of $`A`$, to be the solution of the differential equation: $`\frac{d\mathrm{\Phi }(t)}{dt}=L_{\theta _t}\mathrm{\Phi }(t)`$ of the form: $`\mathrm{\Phi }(t)=\mathrm{𝙸𝚍}__A+t\mathrm{\Phi }_1+t^2\mathrm{\Phi }_2+\mathrm{}.`$ It follows from the construction that $`\mathrm{\Phi }(t)^{}\omega _t=\omega _0`$, see e.g. \[GS\] for more details. Note further that the series $`\mathrm{\Phi }(t)`$ above has only finitely many terms in any given grade degree $`p0`$, i.e. terms that shift the grading on $`A`$ by $`p`$. In particular, setting $`t=1`$ in this series gives a well-defined element of $`\mathrm{\Phi }(1)\mathrm{𝙰𝚞𝚝}_{}(\widehat{A})`$ and we get: $`\mathrm{\Phi }(1)^{}\omega _{_{t=1}}=\omega _0`$. But $`\omega _{_{t=1}}=\omega `$, and part (ii) follows.∎
Because of this result, there is no loss of generality in considering only degree zero symplectic 2-forms $`\omega \mathrm{𝖣𝖱}^2A`$, i.e., such that $`\omega =\omega _0`$. Fix such an $`\omega `$, that is fix $`(E,\omega )`$, a symplectic vector space. Imitating the strategy used in §2 it is possible to define a Lie bracket on the vector space $`𝖱(A)`$. We prefer however to give the following direct explicit construction of this bracket similar to formula (1.4) in the associative case.
For each $`i,j1`$, let $`:\mu \nu \mu (1,\mathrm{},1,\nu ),`$ denote the operad-composition map: $`𝒫(i)𝒫(j)𝒫(1)\mathrm{}𝒫(1)𝒫(i)𝒫(j)𝒫(i+j1),`$ where $`1\mathrm{𝕜}=𝒫(1)`$, see \[GeK, Theorem 2.2(2)\]. We now change the notation and write: $`𝖱^{}(A)=_i𝖱^i(A),`$ where $`𝖱^i(A)`$, previously denoted by $`𝖱(A)_{(i)}`$, is the graded component with respect to the grading induced by one on $`A`$. Also, let $`\mathrm{𝖲𝗒𝗆}`$ be the ‘symmetrisation map’, the projection to $`𝕊_n`$-coinvariants. For each $`i,j1`$, we define a bilinear pairing $`\{,\}{}_{\omega }{}^{}:𝖱^i(A)𝖱^j(A)𝖱^{i+j1}(A)`$ as the following composition
$`𝖱^i(A)𝖱^j(A)=`$ $`\left(𝒫(i)_{𝕊_{i+1}}E^{i+1}\right){\displaystyle \left(𝒫(j)_{𝕊_{j+1}}E^{j+1}\right)}`$
$`\left(𝒫(i)𝒫(j)E^{i+j+2}\right)_{𝕊_{i+1}\times 𝕊_{j+1}}\stackrel{}{}`$
$`\left(𝒫(i+j1)E^{i+j+2}\right)_{𝕊_{i+1}\times 𝕊_{j+1}}\stackrel{\mathrm{𝖲𝗒𝗆}}{}`$
$`\left(𝒫(i+j1)_{𝕊_{i+j}}E^{i+j}\right){\displaystyle E^2}\stackrel{\mathrm{𝗂𝖽}\omega }{}`$
$`𝒫(i+j1)_{_{𝕊_{i+j}}}E^{i+j}=𝖱^{i+j1}(A).`$
An appropriate modification of the proof of Theorem 2.5, or a direct calculation, yields
###### Proposition 6.3.
The bracket $`\{,\}`$ makes $`𝖱^1(A)`$ into a graded Lie algebra. ∎
Let $`\mathrm{𝙳𝚎𝚛}__𝒫(A,\omega )`$ denote the Lie subalgebra in $`\mathrm{𝙳𝚎𝚛}__𝒫A`$ formed by all derivations $`\theta \mathrm{𝙳𝚎𝚛}__𝒫A`$ such that $`L_\theta \omega =0`$. Since $`\omega =\omega _0`$, this is equivalent to the requirement that the degree zero component $`d\theta :A_1A_1`$ induces an endomorphism of $`EE`$ that annihilates $`\omega ^{}EE`$. Using the same argument as in §§2-3, one proves the following two results
###### Lemma 6.4.
The assignment: $`\theta i_\theta \omega `$ gives graded vector space isomorphisms:
$$\mathrm{𝙳𝚎𝚛}__𝒫^{^{_{}}}(A)\stackrel{}{}\mathrm{𝖣𝖱}^1(A^{^{_{}}})\text{and}\mathrm{𝙳𝚎𝚛}__𝒫^{^{_{}}}(A,\omega )\stackrel{}{}\mathrm{𝖣𝖱}^1(A^{^{_{}}})_{_{\mathrm{𝖼𝗅𝗈𝗌𝖾𝖽}}}.$$
$`\mathrm{}`$
###### Proposition 6.5.
There is a canonical graded Lie algebra central extension:
$$0\mathrm{𝕜}𝖱^1(A)\mathrm{𝙳𝚎𝚛}__𝒫^{^{_{}}}(A,\omega )\mathrm{\hspace{0.17em}\hspace{0.17em}0}.$$
$`\mathrm{}`$
We call a pair $`(𝒮,\mathrm{𝗍𝗋})`$, where $`𝒮`$ is a $`𝒫`$-algebra and $`\mathrm{𝗍𝗋}`$ is a symmetric non-degenerate invariant bilinear form $`\mathrm{𝗍𝗋}:𝒮𝒮\mathrm{𝕜}`$, a symmetric $`𝒫`$-algebra. Any such bilinear form is determined, cf. \[GeK\], by a linear function $`\mathrm{𝗍𝗋}:𝖱(𝒮)\mathrm{𝕜},bb^{}\mathrm{𝗍𝗋}(bb^{})=\mathrm{𝗍𝗋}(b,b^{})`$.
From now on, fix a finite-dimensional symmetric $`𝒫`$-algebra $`(𝒮,\mathrm{𝗍𝗋})`$. Let $`\mathrm{𝙰𝚞𝚝}(𝒮,\mathrm{𝗍𝗋})`$ denote the algebraic group of automorphisms of the $`𝒫`$-algebra $`𝒮`$ that preserve the bilinear form $`\mathrm{𝗍𝗋}`$. The corresponding Lie algebra $`\mathrm{𝙳𝚎𝚛}__𝒫(𝒮,\mathrm{𝗍𝗋})`$ is formed by all the derivations $`\theta \mathrm{𝙳𝚎𝚛}__𝒫𝒮`$ such that, for any $`b,b^{}𝒮`$, one has: $`\mathrm{𝗍𝗋}(\theta (b),b^{})+\mathrm{𝗍𝗋}(b,\theta (b^{}))=0.`$
Representation functor. For any finitely generated $`𝒫`$-algebra $`A`$, the set $`\mathrm{𝙷𝚘𝚖}_{_{𝒫\text{-}\mathrm{𝖺𝗅𝗀}}}(A,𝒮)`$ has the natural structure of a finite dimensional affine algebraic variety, acted on by the algebraic group $`\mathrm{𝙰𝚞𝚝}(𝒮,\mathrm{𝗍𝗋})`$. We put $`\mathrm{𝖱𝖾𝗉}(A,𝒮):=\mathrm{𝕜}[\mathrm{𝙷𝚘𝚖}_{_{𝒫\text{-}\mathrm{𝖺𝗅𝗀}}}(A,𝒮)].`$
Let now $`(E,\omega )`$ be a finite dimensional symplectic vector space, and $`A=T__𝒫E`$, the free $`𝒫`$-algebra on $`E`$. Then we clearly have: $`\mathrm{𝙷𝚘𝚖}_{_{𝒫\text{-}\mathrm{𝖺𝗅𝗀}}}(A,𝒮)=\mathrm{𝙷𝚘𝚖}__\mathrm{𝕜}(E,𝒮)=E^{}__\mathrm{𝕜}𝒮`$, is a finite dimensional $`\mathrm{𝕜}`$-vector space. The symplectic 2-form $`\omega `$ on $`E`$ gives rise, as in §3, to the symplectic 2-form $`\omega _{_{\mathrm{𝖱𝖾𝗉}}}:=\omega ^{}\mathrm{𝗍𝗋}`$ on $`\mathrm{𝙷𝚘𝚖}__\mathrm{𝕜}(E,𝒮)=E^{}__\mathrm{𝕜}𝒮`$. The action of the group $`\mathrm{𝙰𝚞𝚝}(𝒮,\mathrm{𝗍𝗋})`$ on $`\mathrm{𝙷𝚘𝚖}_{_{𝒫\text{-}\mathrm{𝖺𝗅𝗀}}}(A,𝒮)`$ preserves this symplectic form and is, moreover, Hamiltonian. In other words, the vector field on $`E^{}__\mathrm{𝕜}𝒮`$ arising from a derivation $`\theta \mathrm{𝙳𝚎𝚛}__𝒫(𝒮,\mathrm{𝗍𝗋})`$ is induced by an $`\mathrm{𝙰𝚞𝚝}(𝒮,\mathrm{𝗍𝗋})`$-invariant Hamiltonian function $`H__\theta \mathrm{𝖱𝖾𝗉}(A,𝒮)`$. Explicitly, the function $`H__\theta `$ is given by the following quadratic polynomial on $`E^{}__\mathrm{𝕜}𝒮`$:
$$H__\theta :_k\stackrel{ˇ}{x}_ks_k_{i<j}\omega ^{}(\stackrel{ˇ}{x}_i,\stackrel{ˇ}{x}_j)\mathrm{𝗍𝗋}(\theta (s_i),s_j),x_lE^{},s_l𝒮,l=i,j,k.$$
Write $`\mathrm{𝖱𝖾𝗉}(A,𝒮)^{\mathrm{𝙰𝚞𝚝}(𝒮,\mathrm{𝗍𝗋})}`$ for the $`\mathrm{𝕜}`$-algebra of $`\mathrm{𝙰𝚞𝚝}(𝒮,\mathrm{𝗍𝗋})`$-invariant polynomial functions on the $`\mathrm{𝕜}`$-vector space $`\mathrm{𝙷𝚘𝚖}__\mathrm{𝕜}(E,𝒮)`$. The symplectic form $`\omega _{_{\mathrm{𝖱𝖾𝗉}}}=\omega \mathrm{𝗍𝗋}`$ makes $`\mathrm{𝖱𝖾𝗉}(A,𝒮)^{\mathrm{𝙰𝚞𝚝}(𝒮,\mathrm{𝗍𝗋})}`$ into a Poisson algebra. We have the standard Lie algebra central extension:
$$0\mathrm{𝕜}\mathrm{𝖱𝖾𝗉}(A,𝒮)^{\mathrm{𝙰𝚞𝚝}(𝒮,\mathrm{𝗍𝗋})}\stackrel{\delta }{}\mathrm{𝙳𝚎𝚛}_{_{\omega _{_{\mathrm{𝖱𝖾𝗉}}}}}\left(\mathrm{𝖱𝖾𝗉}(A,𝒮)^{\mathrm{𝙰𝚞𝚝}(𝒮,\mathrm{𝗍𝗋})}\right)\mathrm{\hspace{0.17em}\hspace{0.17em}0},$$
(6.6)
where $`\mathrm{𝙳𝚎𝚛}_{_{\omega _{_{\mathrm{𝖱𝖾𝗉}}}}}\left(\mathrm{𝖱𝖾𝗉}(A,𝒮)^{\mathrm{𝙰𝚞𝚝}(𝒮,\mathrm{𝗍𝗋})}\right)`$ stands for the Lie algebra of derivations of the commutative algebra $`\mathrm{𝖱𝖾𝗉}(A,𝒮)^{\mathrm{𝙰𝚞𝚝}(𝒮,\mathrm{𝗍𝗋})}`$ respecting the Poisson bracket. It is straightforward to check that the assignment: $`\theta H__\theta `$ gives a Lie algebra splitting: $`\mathrm{𝙳𝚎𝚛}_{_{\omega _{_{\mathrm{𝖱𝖾𝗉}}}}}\left(\mathrm{𝖱𝖾𝗉}(A,𝒮)^{\mathrm{𝙰𝚞𝚝}(𝒮,\mathrm{𝗍𝗋})}\right)\mathrm{𝖱𝖾𝗉}(A,𝒮)^{\mathrm{𝙰𝚞𝚝}(𝒮,\mathrm{𝗍𝗋})}`$ of the surjective morphism $`\delta `$ in the exact sequence above.
Observe next that the ‘infinite dimensional’ group $`\mathrm{𝙰𝚞𝚝}(A)`$ acts naturally on $`\mathrm{𝙷𝚘𝚖}_{_{𝒫\text{-}\mathrm{𝖺𝗅𝗀}}}(A,𝒮)`$. This action commutes with that of the group $`\mathrm{𝙰𝚞𝚝}(𝒮,\mathrm{𝗍𝗋})`$, preserves the symplectic form $`\omega _{_{\mathrm{𝖱𝖾𝗉}}}`$, but it is not Hamiltonian, in general. That means that the induced Lie algebra morphism $`\xi :\mathrm{𝙳𝚎𝚛}__𝒫A\mathrm{𝙳𝚎𝚛}_{_{\omega _{_{\mathrm{𝖱𝖾𝗉}}}}}\left(\mathrm{𝖱𝖾𝗉}(A,𝒮)^{\mathrm{𝙰𝚞𝚝}(𝒮,\mathrm{𝗍𝗋})}\right)`$ cannot be lifted, in general, to a Lie algebra morphism: $`\mathrm{𝙳𝚎𝚛}__𝒫A\mathrm{𝖱𝖾𝗉}(A,𝒮)^{\mathrm{𝙰𝚞𝚝}(𝒮,\mathrm{𝗍𝗋})},`$ see (6.6).
The following result shows that the $`\mathrm{𝙰𝚞𝚝}(A)`$-action becomes Hamiltonian after a 1-dimensional central extension. The result below agrees also with the philosophy advocated in \[KR\], saying that, for any finite-dimensional (symmetric) $`𝒫`$-algebra $`𝒮`$, ‘functions’ on the non-commutative space corresponding to a $`𝒫`$-algebra $`A`$ should go into genuine regular functions on the affine algebraic variety $`\mathrm{𝙷𝚘𝚖}_{_{𝒫\text{-}\mathrm{𝖺𝗅𝗀}}}(A,𝒮)`$.
###### Theorem 6.7.
There is a natural Lie algebra homomorphism $`\phi :𝖱(A)\mathrm{𝖱𝖾𝗉}(A,𝒮)^{\mathrm{𝙰𝚞𝚝}(𝒮,\mathrm{𝗍𝗋})}`$ making the following diagram commute:
Proof. Very similar to the proof of Proposition 3.4.∎
Department of Mathematics, University of Chicago, Chicago IL 60637, USA; ginzburg@math.uchicago.edu
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# An exact-diagonalization study of rare events in disordered conductors
## Abstract
We determine the statistical properties of wave functions in disordered quantum systems by exact diagonalization of one-, two- and quasi-one dimensional tight-binding Hamiltonians. In the quasi-one dimensional case we find that the tails of the distribution of wave-function amplitudes are described by the non-linear $`\sigma `$-model. In two dimensions, the tails of the distribution function are consistent with a recent prediction based on a direct optimal fluctuation method.
It is well established that disordered quantum systems in the metallic regime (i.e., in the limit of weak disorder) and highly excited classically chaotic quantum systems exhibit universal quantum fluctuations that can be described by random matrix theory (RMT): statistical properties, on the scale of the mean level spacing, of eigenvalues, eigenfunctions, and matrix elements are universal, i.e., they do not depend on the microscopic details of the systems under consideration .
However, in ballistic, classically chaotic quantum systems, non-hyperbolic phase-space structures may lead to deviations from universal RMT statistics . Similarly fluctuations in disordered, classically diffusive quantum systems may deviate considerably from the RMT predictions due to increased localization. This effect is naturally very significant in the tails of distribution functions (corresponding to rare events) of wave-function amplitudes , of the local density of states , of inverse participation ratios and of NMR line shapes . In all of these cases (with the exception of Ref. which deals with one-dimensional (1D) systems), the distribution functions have been calculated using the non-linear $`\sigma `$-model (NLSM). Very recently, this approach has been extended to ballistic systems (see also ).
In Ref. a direct optimal fluctuation method was used to calculate the tails of distributions of current relaxation times and wave-function amplitudes; and predictions differing from were put forward. This led the authors of to question the suitability of the NLSM to describe rare events in disordered conductors.
It is thus of great interest to test the predictions of and against results of independent calculations. In this letter, we have determined distribution functions of wave-function amplitudes by exact diagonalization of 1D, 2D and quasi-1D tight-binding Hamiltonians; in this case rare events correspond to unusually high splashes of wave-function amplitudes. We note that wave-function amplitude distributions can be measured in micro-wave experiments .
We use the Anderson model of localization which is a tight-binding model on a $`d`$-dimensional hyper-cubic lattice
$$\widehat{H}=\underset{𝒓,𝒓^{}}{}t_{𝒓𝒓^{}}^{}c_𝒓^{}c_𝒓^{}^{}+\underset{𝒓}{}\upsilon _𝒓^{}c_𝒓^{}c_𝒓^{}.$$
(1)
Here $`𝒓=(x,y,\mathrm{})`$ denotes sites on the lattice, $`c_𝒓^{}`$ and $`c_𝒓`$ are the usual creation and annihilation operators, the hopping amplitudes are $`t_{𝒓𝒓^{}}=1`$ for nearest neighbour sites and zero otherwise. The on-site potential $`\upsilon _𝒓`$ is taken to be uncorrelated white noise, with zero mean and variance $`\upsilon _𝒓\upsilon _𝒓^{}=\delta _{𝒓𝒓^{}}W^2/12`$. The parameter $`W`$ characterizes the disorder strength. As is well-known (see for instance ), the eigenvalues $`E_j`$ and eigenfunctions $`\psi _j(𝒓)`$ of this Hamiltonian, in the metallic regime, exhibit fluctuations described by RMT. In this case, Dyson’s Gaussian orthogonal ensemble is appropriate. When the matrix elements $`t_{𝒓𝒓^{}}`$ of (1) are given an appropriate complex phase factor, Dyson’s unitary ensemble applies. We refer to these two cases by assigning, as usual, the parameter $`\beta =1`$ to the former and $`\beta =2`$ to the latter. The metallic regime is characterized by $`g1`$ where $`g=2\pi \nu VDL^2`$ is the dimensionless conductance (we take $`\mathrm{}=1`$). Here $`\nu =1/(V\mathrm{\Delta })`$, $`\mathrm{\Delta }`$ is the mean level spacing and $`V`$ the volume. $`D=v_\mathrm{F}^2\tau /d`$ is the diffusion constant, $`\tau `$ the mean free time and $`v_\mathrm{F}`$ the Fermi velocity. Four length scales are important: the lattice spacing $`a`$, the linear extension $`L`$, the localization length $`\xi `$ and the mean free path $`\mathrm{}=v_\mathrm{F}\tau `$.
By diagonalizing the Hamiltonian $`\widehat{H}`$ using a modified Lanczos algorithm , we have determined the distribution function
$$f_\beta (E,𝒓;t)=\mathrm{\Delta }\underset{j}{}\delta (t|\psi _j(𝒓)|^2V)g_\eta (EE_j)_W.$$
(2)
Here $`\mathrm{}_W`$ denotes an average over disorder realisations. The wave functions are normalized so that $`|\psi _j(𝒓)|^2_W=V^1`$ and $`g_\eta (E)`$ is a window function of width $`\eta `$, centered around $`E=0`$ and normalized to unity. In the following we describe the results of our calculations and compare them to the predictions of Refs. .
1D case. The eigenstates in a disordered chain are localized with localization length $`\xi =4\mathrm{}`$. According to Ref. the distribution of wave-function amplitudes in a disordered chain of length $`L`$ is
$$f(E;t)\frac{\xi }{Lt}\mathrm{exp}\left(\frac{t\xi }{L}\right)$$
(3)
for $`L\mathrm{}`$, independent of $`x`$ and $`\beta `$. Our results for $`f(E,x;t)_x`$ ($`\mathrm{}_x`$ denotes an average over the $`x`$ coordinate) in Fig. 1 show very good agreement with Eq. (3) for large $`L/\xi `$. The deviations at small $`t`$ for $`L/\xi =4.76`$ are due to the fact that Eq. (3) is only valid asymptotically for large $`L`$. It does not take into account that in a system of finite length $`L`$, the smallest amplitude of a normalized, exponentially decaying wave function is of the order of $`t_\mathrm{c}(L/\xi )\mathrm{exp}(L/\xi )`$. This cut-off is shown in Fig. 1.
Quasi-1D case. In this case, as was shown in Refs. , the NLSM can be solved exactly for the distribution function $`f_\beta (E,x;t)`$, using a transfer matrix approach . The result is
$$f_1(E,x;t)=\frac{2\sqrt{2}}{\pi \sqrt{t}}\frac{\mathrm{d}^2}{\mathrm{d}t^2}_0^{\mathrm{}}\frac{\mathrm{d}z}{\sqrt{z}}Y(z+t/2),$$
(4)
$$f_2(E,x;t)=\frac{\mathrm{d}^2}{\mathrm{d}t^2}Y(t).$$
(5)
Here $`Y(z)=𝒲(z\xi /L,x/\xi )𝒲(z\xi /L,(Lx)/\xi )`$ and $`𝒲(z,\tau )`$ obeys the differential equation
$$\frac{}{\tau }𝒲(z,\tau )=\left(z^2\frac{^2}{z^2}z\right)𝒲(z,\tau )$$
(6)
with initial condition $`𝒲(z,0)=1`$. The function $`𝒲(z,t)`$ may be determined in terms of an eigenfunction expansion of the operator $`z^2_z^2z`$ . The change of the body and the tails of distribution $`f_\beta (E,x;t)`$ due to increasing localization may thus be parameterized by a single parameter which we define to be $`X(\beta /2)L/\xi `$. Here $`\xi \beta \pi \nu DS`$ where $`S`$ is the cross-section of the wire. Thus $`X`$ does not depend on $`\beta `$. In the metallic regime (where $`X0`$) $`Y(z)\mathrm{exp}(z)`$ which leads to the usual RMT results $`f_1^{(0)}(t)=\mathrm{exp}(t/2)/\sqrt{2\pi t}`$ and $`f_2^{(0)}(t)=\mathrm{exp}(t)`$. The former distribution ($`\beta =1`$) is often referred to as the Porter-Thomas distribution . For increasing localization (finite but still small $`X`$), the $`𝒪(X)`$-corrections to the body of the distribution function $`f_\beta ^{(0)}`$ are obtained by expanding $`Y\mathrm{exp}(t)[1+\beta ^1t^2P(x,x;0)]`$ where $`P(x,x^{};\omega )`$ is the one-dimensional diffusion propagator. The result is $`f_\beta (E,x;t)=f_\beta ^{(0)}(t)[1+\delta f_\beta (E,x;t)]`$ with
$`\delta f_\beta (E,x;t)`$ $``$ $`P(x,x;0)\{\begin{array}{cc}3/43t/2+t^2/4\hfill & \text{for}\beta =1,\hfill \\ 12t+t^2/2\hfill & \text{for}\beta =2,\hfill \end{array}`$ (9)
valid for $`tX^{1/2}`$. In the tails ($`tX^1>1`$) of $`f_\beta (E,x;t)`$, Eqs. (4) and (5) simplify to
$$f_\beta (E,x;t)A_\beta (x,X)\mathrm{exp}(2\beta \sqrt{t/X}).$$
(10)
This result may also be obtained within a saddle-point approximation to the NLSM . The prefactors $`A_\beta (x,X)`$ for $`\beta =1,2`$ are given in .
According to Refs. and the NLSM applies provided the following conditions are satisfied:
$$1k_\mathrm{F}\mathrm{}k_\mathrm{F}^2Sk_\mathrm{F}L$$
(11)
where $`k_\mathrm{F}`$ is the Fermi wave vector and $`S`$ is the cross section of the wire. The first condition ensures that disorder is sufficiently weak. The second condition implies that apart from the sample geometry, all other properties are essentially 3D . Due to the third condition the return probability is dominated by diffusive contributions . When $`k_\mathrm{F}=𝒪(a^1)`$, Eq. (11) corresponds to $`1\mathrm{}/aML/a`$, where $`M=k_\mathrm{F}^2S`$ is the number of channels. Furthermore, in the metallic regime, one must have $`\mathrm{}L\xi `$. Since $`\xi M\mathrm{}`$, this implies
$$1L/\mathrm{}M.$$
(12)
In a finite system, the conditions (11) and (12) are not easily met simultaneously. We have performed exact diagonalizations for $`128\times 4\times 4`$ and $`128\times 8\times 8`$ lattices, using open boundary conditions (BC) in the longitudinal direction. In this case, $`P(x,x;0)=2X[1/3x(Lx)/L^2]`$.
The results of our calculations are summarized in Figs. 2 and 3. Figure 2 shows $`\delta f_\beta (E,x;t)_x`$ in comparison with Eqs. (4),(5) and (9). We observe very good agreement. The value of $`PP(x,x;0)_x`$ should be independent of $`\beta `$. As can be seen in Fig. 2, the value of $`P`$ does somewhat change with $`\beta `$, albeit weakly. For narrower wires ($`128\times 4\times 4`$) we have observed that the ratio $`P_1/P_2`$ (determined by fitting $`PP_\beta `$ independently for $`\beta =1,2`$) becomes very small for small values of $`W`$ (corresponding to $`X\stackrel{<}{}0.1`$) while it approaches unity for large values of $`W`$. A possible explanation for this deviation would be that for small $`M`$ and small $`X`$, the condition (11) is no longer satisfied since $`\mathrm{}/a\stackrel{>}{}M`$. Surprisingly, the form of the deviations is still very well described by Eq. (9) (not shown).
Figure 3 shows the tails of $`f_\beta (E,x;t)`$ for weak disorder ($`X\stackrel{<}{}1`$) in comparison with Eqs. (4), (5) and (10). Since for very small values of $`X`$ the tails decay so fast that we cannot reliably calculate them, we decreased the wire cross section and increased the value of $`W`$ in Fig. 3, thus increasing $`X`$. The quoted values of $`X`$ were obtained by fitting Eqs. (4) and (5). The values thus determined differ somewhat between $`\beta =1`$ and $`2`$ (see Fig. 3) and this difference depends on the choice of $`E,\eta `$ and $`W`$.
In summary we conclude that non-universal deviations from RMT statistics in quasi-1D wires are very well described by a NLSM not only in the body (see Fig. 2) but notably also in the tails (see Fig. 3) of the distribution $`f_\beta (E,x;t)`$.
2D case. In this case, according to Ref. , corrections to $`f_\beta ^{(0)}`$ are still given by Eq. (9), but now $`P=P(𝒓,𝒓;0)_𝒓`$ where $`P(𝒓,𝒓^{};\omega )`$ is the 2D diffusion propagator. For the tails of the distributions, the result of the NLSM is within a saddle-point approximation
$$f_\beta (E,𝒓;t)\mathrm{exp}\left[C_\beta (\text{ln}t)^2\right]$$
(13)
with
$$C_\beta =\beta \pi g/[4\text{ln}(L/\mathrm{})].$$
(14)
Note that according to (14) the decay in the tails of Eq. (13) depends on $`\beta `$, as in the quasi-1D case \[Eq. (10)\].
Recently, in Ref. a different approach (a direct optimal fluctuation method ) was used to calculate the tails of $`f_\beta (E,𝒓;t)`$. According to Ref. , the tails of the distribution function are given by Eq. (13) but with $`C_\beta `$ replaced by
$$C=\pi g/[2\text{ln}(L/r_0)]$$
(15)
\[with $`r_0=𝒪(k_\mathrm{F}^1)`$\] which differs in two respects from the prediction of the NLSM: First, $`\mathrm{}`$ in $`C_\beta `$ is replaced by $`r_0`$ in (15). Second, there is no $`\beta `$-dependence.
We have diagonalized the Hamiltonian (1) on a $`100\times 100`$ lattice. Figure 4 shows corrections to $`f_\beta ^{(0)}`$ for weak disorder. We find that the form of the deviations is very well described by Eq. (9). However, the values of $`P_\beta `$ obtained for $`\beta =1,2`$ differ by a factor $`<1/2`$. A possible explanation for this deviation might be the following : In the ballistic regime, $`P`$ is no longer given by the diffusion propagator but may be dominated by a single-scattering expression which involves an additional factor $`\beta /2`$ and thus $`P_1/P_2<1`$. It would be tempting to deduce from this that in our case ballistic effects are important. However, this does not explain $`P_1/P_2<1/2`$. Numerical results (albeit for rather small systems) indicate that in 3D, $`P_1/P_21/2`$ for the parameters chosen in .
Figure 5 shows the tails of the distribution functions. The tails are consistent with an $`\mathrm{exp}[C(\mathrm{ln}t)^2]`$ decay as predicted by Eq. (13). We have thus verified that corrections to RMT distributions in 2D systems do give rise to log-normal tails. Our results suggest that the prefactor in the exponent does not depend on $`\beta `$. This result is consistent with Eq. (15). We have independently calculated the dimensionless conductance $`g`$ using the usual linear response expression. We find that $`g`$ is independent of $`\beta `$ and $`gW^2`$, as expected (inset of Fig. 5).
The inset of Fig. 5 shows that $`C`$ increases with decreasing disorder strength, as it should, albeit slower than $`W^2`$. The increase of $`C`$ for decreasing $`W`$ is underestimated, because for weak disorder, the tails of the distributions have not reached the asymptotic regime.
In summary we have reported on a study of rare events in disordered conductors, by diagonalizing the tight-binding Hamiltonian (1) and analyzing the probability of rare splashes of high wave-function amplitudes. Our 1D results agree with those of . In the quasi-1D case, we have compared our data to an exact solution of the NLSM, and to a saddle-point approximation . We observe very good agreement between our results and those of the calculations based on the NLSM and thus conclude that the NLSM provides a quantitative description of rare events in quasi-1D disordered conductors. In 2D systems, corrections to the body of the distribution functions are well described by results based on the NLSM, with a modified prefactor $`P_\beta `$. Moreover, we could verify that the tails of the distribution function in the vicinity of the metallic regime are log-normal. Thus our numerical investigations, which are complementary to the analytical predictions, corroborate the overall picture suggested in Refs. for quasi-1D and 2D systems. The coefficient describing the tails of wave-function distributions in 2D systems turns out to be independent of $`\beta `$ for the parameters considered in this paper, as opposed to the quasi-1D case. This is consistent with the prediction of the direct optimal fluctuation method (Ref. ).
Acknowledgment. This work was supported by the DFG under project C5/SFB393.
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# Instanton-induced Effects in QCD High-Energy Scattering
## I Introduction
Significant progress reached in the realm of non-perturbative QCD has been mostly related with approaches based on the Euclidean formulation of the theory: numerical simulations using lattice gauge theory, instantons, monopoles, etc. By now, we know a great deal about the important or even dominant role of instanton-induced effects for correlation functions in a variety of hadronic channels, hadronic wave functions and form-factors, for a review see . Unfortunately so far many of those results have not been translated to Minkowski space, a crucial step for understanding hadronic high-energy processes. It is however clear that there must be a very general and direct relationship between the hadronic substructure and the details of high energy reactions. Indeed, the non-perturbative modification of QCD vacuum fields induced by the valence quarks studied in Euclidean space-time should look like parton correlations in the transverse plane in a boosted frame. Many known features of partonic distributions, including spin and flavor of the sea quarks, point to their non-perturbative origin. Many more features (like $`fluctuations`$ of these cross sections and $`correlations`$ in the parton positions in the transverse plane which we briefly discuss at the end of the paper) are still to be studied in details.
The first systematic step towards a semi-classical but non-perturbative formulation of high-energy scattering in QCD was suggested by Nachtmann , who has related the scattering amplitude to expectations of pairs of Wilson lines. Semi-classical expressions with a similar pair of Wilson lines for DIS structure functions were also proposed by Muller : in contrast to their traditional interpretation as partonic densities, they were treated as cross sections for targets penetrated by small dipole-like probes at high energy. One systematic way to use these semi-classical expressions is to go back to the perturbative domain and try to improve on the diagrammatic approaches (like the celebrated BFKL re-summation): see e.g. calculations of the anomalous dimension of the cross-singularity between two Wilson lines or the analysis of the path-exponents in .
The approach we will follow in this paper is different: the Wilson lines in question are evaluated semi-classically using instantons. In order to be able to do so, one should start in Euclidean space-time, where those solutions are the saddle points of the functional integrals. The results are then analytically continued back to Minkowski space. Although it was not done before in this form, there are similar approaches in the perturbative context (e.g. and references therein). Another methodically close approach to our analysis is where recent progress on the non-perturbative dynamics in N=4 SYM theory was used. In particular, the AdS/CFT correspondence have been used to evaluate the partonic cross section geometrically, using a deformed string in the curved anti-de-Sitter space.
The instanton-induced processes to be considered in this work are either $`elastic`$ scattering of partons, or quasi-elastic ones, with color transfer between them. They are very different from (and should not be confused with) multi-quanta production processes originally discussed in electroweak theory in connection with baryon number violation and later in QCD in connection to DIS . Such phenomena, associated with small-size instantons, are easier to evaluate and also they should lead to much more spectacular events. However, those lead to much smaller cross sections in comparison to the processes to be discussed below.
In this paper, we will not aim at a development of a realistic model for high-energy hadronic reactions based on instanton physics. Instead, we will answer few questions of principle, such as : Is it possible to assess non-perturbatively scattering amplitudes using the Euclidean formulation of the theory? How is the analytical continuation enforced on the non-perturbative amplitudes? What is the magnitude of the instanton induced effects in comparison to the perturbative effects in the scattering of near-forward high energy partons?
In section 2 we review the perturbative effects on the dipole-dipole potential, including the derivation of a renormalization group solution that can be tested using QCD lattice simulations. In section 3 we extend the perturbative analysis in Euclidean space to the case of scattering between two quarks and two dipoles. Particular issues regarding the analytical continuation of the perturbative results to Minkowski space are discussed. In section 4, we discuss the effects of instantons on the static potentials for quarks and dipoles. At large distances the results resemble perturbation theory apart from the large classical enhancement of $`(8\pi ^2/g^2)^210^2`$, which is partially compensated by the diluteness factor $`n_0\rho _0^4(1/3)^4`$ of the instantons in the vacuum. In section 5, we calculate the scattering amplitudes for quarks and dipoles in the one-instanton approximation. The color preserving part of the amplitude is real and vanishes at high energy. The color exchange part is real but finite at high energy, thereby contributing to the near-forward inelastic scattering or re-scattering of partons. In section 6, we extend our discussion to two-instantons. We found that for two quarks the cross section is of the order of $`\sigma _{qq}(n_0\rho _0^4)^2\rho _0^2`$, while for two dipoles it is further suppressed $`\sigma _{dd}\sigma _{qq}(d_1^2d_2^2/\rho _0^4)`$. These results are supported by our calculations. In section 7, we discuss the possible role of instantons in cross-section fluctuations. Our conclusions and recommendations are in section 8.
## II Perturbative Analysis of Potentials
### A Dipole- dipole potential
We start with the simplest analysis in Euclidean space, in which the perturbative expansion of two Euclidean Wilson lines leads to the well-known result for the potential between static charges. Indeed, by expanding two Wilson lines to first order in the gauge-coupling $`g`$, using the Euclidean propagator $`A(x)A(y)1/(xy)^2`$ with $`x,y`$ located on two parallel but straight lines, and finally integrating over the relative time $`x_0y_0`$, we readily obtain the Coulomb potential $`V(R)\alpha _s(R)/R`$.
Now, consider the case of the interaction between two color neutral objects, such as two static color dipoles. The simplest perturbative process in this case includes double photon/gluon exchange. The problem was solved in QED by Casimir and Polder , who have shown that the potential at large distances R is
$`V(R)={\displaystyle \frac{\alpha _1\alpha _2}{R^7}}`$ (1)
where the polarizabilities $`\alpha _{1,2}`$ are of the order of $`\alpha _{1,2}\tau _0d^2`$, $`d`$ is the dipole size and $`\tau _0`$ is some characteristic time (see below). This result differs from the Van-der-Waals potential $`1/R^6`$ (valid at smaller R) because of the time delay effects. These observations were generalized to perturbative QCD in .
The Euclidean approach leads to the 7-th power of R in a simple way, provided that the following conditions are satisfied: (i) $`d_{1,2}R`$ which justifies the dipole approximation and identifies the relevant field operators $`(\stackrel{}{d}\stackrel{}{E})^2`$; (ii) $`both`$ exchanged photons (or gluons) are emitted and absorbed at close $`x_0`$ and $`y_0`$ times. As a result, the perturbative field correlator $`<E^2(x)E^2(y)>1/(xy)^8=1/(R^2+\tau ^2)^4`$, once integrated over the relative time $`\tau =x_0y_0`$, leads the result $`1/R^7`$.
The condition (ii) can be understood for complex systems like atoms or hadrons in the following way: the first dipole emission excites the system from (usually an S-wave) ground state to (usually a P-wave) excited state, while the second dipole emission returns the system back. The energies of the intermediate state sets the characteristic life-time $`\tau _01/(E_PE_S)`$.
However for static dipoles the situation is different in QED and QCD. In QED the emission times of two exchanged quanta are independent, but in QCD they are not. Even a static dipole can change its color degrees of freedom. Because different total color states of the dipole have different energies, thanks to the Coulomb interaction, we again have excited intermediate states. Therefore the characteristic time is determined by the difference in Coulomb energy between the singlet and octet states
$`1/\tau _0=\mathrm{\Delta }E=(3\alpha _s/2)/d.`$ (2)
Although the dipoles may be $`small`$ $`dR`$, this time may still be long because in the perturbative domain the coupling constant is small $`g^2(d)1`$. As a result, there are two different regimes, when the distance R is large (i) $`R\tau _0`$ or small (ii) $`R\tau _0`$. In the former case again the power is 7 and the polarizability is<sup>1</sup><sup>1</sup>1Amusingly, the result is just the volume of a sphere of radius d, from which the perturbative coupling constant $`g`$ dropped out. $`\alpha =4\pi d^3/3`$. The latter case is the Van-der-Walls domain.
### B RGE analysis of the dipole-dipole potential
On general grounds, the potential between two interacting dipoles can be shown to obey the following equation
$$\alpha _s\frac{𝒱(b)}{\alpha _s}=\frac{1}{2}d^3x\mathrm{Tr}F^2(x)_b$$
(3)
where $`\alpha _s`$ is the QCD running coupling and the averaging in (3) is carried in the presence of the two static dipoles a distance $`b`$ apart. Generically,
$$𝒱(b)𝒱(b,a,\mu ,\alpha _s)\mu (\mu a)^\kappa F(\mu b,\alpha _s)$$
(4)
where $`\mu `$ is the renormalization scale. Hence,
$$\frac{𝒱}{\alpha _s}=\frac{1}{\beta }\left((\kappa +1)𝒱+b\frac{𝒱}{b}\right)$$
(5)
with $`\beta =d\alpha _s/d\mathrm{ln}\mu `$ is the QCD beta function. Inserting (5) into (3) yields
$$(\kappa +1)𝒱+b\frac{𝒱}{b}=\frac{\beta }{2\alpha _s}d^3xTrF^2(x)_b,$$
(6)
which is the RGE equation satisfied by the dipole-dipole potential. At large separations we may assume the dipole-dipole potential in quenched QCD to follow like a power law, i.e.
$$𝒱(b)\mu (\mu a)^\kappa (\mu b)^\gamma ,$$
(7)
turning (6) into an algebraic equation
$$(1+\gamma +\kappa )𝒱(b)=\frac{\beta }{2\alpha _s}d^3x\mathrm{Tr}F^2(x)_b.$$
(8)
Alternatively, the potential between two dipoles is a measure of the energy density in the presence of two-dipoles
$$𝒱(b)=d^3x\mathrm{\Theta }_{00}(x)_b.$$
(9)
The combination of the RGE equation (8) and the definition (9) yields a constraint between the exponents $`\kappa `$ and $`\gamma `$ in (7) asymptotically, namely
$$\gamma =1\kappa +\frac{\beta }{\alpha _s}\frac{1R}{1+R}$$
(10)
with
$$R=\frac{d^3xB^2(x)_b}{d^3xE^2(x)_b}$$
(11)
a measure of the magnetic-to-electric ratio in the configuration composed of two static dipoles a distance $`b`$ away from each other. For a self-dual field $`R=1`$ and $`\gamma =1\kappa `$ if the asymptotic (7) is assumed.
## III Perturbative Scattering In Euclidean Geometry
### A Quark-Quark scattering
Generically, we will refer to quark-quark scattering as
$`Q_A(p_1)+Q_B(p_2)Q_C(k_1)+Q_D(k_2)`$ (12)
We denote by $`AB`$ and $`CD`$ respectively, the incoming and outgoing color and spin of the quarks (polarization for gluons). Using the eikonal approximation and LSZ reduction, the scattering amplitude $`𝒯`$ for quark-quark scattering reads
$`𝒯_{AB,CD}(s,t)2is{\displaystyle d^2be^{iq_{}b}}`$ (13)
$`\times <(𝐖_1(b)\mathrm{𝟏})_{AC}(𝐖_2(0)\mathrm{𝟏})_{BD}>`$ (14)
where as usual $`s=(p_1+p_2)^2`$, $`t=(p_1k_1)^2`$, $`s+t+u=4m^2`$ and
$`𝐖_{1,2}(b)=𝐏_c\mathrm{exp}\left(ig{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑\tau A(b+v_{1,2}\tau )v_{1,2}\right)`$ (15)
The 2-dimensional integral in (14) is over the impact parameter $`b`$ with $`t=q_{}^2`$, and the averaging is over the gauge configurations using the QCD action. The color bearing amplitude (14) allows for scattering into a singlet or an octet configuration, i.e.
$`𝒯=𝒯_1\mathbf{\hspace{0.17em}\hspace{0.17em}1}\mathrm{𝟏}+𝒯_{N_c^21}(\tau ^a\tau ^a)`$ (16)
following the decomposition $`N_cN_c=1(N_c^21)`$. For gluon-gluon scattering the lines are doubled in color space (adjoint representation) and further gauge-invariant contractions are possible. For quark-quark scattering the singlet exchange in t-channel is 0<sup>+</sup> (pomeron) while for quark-antiquark it is 0<sup>-</sup> (odderon) as the two differ by charge conjugation.
A quark with large momentum $`p`$ travels on a straight line with 4-velocity $`\dot{x}=v=p/m`$ and $`v^2=1`$. In the the eikonal approximation an ordinary quark transmutes to a scalar quark. The argument applies to any charged particle in a background gluon field, with the following amendments: for anti-quarks the 4-velocity $`v`$ is reversed in the Wilson line and for gluons the Wilson lines are in the adjoint representation. Quark-quark scattering can be also extended to quark-antiquark, gluon-gluon or scalar-scalar scattering. For quark-antiquark scattering the elastic amplitude dominates at large $`\sqrt{s}`$ since the annihilation part is down by $`\sqrt{t/s}`$.
It can be described in Minkowski geometry in the CM frame with $`p_1/m=(\mathrm{cosh}\gamma /2,\mathrm{sinh}\gamma /2,0_{})`$ and $`p_2/m=(\mathrm{cosh}\gamma /2,\mathrm{sinh}\gamma /2,0_{})`$ with the rapidity $`\gamma `$ defined through $`\mathrm{cosh}\gamma /2=\sqrt{s}/2m`$. For $`sm^2`$ the rapidity gap between the receding scatterers become large with $`\gamma \mathrm{log}(s/m^2)`$. The momentum transfer between the scatterers is $`q=p_1k_1`$, with $`q_0q_3t/\sqrt{s}`$ and $`q_{}^2=tu/(s4m^2)t`$. Hence $`q=(0,0,q_{})`$ with $`q^2=q_{}^2=t`$. Although the partons or dipoles change their velocities after scattering, this change is small for $`st`$. This is the kinematical assumption behind the use of the eikonal approximation.
In Euclidean geometry, the kinematics is fixed by noting that the Lorenz contraction factor translates to
$$\mathrm{cosh}\gamma =\frac{1}{\sqrt{1v^2}}=\frac{s}{2m^2}1\mathrm{cos}\theta .$$
(17)
Scattering at high-energy in Minkowski geometry follows from scattering in Euclidean geometry by analytically continuing $`\theta i\gamma `$ in the regime $`\gamma \mathrm{log}(s/m^2)1`$ . It is sufficient to analyze the scattering for $`p_1/m=(1,0,0_{})`$, $`p_2/m=(\mathrm{cos}\theta ,\mathrm{sin}\theta ,\mathrm{\hspace{0.17em}0}_{})`$, $`q=(0,0,q_{})`$ and $`b=(0,0,b_{})`$. The Minkowski scattering amplitude at high-energy can be altogether continued to Euclidean geometry through
$`𝒯_{AB,CD}(\theta ,q)4m^2\mathrm{sin}\theta {\displaystyle d^2be^{iq_{}b}}`$ (18)
$`\times <(𝐖(\theta ,b)\mathrm{𝟏})_{AC}(𝐖(0,0)\mathrm{𝟏})_{BD}>`$ (19)
where
$`𝐖(b,\theta )=𝐏_c\mathrm{exp}\left(ig{\displaystyle _\theta }𝑑\tau A(b+v\tau )v\right)`$ (20)
with $`v=p/m`$. The line integral in (20) is over a straight-line sloped at an angle $`\theta `$ away from the vertical.
In QCD perturbation theory, different time-ordering contributions to quark-quark scattering are shown in Fig. 1 to order $`g^2`$. They contribute to the T-matrix as $`𝒯=2𝒯_1+2𝒯_2`$ with ($`T\mathrm{}`$<sup>2</sup><sup>2</sup>2The color factors can be restored trivially.
$`𝒯_1(\theta ,b)=`$ $`{\displaystyle \frac{g^2}{4\pi ^2}}{\displaystyle _0^T}𝑑\tau _1{\displaystyle _0^T}𝑑\tau _2`$ (22)
$`\times {\displaystyle \frac{\mathrm{cos}\theta }{(\tau _1\tau _2\mathrm{cos}\theta )^2+\tau _2^2\mathrm{sin}^2\theta +b^2}}`$
$`=`$ $`{\displaystyle \frac{\theta }{\mathrm{tan}\theta }}{\displaystyle \frac{g^2}{4\pi ^2}}\mathrm{log}\left({\displaystyle \frac{T}{b}}\right).`$ (23)
and
$`𝒯_2(\theta ,b)=`$ $`{\displaystyle \frac{g^2}{4\pi ^2}}{\displaystyle _0^T}𝑑\tau _1{\displaystyle _T^0}𝑑\tau _2`$ (25)
$`\times {\displaystyle \frac{\mathrm{cos}\theta }{(\tau _1\tau _2\mathrm{cos}\theta )^2+\tau _2^2\mathrm{sin}^2\theta +b^2}}`$
$`=`$ $`{\displaystyle \frac{(\pi \theta )}{\mathrm{tan}\theta }}{\displaystyle \frac{g^2}{4\pi ^2}}\mathrm{log}\left({\displaystyle \frac{T}{b}}\right).`$ (26)
with $`𝒯_2(\theta ,b)=𝒯_1(\pi \theta ,b)`$ as expected from geometry<sup>3</sup><sup>3</sup>3The reader may be puzzled by why we are emphasizing this simple point. We note that for more involved multi-gluon processes this cancellation is spoiled by color factors and powers of the angle survive in the answer: after the analytic continuation to Minkowski space these powers become powers of rapidity. They exponentiate and produce powers of the collision energy characteristic of Reggeon behavior (to be described elsewhere).. We note that the overall linear dependence in $`\theta `$ reflects on the range of the gluon exchanged in rapidity space caused by our ordering in time. This dependence becomes $`\theta +(\pi \theta )=\pi `$ in the sum $`𝒯`$, i.e.
$$𝒯(\theta ,b)=\frac{g^2}{2\pi ^2}\frac{\pi }{\mathrm{tan}\theta }\mathrm{log}\left(\frac{T}{b}\right)$$
(27)
as the ordering is unrestricted between $`0`$ and $`\pi `$. All gluons between the spatial distance $`b`$ and $`T`$ are also exchanged, hence the infrared sensitivity of the quark-quark scattering amplitude in perturbation theory. This sensitivity drops from the cross section (see below). We note that the order $`g^2`$ contribution to (27) is of order $`s^0`$ after analytical continuation, in agreement with the general energy-spin assignment for vector exchange. We recall that the expected behavior is $`s^{J1}`$ for a spin-J exchange.
The contribution of (27) to $`𝒯`$ follows after integrating over the impact parameter $`b`$. The result in Euclidean geometry is
$`𝒯(\theta ,q)=`$ $`4m^2\mathrm{sin}\theta {\displaystyle d^2be^{iqb}𝒯(\theta ,b)}`$ (28)
$`=`$ $`\mathrm{cos}\theta {\displaystyle \frac{g^2}{2}}{\displaystyle \frac{4m^2}{q^2}}{\displaystyle _0^{\mathrm{}}}𝑑xJ_0(x)\mathrm{log}x.`$ (29)
which can be translated into Minkowski geometry by analytical continuation through $`\theta i\gamma `$ with $`q^2=t`$. In both geometries, $`𝒯`$ is purely real and divergent as $`t0`$, leading to a differential cross section of the order of $`d\sigma /dtg^4/t^2`$ with a corresponding divergent Coulomb cross section $`\sigma g^4/(t_{min})`$. In perturbation theory, the $`𝒯`$ matrix acquires absorptive parts and turns complex to higher-order, i.e. $`𝒯=g^2/t+ig^4/t+\mathrm{}`$. The Euclidean perturbative analysis can be carried out to higher orders as well, in close analogy with analytically continued Feynman diagrams .
### B Dipole-Dipole Scattering
We now consider dipole-dipole scattering
$`D_A(p_1)+D_B(p_2)D_C(k_1)+D_D(k_2)`$ (30)
emphasizing its color degrees of freedom. For simplicity we will assume both dipoles to have sizes $`d`$, and (in this section) average over their orientations. For pedagogical reasons, we start with a “naive” Euclidean approach at large impact parameter b, analogous to the calculation of the dipole-dipole potential above. This would be shown to lead to an incorrect answer for the high energy scattering amplitude. The reason will be given below along with the correct answer.
We will assume that the impact parameter $`b`$ is large in comparison to the typical time characteristic of the Coulomb interaction inside the dipole, i.e. $`b\tau _0d/g^2`$. In the elastic dipole-dipole amplitude the dipoles remain color-neutral, and may argue that the leading order is 2-gluon dominated. In analogy to the potential, one may rely on the Coulomb interaction inside the dipole to write the dipole-dipole effective vertex in the form
$$S_{eff}=\alpha _E_{\mathrm{}}^{\mathrm{}}𝑑\tau \dot{x}_\mu \dot{x}_\nu F_{\mu \alpha }^aF_{\nu \alpha }^a(x)$$
(31)
where the electric polarizability $`\alpha _E(gd)^2/`$ with $`g^2/d`$ its Rydberg energy . (Higher order operators are suppressed by powers of the dipole size $`d`$.)
In leading order in the dipole size, the scattering amplitude then reduces to
$`𝒯(\theta ,b)`$ $`\alpha _{E}^{}{}_{}{}^{2}\dot{x}_{1\mu }\dot{x}_{1\nu }\dot{x}_{2\lambda }\dot{x}_{2\sigma }`$ (33)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}d\tau _1d\tau _2<F_{\mu \alpha }^aF_{\nu \alpha }^a(x_1)F_{\lambda \beta }^bF_{\sigma \beta }^b(x_2)>,`$
with $`x_1=v_1\tau _1`$ and $`x_2=v_2\tau _2+b`$. The last expectation value can be unwound using free field theory to obtain
$$𝒯(\theta ,b)\frac{(N_c^21)}{\pi ^3}\frac{\alpha _E^2}{b^6}\left(\frac{11}{25}\frac{1}{\mathrm{sin}\theta }+\frac{8}{5}\frac{\mathrm{cos}^2\theta }{\mathrm{sin}\theta }\right).$$
(34)
We note that the result (34) diverges as $`\theta 0`$. For the case $`\theta =0`$, we obtain the Casimir-Polder-type amplitude
$$𝒯(0,b)\frac{(N_c^21)}{\pi ^3}\frac{T\alpha _E^2}{b^7}\frac{23}{8}$$
(35)
with $`T\mathrm{}`$, which differs from the $`\theta 0`$ by the occurrence of the infrared sensitive factor $`T/b`$.
The analytical continuation of (34) to Minkowski space shows that the first contribution is of order $`1/s`$, while the second contribution is of order $`s`$. This implies that the total cross section is unbound, i.e. $`\sigma s`$, which is clearly incorrect. Indeed, on physical grounds the total cross section should be constant at large $`s`$. In Minkowski space it is easy to understand what went wrong. The electric field of a boosted dipole looks like a Lorenz contracted disk with a very small longitudinal width $`b/\mathrm{ch}yb`$. Clearly, at high energy the interaction time of two dipoles is of this order of magnitude, which is much shorter than the Coulomb time $`\tau _0`$. During this short time, the color rotation induced by the Coulomb interaction can be ignored. Therefore, the use of (31) in the form of a local 2-gluon exchange is incorrect <sup>4</sup><sup>4</sup>4Note that the result (35) is not based on this approximation, and therefore is still valid.. This point is actually missed in the Euclidean formulation as the Lorenz factor is $`\mathrm{cos}\theta 1`$. Although any particular integral can be analytically continued from Euclidean to Minkowski space, kinematical approximations can only be inferred from the Minkowski domain where all parameters have their physical values. This will be understood throughout.
So ignoring the Coulomb interaction and using the eikonal approximation, LSZ reduction and the analytical continuation discussed above, we can write the dipole-dipole scattering amplitude $`𝒯`$ in Euclidean geometry similarly to (19) with
$`𝐖(\theta ,b)={\displaystyle \frac{1}{N_c}}\mathrm{Tr}\left(𝐏_c\mathrm{exp}\left(ig{\displaystyle _{𝒞_\theta }}𝑑\tau A(x)v\right)\right)`$ (36)
where $`x`$ is an element of $`𝒞_\theta `$. In Euclidean geometry $`𝒞_\theta `$ is a closed rectangular loop of width $`d`$, that is slopped at an angle $`\theta `$ with respect to the vertical direction. To leading order in the dipole-interaction, $`𝒯`$ can be assessed by expanding each Wilson-line (36) in powers of $`g`$, and treating the resulting 2-gluon correlations perturbatively. The result is
$`𝒯(\theta ,b){\displaystyle \frac{N_c^21}{N_c^2}}{\displaystyle \frac{(gd)^4}{32\pi ^2}}{\displaystyle \frac{\mathrm{cotan}^2\theta }{b^4}}`$ (37)
for two identical dipoles $`d_1=d_2=d`$ with polarizations along the impact parameter $`b`$. For small size dipoles, (37) is the dominant contribution to the scattering amplitude. The analytical continuation shows that $`\mathrm{cotan}\theta i`$, leading to a finite total cross section as expected.
## IV Instanton effects on the potentials
### A Generalities
Instantons are self-dual solutions to the classical Yang-Mills equations in vacuum originally discovered in ref. . They are classical paths describing tunneling between topologically inequivalent vacua of the gauge theory. In QCD, instantons were argued to be responsible for observable phenomena such as the resolution of the U(1) problem (large $`\eta ^{}`$ mass) and the spontaneous breaking of chiral symmetry . The interacting instanton liquid model (IILM) has been shown to reproduce multiple correlation functions, including hadronic spectra and coupling constants (for a review see ).
Instantons are also commonly used in other gauge theories, especially in supersymmetric gauge theories where supersymmetry makes their effects dominant in the non-perturbative regime. Indeed, some exact results (such as the effective low energy Lagrangian for $`N`$=2 supersymmetric theories derived by Seiberg and Witten and also the AdS/CFT correspondence suggested by Maldacena <sup>5</sup><sup>5</sup>5In fact, the 5-dimensional anti-de Sitter space emerges from the space of the instanton collective coordinates (the center position and size $`d^4zd\rho /\rho ^5`$) which will be extensively used for averaging below. for the N=4 super-conformal theory) can be exactly reproduced using exclusively the instanton calculus developed in .
For the purpose of this paper the topology of instantons is not important: heavy quarks do not interact with fermionic zero modes, and high energy quarks for all purposes behave as heavy quarks. What is important instead is the following technical point: in the instanton field the path-ordered exponents can be evaluated $`analytically`$, since the color phase rotations take place around the same axis for a fixed path (the instanton is a hedgehog in color-space). The self-duality of the instanton field will also have an effect on some of our results. Once a path-ordered exponent is evaluated in the one-instanton field, the vacuum averages follow through the instanton ensemble average representing the QCD vacuum (dilute phase). This includes averaging over the instanton center-position $`z_\mu `$ and size $`\rho `$. Specifically, we will use the measure
$`dn=d\rho d^4z{\displaystyle \frac{D(\rho )}{\rho ^5}}`$ (38)
for both instantons and anti-instantons. The integral over $`z`$ can be sometimes carried out analytically, but most of the times will be done numerically. The understanding of the instanton size distribution $`D(\rho )`$ remains an open problem. Naive semi-classical results suggest :
$`D_0(\rho )`$ $`C_{N_c}({\displaystyle \frac{8\pi ^2}{g^2(\rho )}})^{2N_c}\mathrm{exp}({\displaystyle \frac{8\pi ^2}{g^2(\rho )}})`$ (39)
$``$ $`(\rho \mathrm{\Lambda })^{(11/3)N_c(2/3)N_f}`$ (40)
where $`C_{N_c}`$ is a constant depending on the number of colors $`N_c`$. We have used the asymptotic freedom formula in the exponent to show that this density dramatically grows with the instanton size $`\rho `$. However in the true QCD vacuum instantons and antiinstantons interact with each other and othe quantum fields, so that the real function $`D(\rho )`$ deviates from the semiclassical one for large sizes.
For qualitative estimates we will often use parameters of the instanton liquid model , which assumes that all instantons have the same size
$`dn(\rho )=n_0d^4zd\rho \delta (\rho \rho _0)`$ (41)
where $`n_0`$ is the total instanton (plus anti-instanton) with a typical radius $`\rho _0`$, i.e.
$`n_01fm^4;\rho _01/3fm`$ (42)
These values were deduced from phenomenological data extracted from the QCD sum rules, the topological succeptibility and the chiral condensate long before direct lattice data became available. In Fig. 2 we show a sample of such lattice measurements, together with the parameterization for the instanton suppression suggested in . Specifically,
$`dn(\rho )=dn_0(\rho )e^{2\pi \sigma \rho ^2}`$ (43)
can be used for averaging in any integral over the instanton density. Typically, the string tension $`\sigma (0.440MeV)^2`$, so that
$`<\rho ^2>(0.28fm)^2;`$ (44)
$`<\rho ^4>(0.31fm)^4;`$ (45)
$`<\rho ^5>(0.32fm)^5,`$ (46)
which shows that the difference between the realistic averages and simple powers of $`\rho _0`$ are relatively small. We will ignore these differences below.
In the analysis to follow, the parameters capturing the instanton physics will appear as two dimensionless quantities: (i) a small diluteness parameter and (ii) a large action of an instanton (per $`\mathrm{}`$):
$`n_0\rho _0^4\left({\displaystyle \frac{1}{3}}\right)^4S_0={\displaystyle \frac{8\pi ^2}{g^2(\rho _0)}}(1015)`$ (47)
The small factor is a penalty for finding the instanton, and the large factor is a classical enhancement relative to perturbation theory. Their interplay would cause particular effects to be parametrically large or small.
### B Static quarks
At the one instanton level, the various potentials for a static quark-antiquark potential have been assessed long ago , including the spin-dependent part. We will briefly review this assessment for completeness. We recall that the various components of the potential follow from the rectangular $`T\times R`$ Wilson loop
$`V(R)={\displaystyle \frac{1}{T}}\underset{T\mathrm{}}{lim}\mathrm{ln}<𝐖(T,R)>`$ (48)
evaluated in a classical instanton field, after averaging over the instanton position. In the Wilson loop, the path-ordered exponents $`Pexp(igA_\mu 𝑑x_\mu )`$ can be evaluated analytically as the instanton locks the color orientation to space. Indeed, the static potentials involve $`A_0^a\eta _{0,\nu }^a(xz)_\nu (xz)^a`$ where $`(\stackrel{}{x}\stackrel{}{z})`$ refers to the distance between the quark position and the 3-d coordinate of the instanton center<sup>6</sup><sup>6</sup>6The time position of the instanton $`z_4`$ is irrelevant.. The resulting color rotation angle $`\alpha `$ and the unit vector $`n^a`$ around which the rotation takes place are defined through
$`𝐖=`$ $`\mathrm{exp}\left(i\pi {\displaystyle \frac{\tau ^a(z_ar_a)}{((r_az_a)^2+\rho ^2)^{1/2}}}\right)`$ (49)
$`=`$ $`\mathrm{exp}(i\pi \tau ^an_a\alpha ).`$ (50)
If all relevant distances are comparable, $`|r_az_a|\rho _0`$, the rotation angle is O(1), showing that the expansion in field strength is in general not justified. For a small-size dipole, the potential is small $`V(R0)R^2`$, since the path-ordered lines in $`𝐖`$ are close enough to cause partial cancellation. However when $`R\rho _0`$, and both path-ordered lines happen to be on the opposite sides of the instanton center, the color rotations on both lines adds up and the potential becomes roughly linear in R and more sizable. Finally, when the dipole is too large, the potential saturates<sup>7</sup><sup>7</sup>7In ref. one of us has noticed that this behavior is surprisingly similar to that experimentally observed in deep inelastic scattering, if $`Q^2`$ dependence of structure functions is treated as dependence of the cross section on the dipole size.
The quark-antiquark potential calculated in can be expressed as
$`V(R)={\displaystyle 𝑑n(\rho )\rho ^3F(R/\rho )}`$ (51)
where the dimensionless function F is defined as
$`F={\displaystyle \frac{d^3z}{N_c\rho ^3}\mathrm{Tr}(1𝐖_1𝐖_2^{})}.`$ (52)
The trace-part of the integrand is
$`2\left(1\mathrm{cos}\alpha _1\mathrm{cos}\alpha _2\stackrel{}{n}_1\stackrel{}{n}_2\mathrm{sin}\alpha _1\mathrm{sin}\alpha _2\right).`$ (53)
where the angles $`\alpha _i`$ and vectors $`n_i^a`$ are defined in (49). This function is shown in Fig. 3 a. In order to emphasize the small-R “dipole limit” $`V(R)R^2`$ (to be important for what follows), we have also plotted the ratio of this function to its dipole limit in Fig. 3 b. One can see that the dipole approximation has an unexpectedly large range of applicability: this ratio does not change appreciably (less than 25%) till $`R\rho _0`$. One may expect similar accuracy of the dipole approximation in other applications to be discussed.
For large R the potential goes to a constant plus a Coulomb term
$`V(R\mathrm{})=37{\displaystyle \frac{d\rho }{\rho ^2}D(\rho )}{\displaystyle \frac{4\pi ^3}{3R}}{\displaystyle \frac{d\rho }{\rho }}D(\rho )+\mathrm{}`$ (54)
which can be interpreted as the instanton contribution to the $`mass`$ and $`charge`$ renormalization, respectively. It is instructive to compare the magnitude of the latter to the perturbative potential, through
$`{\displaystyle \frac{V_{inst}}{V_{pert}}}={\displaystyle \frac{\pi ^2}{2}}(n_0\rho _0^4)({\displaystyle \frac{8\pi ^2}{g^2(\rho _0)}}),`$ (55)
with $`V_{pert}=4\alpha _s/3R`$. The ratio is the product of the diluteness parameter <sup>8</sup><sup>8</sup>8 The coefficient in front of $`\pi ^2\rho ^4/2`$ happens to be the volume of a 4-sphere. (the fraction of space-time occupied by instantons) times the classical enhancement through the instanton action (per $`\mathrm{}`$). Using the phenomenological parameters discussed above, we observe that the diluteness is compensated by the classical enhancement, so that the instanton corrections at $`R\rho _0`$ are actually comparable to the perturbative Coulomb effect.
However, instantons are not the only non-perturbative effects contributing to the static quark-antiquark potential. At large $`R`$ confinement in the form of a QCD string with $`V_{conf}\sigma R`$ dominates. In fact, already for $`R\rho _00.3fm`$ confinement is dominant, with the instanton-induced potential accounting for only 10-15% <sup>9</sup><sup>9</sup>9The claim made in , that instanton effects account for the confining potential is incorrect.. For a detailed study of these issues at the multi-instanton level, one can consult refs for a numerical analysis and for analytical results.
### C Static dipoles
Unlike the quark-antiquark potential, the dipole-dipole potential is insensitive to confinement, and the instanton-induced interaction may be easier to identify. In the latter case, we will consider two cases where the characteristic time within the dipole is either (i) $`short`$ $`\tau _0d/g^2\rho `$ or (ii) $`long`$ $`\tau _0\rho _0`$ in comparison to the instanton size. These two cases translate to a magnitude of the dipole field $`A_0g/d`$ which is large (i) or small (ii) in comparison to that of the instanton field $`A_\mu 1/g\rho `$.
In the case (i) the static potential can be written in terms of the polarizabilities, and the correlator of gluo-electric fields
$`V(R)=\alpha _1\alpha _2{\displaystyle 𝑑\tau <\stackrel{}{E}^2(\tau ,R)\stackrel{}{E}^2(0,0)>}`$ (56)
This field strength correlator can be evaluated by substituting the expression for the instanton field
$`\stackrel{}{E}^2(x)=\stackrel{}{B}^2(x)={\displaystyle \frac{96\rho ^4}{g^2}}{\displaystyle \frac{1}{((xz)^2+\rho ^2)^4}}`$ (57)
The averaging of the correlator over the location of the instanton position $`z`$ can carried out analytically
$`<(gG_{\mu \nu }^a(x))^2(gG_{\mu \nu }^a(0))^2>=`$ (58)
$`{\displaystyle \frac{384g^4}{\pi ^4x^8}}+(n_0\rho _0^4)\mathrm{\Pi }_{inst}(x/\rho )/\rho ^8,`$ (59)
where the last term was added to account for the perturbative contribution. The dimensionless function describing the instanton contribution is
$`\mathrm{\Pi }_{inst}(y)=`$ $`{\displaystyle \frac{12288\pi ^2}{y^6(y^2+4)^5}}`$ (63)
$`\times (y^8+28y^694y^4160y^2120`$
$`+{\displaystyle \frac{240}{y\sqrt{y^2+4}}}`$
$`\times (y^6+2y^4+3y^2+2)\mathrm{arcsinh}(y/2)).`$
Its behavior is shown in Fig.4a. Its ratio to the perturbative contribution to the same correlator (for $`g=2`$ or $`\alpha _s=0.32`$) is shown in fig.4b. As expected, it is small at small distances $`x\rho _0`$. At large distances, the instanton-induced contribution has the same behavior $`\mathrm{\Pi }_{inst}1/R^8`$, as the perturbative one. Furthermore, the ratio of the two is about 30, much more than the “instanton-induced charge renormalization” (55) we discussed in the preceding subsection. About the same is found in the potentials themselves (the correlator integrated over the time difference) as shown in Fig. 4 d. The perturbative behavior is dominated by $`two`$ gluons rather than one, and therefore the instanton effect occurs with a classical enhancement $`squared`$:
$`{\displaystyle \frac{V_{inst}}{V_{pert}}}(n_0\rho _0^4)\left({\displaystyle \frac{8\pi ^2}{g^2(\rho _0)}}\right)^2`$ (64)
This feature implies that instanton effects are much more important for dipole-dipole interactions at $`R\rho _00.3`$ fm than the perturbative Casimir-Polder effects. We will argue below that this is generic for all processes demanding multi-gluon exchanges, and that instanton-induced processes can become dominant in this case.
In the case ii), the dipoles can be considered quasi-static in time, $`\tau _0d/g^2\rho _0`$, and the time evolution of the color degrees of freedom due to the Coulomb interaction can be ignored. In other words, the dynamics is driven entirely by the instanton field. The potential between two dipoles is now
$`V_{dd}(R)={\displaystyle 𝑑n(\rho )\rho ^3F_{dd}(R/\rho )}`$ (65)
with
$`F_{dd}={\displaystyle \frac{d^3z}{N_c\rho ^3}(1\mathrm{Tr}𝐖_1\mathrm{Tr}𝐖_2)}`$ (66)
Here W are rectangular Wilson $`loops`$ for each dipole, traced separately. Averaging over the instanton position can be done numerically. The results are shown in Fig.5. The outcome is proportional to $`d_1^2d_2^2`$ (dipole-moments) rather than $`\alpha _1\alpha _2`$ (electric polarizabilities), when $`d`$ is reasonably small in comparison to $`\rho _0`$. The large distance potential is few % that of $`V(R)d_1^2d_2^2\rho _0^2/R^7`$. Note that it is larger than the perturbative one since $`\rho _0^2`$ is assumed to be much larger than $`d_1d_2`$, but both answers have the same (zeroth) power of g.
In general the dipole-dipole potential cannot be approximated by the correlator of scalars $`E^2`$, as can be checked through its dependence on the relative orientation of the dipoles. Even in the dipole (quadratic) approximations for sufficiently small dipoles ($`d_i\rho `$) one can define 4 invariant functions for the dipole-dipole interaction
$`V(R)=`$ $`d_1^id_1^jd_2^ld_2^m(A(R)\delta _{ij}\delta _{lm}+`$ (69)
$`{\displaystyle \frac{1}{2}}B(R)(n^in^j\delta _{lm}+n^ln^m\delta _{ij})+`$
$`C(R)n^in^l\delta _{jm}+D(R)n^in^jn^ln^m)`$
The first function A(R) accounts for the spin-zero gluonic operator $`\stackrel{}{E}^2`$ discussed at the beginning of this subsection. However, as one can see from fig(5), other functions also contribute. In (a) we compare the xx orientation (or A+B+C+D) with the xy one (or A+B/2) and see a clear difference. In (b) we note the dependence on the rotation angle for one of the dipoles, which shows a clear $`\mathrm{cos}^2\theta `$ behavior expected from the expression above.
## V One-Instanton effect on scattering
### A Quark-Quark scattering
Our first step now is the generalization of (49) to an arbitrary orientation $`\theta `$ of the Wilson line. The analytical continuation to Minkowski space follows from $`\theta iy`$ with $`y`$ identified as the rapidity difference between the receding partons. The untraced and tilted Wilson line in the one-instanton background reads
$$𝐖(\theta ,b)=\mathrm{cos}\alpha i\tau \widehat{n}\mathrm{sin}\alpha $$
(70)
where
$$n^a=𝐑^{ab}\eta _{\mu \nu }^b\dot{x}_\mu (zb)_\nu =𝐑^{ab}𝐧^b$$
(71)
and $`\alpha =\pi \gamma /\sqrt{\gamma ^2+\rho ^2}`$ with
$`\gamma ^2=`$ $`nn=𝐧𝐧`$ (72)
$`=`$ $`(z_4\mathrm{sin}\theta z_3\mathrm{cos}\theta )^2+(bz_{})^2.`$ (73)
The one-instanton contribution to the untraced QQ-scattering amplitude follows from the following correlator
$`𝐖_{AC}(\theta ,b)𝐖_{BD}(0,0)n_0{\displaystyle d^4z}`$ (74)
$`\times (\mathrm{cos}\alpha \mathrm{cos}\underset{¯}{\alpha }\mathbf{\hspace{0.17em}\hspace{0.17em}1}_{AC}\mathbf{\hspace{0.17em}1}_{BD}`$ (75)
$`{\displaystyle \frac{1}{N_c^21}}\widehat{𝐧}\underset{¯}{\overset{^}{𝐧}}\mathrm{sin}\alpha \mathrm{sin}\underset{¯}{\alpha }(\tau ^a)_{AC}(\tau ^a)_{BD}),`$ (76)
where the (under) bar notation means the same as the corresponding un-bar one with $`\theta =0`$ and $`b=0`$.
Furthermore,
$`{\displaystyle \frac{1}{N_c}}\mathrm{Tr}\left(𝐖(\theta ,b)𝐖(0,0)\right)=`$ (77)
$`{\displaystyle \frac{2n_0}{N_c}}{\displaystyle d^4z\left(\mathrm{cos}\alpha \mathrm{cos}\underset{¯}{\alpha }\widehat{𝐧}\underset{¯}{\overset{^}{𝐧}}\mathrm{sin}\alpha \mathrm{sin}\underset{¯}{\alpha }\right)}.`$ (78)
The integrand in (78) can be simplified by changing variable $`(z_4\mathrm{sin}\theta z_3\mathrm{cos}\theta )z_4`$ and dropping the terms that vanish under the z-integration. Hence
$`{\displaystyle \frac{1}{N_c}}\mathrm{Tr}\left(𝐖(\theta ,b)𝐖(0,0)\right)={\displaystyle \frac{2n_0}{N_c}}{\displaystyle d^4z}`$ (79)
$`\left({\displaystyle \frac{1}{\mathrm{sin}\theta }}\mathrm{cos}\stackrel{~}{\alpha }\mathrm{cos}\underset{¯}{\overset{~}{\alpha }}{\displaystyle \frac{1}{\mathrm{tan}\theta }}\mathrm{sin}\stackrel{~}{\alpha }\mathrm{sin}\underset{¯}{\overset{~}{\alpha }}{\displaystyle \frac{z_{}^2z_{}b}{\stackrel{~}{\gamma }\underset{¯}{\overset{~}{\gamma }}}}\right).`$ (80)
The tilde parameters follow from the un-tilde ones by setting $`\theta =\pi /2`$. We note that $`\underset{¯}{\overset{~}{\gamma }}=\underset{¯}{\gamma }=|\stackrel{}{z}|`$. After analytical continuation, the first term produces the elastic amplitude which decays as $`1/s`$ with the energy. The second term corresponds to the color-changing amplitude. It is of order $`s^0`$ and dominates at high energy. Specifically
$`{\displaystyle \frac{1}{N_c}}\mathrm{Tr}\left(𝐖(\theta ,b)𝐖(0,0)\right)=`$ (81)
$`{\displaystyle \frac{2n_0}{N_c}}\left({\displaystyle \frac{1}{\mathrm{sin}\theta }}F_{cc}(b/\rho _0){\displaystyle \frac{1}{\mathrm{tan}\theta }}F_{ss}(b/\rho _0)\right).`$ (82)
We show in Fig.6 the numerical behavior of the two contributions in (82). Note that the second function (which describes color-inelastic collisions and survives in the high energy limit) changes sign, before decreasing as a power law to zero at large b.
### B Dipole-Dipole and Multi-parton Scattering
One can directly generalize the calculation of the quark-quark scattering amplitude to that of any number of partons. For that, we assume that they all move with high energy in some reference frame and opposite direction: in Euclidean space those would propagate along two directions, with parton numbers $`N_1`$ and $`N_2`$ respectively. Any one of them, passing through the instanton field, is rotated in color space by a different angle $`\alpha _i`$ around a different axis $`\stackrel{}{n}_i`$, depending on the shortest distance between its path and the instanton center. Integration over all possible color orientations of the instanton leads then to global color conservation.
Before discussing specific cases in details, let us make a general qualitative statement about such processes. We have found in the previous section that (the color-changing) quark-quark instanton-induced scattering has a finite high energy limit. For perturbative n-gluon exchange a factor of $`\alpha _s^n`$ is paid, while for an instanton mediated scattering a factor of $`n_0\rho _0^4`$ is paid (the price to find the instanton at the right place), no matter how many partons participate. Since the instanton vacuum is dilute, the one-gluon mediated process dominates the instanton one. However, the situation dramatically changes for two or more gluon exchanges: the instanton-induced amplitude is about the same for any number of partons, provided that all of them pass at a distance $`\rho _0`$ from the instanton center.
Now, consider a dipole configuration of size $`d`$ chosen in the transverse plane of a $`\overline{q}q`$ located on a straight-line sloped at an angle $`\theta `$ in Euclidean space. Let $`AA`$ be the initial color of the dipole and $`CD`$ its final color. The Wilson loop with open color for the dipole configuration in the one-instanton background is
$`𝒲_{AA}^{CD}(\theta ,b)=\mathrm{cos}\alpha _{}\mathrm{cos}\alpha _+\mathbf{\hspace{0.17em}1}_{CD}`$ (83)
$`+i\mathrm{cos}\alpha _{}\mathrm{sin}\alpha _+𝐑^{ab}\widehat{𝐧}_+^b(\tau ^a)_{DC}`$ (84)
$`i\mathrm{sin}\alpha _{}\mathrm{cos}\alpha _+𝐑^{ab}\widehat{𝐧}_{}^b(\tau ^a)_{DC}`$ (85)
$`+\mathrm{sin}\alpha _{}\mathrm{sin}\alpha _+𝐑^{ab}𝐑^{cd}\widehat{𝐧}_{}^b\widehat{𝐧}_+^d(\tau ^c\tau ^a)_{DC}.`$ (86)
We have defined
$`\alpha _\pm =`$ $`{\displaystyle \frac{\pi \gamma _\pm }{\sqrt{\gamma _\pm +\rho ^2}}}`$ (87)
$`\gamma _\pm ^2=`$ $`(z_4\mathrm{sin}\theta z_3\mathrm{cos}\theta )^2+(z_{}b\pm {\displaystyle \frac{d}{2}})^2`$ (88)
$`𝐧_+𝐧_{}=`$ $`\left((z_4\mathrm{sin}\theta z_3\mathrm{cos}\theta )^2+(bz_{})^2{\displaystyle \frac{d^2}{4}}\right)`$ (89)
with $`𝐧_\pm 𝐧_\pm =\gamma _\pm ^2`$. The scattering amplitude of an initial dipole through an instanton after averaging over the global color orientation $`𝐑`$ is
$`{\displaystyle \frac{2}{N_c}}\left(\mathrm{cos}\alpha _{}\mathrm{cos}\alpha _++\widehat{𝐧}_{}\widehat{𝐧}_+\mathrm{sin}\alpha _{}\mathrm{sin}\alpha _+\right)\mathbf{\hspace{0.17em}1}_{CD}`$ (90)
which reduces to the color-singlet channel. Specifically,
$`𝒲(\theta ,b)={\displaystyle \frac{2}{N_c}}\left(\mathrm{cos}\alpha _{}\mathrm{cos}\alpha _++\widehat{𝐧}_{}\widehat{𝐧}_+\mathrm{sin}\alpha _{}\mathrm{sin}\alpha _+\right).`$ (91)
(92)
The $`\theta `$ dependence in (90-92) can be readily eliminated by carrying the integration over the instanton position $`z`$ through the same change of variable discussed in the quark-quark scattering, resulting in an amplitude that depends only on $`1/\mathrm{sin}\theta `$. In Minkowski space this translates to $`1/s`$ which vanishes at high energy. Indeed, the dipole-dipole scattering amplitude through a single instanton is
$`𝐖(\theta ,b)𝐖(0,0){\displaystyle \frac{n_0}{\mathrm{sin}\theta }}{\displaystyle d^4z\stackrel{~}{𝒲}(\theta ,b)\stackrel{~}{𝒲}(0,0)}`$ (93)
where $`\stackrel{~}{𝒲}`$ follows from $`𝒲`$ by setting $`\theta =\pi /2`$. Note that in this case $`𝒲(0,0)=\stackrel{~}{𝒲}(0,0)`$.
It is clear from (86) that while scattering through an instanton, the dipole has to flip-color to keep track of the velocity of the quarks in the dipole. The process is color-inelastic and therefore only contributes to the inelastic amplitude to first order in the instanton density $`n_0`$, and to the elastic amplitude to second order in the instanton density, a situation reminiscent of one- and two-gluon exchange.
The dipole-dipole scattering amplitude with open-color in the final state can be constructed by using two dipole configurations as given by (86) with a relative angle $`\theta `$. After averaging over the instanton color-orientations we obtain
$`𝒲_{AA}^{CD}(\theta ,b)𝒲_{A^{}A^{}}^{C^{}D^{}}(0,0)=`$ (94)
$`{\displaystyle \frac{2}{N_c}}𝒲_1\mathbf{\hspace{0.17em}\hspace{0.17em}1}_{CD}\mathbf{\hspace{0.17em}1}_{C^{}D^{}}+{\displaystyle \frac{1}{N_c^21}}𝒲_{N_c^21}(\tau ^a)_{DC}(\tau ^a)_{D^{}C^{}},`$ (95)
(96)
with the singlet part
$`𝒲_1=\mathrm{cos}\alpha _{}\mathrm{cos}\alpha _+\mathrm{cos}\underset{¯}{\alpha }_{}\mathrm{cos}\underset{¯}{\alpha }_+`$ (97)
$`+𝐧_{}𝐧_+\underset{¯}{𝐧}_{}\underset{¯}{𝐧}_+\mathrm{sin}\alpha _{}\mathrm{sin}\alpha _+\mathrm{sin}\underset{¯}{\alpha }_{}\mathrm{sin}\underset{¯}{\alpha }_+`$ (98)
$`+\underset{¯}{𝐧}_{}\underset{¯}{𝐧}_+\mathrm{cos}\alpha _{}\mathrm{cos}\alpha _+\mathrm{sin}\underset{¯}{\alpha }_{}\mathrm{sin}\underset{¯}{\alpha }_+`$ (99)
$`+𝐧_{}𝐧_+\mathrm{sin}\alpha _{}\mathrm{sin}\alpha _+\mathrm{cos}\underset{¯}{\alpha }_{}\mathrm{cos}\underset{¯}{\alpha }_+,`$ (100)
and the octet part
$`𝒲_{N_c^21}=\mathrm{cos}\alpha _{}\mathrm{sin}\alpha _+\mathrm{cos}\underset{¯}{\alpha }_{}\mathrm{sin}\underset{¯}{\alpha }_+𝐧_+\underset{¯}{𝐧}_+`$ (101)
$`\mathrm{sin}\alpha _{}\mathrm{cos}\alpha _+\mathrm{sin}\underset{¯}{\alpha }_{}\mathrm{cos}\underset{¯}{\alpha }_+𝐧_{}\underset{¯}{𝐧}_{}`$ (102)
$`+\mathrm{cos}\alpha _{}\mathrm{sin}\alpha _+\mathrm{sin}\underset{¯}{\alpha }_{}\mathrm{cos}\underset{¯}{\alpha }_+𝐧_+\underset{¯}{𝐧}_{}`$ (103)
$`+\mathrm{sin}\alpha _{}\mathrm{cos}\alpha _+\mathrm{cos}\underset{¯}{\alpha }_{}\mathrm{sin}\underset{¯}{\alpha }_+𝐧_{}\underset{¯}{𝐧}_+`$ (104)
$`\mathrm{sin}\alpha _{}\mathrm{sin}\alpha _+\mathrm{sin}\underset{¯}{\alpha }_{}\mathrm{sin}\underset{¯}{\alpha }_+`$ (105)
$`\times (𝐧_{}\underset{¯}{𝐧}_{}𝐧_+\underset{¯}{𝐧}_+𝐧_{}\underset{¯}{𝐧}_+𝐧_+\underset{¯}{𝐧}_{}).`$ (106)
As in the case of quark-quark scattering, the (color) elastic dipole-dipole amplitude scales as $`1/\mathrm{sin}\theta `$ and vanishes at high energy after analytical continuation. However the (color) inelastic part of the amplitude is not. After performing the shift of variables described before, the $`\theta `$ dependence drops from all the angles $`\alpha `$. There is a remaining $`\theta `$ dependence in the four combinations $`𝐧\underset{¯}{𝐧}`$. In general, the $`\theta `$ dependence in the latter is linear in $`\mathrm{sin}\theta `$ or $`\mathrm{cos}\theta `$, and one may worry that the last term in (106) may involve higher powers of the trigonometric functions, which would yield to an unphysical cross section growing as $`s`$ after analytical continuation. We have checked that this is not the case, since
$`𝐧_{}\underset{¯}{𝐧}_{}𝐧_+\underset{¯}{𝐧}_+𝐧_{}\underset{¯}{𝐧}_+𝐧_+\underset{¯}{𝐧}_{}d^2\left(z_2^2\mathrm{cos}\theta z_3z_4^{}\right)`$ (107)
where $`z_4^{}`$ is the new $`z_4`$ after the change of variable. Moreover, the $`\mathrm{cos}\theta `$ term drops in the integral over z (odd under $`z_3z_3`$), making this contribution to (106) subleading at high-energy after analytical continuation <sup>10</sup><sup>10</sup>10This cancellation is not generic. Indeed, the square of this contribution would be leading.. Finally, we note that all $`\mathrm{sin}\theta `$ contributions in (106) drop following similar parity considerations. As a result, the pertinent octet contribution to the scattering amplitude is proportional to $`\mathrm{cotan}\theta `$ which is $`1/i\mathrm{tan}y=1/iv`$ after analytical continuation.
We have assessed numerically the function
$`F_{N_c^21}({\displaystyle \frac{b}{\rho }}_0,{\displaystyle \frac{d}{\rho }}_0)={\displaystyle \frac{n_0}{\mathrm{cos}\theta }}{\displaystyle d^4z𝒲_{N_c^21}}`$ (108)
which is shown in Fig. 7 for different dipole sizes. We find that the dipole approximation scaling $`F_{N_c^21}d^2`$ works well, even for sizes as large as the instanton size $`d=\rho _0`$.
## VI Two-instanton effect
We have shown above that the instanton contribution at large $`s`$ but small $`t`$ behaves in a way similar to one-gluon exchange: only color-inelastic channels survive in the high energy limit. This means that the contribution to the total cross section appears in the amplitude squared, leading naturally to the concept of two-instanton exchange. The latter contribution to each Wilson-line is more involved. To streamline the discussion we will present the analysis of the two instanton contribution to the differential cross-section of quark-quark scattering at high energy. Similar considerations apply to dipole-dipole scattering as we briefly mention at the end of this section. Indeed, for the quark-quark scattering, unitarity implies that the two-instanton contribution to the differential cross section is
$$\frac{d\sigma }{dt}\frac{1}{s^2}\underset{CD}{}\left|𝒯_{AC}^{BD}\right|^2,$$
(109)
with the averaging over the initial colors $`A,B`$ understood. Inserting (14) after the substitution (76), we obtain
$$\frac{d\sigma }{dt}\left(\frac{4n_0}{N_c}\right)^2𝑑b𝑑b^{}e^{iq(bb^{})}\left(𝐉+\frac{1}{(N_c^21)}𝐊\right)$$
(110)
with
$`𝐉=`$ $`{\displaystyle d^4z(\mathrm{cos}\alpha 1)(\mathrm{cos}\underset{¯}{\alpha }1)}`$ (112)
$`\times {\displaystyle }d^4z^{}(\mathrm{cos}\alpha ^{}1)(\mathrm{cos}\underset{¯}{\alpha }^{}1)`$
$`𝐊=`$ $`{\displaystyle d^4z\widehat{𝐧}\underset{¯}{\overset{^}{𝐧}}\mathrm{sin}\alpha \mathrm{sin}\underset{¯}{\alpha }}`$ (114)
$`\times {\displaystyle }d^4z^{}\widehat{𝐧}^{}\underset{¯}{\overset{^}{𝐧}}^{}\mathrm{sin}\alpha ^{}\mathrm{sin}\underset{¯}{\alpha }^{}.`$
The primed variables follow from the unprimed ones through the substitution $`z,bz^{},b^{}`$. For large $`\sqrt{s}`$, $`𝐉(1F_{cc})(1F_{cc}^{})/s^2`$ <sup>11</sup><sup>11</sup>11Up to self-energies. and $`𝐊=F_{ss}F_{ss}^{}`$, so that
$`{\displaystyle \frac{d\sigma }{dt}}{\displaystyle \frac{16n_0^2}{N_c^2(N_c^21)}}\left|{\displaystyle 𝑑be^{iqb}F_{ss}\left(\frac{b}{\rho _0}\right)}\right|^2.`$ (115)
In particular, the forward scattering amplitude in the two-instanton approximation is
$`\sigma (t=0)`$ $`{\displaystyle \frac{16n_0^2}{N_c^2(N_c^21)}}`$ (117)
$`\times {\displaystyle _0^{\mathrm{}}}dq_{}^2|{\displaystyle }dbe^{iqb}F_{ss}\left({\displaystyle \frac{b}{\rho _0}}\right)|^2,`$
which is finite at large $`\sqrt{s}`$. Hence, for forward scattering partons in the instanton vacuum model, we have
$`\sigma _{qq}(n_0\rho _0^4)^2\rho _0^2`$ (118)
Clearly, the present analysis generalizes to the dipole-dipole scattering amplitude by using (96) instead of (76) and proceeding as before. The outcome is a finite scattering cross section,
$`\sigma (t=0){\displaystyle \frac{4n_0^2}{(N_c^21)}}`$ (119)
$`\times {\displaystyle _0^{\mathrm{}}}dq_{}^2|{\displaystyle }dbe^{iqb}F_{N_c^21}({\displaystyle \frac{b}{\rho _0}},{\displaystyle \frac{d}{\rho _0}})|^2.`$ (120)
Generically, the dipole-dipole cross section relates to the quark-quark cross-section in the forward direction through
$`\sigma _{dd}\sigma _{qq}{\displaystyle \frac{(d_1d_2)^2}{\rho _0^4}}.`$ (121)
It is instructive to compare our instanton results to those developed by Dosch and collaborators in the context of the stochastic vacuum model (SVM). In brief, in the SVM model the Wilson-lines are expanded in powers of the field-strength using a non-Abelian form of Stokes theorem in the Gaussian approximation. A typical hadronic cross section in the SVM model is
$`\sigma <(gG)^2>^2a^{10}𝐅(R_h/a)`$ (122)
where the first factor is the “gluon condensate”, $`a`$ is a fitted correlation length, F is some dimensionless function depending on the hadronic radius $`R_h`$. Although our assumptions and those of are very different regarding the character of the vacuum state, it is amusing to note the agreement between (121) and (122). Indeed, the correlation length $`a`$ of the SVM model is related (and in fact numerically close) to our instanton radius $`\rho _01/3`$ fm, while the gluon condensate $`<(gG)^2>`$ of the SVM model is simply proportional to the instanton density $`n_0`$ in the instanton model.
The most significant difference between these two approaches apart from their dynamical content and the way we have carried the analytical continuation, is the fact that we do not expand in field strength. In fact, in the instanton model there is no parameter which would allow to do so for strong instanton fields. This difference is rather important as it is on it that our conclusion regarding multiple color exchanges is based. (In the SVM model with Wick-theorem-like decomposition, those would be just products of single exchanges, like in pQCD.)
## VII Cross section fluctuations
In so far, we have considered the $`average`$ value of the cross section for a parton in a state of unit probability. However, partons and in general hadrons, are complex quantum mechanical states <sup>12</sup><sup>12</sup>12A truly elementary particle may have only one state and non-fluctuating cross section: it may have diffraction but no inelastic diffraction.. Hence, the quantum system is characterized by some amplitude of probability through its wave function, and its corresponding scattering cross section is probabilistic with a probability distribution $`P(\sigma )`$. This idea was originally suggested by Good and Walker , who emphasized that inelastic diffraction is a way to quantify this distribution via the second moment $`\mathrm{\Delta }\sigma ^2=<(\sigma ^2<\sigma >)^2>`$.
The extraction of this and the next (cubic) moment for the pion and the nucleon using available data has been carried out years later allowing for a reconstruction of the distribution $`P(\sigma )`$. A striking aspect of these results is that the nucleon fluctuations are large and comparable to the pion fluctuations. This outcome does not fit with the constituent quark model where the pion is a 2-body system, and the nucleon is a 3-body system, with more degrees of freedom. One of us had already noticed that this can be a further indication for strongly correlated scalar diquarks in a nucleon. An experimental test for this idea is to measure cross section fluctuations for a decuplet baryon such as $`\mathrm{\Omega }^{}`$. In the latter there are no diquarks, and smaller fluctuations (typical of a 3-body state) are expected. Another aspect of these fluctuations worth mentioning here is that they seem to be maximal for $`\sqrt{s}100`$ GeV, decreasing at very large energies. It supports well the idea that the “most fluctuating” partons are at $`x10^2`$, while at much smaller x one basically approaches a non-fluctuating black disk.
Although in the present paper we have limited our discussion to issues of methodology, it is worth pointing out that the present concept of fluctuations in cross section can be used to discriminate between the instanton effects herein described and other descriptions based either on perturbative multi-gluon exchange or non-perturbative vacuum structures.
Indeed, the standard multi-photon exchange in QED leads to an (eikonalized) exponential scattering amplitude, with Poisson-like fluctuations. If the mean-number of quanta exchanged $`<n>1`$ (e.g. for heavy ions with large Z$`1/\alpha `$), the distribution becomes narrow and we approach a classical limit, with weakly fluctuating scattering. Modulo color factors, the same conclusion applies to multi-gluon exchange in QCD. In contrast, the instanton-induced effects have completely different statistical properties. The field of the instanton itself is classical, hence coherent. However, the distribution over the instanton size and position is quantum (in contrast to the Coulomb field of the ion just mentioned), thereby leading to cross section fluctuations. The latter are further enhanced by the $`diluteness`$ of the instanton ensemble: the quark may appear very black, provided a tunneling event happens to be close to it, and rather transparent otherwise. As noticed already in , quarks are “twinkling” objects, as the associated gauge/quark fields are strongly fluctuating.
To quantify some of these statements we show in Fig.(8) how such distribution looks like. We plot $`|F_{ss}(b=1)|^2`$, at fixed impact parameter b=$`\rho _0`$. The distribution corresponds to instantons being homogeneously distributed in the 4d sphere around the center of the collision point, with a radius $`R_s2.2\rho _0`$ such that $`\pi ^2R_s^4/2=1/n_0`$, or in a smaller sphere within $`R<\rho _0`$. However the resulting amplitude is highly inhomogeneous, with a large peak at small amplitude, and a long tail at large amplitude. Comparing the solid and dashed curves, one can see that the latter is due to instantons sitting near the center of the system.
## VIII Conclusions and Outlook
### A Conclusions
Several new instanton-generated phenomena have been studied in this work: static potentials for color dipoles, and high energy quark-quark and dipole-dipole scattering. The nature of the instanton effects makes their contribution to these processes different from the contribution expected in perturbation theory.
Overall, the magnitude of the instanton contribution is governed by two competing factors: (i) a diluteness factor $`n_0\rho _0^41`$ reflecting the fact that their density in the QCD vacuum is small ($`n_0\rho _0^41`$), and (ii) a classical enhancement factor, the instanton action which is large ($`S_0/\mathrm{}101)`$. Naturally, the more partons are involved in a particular process, the more powers of $`\alpha _s`$ appear in the perturbative result for a particular process. This penalty does not apply to the instanton contribution. One way to quantify this difference is to note that the ratio of the instanton-to-perturbative contributions contains a power of the classical enhancement parameter, and this power grows with the number of partons involved. Typically the first power due to the classical instanton enhancement cannot really compensate for the small diluteness of the instantons in the vacuum. However, the second power is already sufficient to make the instanton effects larger than the perturbative ones as we have now established for the potentials. Indeed, the dipole-dipole instanton-induced potential exceeds significantly (by a factor $``$ 25) the perturbative contribution for distances $`R>\rho _0`$.
Based on these ideas, we have extended the analysis to near-forward parton-parton scattering amplitudes, treating in details the case of quark-quark and dipole-dipole scattering. Key to our analysis was the concept of analytical continuation in the rapidity variable, which we have applied to both the perturbative and instanton analysis for comparison.
In the perturbative analysis, one- and two-gluon exchanges differ fundamentally in the sense that the former is color-changing (inelastic), while the latter is color-preserving (elastic). Indeed, the two-gluon exchange mechanism constitutes the starting ground for the soft pomeron approach to dipole-dipole scattering. Since the instantons can be viewed as multi-gluon configurations (classical fields), we have suggested that they maybe a viable starting point to analyze soft parton-parton scatterings. We have shown that the instanton-induced amplitudes involve also color-elastic and color-inelastic channels. After analytical continuation, the one-instanton contribution to the color-elastic channel is purely real and vanishes as $`1/\sqrt{s}`$ (much like a scalar exchange). In other words, in this work a single instanton is not “cut”, its multi-gluon content is not used. Instantons contribute to soft parton-parton scattering like the t-channel gluons mostly through color exchange channels, or through re-scattering in the elastic channel. The leading instanton contribution involves a two-instanton-prong channel, and yields a finite elastic parton-parton scattering amplitude after analytical continuation in rapidity space. Our result is reminiscent of the one reached in the stochastic vacuum model , although our assumptions and methodology are different.
### B Outlook
The results we have derived were achieved in Euclidean space prior to our pertinent analytical continuation. Therefore, they are testable from first principles by repeating our analysis using instead lattice QCD simulations. Indeed, the non-perturbative dipole-dipole forces could be studied. In contrast to the quark-antiquark potential and to the best of our knowledge, those forces have not been investigated on the lattice. Also, the various scattering amplitudes discussed in the present work can and should be looked at, leading to multi-parton amplitudes as we have qualitatively discussed. Note, that not only the potentials and scattering amplitudes themselves can be derived, but the degree of their correlation with the presence of instantons in the underlying configurations can be revealed as well, using lattice techniques such as “cooling” and alike to help discriminate instantons by their topological charge.
Regarding the applications of our results, we admit that there remains a significant distance to the description of real hadronic processes. Although we hope to cover further phenomenological applications elsewhere, we still would like to comment on two broad but important dynamical issues: (i) the mechanism of color rearrangements in high energy collisions and (ii) the issue of hadronic substructure of the non-perturbative effects in the hadronic wave functions.
It is generally accepted that high energy hadronic processes can be split into three stages: (i) formation of hadronic wave function (to which we turn later); (ii) color re-arrangements of partons in a collision; and (iii) decay of the arising system into multi-hadron final states. It is further believed that at stage (iii) color flux tubes (QCD strings) are formed with basically $`unit`$ probability<sup>13</sup><sup>13</sup>13The fluxes are described by multiple phenomenological models/codes, e.g. the Lund model., so that one can ignore them in the calculation of the cross section. Such assumption is implied in any perturbative approach (such as the Low-Nussinov gluon-exchange model ), and we assume the same is true for instanton-induced color exchanges as well.
Our main suggestion for further work is that although the instanton-induced mechanism yields relatively small cross sections, this mechanism is likely to dominate over events with multiple color rearrangements. Is there experimental evidence for this assertion in high-energy hadronic collisions? An answer is provided by Fig. 9 (taken from ), which shows a (specially normalized) compilation of multiplicity distributions in $`\overline{p}p`$ collisions at various energies. The data shows that there is indeed (at least) two components: (i) one, with the cross section $`\sigma _1(s)`$ and standard KNO distribution (well known from lower energy pp collisions), as indicated by the solid curve; and (ii) another one with a different cross section $`\sigma _2(s)`$ and much higher multiplicity. Ascribing the main peak at $`N/<N_1>0.8`$ to a $`single`$ color rearrangement reaction (=2 QCD strings formed), one can conclude that at the highest energy $`\sqrt{s}=1800GeV`$ the multiplicity seen may amount up to 10 strings.
The existence of the second component with double and triple multiplicity was anticipated by people doing Regge theory decades ago, in the form of multi-pomeron exchanges. However, the multiplicity data shown in Fig.9 do not really fit well into this description. The second components simply does not look as iterations of the first one. There are no separate peaks and, more importantly, the s-dependence is completely different. The first component is in fact consistent with the approximation used above, namely asymptotically constant cross section (zero pomeron intercept), while the latter grows with $`\sqrt{s}`$ very strongly.
Attempts to solve this puzzle in pQCD, by summing ladder-type diagrams in leading log(x) approximation are well known , and they do indeed produce strongly growing cross sections and multi-parton states. So, the second component may well be due to those perturbative processes. Non-perturbative approaches (aiming mostly at the “soft pomeron” or the first component discussed) have also been tried, from old fashion multi-peripheral hadronic models (e.g. recent work ) to mixed gluon-hadron ladders . Unfortunately, none of these approaches have lead so far to a quantitative theory.
Results/estimates made in this work lead to the conclusion, that instanton-induced color exchanges should dominate over pQCD t-channel gluons starting from the $`double`$ exchange amplitudes. It is therefore logical to conjecture, that the second high-multiplicity component of pp collisions may be generated by this mechanism. That would explain why multiple-string events are not just iteration of the first component, and even consistent with where the transition appears to be. Needless to say that much more work is still needed for a further test of this conjecture.
Acknowledgements We thank Maciek Nowak for discussions. This work was supported in parts by the US-DOE grant DE-FG-88ER40388.
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# Basis-independent analysis of the sneutrino sector in R-parity violating supersymmetry
## I Introduction
In low-energy supersymmetric extensions of the Standard Model, lepton and baryon number conservation are not automatically respected by the most general set of renormalizable interactions . Nevertheless, experimental observations imply that lepton number violating effects, if they exist, must be rather small. Moreover, baryon number violation, if present, must be consistent with the observed stability of the proton. If one wants to enforce lepton and baryon number conservation, it is sufficient to impose one extra discrete symmetry. In the minimal supersymmetric extension of the Standard Model (MSSM), a multiplicative symmetry called R-parity is introduced , such that the R quantum number of an MSSM field of spin $`S`$, baryon number $`B`$ and lepton number $`L`$ is given by $`(1)^{[3(BL)+2S]}`$. By introducing $`BL`$ conservation modulo 2, one eliminates all dimension-four lepton number and baryon number-violating interactions.
The observation of neutrino mixing effects in solar and atmospheric neutrinos suggest that lepton-number is not an exact global symmetry of the low-energy theory. One can develop a supersymmetric model of neutrino masses that generalizes the see-saw mechanism while maintaining R-parity as a good symmetry (where lepton number is violated by two units). In this paper, we consider the alternative possibility that neutrino masses and mixing arise in a theory of R-parity violation, in which lepton number is violated by one unit . In the most general R-parity-violating (RPV) model, both $`B`$ and $`L`$ are violated. However, it is difficult to relax both lepton and baryon number conservation in the low-energy theory without generating a proton decay rate many orders of magnitude above the present bounds. It is possible to enforce baryon number conservation, while allowing for lepton number violating interactions by imposing a discrete baryon $`𝐙_\mathrm{𝟑}`$ symmetry on the low-energy theory , in place of the standard $`𝐙_\mathrm{𝟐}`$ R-parity. Henceforth, we consider R-parity-violating low-energy supersymmetry with an unbroken discrete baryon $`𝐙_\mathrm{𝟑}`$ symmetry. This model exhibits lepton-number-violating phenomena such as neutrino masses, sneutrino/antisneutrino mixing, and lepton-number violating decays.
In RPV low-energy supersymmetry, there is no quantum number that distinguishes the lepton supermultiplets $`\widehat{L}_m`$ and the down-type Higgs supermultiplet $`\widehat{H}_D`$ ($`m`$ is a generation label that runs from 1 to $`n_g=3`$). Each supermultiplet transforms as a $`Y=1`$ weak doublet under the electroweak gauge group. It is therefore convenient to denote $`\widehat{L}_0\widehat{H}_D`$ and unify the four supermultiplets by one symbol $`\widehat{L}_\alpha `$ ($`\alpha =0,1,\mathrm{},n_g`$). Then, the relevant terms in the (renormalizable) superpotential are
$$W=ϵ_{ij}\left[\mu _\alpha \widehat{L}_\alpha ^i\widehat{H}_U^j+\frac{1}{2}\lambda _{\alpha \beta m}\widehat{L}_\alpha ^i\widehat{L}_\beta ^j\widehat{E}_m+\lambda _{\alpha nm}^{}\widehat{L}_\alpha ^i\widehat{Q}_n^j\widehat{D}_m\right],$$
(1)
where $`\widehat{H}_U`$ is the up-type Higgs supermultiplet, the $`\widehat{Q}_n`$ are doublet quark supermultiplets, the $`\widehat{D}_m`$ are singlet down-type quark supermultiplets and the $`\widehat{E}_m`$ are the singlet charged lepton supermultiplets. Note that $`\mu _\alpha `$ and $`\lambda _{\alpha nm}^{}`$ are vectors and $`\lambda _{\alpha \beta m}`$ is an antisymmetric matrix in the generalized lepton flavor space.
Next, the soft-supersymmetry-breaking terms are also generalized in similar way. The relevant terms are
$$V_{\mathrm{soft}}=(M_{\stackrel{~}{L}}^2)_{\alpha \beta }\stackrel{~}{L}_\alpha ^i\stackrel{~}{L}_\beta ^i(ϵ_{ij}b_\alpha \stackrel{~}{L}_\alpha ^iH_U^j+\mathrm{h}.\mathrm{c}.)+ϵ_{ij}[\frac{1}{2}a_{\alpha \beta m}\stackrel{~}{L}_\alpha ^i\stackrel{~}{L}_\beta ^j\stackrel{~}{E}_m+a_{\alpha nm}^{}\stackrel{~}{L}_\alpha ^i\stackrel{~}{Q}_n^j\stackrel{~}{D}_m+\mathrm{h}.\mathrm{c}.],$$
(2)
where the fields appearing in eq. (2) are the scalar partners of the superfields that appear in eq. (1). Here, $`b_\alpha `$ and $`a_{\alpha nm}^{}`$ are vectors, $`a_{\alpha \beta m}`$ is an antisymmetric matrix and $`(M_{\stackrel{~}{L}}^2)_{\alpha \beta }`$ is a Hermitian matrix in the generalized lepton flavor space.
When the scalar potential is minimized (see Section II), one finds a vacuum expectation value for the neutral scalar fields: $`\stackrel{~}{L}_\alpha =v_\alpha /\sqrt{2}`$ and $`H_U=v_u/\sqrt{2}`$. To make contact with the usual notation of the MSSM, we define the length of the vector $`v_\alpha `$ by $`v_d(v_\alpha v_\alpha )^{1/2}`$ and $`\mathrm{tan}\beta v_u/v_d`$. The mass of the $`W`$ boson constrains the value $`v^2v_u^2+v_d^2=(246\mathrm{GeV})^2`$.
So far, there is no distinction between the neutral Higgs bosons and neutral sleptons. Nevertheless, we know that RPV-interactions, if present, must be small. It is tempting to choose a particular convention corresponding to a specific choice of basis in the generalized lepton flavor space. For example, one can choose to define the down-type Higgs multiplet such that $`H_D\stackrel{~}{L}_0=v_d/\sqrt{2}`$ and $`\stackrel{~}{L}_m=0`$. This means that we let the dynamics (which determines the direction of the vacuum expectation value in the generalized lepton flavor space) choose the definition of the down-type Higgs field. In this basis, all the RPV-parameters are well defined and must be small to satisfy phenomenological constraints.
Nevertheless, the above convention is only one possible basis choice. Other conventions are equally sensible. For example, one could choose a second basis where $`\mu _m=0`$ and a third basis where $`b_m=0`$. In each case, the corresponding RPV-parameters are small. But comparing results obtained in different bases requires some care. Moreover, it is often desirable to study the evolution of couplings from some high (unification) scale to the low-energy (electroweak) scale. The renormalization group equations for the RPV-parameters is not basis preserving. That is, a particular basis choice at the high energy scale will lead to some complicated effective basis choice at the low-energy scale.
The problems described above can be ameliorated by avoiding basis-specific definitions of parameters. The challenge of such an approach is to determine a set of basis-independent RPV parameters, in the spirit of the Jarlskog invariant which characterizes the strength of CP-violation in the Standard Model . Such an approach has been applied to RPV models in the past, where neutrino masses , early universe physics and the Higgs sector were studied. It is instructive to examine the neutrino spectrum of the RPV model. At tree level, one neutrino become massive due to the RPV mixing of the neutrinos and the neutralinos. The other $`n_g1`$ neutrinos remain massless at tree-level, although they can acquire smaller radiative masses at one-loop. To first order in the small RPV-parameter the basis independent formula can be written in the following form :
$$m_\nu =\frac{m_Z^2\mu M_{\stackrel{~}{\gamma }}\mathrm{cos}^2\beta }{m_Z^2M_{\stackrel{~}{\gamma }}\mathrm{sin}2\beta M_1M_2\mu }|\widehat{v}\times \widehat{\mu }|^2,$$
(3)
where $`M_{\stackrel{~}{\gamma }}\mathrm{cos}^2\theta _WM_1+\mathrm{sin}^2\theta _WM_2`$ depends on gaugino mass parameters $`M_1`$ and $`M_2`$. In eq. (3), $`\widehat{v}`$ and $`\widehat{\mu }`$ are unit vectors in the $`v_\alpha `$ and $`\mu _\alpha `$ directions, respectively. It is convenient to introduce the notation of the cross-product of two vectors. Although the cross-product technically exists only in three-dimensions, the dot product of two cross-products can be expressed as a product of dot-products
$$(a\times b)(c\times d)=(ac)(bd)(ad)(bc),$$
(4)
which exists in any number of dimensions. This notation is useful, since any expression that involves the cross-product of two vectors vanishes if the corresponding vectors are parallel. This provides a nice geometrical characterization of the small RPV-parameters of the model. For example, $`|\widehat{v}\times \widehat{\mu }|^2=\mathrm{sin}^2\xi `$ where $`\xi `$ is the angle between $`\widehat{v}`$ and $`\widehat{\mu }`$. Thus, in eq. (3), $`|\widehat{v}\times \widehat{\mu }|^2`$ is the small RPV-parameter, while the prefactor can be computed in the R-parity-conserving (RPC) limit of the model.
In this paper, we focus on a basis-independent description of the RPV-parameters that govern the sneutrino spectrum. The model possesses lepton-number-violating $`\mathrm{\Delta }L=1`$ interactions that give rise to the mixing of sleptons and Higgs bosons.<sup>a</sup><sup>a</sup>aThese interactions modify the phenomenology of the charged and neutral scalars (relative to the RPC limit), as discussed in ref. . These interactions also generate $`\mathrm{\Delta }L=2`$ effective operators that give rise to sneutrino/antisneutrino mixing . In this case, the sneutrino ($`\stackrel{~}{\nu }`$) and antisneutrino ($`\overline{\stackrel{~}{\nu }}`$), which are eigenstates of lepton number, are no longer mass eigenstates. The mass eigenstates are superpositions of $`\stackrel{~}{\nu }`$ and $`\overline{\stackrel{~}{\nu }}`$, and sneutrino mixing effects can lead to a phenomenology analogous to that of $`K`$$`\overline{K}`$ and $`B`$$`\overline{B}`$ mixing . The mass splitting between the two sneutrino mass eigenstates is related to the magnitude of lepton number violation, which is typically characterized by the size of neutrino masses . As a result, the sneutrino/antisneutrino mass splitting is expected generally to be very small. Yet, it can be detected in many cases, if one is able to observe the lepton number oscillation .
In contrast to the neutrino sector (where only one neutrino mass eigenstate acquires a tree-level mass), in general all sneutrinos/antisneutrino pairs are split in mass at tree level.<sup>b</sup><sup>b</sup>bThis result is a consequence of the fact that in the RPC limit, neutrinos are massless and hence degenerate, whereas the RPC sneutrino masses are in general non-degenerate. See Section V and ref. for further discussions of this point. For simplicity, we consider the case of a CP-conserving scalar sector. In the RPC limit, the CP-even scalar sector consists of two Higgs scalars ($`h^0`$ and $`H^0`$, with $`m_{h^0}<m_{H^0}`$) and $`n_g`$ generations of CP-even sneutrinos $`(\stackrel{~}{\nu }_+)_m`$, while the CP-odd scalar sector consists of the Higgs scalar, $`A^0`$, the Goldstone boson (which is absorbed by the $`Z`$), and $`n_g`$ generations of CP-odd sneutrinos $`(\stackrel{~}{\nu }_{})_m`$. Here, we have implicitly chosen a flavor basis in which the sneutrinos are mass eigenstates. Moreover, the $`(\stackrel{~}{\nu }_\pm )_m`$ are mass degenerate (separately for each $`m`$), so that the standard practice is to define eigenstates of lepton number: $`\stackrel{~}{\nu }_m[(\stackrel{~}{\nu }_+)_m+i(\stackrel{~}{\nu }_{})_m]/\sqrt{2}`$ and $`\overline{\stackrel{~}{\nu }}_m\stackrel{~}{\nu }_m^{}`$. When R-parity is violated, the sneutrinos in each CP-sector mix with the corresponding Higgs scalars, and the mass degeneracy of $`(\stackrel{~}{\nu }_+)_m`$ and $`(\stackrel{~}{\nu }_{})_m`$ is broken. In ref. we computed the mass-splitting in a special basis where $`v_m=0`$ and the matrix $`(M_{\stackrel{~}{\nu }\stackrel{~}{\nu }^{}}^2)_{ij}`$ \[which is the $`3\times 3`$ block sub-matrix of $`(M_{\stackrel{~}{\nu }\stackrel{~}{\nu }^{}}^2)_{\alpha \beta }`$ defined in eq. (7)\] is diagonal. In this basis, we identified the $`b_m`$ as the relevant small RPV-parameters. To leading order in $`b_m^2`$,
$$(\mathrm{\Delta }m_{\stackrel{~}{\nu }}^2)_m=\frac{4b_m^2m_Z^2m_{\stackrel{~}{\nu }_m}^2\mathrm{sin}^2\beta }{(m_H^2m_{\stackrel{~}{\nu }_m}^2)(m_h^2m_{\stackrel{~}{\nu }_m}^2)(m_A^2m_{\stackrel{~}{\nu }_m}^2)},$$
(5)
where $`(\mathrm{\Delta }m_{\stackrel{~}{\nu }}^2)_m(m_{\stackrel{~}{\nu }_+}^2)_m(m_{\stackrel{~}{\nu }_{}}^2)_m`$. As in the neutrino case described above, we may evaluate the prefactor that multiplies $`b_m^2`$ in the RPC limit. In deriving eq. (5), it was assumed that all RPC Higgs and sneutrino masses are all distinct. If degeneracies exist, the above formula must be modified.
The goal of this paper is to reanalyze the sneutrino mass spectrum in a basis-independent formalism. We identify the small RPV-parameters that govern the sneutrino/antisneutrino mass splittings. Our technique will also allow us to generalize the analysis to treat the case of scalar mass degeneracies. In Section II, we derive a convenient form for the CP-even and CP-odd scalar squared-mass matrices. We compute the sneutrino/antisneutrino squared-mass difference in the case of one sneutrino flavor in Section III. In Section IV, we generalize to an arbitrary number of generations, and exhibit explicit formulae for the two and three generation cases. The latter results assume that in the RPC limit, there are no degeneracies among different sneutrino flavors. The degenerate case is treated in Section V. A discussion of our results and conclusions are presented in Section VI. Details of our computations are provided in six appendices.
## II Minimum condition and basic equations
We begin our analysis by collecting the relevant formulae given in ref. . We assume that the scalar sector is CP-conserving,<sup>c</sup><sup>c</sup>cIn the MSSM, the Higgs sector is automatically CP-conserving at tree-level, since all phases can be removed by suitable redefinitions of the fields. In the RPV model, new phases enter through $`(M_{\stackrel{~}{L}}^2)_{\alpha \beta }`$, $`b_\alpha `$ and $`\mu _\alpha `$, which cannot all be simultaneously removed in the general case. which implies that the scalar fields can be defined such that $`M_{\stackrel{~}{L}}^2`$ is a real symmetric matrix and $`b_\alpha `$ and $`\mu _\alpha `$ are real. The vacuum expectation value $`L_\alpha v_\alpha /\sqrt{2}`$ is determined by minimizing the scalar potential. With the assumption of CP-conservation, one can separate out the scalar potential for the CP-even and CP-odd sector, $`V=V_{\mathrm{even}}+V_{\mathrm{odd}}`$. Then, $`v_\alpha `$ is determined by minimizing $`V_{\mathrm{even}}`$, and the resulting condition is given by:
$$(M_{\stackrel{~}{\nu }\stackrel{~}{\nu }^{}}^2)_{\alpha \beta }v_\beta =v_ub_\alpha ,$$
(6)
where
$$(M_{\stackrel{~}{\nu }\stackrel{~}{\nu }^{}}^2)_{\alpha \beta }(M_{\stackrel{~}{L}}^2)_{\alpha \beta }+\mu _\alpha \mu _\beta \frac{1}{8}(g^2+g^2)(v_u^2v_d^2)\delta _{\alpha \beta }.$$
(7)
Note that eq. (6) determines both the size of $`v_\alpha `$ and its direction. In a perturbative treatment of RPV terms, we can use the RPC value for the squared magnitude $`v_d^2_\alpha v_\alpha v_\alpha `$, and then use eq. (6) to determine the direction of $`v_\alpha `$ in the generalized lepton flavor space.
We next separate the scalar squared-mass matrices into CP-odd and CP-even blocks. In the $`H_U`$$`\stackrel{~}{L}_\alpha `$ basis, the CP-odd squared-mass matrix is given by
$$M_{\mathrm{odd}}^2=\left(\begin{array}{cc}b_\rho v_\rho /v_u& b_\beta \\ b_\alpha & (M_{\stackrel{~}{\nu }\stackrel{~}{\nu }^{}}^2)_{\alpha \beta }\end{array}\right),$$
(8)
while the CP-even squared-mass matrix is given by
$$M_{\mathrm{even}}^2=\left(\begin{array}{cc}\frac{1}{4}(g^2+g^2)v_u^2+b_\rho v_\rho /v_u& \frac{1}{4}(g^2+g^2)v_uv_\beta b_\beta \\ \frac{1}{4}(g^2+g^2)v_uv_\alpha b_\alpha & \frac{1}{4}(g^2+g^2)v_\alpha v_\beta +(M_{\stackrel{~}{\nu }\stackrel{~}{\nu }^{}}^2)_{\alpha \beta }\end{array}\right),$$
(9)
where $`(M_{\stackrel{~}{\nu }\stackrel{~}{\nu }^{}}^2)_{\alpha \beta }`$ is defined in eq. (7).
To compute the squared-mass differences of the corresponding CP-even and CP-odd sneutrinos, we must diagonalize both of the above matrices. We wish to employ a perturbative procedure by identifying the small RPV-parameters, without resorting to a specific choice of basis. Our strategy is to recast the two scalar squared-mass matrices in a more convenient form.
First, consider the CP-odd squared mass matrix \[eq. (8)\]. Note that the vector $`(v_u,v_\beta )`$ is an eigenvector of $`M_{\mathrm{odd}}^2`$ with zero eigenvalue; this is the Goldstone boson that is absorbed by the $`Z`$. We can remove the Goldstone boson by introducing the following orthogonal $`(n_g+2)\times (n_g+2)`$ matrix
$$U_o=\left(\begin{array}{ccc}v_u/v& v_d/v& 0\\ v_\beta /v& v_uv_\beta /(v_dv)& X_{\beta i}\end{array}\right),$$
(10)
where $`v(v_u^2+v_d^2)^{1/2}`$. Note that the index $`i`$ runs from 1 to $`n_g`$; thus $`X_{\alpha i}`$ is an $`(n_g+1)\times n_g`$ matrix. The orthogonality of $`U_o`$ implies that each column of $`U_o`$ is a real unit vector and different columns are orthogonal. In addition, the set $`\{v_\beta /v,X_{\beta i}\}`$ forms an orthonormal set of vectors in an $`(n_g+1)`$-dimensional vector space. It follows that:
$`v_\alpha X_{\alpha i}`$ $`=`$ $`0,`$ (11)
$`X_{\alpha i}X_{\alpha j}`$ $`=`$ $`\delta _{ij},`$ (12)
$`X_{\alpha i}X_{\beta i}`$ $`=`$ $`\delta _{\alpha \beta }{\displaystyle \frac{v_\alpha v_\beta }{v_d^2}}.`$ (13)
In our computations, no explicit realization of the $`X_{i\alpha }`$ will be required. A simple computation yields:
$$U_o^TM_{\mathrm{odd}}^2U_o=\left(\begin{array}{cc}0& 0_\beta \\ 0_\alpha & (\stackrel{~}{M}_{\mathrm{odd}}^2)_{\alpha \beta }\end{array}\right),$$
(14)
where $`0_\beta `$ \[$`0_\alpha `$\] is a row \[column\] matrix of zeros and<sup>d</sup><sup>d</sup>dWe define $`X_{i\alpha }`$ to be the transpose of the matrix $`X_{\alpha i}`$. When no ambiguity arises, we will not explicitly exhibit the transpose symbol (superscript $`T`$).
$$\stackrel{~}{M}_{\mathrm{odd}}^2=\left(\begin{array}{cc}v^2(vb)/(v_uv_d^2)& vb_\beta X_{\beta i}/v_d\\ vX_{j\alpha }b_\alpha /v_d& X_{j\alpha }(M_{\stackrel{~}{\nu }\stackrel{~}{\nu }^{}}^2)_{\alpha \beta }X_{\beta i}\end{array}\right),$$
(15)
where $`vbv_\alpha b_\alpha `$. The eigenstates of $`\stackrel{~}{M}_{\mathrm{odd}}`$ correspond to the CP-odd Higgs boson $`A^0`$ and $`n_g`$ generations of CP-odd sneutrinos.
It turns out that it is also convenient to rotate the CP-even squared-mass matrix, but by a slightly different orthogonal transformation. In the limit of $`m_Z=0`$, the CP-even squared-mass matrix also possesses a Goldstone boson, which we can explicitly isolate. Comparing the CP-odd and CP-even cases, we see that when $`g=g^{}=0`$ the two matrices are related by $`b_\alpha b_\alpha `$ and $`v_uv_u`$. Thus, if we introduce $`U_eU_o(v_uv_u)`$ and define $`\stackrel{~}{M}_{\mathrm{even}}^2U_e^TM_{\mathrm{even}}^2U_e`$, then
$$\stackrel{~}{M}_{\mathrm{even}}^2=\left(\begin{array}{ccc}m_Z^2\mathrm{cos}^22\beta & m_Z^2\mathrm{cos}2\beta \mathrm{sin}2\beta & 0\\ m_Z^2\mathrm{cos}2\beta \mathrm{sin}2\beta & m_Z^2\mathrm{sin}^22\beta +v^2(vb)/(v_uv_d^2)& vb_\beta X_{\beta i}/v_d\\ 0& vX_{j\alpha }b_\alpha /v_d& X_{j\alpha }(M_{\stackrel{~}{\nu }\stackrel{~}{\nu }^{}}^2)_{\alpha \beta }X_{\beta i}\end{array}\right),$$
(16)
where we used $`m_Z^2=\frac{1}{4}(g^2+g^2)v^2`$ and $`\mathrm{tan}\beta v_u/v_d`$. The eigenstates of $`\stackrel{~}{M}_{\mathrm{even}}^2`$ correspond to the CP-even Higgs bosons $`h^0`$ and $`H^0`$ and $`n_g`$ generations of CP-even sneutrinos.
In the RPC limit, one can choose a basis in which $`v_m=b_m=X_{0m}=0`$. It follows that the sneutrino and Higgs mass matrices decouple and one recovers the known RPC result. A basis-independent characterization of the RPC limit in the scalar (sneutrino/Higgs) sector is the condition that the vectors $`b_\beta `$ and $`v_\beta `$ are aligned.<sup>e</sup><sup>e</sup>eFor example, if one chooses the basis where $`v_m=0`$, then by eq. (11), $`X_{0m}=0`$. In this basis, the $`b_m`$ are the small RPV-parameters. Equivalently, by using the transformed mass matrices given above, it is clear that the quantities
$$B_i\frac{vb_\beta X_{\beta i}}{v_d}$$
(17)
can be identified as $`n_g`$ basis-independent parameters that vanish in the RPC limit, and thus provide good candidates for the small quantities that can be used in a perturbative expansion. If $`b_\beta `$ and $`v_\beta `$ are aligned, then eq. (11) implies that the $`B_i=0`$ and we are back to the RPC limit.
Although the $`B_i`$ provide a basis-independent set of small RPV-parameters, the explicit dependence on $`X_{\alpha i}`$ is inconvenient. Clearly, it is preferable to re-express the $`B_i`$ directly in terms of the original model parameters. In the following sections, we will exhibit this procedure in the one-generation case and generalize it to the multi-generation case. To this end, it is convenient to introduce a set of new RPV-parameters. In the case of $`n_g3`$, only one new vector is required for our final results:
$$c_\alpha \frac{(M_{\stackrel{~}{\nu }\stackrel{~}{\nu }^{}}^2)_{\alpha \beta }b_\beta }{b^2}.$$
(18)
For $`n_g>3`$, further vectors are required. It is convenient to introduce a series of vectors
$$c_\alpha ^{(n+1)}(M_{\stackrel{~}{\nu }\stackrel{~}{\nu }^{}}^2)_{\alpha \beta }c_\beta ^{(n)},$$
(19)
where $`c_\alpha ^{(1)}c_\alpha `$. Clearly, the maximal number of linearly independent vectors (along with $`b`$ and $`v`$) is $`n=n_g1`$, although not all of these will appear in our final results.
To simplify the presentation we introduce the following shorthand for the elements of the transformed CP-even and CP-odd squared-mass matrices:
$$\stackrel{~}{M}_{\mathrm{odd}}^2=\left(\begin{array}{cc}A& B_i\\ B_j& C_{ij}\end{array}\right),$$
(20)
$$\stackrel{~}{M}_{\mathrm{even}}^2=\left(\begin{array}{ccc}D& E& 0\\ E& F+A& B_i\\ 0& B_j& C_{ij}\end{array}\right).$$
(21)
Our strategy is to employ a first-order perturbative analysis to diagonalize these matrices, taking the small parameters to be the $`B_i`$ \[eq. (17)\]. At zeroth order (setting the $`B_i`$ to zero), the above matrices can be evaluated in the RPC limit. In particular, $`A=m_A`$, the eigenvalues of $`\left(\genfrac{}{}{0pt}{}{D}{E}\genfrac{}{}{0pt}{}{E}{F+A}\right)`$ are the CP-even Higgs squared-masses $`m_{h^0}^2`$ and $`m_{H^0}^2`$, and $`C_{ij}`$ is the RPC sneutrino squared-mass matrix. In order for non-degenerate perturbation theory to be valid, we henceforth assume that none of the eigenvalues of $`C_{ij}`$ is (approximately) equal to any of the neutral Higgs boson squared-masses.<sup>f</sup><sup>f</sup>fRelaxing this assumption would require one to use degenerate perturbation theory. The resulting analysis would be more involved and not very illuminating, so we spare the reader by omitting the consideration of this possibility. This is not a serious restriction, since there is no natural choice of supersymmetric parameters that would guarantee such a (near) degeneracy. In particular, the Higgs mass matrices arise in part from the mixing of $`H_D`$ and $`H_U`$, which is governed by the parameter $`b_0`$ and not related to the generation of sneutrino masses.
The result of the calculation outlined above will be basis-independent expressions for the sneutrino/antisneutrino squared-mass splittings. We first illustrate the method for one sneutrino generation in Section III, and then generalize it to the multi-generational case in Section IV, assuming that the eigenvalues of $`C_{ij}`$ are non-degenerate. The case where some of the eigenvalues of $`C_{ij}`$ are degenerate will be treated in Section V.
## III The case of one generation
In the one generation case we can drop the Roman indices from $`X`$, $`B`$ and $`C`$, namely we define $`X_\alpha X_{\alpha 1}`$, $`BB_1`$ and $`CC_{11}`$. To zeroth order in $`B`$, the CP-even and CP-odd sneutrino squared-masses are equal to $`C`$. The corrections can be calculated perturbatively.
The eigenvalue equation for the CP-odd squared-mass matrix, $`det(\stackrel{~}{M}_{\mathrm{odd}}^2\lambda I)=0`$ reads
$$(A\lambda )(C\lambda )B^2=0.$$
(22)
For $`B`$ small, $`\lambda =C+𝒪(B^2)`$, so we can take $`\lambda C=aB^2`$ and solve for $`a`$. Thus, to first order in $`B^2`$, the squared mass of the CP-odd sneutrino is
$$m_{\mathrm{odd}}^2=C\frac{B^2}{AC}.$$
(23)
A similar analysis for the CP-even squared-mass matrix yields the following result for the squared-mass of the CP-even sneutrino:
$$m_{\mathrm{even}}^2=C\frac{B^2(DC)}{(F+AC)(DC)E^2}.$$
(24)
The squared-mass splitting, $`\mathrm{\Delta }m_{\stackrel{~}{\nu }}^2=m_{\mathrm{even}}^2m_{\mathrm{odd}}^2`$, is given by
$$\mathrm{\Delta }m_{\stackrel{~}{\nu }}^2=\frac{FCB^2}{(AC)[(F+AC)(DC)E^2]},$$
(25)
where we have used the fact that $`FD=E^2`$. Note that the only small parameter in eq. (25) is $`B^2`$. We can therefore use the RPC values for the prefactor that multiplies $`B^2`$. For example, as noted above, $`C=m_{\stackrel{~}{\nu }}^2`$ is the RPC sneutrino squared-mass.
Although $`B=vb_\alpha X_\alpha /v_d^2`$ is expressed in a basis-independent manner, the explicit dependence on $`X_\alpha `$ is inconvenient. Clearly, it is preferable to re-express $`B`$ directly in terms of the original model parameters. To this end, note that the orthogonality condition \[eq. (11)\] implies that $`X_\alpha `$ and $`v_\alpha `$ are orthogonal, and thus a dot product of any vector with $`X`$ is equivalent to a cross product with $`v`$. Using eq. (A3) for $`B^2`$ and the RPC values for the other parameters of eq. (25), we end up with
$$\mathrm{\Delta }m_{\stackrel{~}{\nu }}^2=\frac{4b^2m_Z^2m_{\stackrel{~}{\nu }}^2\mathrm{sin}^2\beta }{(m_{H^0}^2m_{\stackrel{~}{\nu }}^2)(m_{h^0}^2m_{\stackrel{~}{\nu }}^2)(m_{A^0}^2m_{\stackrel{~}{\nu }}^2)}|\widehat{v}\times \widehat{b}|^2,$$
(26)
where $`\widehat{b}`$ is a unit vector in the $`b_\alpha `$ direction and the square of the cross-product is formally defined according to eq. (4). It is easy to check \[see Appendix B\] that in the special basis where $`v_1=0`$, the basis-independent result above \[eq. (26)\] reduces to the basis-dependent result quoted in eq. (5).
The basis-independent result obtain in eq. (26) is still not in optimal form, since it depends on $`v`$, which is a derived quantity that requires one to determine the minimum of the scalar potential \[eq. (6)\]. However, we can employ the vector $`c_\alpha `$ \[defined in eq. (18)\] to our advantage by noting that in the $`n_g=1`$ case \[see eq. (C7)\],
$$|b\times c|^2=m_{\stackrel{~}{\nu }}^4|\widehat{v}\times \widehat{b}|^2.$$
(27)
Consequently, we can express the sneutrino squared-mass splitting in the one-generation case directly in terms of fundamental parameters of the RPV-Lagrangian in a completely basis-independent form.
## IV The case of an arbitrary number of generations
In this section, we obtain results for an arbitrary number of generations. We then explicitly exhibit the corresponding results for $`n_g=2`$ and 3 generations. The eigenvalue equation for the CP-odd scalar squared-mass matrix is
$$(A\lambda )det(C\lambda I)+Y^{(N)}(\lambda )=0,$$
(28)
where $`I`$ is the $`N\times N`$ unit matrix, $`Nn_g`$, and
$$Y^{(N)}(\lambda )B_i\text{cof}\left[(\stackrel{~}{M}_{\mathrm{odd}}^2\lambda I)_{0i}\right],$$
(29)
where the sum over the repeated index $`i`$ is assumed implicitly. As usual, the cofactor is defined as $`\text{cof}[A_{ij}]=(1)^{i+j}det\stackrel{~}{A}(i,j)`$ where $`\stackrel{~}{A}(i,j)`$ is the matrix $`A`$ whose $`i`$th row and $`j`$th column are removed. In the special case of a one-dimensional matrix, we can define $`\text{cof}[A_{11}]=1`$.
Let $`\lambda _m^{(0)}`$ ($`m=1,2,\mathrm{},N`$) be the roots of eq. (28) to zeroth order in the $`B_i`$, namely $`det(C\lambda _m^{(0)}I)=0`$. For small $`B_i`$, we insert $`\lambda _m=\lambda _m^{(0)}+(\delta \lambda _m)_{\mathrm{odd}}`$ into eq. (28). Working to the lowest non-trivial order in the $`B_i`$, we make use of eq. (D2) to obtain
$$(\delta \lambda _m)_{\mathrm{odd}}=\frac{Y^{(N)}(\lambda _m)}{(A\lambda _m)det^{}(C\lambda _mI)},$$
(30)
where $`det^{}A`$ is the product of all the non-zero eigenvalues of $`A`$. In this analysis, we assume that there are no degenerate eigenvalues (the degenerate case will be considered in Section V); hence
$$\stackrel{}{det}(C\lambda _mI)=\underset{im}{}(\lambda _i\lambda _m).$$
(31)
Note that we do not distinguish between $`\lambda _m^{(0)}`$ and $`\lambda _m`$ in eq. (30). Since the $`B_iB_j`$ are the small parameters, any distinction between the two estimates for $`\lambda _m`$ would yield a result that is higher order in the product of the $`B_i`$.
By a similar technique, we may solve the eigenvalue equation for the CP-even scalar squared-mass matrix. Noting that $`\lambda _m^{(0)}`$ is the same in both the CP-odd and CP-even squared-mass computations, we can write $`\lambda _m=\lambda _m^{(0)}+(\delta \lambda _m)_{\mathrm{even}}`$. The end result is
$$(\delta \lambda _m)_{\mathrm{even}}=\frac{Y^{(N)}(\lambda _m)(D\lambda _m)}{[(D\lambda _m)(F+A\lambda _m)E^2]det^{}(C\lambda _mI)}.$$
(32)
We may evaluate the denominator of the above expression in the RPC limit (where $`B^2=0`$). In this limit,
$`(D\lambda _m)(F+A\lambda _m)E^2`$ $`=`$ $`(m_h^2\lambda _m)(m_H^2\lambda _m),`$ (33)
$`A\lambda _m`$ $`=`$ $`m_A^2\lambda _m.`$ (34)
The squared-mass difference of the $`m`$th sneutrino/antisneutrino pair is denoted by $`\mathrm{\Delta }m_{\stackrel{~}{\nu }_m}^2=(m_{\mathrm{even}}^2)_m(m_{\mathrm{odd}}^2)_m`$. Plugging in the results of eqs. (30) and (32), we obtain
$$\mathrm{\Delta }m_{\stackrel{~}{\nu }_m}^2=\frac{Y^{(N)}(\lambda _m)}{det^{}(C\lambda _mI)}\left[\frac{D\lambda _m}{(m_h^2\lambda _m)(m_H^2\lambda _m)}\frac{1}{(m_A^2\lambda _m)}\right].$$
(35)
We may further simplify this result by employing RPC values for any expression that multiplies a term of order $`B_iB_j`$. In particular, we can make use of the well known tree-level MSSM Higgs results: $`m_{h^0}^2+m_{H^0}^2=m_{A^0}^2+m_Z^2`$ and $`m_{H^0}^2m_{h^0}^2=m_{A^0}^2m_Z^2\mathrm{cos}^22\beta `$. Moreover, we may take $`\lambda _m=m_{\stackrel{~}{\nu }_m}^2`$. The end result is:
$$\mathrm{\Delta }m_{\stackrel{~}{\nu }_m}^2=\frac{m_{\stackrel{~}{\nu }_m}^2m_Z^2\mathrm{sin}^22\beta Y^{(N)}(m_{\stackrel{~}{\nu }_m}^2)}{(m_A^2m_{\stackrel{~}{\nu }_m}^2)(m_h^2m_{\stackrel{~}{\nu }_m}^2)(m_H^2m_{\stackrel{~}{\nu }_m}^2)_{im}(m_{\stackrel{~}{\nu }_i}^2m_{\stackrel{~}{\nu }_m}^2)}.$$
(36)
The small RPV parameters that govern the above expression has been completely isolated into $`Y^{(N)}`$. One additional consequence of this result is a simple sum rule that holds for an appropriately weighted sum of sneutrino squared-mass differences. The sum rule and its derivation is given in Appendix E.
We next derive a method for computing $`Y^{(N)}(\lambda )`$. First, we evaluate $`Y^{(N)}`$ for $`\lambda =0`$. From eq. (29), it is straightforward to evaluate $`Y^{(N)}(0)=B_i\text{cof}\left[(\stackrel{~}{M}_{\mathrm{odd}}^2)_{0i}\right]`$. Using eq. (20),
$$Y^{(N)}(0)=B_iB_j\text{cof}\left[C_{ij}\right].$$
(37)
Note that for $`N=1`$, $`Y^{(1)}(\lambda )=B^2`$, independent of the value of $`\lambda `$. One can extend eq. (37) for arbitrary $`\lambda `$. We have found the following recursion relation:
$$Y^{(N)}(\lambda )=B_iB_j\text{cof}\left[C_{ij}\right]\lambda Y^{(N1)}(\lambda ),$$
(38)
with $`Y^{(1)}(\lambda )=B^2`$. The proper use of this equation requires some care. One must first express $`Y^{(N1)}`$ covariantly in terms of the $`(N1)`$-dimensional vector $`B_i`$ and $`(N1)\times (N1)`$ dimensional matrix $`C_{ij}`$. Then, the term $`Y^{(N1)}`$ that appears in eq. (38) is given by precisely the same expression \[obtained in the $`(N1)`$-dimensional case\], but with $`B_i`$ and $`C_{ij}`$ now $`N`$-dimensional objects. For example,
$$Y^{(2)}(\lambda )=Y^{(2)}(0)+\lambda B^2,$$
(39)
but in eq. (39), $`B^2=_{i=1}^NB_iB_i`$, with $`N=2`$.
The solution to the recursion relation \[eq. (38)\] is
$$Y^{(N)}(\lambda )=\underset{k=0}{\overset{N1}{}}(1)^k\lambda ^kY^{(Nk)}(0).$$
(40)
Again, we emphasize that the $`Y^{(Nk)}`$ are first obtained by an $`(Nk)`$-dimensional computation. Once these terms are expressed covariantly in terms of $`B_i`$ and $`C_{ij}`$, the resulting expressions for $`Y^{(Nk)}`$ may be used in eq. (40), with $`B_i`$ and $`C_{ij}`$ promoted to full $`N`$-dimensional objects. Thus, for each extra generation, we need only calculate one new invariant, $`Y^{(N)}(0)`$.
We illustrate the general formulae above for the cases of $`N=1`$, 2 and 3 generations. The case of $`N=1`$ is trivial. Here $`Y^{(1)}(0)=B^2`$. Using eq. (A3), we quickly recover the result of eq. (26). For the case of $`N=2`$, we use eqs. (37) and (40) to obtain:
$$Y^{(2)}(\lambda )=B^2\left[\lambda \text{Tr}(C)\right]+B_iB_jC_{ij}.$$
(41)
Using the results of Appendix A we can express $`Y^{(2)}`$ directly in terms of the model parameters:
$$Y^{(2)}(\lambda )=\frac{1}{v_d^2\mathrm{cos}^2\beta }\left\{|v\times b|^2[\lambda \text{Tr}(M_{\stackrel{~}{\nu }\stackrel{~}{\nu }^{}}^2)]+b^2(v\times b)(v\times c)\right\}.$$
(42)
The final result for the sneutrino squared-mass splittings in the two generation case is
$$\mathrm{\Delta }m_{\stackrel{~}{\nu }_m}^2=\frac{4m_{\stackrel{~}{\nu }_m}^2m_Z^2\mathrm{tan}^2\beta \left\{|v\times b|^2[m_{\stackrel{~}{\nu }_m}^2\text{Tr}(M_{\stackrel{~}{\nu }\stackrel{~}{\nu }^{}}^2)]+b^2(v\times b)(v\times c)\right\}}{v^2(m_{\stackrel{~}{\nu }_n}^2m_{\stackrel{~}{\nu }_m}^2)(m_A^2m_{\stackrel{~}{\nu }_m}^2)(m_h^2m_{\stackrel{~}{\nu }_m}^2)(m_H^2m_{\stackrel{~}{\nu }_m}^2)},$$
(43)
where $`nm`$ and we have put $`\lambda _m=m_{\stackrel{~}{\nu }_m}^2`$. We may evaluate $`\text{Tr}(M_{\stackrel{~}{\nu }\stackrel{~}{\nu }^{}}^2)`$ in the RPC limit:
$$\text{Tr}(M_{\stackrel{~}{\nu }\stackrel{~}{\nu }^{}}^2)=|b|\mathrm{tan}\beta +m_{\stackrel{~}{\nu }_1}^2+m_{\stackrel{~}{\nu }_2}^2,$$
(44)
where $`|b|(b_\alpha b_\alpha )^{1/2}`$. In Appendix B, we verify that eqs. (5) and (43) agree in the special basis.
As in the previous section, we note that eq. (43) depends on the derived quantity $`v`$. At the expense of a somewhat more complex result, we can re-express eq. (43) in terms of the vectors $`b`$, $`c`$, and a new vector $`c^{(2)}`$ introduced in eq. (19), as shown in Appendix C.
For $`N=3`$ generations, the new invariant that arises is again obtained from eq. (37):
$$Y^{(3)}(0)=\frac{1}{2}B^2\left[\text{Tr}(C^2)[\text{Tr}(C)]^2\right]+B_iB_jC_{ij}\text{Tr}(C)B_iB_jC_{ik}C_{kj}.$$
(45)
Following the procedure outlined above, eqs. (38) and (41) yield:
$$Y^{(3)}(\lambda )=Y^{(3)}(0)\lambda Y^{(2)}(\lambda ),$$
(46)
where $`Y^{(2)}(\lambda )`$ is given by eq. (42), with $`v`$, $`b`$ and $`c`$ promoted to three-dimensional vectors and $`M_{\stackrel{~}{\nu }\stackrel{~}{\nu }^{}}^2`$ promoted to a $`3\times 3`$ matrix. An explicit evaluation of $`Y^{(3)}(0)`$ is given in eq. (A22). Inserting the corresponding results into eq. (36), we end up with the sneutrino squared-mass splittings in the three generation case:
$`\mathrm{\Delta }m_{\stackrel{~}{\nu }_m}^2`$ $`=`$ $`{\displaystyle \frac{4m_{\stackrel{~}{\nu }_m}^2m_Z^2\mathrm{tan}^2\beta }{v^2(m_{\stackrel{~}{\nu }_n}^2m_{\stackrel{~}{\nu }_m}^2)(m_{\stackrel{~}{\nu }_k}^2m_{\stackrel{~}{\nu }_m}^2)(m_A^2m_{\stackrel{~}{\nu }_m}^2)(m_h^2m_{\stackrel{~}{\nu }_m}^2)(m_H^2m_{\stackrel{~}{\nu }_m}^2)}}`$ (49)
$`\times \{|v\times b|^2[m_{\stackrel{~}{\nu }_m}^4m_{\stackrel{~}{\nu }_m}^2\text{Tr}(M_{\stackrel{~}{\nu }\stackrel{~}{\nu }^{}}^2)\frac{1}{2}[\text{Tr}(M_{\stackrel{~}{\nu }\stackrel{~}{\nu }^{}}^4)[\text{Tr}(M_{\stackrel{~}{\nu }\stackrel{~}{\nu }^{}}^2)]^2]`$
$`+b^2(v\times b)(v\times c)[m_{\stackrel{~}{\nu }_m}^2bc\text{Tr}(M_{\stackrel{~}{\nu }\stackrel{~}{\nu }^{}}^2)]+b^2v_d^2|b\times c|^2\},`$
where $`nkm`$. The traces in the RPC limit are given by:
$`\text{Tr}(M_{\stackrel{~}{\nu }\stackrel{~}{\nu }^{}}^2)`$ $`=`$ $`|b|\mathrm{tan}\beta +{\displaystyle \underset{k=1}{\overset{3}{}}}m_{\stackrel{~}{\nu }_k}^2,`$ (50)
$`\text{Tr}(M_{\stackrel{~}{\nu }\stackrel{~}{\nu }^{}}^4)`$ $`=`$ $`b^2\mathrm{tan}^2\beta +{\displaystyle \underset{k=1}{\overset{3}{}}}m_{\stackrel{~}{\nu }_k}^4.`$ (51)
Again, we can check that in the special basis \[see Appendix B\], eqs. (5) and (49) agree. The extension of these results to four and more generations is straightforward.
From the results above, we learn that if the vectors $`v`$ and $`b`$ are parallel, then $`\mathrm{\Delta }m_{\stackrel{~}{\nu }_m}^2=0`$ for all $`m`$. This is obvious for the one and two generation cases, since we may then put $`v\times b=0`$ in eqs. (26) and (43). In the three-generation case, it is sufficient to note that the vectors $`b`$ and $`c`$ are also parallel; it then again follows \[see eq. (49)\] that $`\mathrm{\Delta }m_{\stackrel{~}{\nu }_m}^2=0`$. In fact, if $`v`$ and $`b`$ are parallel, then all the vectors $`c^{(n)}`$ \[eq. (19)\] are simultaneously parallel to $`v`$. As a result, $`\mathrm{\Delta }m_{\stackrel{~}{\nu }_m}^2=0`$ for all $`m`$ in the case of an arbitrary number of generations. Conversely, if $`v`$ and $`b`$ are not parallel, then there must exist at least one sneutrino/anti-sneutrino pair that is split in mass, as a consequence of the sum rule derived in eq. (E1).
## V The Degenerate Case
So far we have assumed that the sneutrinos are non-degenerate. We expect this assumption to hold in any realistic model, since the sneutrino/antisneutrino mass splittings are of order the neutrino masses. Thus in order for the result of the previous section not to hold, the flavor degeneracy has to be very good, namely the mass splitting between different sneutrino flavors should be much smaller than the neutrino mass. In any realistic model, we do not expect such a high degree of degeneracy. Even in models where supersymmetry breaking is flavor blind, a mass-splitting between sneutrino flavors will be generated via renormalization group (RG) evolution (from the scale of primordial supersymmetry breaking to the electroweak scale) that is proportional to the corresponding charged lepton masses. Even a very small amount of running is sufficient to generate a mass splitting that is many orders of magnitude larger than the neutrino mass.
Nevertheless, as a mathematical exercise and for completeness, we generalize the results of the previous sections to the case of degenerate sneutrinos. First we give a basis-dependent argument that explains how one can obtain the sneutrino squared-mass splittings in the degenerate case from the results already obtained in the non-degenerate case without any additional calculation. A basis-independent proof is relegated to Appendix F.
Consider a case with $`n_f`$ sneutrinos, of which $`n_d`$ sneutrinos are degenerate in mass (where $`2n_dn_f`$) in the RPC limit. Consider the $`n_d`$ degenerate sneutrinos and their corresponding antisneutrinos. Of these, $`n_d1`$ sneutrino/antisneutrino pairs remain degenerate when RPV effects are included, while one pair is spilt in mass.<sup>g</sup><sup>g</sup>gThis corrects a misstatement made at the end of Section III in ref. . In total, $`n_fn_d+1`$ sneutrino/antisneutrino pairs are split in mass. The corresponding squared-mass differences are then given by eq. (36) for the $`(n_fn_d+1)`$-generation case, but with all vectors and tensors appearing in the formula promoted to $`n_f`$ dimensions. The proof of this assertion is as follows. For the case of $`n_d`$ degenerate sneutrinos, the matrix $`C`$ \[that appears in eqs. (20) and (21)\] has $`n_d`$ degenerate eigenvalues. Thus, we are free to make arbitrary rotations within the $`n_d`$ dimensional subspace corresponding to the degenerate states. By a suitable rotation, we can choose of basis in which only one of the $`B_i`$ within the degenerate subspace is non-zero. In this basis the CP-odd and the CP-even squared-mass matrices \[eqs. (20) and (21)\] separate into $`(n_d1)`$ and $`(n_fn_d+1)`$-dimensional blocks. Clearly, the sneutrino eigenvalues in the corresponding $`(n_d1)`$-dimensional blocks are not affected by the presence of RPV terms, while the $`(n_fn_d+1)`$-dimensional block can be treated by the methods of Section IV.
Further generalizations, where more than one set of sneutrinos are each separately degenerate, can also be studied. The procedure for computing the resulting sneutrino squared-mass differences is now clear, so we shall not elaborate further.
## VI Discussion
This paper provides formulae for the sneutrino/antisneutrino squared-mass differences at tree-level in terms of basis-independent R-parity-violating (RPV) quantities. In contrast to the neutrino sector, where only one tree-level neutrino mass is generated by RPV-effects, we expect that all sneutrino/antisneutrino squared-mass differences are generated at tree-level with roughly the same order of magnitude. The sneutrino/antisneutrino mass difference is expected to be of the same order of magnitude as the (tree-level) neutrino mass. However, these quantities could be significantly different, as they depend on independent RPV-parameters. One can also analyze the case of degenerate masses for different sneutrino flavors; although this case can only arise as a result of a high degree of fine-tuning of low-energy parameters. The pattern of sneutrino/antisneutrino squared-mass differences would provide some insight into the fundamental origin of lepton flavor at a very high energy scale.
The sneutrino/antisneutrino squared-mass splittings can be explored either directly by observing sneutrino oscillation , or indirectly via its effects on other lepton number violating processes, such as neutrinoless double beta decay and neutrino masses . Moreover, the effects of tree-level sneutrino/antisneutrino squared-mass splittings on neutrino masses are expected to be significant. The neutrino spectrum is determined by the relative size of the different RPV couplings that control three sources of neutrino masses: (i) the tree-level mass, (ii) the sneutrino induced one-loop masses, and (iii) the trilinear RPV induced one-loop masses . Since only one neutrino acquires a tree-level mass, the other two mechanisms are responsible for the masses of the other two neutrinos. In the literature, only the trilinear RPV-induced one-loop masses have been considered in most studies. In refs. and , it is argued that the sneutrino-induced one-loop contributions to the neutrino masses are generically dominant, since the trilinear RPV-induced one-loop masses are additionally suppressed by a factor proportional to the Yukawa coupling squared.
The results of our basis-independent formalism are useful for comparing the two radiative neutrino mass generation mechanisms. In particular, in models in which a theory of flavor determines the structure of the soft-supersymmetry-breaking parameters at some high energy scale, RG-evolution provides the connection between the observed low-energy spectrum and the high-energy values of the fundamental parameters of the theory . Basis-dependent quantities are not renormalization-group invariant; hence the RG-evolution of basis-independent quantities can significantly simplify the analysis. For example, the direction of the vacuum expectation value of the generalized slepton/Higgs scalar field is dynamically generated at each energy scale. Since the model parameters generically depend on the scale, the direction of the vacuum expectation value in the generalized lepton flavor space is scale dependent. Clearly, in the basis-independent approach, such complications are avoided. This will be the subject of a subsequent paper.
A few possible directions for future research are worth noting. First, recall that in this paper, we assumed that CP was conserved in the scalar sector. If CP is violated, the required analysis is more complicated. Instead of diagonalizing separately CP-even and CP-odd squared-mass matrices, one must diagonalize a single squared-mass matrix in which the formerly CP-even and CP-odd states can mix. Then, one must identify the two sneutrino mass eigenstates (in the limit of small RPV couplings). It should be possible to extend the techniques developed in this paper to address this more general case. Second, in exploring the phenomenology of sneutrino interactions (production cross-sections and decay), one can generally assume that RPV-couplings are irrelevant except in the decay of the lightest sneutrino state. In that case, new RPV-couplings enter, in particular the corresponding $`\lambda `$ and $`\lambda ^{}`$ parameters given in eq. (1). In the spirit of this paper, one should also develop a basis-independent formalism to describe the RPV sneutrino decay. We hope to return to some of these issues in a future work.
###### Acknowledgements.
We thank Kiwoon Choi for helpful conversations. YG is supported by the U.S. Department of Energy under contract DE-AC03-76SF00515, and HEH is supported in part by the U.S. Department of Energy under contract DE-FG03-92ER40689.
## A Evaluation of $`𝒀^{\mathbf{(}𝑵\mathbf{)}}\mathbf{(}\mathrm{𝟎}\mathbf{)}`$ for $`𝑵\mathbf{=}\mathrm{𝟏}\mathbf{,}\mathrm{𝟐}`$ and 3
Using eq. (37) for $`Y^{(N)}(0)`$, we provide below the explicit computation for the cases of $`N=1,2`$ and 3. The computation makes use of the definitions of $`B_i`$ and $`C_{ij}`$:
$`B_i`$ $``$ $`{\displaystyle \frac{b_\beta X_{\beta i}}{\mathrm{cos}\beta }},`$ (A1)
$`C_{ij}`$ $``$ $`X_{\alpha j}M_{\alpha \beta }^2X_{\beta i},`$ (A2)
where $`M^2M_{\stackrel{~}{\nu }\stackrel{~}{\nu }^{}}^2`$, and the properties of the $`X_{\alpha i}`$ given in eqs. (11)–(13).
The case of $`N=1`$ is very simple.
$$Y^{(1)}(0)=B^2=\frac{1}{\mathrm{cos}^2\beta }b_\alpha X_{\alpha i}b_\beta X_{\beta i}=\frac{1}{\mathrm{cos}^2\beta }\left[b^2\frac{(bv)^2}{v_d^2}\right]=\frac{b^2}{\mathrm{cos}^2\beta }|\widehat{v}\times \widehat{b}|^2,$$
(A3)
where the product of cross-products, defined in eq. (4), can be used in any number of dimensions.
For the case of $`N=2`$, we compute
$`Y^{(2)}(0)`$ $`=`$ $`B_iB_jC_{ij}B^2\text{Tr}(C)`$ (A4)
$`=`$ $`{\displaystyle \frac{1}{\mathrm{cos}^2\beta }}\left[X_{\alpha i}M_{\alpha \beta }^2X_{\beta j}b_\mu X_{\mu i}b_\nu X_{\nu j}X_{\alpha i}M_{\alpha \beta }^2X_{\beta i}b_\mu X_{\mu j}b_\nu X_{\nu j}\right]`$ (A5)
$`=`$ $`{\displaystyle \frac{1}{v_d^2\mathrm{cos}^2\beta }}\left\{b^2v_d^2(cb)v_ub^2(bv)[b^2v_d^2(bv)^2]\text{Tr}(M^2)\right\}`$ (A6)
$`=`$ $`{\displaystyle \frac{1}{v_d^2\mathrm{cos}^2\beta }}\left\{b^2(v\times b)(c\times v)|v\times b|^2\text{Tr}(M^2)\right\}.`$ (A7)
Note that the resulting expression has simplified considerably after introducing the vector $`c`$ \[defined in eq. (18)\].
The case of $`N=3`$ is more involved.
$$Y^{(3)}(0)=\frac{1}{2}B^2\left[\text{Tr}(C^2)[\text{Tr}(C)]^2\right]+B_iB_jC_{ij}\text{Tr}(C)B_iB_jC_{ik}C_{kj}.$$
(A8)
We calculate separately the two terms above. First, $`B^2`$ is obtained from eq. (A3) and
$`\text{Tr}(C^2)[\text{Tr}(C)]^2`$ $`=`$ $`\left[X_{\alpha i}M_{\alpha \beta }^2X_{\beta j}X_{\mu j}M_{\mu \nu }^2X_{\nu i}X_{\alpha i}M_{\alpha \beta }^2X_{\beta i}X_{\mu j}M_{\mu \nu }^2X_{\nu j}\right]`$ (A9)
$`=`$ $`\text{Tr}(M^4)[\text{Tr}(M^2)]^2{\displaystyle \frac{2v_u}{v_d^2}}\left[v_ub^2(bv)\text{Tr}(M^2)\right].`$ (A10)
Next, we evaluate
$`[B_iB_jC_{ij}\text{Tr}(C)B_iB_jC_{ij}^2]\mathrm{cos}^2\beta `$ (A11)
$`=X_{\alpha i}M_{\alpha \beta }^2X_{\beta i}X_{\mu j}M_{\mu \nu }^2X_{\nu k}b_\sigma X_{\sigma k}b_\rho X_{\rho j}X_{\alpha i}M_{\alpha \beta }^2X_{\beta j}X_{\mu j}M_{\mu \nu }^2X_{\nu k}b_\sigma X_{\sigma k}b_\rho X_{\rho i}`$ (A12)
$`=M_{\mu \nu }^2b_\mu b_\nu \text{Tr}(M^2)M_{\alpha \mu }^2M_{\mu \nu }^2b_\nu b_\alpha {\displaystyle \frac{v_u}{v_d^2}}\left[b^2(bv)\text{Tr}(M^2)(bv)M_{\mu \nu }^2b_\nu b_\alpha \right]`$ (A13)
$`+{\displaystyle \frac{v_u}{v_d^2}}\left[v_ub^4b^2(bv)\text{Tr}(M^2)\right]+{\displaystyle \frac{v_u}{v_d^4}}\left[(bv)^3\text{Tr}(M^2)v_ub^2(bv)^2\right].`$ (A14)
The above result can be simplified further. First, the last two terms can be combined by noting that
$$v_ub^4b^2(bv)\text{Tr}(M^2)+\frac{1}{v_d^2}\left[(bv)^3\text{Tr}(M^2)v_ub^2(bv)^2\right]=\frac{|v\times b|^2}{v_d^2}\left[v_ub^2(bv)\text{Tr}(M^2)\right].$$
(A15)
This term will end up canceling a similar term in eq. (A9).
At this point, it is convenient to re-express some of the terms of eq. (A11) in terms of the vector $`c`$. First, we observe that
$`v_d^2M_{\mu \nu }^2b_\mu b_\nu \text{Tr}(M^2)v_ub^2(bv)\text{Tr}(M^2)`$ $`=`$ $`b^2\left[(bc)v_d^2(bv)(cv)\right]\text{Tr}(M^2)`$ (A16)
$`=`$ $`b^2(v\times c)(v\times b)\text{Tr}(M^2).`$ (A17)
In deriving the above result, we noted that $`M_{\mu \nu }^2b_\nu v_\nu =v_ub^2=b^2(vc)`$ \[using eqs. (6) and (18) and the fact that $`M_{\stackrel{~}{\nu }\stackrel{~}{\nu }^{}}^2`$ is a symmetric matrix\], which implies that
$$v_u=vc.$$
(A18)
Second,
$`v_d^2M_{\alpha \mu }^2M_{\mu \nu }^2b_\nu b_\alpha v_u(bv)M_{\mu \nu }^2b_\nu b_\nu `$ $`=`$ $`b^4c^2v_d^2v_ub^2(bv)(bc)`$ (A19)
$`=`$ $`b^2\left[v_d^2|b\times c|^2+(bc)(v\times b)(v\times c)\right].`$ (A20)
Collecting all of the above results, the final expression is quite compact:
$`Y^{(3)}(0)`$ $`=`$ $`{\displaystyle \frac{1}{v_d^2\mathrm{cos}^2\beta }}\{\frac{1}{2}|v\times b|^2[\text{Tr}(M^4)[\text{Tr}(M^2)]^2]+b^2v_d^2|b\times c|^2`$ (A22)
$`+b^2(v\times c)(v\times b)[\text{Tr}(M^2)(bc)]\}.`$
## B Sneutrino Squared-Mass Splitting Formulae in the Special Basis
We define the special basis in which $`v_m=0`$ (i.e, the neutral scalar vacuum expectation values determines the definition of the down-type Higgs field) and the matrix $`(M_{\stackrel{~}{\nu }\stackrel{~}{\nu }^{}}^2)_{ij}`$ \[which is the $`3\times 3`$ block sub-matrix of $`(M_{\stackrel{~}{\nu }\stackrel{~}{\nu }^{}}^2)_{\alpha \beta }`$ defined in eq. (7)\] is diagonal. As in Appendix A we define $`M^2M_{\stackrel{~}{\nu }\stackrel{~}{\nu }^{}}^2`$.
In the special basis, one can use eqs. (6), (18) and (A18) to obtain the following relations:
$$c_0=\mathrm{tan}\beta ,M_{00}^2=b_0c_0,M_{0i}^2=b_i\mathrm{tan}\beta ,c_i=\frac{b_i(M_{00}^2+M_{ii}^2)}{b^2}.$$
(B1)
Using these results, it follows that
$`|v\times b|^2`$ $`=`$ $`v_d^2{\displaystyle \underset{i}{}}b_i^2,`$ (B2)
$`b^2(v\times b)(v\times c)`$ $`=`$ $`v_d^2\left[M_{00}^2{\displaystyle \underset{i}{}}b_i^2+{\displaystyle \underset{i}{}}M_{ii}^2b_i^2\right].`$ (B3)
We will also need a similar expression for $`|b\times c|^2`$. First, we note that the last relation of eq. (B1) implies
$$c^2=c_0^2+\frac{1}{b^4}\underset{i}{}b_i^2(M_{00}^2+M_{ii}^2)^2,bc=M_{00}^2+\frac{1}{b^2}\underset{i}{}b_i^2(M_{00}^2+M_{ii}^2).$$
(B4)
These results can be used to obtain:
$$|b\times c|^2=b^2c^2(bc)^2=\frac{1}{b^2}\underset{i}{}(b_iM_{ii}^2)^2+𝒪(b_i^4).$$
(B5)
We now turn to the specific cases. For the case of one generation, the basis-independent result is given in eq. (26). In the special basis, eq. (B2) yields
$$b^2|\widehat{v}\times \widehat{b}|^2=b_1^2.$$
(B6)
Inserting this result into eq. (26), we immediately obtain eq. (5).
In the two generation case, the basis-independent result is given in eq. (43). In the special basis, eqs. (B2) and (B3) yield
$`b^2(v\times b)(v\times c)`$ $`=`$ $`v_d^2\left[M_{00}^2(b_1^2+b_2^2)+M_{11}^2b_1^2+M_{22}^2b_2^2\right],`$ (B7)
$`[M_{ii}^2\text{Tr}(M^2)]|v\times b|^2`$ $`=`$ $`v_d^2\left(M_{00}^2+M_{jj}^2\right)(b_1^2+b_2^2),`$ (B8)
for the two cases of $`i=1`$, $`j=2`$ and $`i=2`$, $`j=1`$, respectively. Adding the above two equations, one finds
$$[M_{ii}^2\text{Tr}(M^2)]|v\times b|^2+b^2(v\times b)(v\times c)=(M_{ii}^2M_{jj}^2)b_i^2v_d^2.$$
(B9)
Working to leading order in the RPV-parameters $`b_i^2`$, we may set the diagonal elements of $`M^2`$ to their RPC values, $`M_{ii}^2=m_{\stackrel{~}{\nu }_i}^2`$. Plugging the result into eq. (43), one again recovers eq. (5).
In the three generation case, the basis-independent result is given in eq. (49). Again, it is sufficient to work to leading order in the $`b_i^2`$. Then, one finds that in the special basis,
$`[M_{11}^2]^2M_{11}^2\text{Tr}(M^2)\frac{1}{2}[\text{Tr}(M^4)[\text{Tr}(M^2)]^2]`$ $`=`$ $`M_{00}^2\left[M_{22}^2+M_{33}^2\right]+M_{22}^2M_{33}^2+𝒪(b_i^2),`$ (B10)
$`M_{11}^2bc\text{Tr}(M^2)`$ $`=`$ $`M_{22}^2M_{33}^2+𝒪(b_i^2).`$ (B11)
Using these results and those of eqs. (B2), (B3) and (B5), we end up with
$`\{[M_{11}^2]^2M_{11}^2\text{Tr}(M^2)\frac{1}{2}[\text{Tr}(M^4)[\text{Tr}(M^2)]^2\}|v\times b|^2+b^2v_d^2|b\times c|^2`$ (B12)
$`+b^2(v\times b)(v\times c)\left[M_{11}^2bc\text{Tr}(M^2)\right]=v_d^2b_1^2(M_{11}^2M_{22}^2)(M_{11}^2M_{33}^2).`$ (B13)
Two additional equations can be generated by permuting the indices 1,2 and 3. Finally, setting $`M_{ii}^2=m_{\stackrel{~}{\nu }_i}^2`$ and plugging the result into eq. (49), one confirms eq. (5) for the third time.
## C How to eliminate $`𝒗`$ in favor of other vectors
Our final expressions for the sneutrino squared mass differences depend on basis-independent products of vectors, $`v`$, $`b`$, $`c,\mathrm{}`$, and traces of powers of $`M_{\stackrel{~}{\nu }\stackrel{~}{\nu }^{}}^2`$. However, the vector $`v`$ is not a fundamental parameter of the model, but a derived parameter which arises as a solution to eq. (6). With some manipulation, it is possible to eliminate $`v`$ in favor of the other vectors (which correspond more directly to the fundamental supersymmetric model parameters, namely $`b`$ and a series of vectors obtained by multiplying $`b`$ some number of times by $`M_{\stackrel{~}{\nu }\stackrel{~}{\nu }^{}}^2`$). In this appendix, we illustrate the procedure in the case of the one and two generation models.
In the one generation model, $`M_{\stackrel{~}{\nu }\stackrel{~}{\nu }^{}}^2`$ is a $`2\times 2`$ matrix. Consider an arbitrary $`2\times 2`$ matrix $`A`$ and its characteristic equation $`det(A\lambda I)=0`$. Since any matrix satisfies its own characteristic equation, we obtain<sup>h</sup><sup>h</sup>hWe henceforth suppress the obvious factors of the identity matrix $`I`$.
$$A^2A\text{Tr}(A)+det(A)=0,$$
(C1)
which after multiplication by $`A^1`$ yields
$$A^1=\frac{\text{Tr}(A)A}{det(A)}.$$
(C2)
Using eq. (C2) we can express $`|v\times b|^2`$ in terms of $`|b\times c|^2`$. Let $`AM_{\stackrel{~}{\nu }\stackrel{~}{\nu }^{}}^2`$, and use eqs. (6) and (18) to obtain
$$v=v_uA^1b=\frac{v_u\left[b\text{Tr}(A)b^2c\right]}{det(A)},$$
(C3)
We now substitute eq. (C3) for $`v`$ in $`|v\times b|^2b^2v_d^2(vb)^2`$.
$`b^2v_d^2`$ $`=`$ $`{\displaystyle \frac{b^4v_u^2}{[det(A)]^2}}\left\{[\text{Tr}(A)]^22(bc)\text{Tr}(A)+b^2c^2\right\},`$ (C4)
$`(vb)^2`$ $`=`$ $`{\displaystyle \frac{b^4v_u^2}{[det(A)]^2}}\left\{[\text{Tr}(A)]^22(bc)\text{Tr}(A)+(bc)^2\right\}.`$ (C5)
Subtracting these two equations, we end up with
$$|v\times b|^2=\frac{b^4v_u^2}{[det(A)]^2}|b\times c|^2.$$
(C6)
Since $`|b\times c|^2`$ is the small RPV parameter, we may evaluate $`det(A)`$ in the RPC limit. Using eq. (6) in the RPC limit, $`[det(A)]^2=m_{\stackrel{~}{\nu }}^4b^2\mathrm{tan}^2\beta `$. The end result is
$$|b\times c|^2=m_{\stackrel{~}{\nu }}^4|\widehat{v}\times \widehat{b}|^2.$$
(C7)
In the two generation model, $`M_{\stackrel{~}{\nu }\stackrel{~}{\nu }^{}}^2`$ is a $`3\times 3`$ matrix. The procedure again employs the characteristic equation. For an arbitrary $`3\times 3`$ matrix,
$$A^1=\frac{A^2A\text{Tr}(A)+\text{Sym}_2A}{det(A)},$$
(C8)
where<sup>i</sup><sup>i</sup>iTo prove eq. (C9), simply note that $`\text{Tr}(A^2)=_k\lambda _k^2`$.
$$\text{Sym}_2(A)\underset{i<j}{}\lambda _i\lambda _j=\frac{1}{2}\left\{\left[\text{Tr}(A)\right]^2\text{Tr}(A^2)\right\},$$
(C9)
where $`\lambda _j`$ are the eigenvalues of $`A`$. We can again solve for $`v`$ following the method used in the one-generation case:
$$v=v_uA^1b=\frac{v_u\left[b^2db^2c\text{Tr}(A)+b\text{Sym}_2(A)\right]}{det(A)},$$
(C10)
where we have defined $`dc^{(2)}=Ac`$. After some algebra, we obtain<sup>j</sup><sup>j</sup>jObserve that the result for $`|v\times b|`$ depends on the number of generations, i.e. the dimension of the matrix $`A`$ \[compare eqs. (C6) and (C11)\].
$$|v\times b|^2=\frac{b^4v_u^2}{[det(A)]^2}\left[|b\times d|^2+|b\times c|^2[\text{Tr}(A)]^22(b\times c)(b\times d)\text{Tr}(A)\right],$$
(C11)
and
$`(v\times b)(v\times c)`$ $`=`$ $`{\displaystyle \frac{b^2v_u^2}{[det(A)]^2}}[(d\times b)(d\times c)[b^2\text{Sym}_2(A)]`$ (C13)
$`+b^2(c\times d)(c\times b)\text{Tr}(A)+|b\times c|^2\text{Tr}(A)\text{Sym}_2(A)].`$
It is easy to evaluate the three invariants in the RPC limit:<sup>k</sup><sup>k</sup>kSince $`vc=v_u`$, it follows that in the RPC limit, $`|c|=\mathrm{tan}\beta `$.
$`\text{Tr}(A)`$ $`=`$ $`|b|\mathrm{tan}\beta +m_1^2+m_2^2,`$ (C14)
$`det(A)`$ $`=`$ $`m_1m_2|b|\mathrm{tan}\beta ,`$ (C15)
$`\text{Sym}_2(A)`$ $`=`$ $`m_1^2m_2^2+|b|\mathrm{tan}\beta (m_1^2+m_2)^2,`$ (C16)
where $`m_i^2m_{\stackrel{~}{\nu }_i}^2`$.
Inserting the above results into eq. (43) yields the desired result. Further algebraic manipulations of the resulting expression do not lead to a particularly simple result.
## D Evaluation of $`𝐝𝐞𝐭\mathbf{[}𝑨\mathbf{}\mathbf{(}𝝀_𝒎\mathbf{}\mathit{ϵ}\mathbf{)}𝑰\mathbf{]}`$
Consider a general $`N\times N`$ matrix $`A`$, with eigenvalues $`\lambda _k`$. Then,
$$det(A\lambda I)=\underset{k}{\overset{N}{}}(\lambda _k\lambda ),$$
(D1)
where the product is taken over all $`N`$ eigenvalues (some of which might be degenerate). First, suppose that there are no degenerate eigenvalues. If $`\lambda _m`$ is one of the eigenvalues and $`ϵ1`$, then
$`det[A(\lambda _mϵ)I]`$ $`=`$ $`ϵ{\displaystyle \underset{im}{}}(\lambda _i\lambda _m)+𝒪(ϵ^2)`$ (D2)
$`=`$ $`ϵ\stackrel{}{det}(A\lambda _mI)+𝒪(ϵ^2),`$ (D3)
where $`det^{}M`$ is the product of all the non-zero eigenvalues of $`M`$ \[see eq. (31)\].
The case of degenerate eigenvalues is easily handled. We can still use eq. (31) if it is understood that all terms in which $`\lambda _i`$ is equal to the degenerate eigenvalue are omitted from the product. If $`\lambda _d`$ is an eigenvalue which is $`n_d`$-fold degenerate, then eq. (D2) is generalized to
$$det[A(\lambda _dϵ)I]=ϵ^{n_d}\stackrel{}{det}(A\lambda _dI)+𝒪(ϵ^{n_d+1}).$$
(D4)
In Section IV and Appendix F, we have employed these results with $`ϵ=\delta \lambda _m`$ and $`ϵ=\delta \lambda _d`$, respectively.
## E Sneutrino Squared-Mass Splitting Sum Rules
In the case of $`N`$ sneutrino generations, one can calculate the corresponding sneutrino squared-mass splittings. In the case of non-degenerate tree-level sneutrino masses, the squared-mass splittings were obtained in eq. (36). In the case of degenerate masses, one employs the modified results according to the discussion given in Section V and Appendix F. We then find the following interesting sum rule:
$$\underset{m=1}{\overset{N}{}}\frac{v^2(m_A^2m_{\stackrel{~}{\nu }_m}^2)(m_h^2m_{\stackrel{~}{\nu }_m}^2)(m_H^2m_{\stackrel{~}{\nu }_m}^2)}{4m_Z^2\mathrm{tan}^2\beta }\frac{\mathrm{\Delta }m_{\stackrel{~}{\nu }_m}^2}{m_{\stackrel{~}{\nu }_m}^2}=|v\times b|^2.$$
(E1)
We shall prove this result for the non-degenerate case. Using eqs. (36) and (A3), we see that eq. (E1) is equivalent to the following result:
$$\underset{m=1}{\overset{N}{}}\frac{Y^{(N)}(\lambda _m)}{_{im}(\lambda _i\lambda _m)}=Y^{(1)}(0).$$
(E2)
To prove eq. (E2), we insert the expansion for $`Y^{(N)}(\lambda _m)`$ \[eq. (40)\] into eq. (E2), and make use the following identity:
$$S_{N,k}\underset{m=1}{\overset{N}{}}\frac{\lambda _m^k}{_{im}(\lambda _m\lambda _i)}=\{\begin{array}{cc}0,\hfill & k=0,1,\mathrm{},N2\text{ ,}\hfill \\ 1,\hfill & k=N1\text{ ,}\hfill \end{array}$$
(E3)
where all the $`\lambda _m`$ are assumed to be distinct.<sup>l</sup><sup>l</sup>lNote that in eq. (E3), the sign of the factors $`\lambda _m\lambda _i`$ is reversed compared to eq. (E2). Thus, an extra factor of $`(1)^{N1}`$ is generated which cancels with the corresponding sign in front of $`Y^{(1)}`$ in eq. (40).
Eq. (E3) is established as follows. Let $`f(x)=(x\lambda _1)(x\lambda _2)\mathrm{}(x\lambda _N)`$, where the $`\lambda _m`$ are distinct. Consider the resolution of $`x^{k+1}/f(x)`$ into partial fractions (where $`k`$ is an integer such that $`0kN2`$):
$$\frac{x^{k+1}}{f(x)}=\underset{m=1}{\overset{N}{}}\frac{A_m}{x\lambda _m}.$$
(E4)
Combining denominators, it follows that:
$$x^{k+1}=\underset{m=1}{\overset{N}{}}A_m\underset{im}{}(x\lambda _i).$$
(E5)
The right hand side of eq. (E5) is a polynomial of degree $`N1`$ or less. Since this must be an identity for all $`x`$, we can solve for each coefficient $`A_m`$ separately by setting $`x=\lambda _m`$:
$$A_m=\frac{\lambda _m^{k+1}}{_{im}(\lambda _m\lambda _i)}.$$
(E6)
Inserting this result into eq. (E4) and setting $`x=0`$ yields eq. (E3) for the case of $`0kN2`$. The case of $`k=0`$ where one of the $`\lambda _m`$ vanishes must be treated separately, although it is easy to show that the end result is unchanged. Thus, $`S_{N,k}=0`$ for $`0kN2`$.
To derive eq. (E3) in the case of $`k=N1`$, we set $`k=N2`$ in eq. (E5). On the right hand side of eq. (E5), we note that the term proportional to $`x^{N1}`$ arises simply by setting the $`\lambda _i=0`$. It follows that $`_{m=1}^NA_m=1`$ (for $`k=N2`$) which is precisely equivalent to $`S_{N,N1}=1`$ and the proof is complete.
Finally, we note a useful recursion relation satisfied by the $`S_{N,k}`$. Multiply the $`m`$th term of eq. (E3) by $`(\lambda _m\lambda _{N+1})/(\lambda _m\lambda _{N+1})`$. One immediately deduces that relation:
$$S_{N,k}=S_{N+1,k+1}\lambda _{N+1}S_{N+1,k}.$$
(E7)
The boundary conditions for the recursion relation are: $`S_{N,0}=0`$ for $`N2`$ (which is a consequence of the proof given above), and $`S_{1,0}=1`$.<sup>m</sup><sup>m</sup>mThe condition $`S_{1,0}=1`$ formally defines the sum \[eq. (E3)\] in the case of $`N=1`$. Alternatively, one can check by explicit evaluation that $`S_{2,1}=1`$. Thus, we see that the assigned definition of $`S_{1,0}`$ is consistent. Note that one can similarly define $`S_{1,k}=\lambda _1^k`$. It follows that $`S_{N,k}=0`$ for $`1kN2`$. Choosing $`k=N1`$ in eq. (E7) then yields $`S_{N+1,N}=S_{N,N1}`$; it follows that $`S_{N+1,N}=1`$ for all $`N1`$. Finally, it is easy to increase $`k`$ further. For example, eq. (E7) implies that $`S_{N+1,N+1}=S_{N,N}+\lambda _{N+1}`$. It follows that $`S_{N,N}=_{i=1}^N\lambda _i`$, and so on.
The degenerate case can be treated as a limiting case of the non-degenerate results obtained above. In the final analysis, we find that eq. (E1) applies in general, and serves as a useful check of our results.
## F Basis-Independent Treatment of the Degenerate Case
The degenerate case for $`n_f`$ flavors was treated in Section V. If $`n_d`$ sneutrino/antisneutrino pairs are degenerate in mass in the RPC limit, then when RPV-effects are incorporated, one finds that $`n_d1`$ pairs remain degenerate, while $`n_fn_d+1`$ sneutrino/antisneutrino pairs are split in mass. The squared-mass splittings of the latter can be obtained from the corresponding formulae of the non-degenerate $`n_fn_d+1`$ flavor case.
In this appendix, we briefly sketch the required steps of a proof that generalizes the basis-independent results of Section IV. Consider first the squared-mass matrix of the CP-odd scalars \[eq. (20)\]. The characteristic equation, eq. (28), is still valid in the case of degenerate sneutrinos. First we consider the quantity $`Y^{(N)}(\lambda )`$ \[eq. (29)\]. Suppose that the matrix $`C`$ which appears in $`\stackrel{~}{M}_{\mathrm{odd}}^2`$ has an eigenvalue $`\lambda _d`$ that is $`n_d`$-fold degenerate (with the remaining eigenvalues of $`C`$ distinct). We assert that the following formula holds:
$$Y^{(N)}(\lambda )=(\lambda _d\lambda )^{n_d1}Y_{\mathrm{deg}}^{(Nn_d+1)}(\lambda ),$$
(F1)
where $`Y_{\mathrm{deg}}^{(Nn_d+1)}(\lambda )`$ is obtained as follows. First, one evaluates $`Y^{(Nn_d+1)}(\lambda )`$ as in Section IV, and expresses the result covariantly in terms of the vector $`B_i`$ and the matrix $`C_{ij}`$. Next, these quantities are reinterpreted as $`N`$-dimensional objects. Finally, all traces that appear in the result are replaced by:
$$\text{Tr}^{}C^n\text{Tr}C^n(n_d1)\lambda _d^n.$$
(F2)
Consider first the effect of the RPV terms on the non-degenerate sneutrinos. Then the analysis of Section IV can be used, and we obtain eq. (36) for the squared-mass splitting of sneutrino/antisneutrino pairs. If we now insert eq. (F1) for $`Y^{(N)}(m_{\stackrel{~}{\nu }_m}^2)`$, we see that we obtain a new formula which has the same form as eq. (36), with the following modifications: (i) $`Y^{(N)}(m_{\stackrel{~}{\nu }_m}^2)`$ is replaced by $`Y_{\mathrm{deg}}^{(Nn_d+1)}(m_{\stackrel{~}{\nu }_m}^2)`$; and (ii) the product that appears in the denominator of eq. (36) is modified to $`_{im}^{}(m_{\stackrel{~}{\nu }_i}^2m_{\stackrel{~}{\nu }_m}^2)`$, where the prime indicates that degenerate squared-masses appear only once in the product. To obtain a covariant expression for $`Y_{\mathrm{deg}}^{(Nn_d+1)}(m_{\stackrel{~}{\nu }_m}^2)`$, we first obtain the expression of $`Y^{(Nn_d+1)}(m_{\stackrel{~}{\nu }_m}^2)`$ in terms of the various vectors ($`v`$, $`b`$, $`c,\mathrm{}`$) and traces of powers of $`M^2`$ using the results of Section IV and Appendix A. The resulting expression can then be used for $`Y_{\mathrm{deg}}^{(Nn_d+1)}(m_{\stackrel{~}{\nu }_m}^2)`$ by replacing $`\text{Tr}M^{2n}`$ with $`\text{Tr}^{}M^{2n}`$ \[the latter is defined by replacing $`C`$ with $`M^2`$ in eq. (F2)\] and interpreting all the vectors and matrices as $`N`$-dimensional objects.
Finally, consider the effect of the RPV terms on the $`n_d`$ degenerate sneutrinos. Now, we must return to eq. (28) and insert $`\lambda =\lambda _d^{(0)}+(\delta \lambda _d)_{\mathrm{odd}}`$. Working to the lowest non-trivial order in the $`B_i`$, we make use of eqs. (F1) and (D4) to obtain
$$\left[(\delta \lambda _d)_{\mathrm{odd}}\right]^{n_d}=\frac{Y_{\mathrm{deg}}^{(Nn_d+1)}(\lambda _d)\left[(\delta \lambda _d)_{\mathrm{odd}}\right]^{n_d1}}{(A\lambda _d)det^{}(C\lambda _dI)}.$$
(F3)
The solution to this equation has $`n_d1`$ degenerate solution, $`(\delta \lambda _d)_{\mathrm{odd}}=0`$, and one non-degenerate solution for $`(\delta \lambda _d)_{\mathrm{odd}}`$ which has the same form as eq. (30) for the $`(Nn_d+1)`$-dimensional problem. As described above, we can make use of the relevant covariant expressions obtained in Section IV and Appendix A by replacing Tr with $`\text{Tr}^{}`$ and interpreting all the vectors and matrices as $`N`$-dimensional objects. The end result is that $`n_d1`$ sneutrino/antisneutrino pairs remain degenerate, while one of the original degenerate pairs is split according to the $`Nn_d+1`$-dimensional version of eq. (36) \[with all vectors and tensors promoted to $`N`$-dimensional objects\]. One can also check that the sum rule obtained in Appendix E for sneutrino squared-mass differences (appropriately weighted) applies even when there are degenerate sneutrino masses.
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# 1 Introduction
## 1 Introduction
The Thirring model has been investigated by many people. It is well-known that the Thirring model is an exactly solvable quantum field theory model in (1+1) dimensions . An extensive investigation of the model was given by Hagen and Klaiber . Hagen introduced an external field and gave the general solution of the Thirring model. Klaiber analyzed the Thirring model and found the operator solutions which are expressed in terms of a free massless Dirac field. He constructed the solution to fulfill the positive definiteness. On the other hand, Nakanishi expressed the solution in terms of the free massless bosonic field . He asserted that all Heisenberg operators should be expressed in terms of asymptotic fields from the standpoint of the general principle of quantum field theory. In the present paper, we use the bosonic expression (bosonization) .
One of the methods for solving the quantum field theory is to determine the Operator Product Expansion (OPE) . In this formalism, the short distance behavior for products of the two local fields is important. For some case, we can determine the OPE exactly, e.g. (1+1)-dimensional Conformal Field Theory . Concerning the operator products of the quantum field, there are difficulties with respect to the current regularization. In most cases, the current is defined by the limiting procedures as $`\overline{\psi }(x+ϵ)\gamma _\mu \psi (x)`$. In the Thirring model, there are several definitions of the current. For example, the Schwinger current is defined by limiting from spacelike direction only. The Johnson current is defined by limiting from not only spacelike direction but also timelike direction symmetrically as
$$j^\mu (x)=\frac{1}{2}\left[j^\mu (x;ϵ)+j^\mu (x;\stackrel{~}{ϵ})\right]_{ϵ,\stackrel{~}{ϵ}0},$$
(1)
where $`ϵ`$ and $`\stackrel{~}{ϵ}`$ are a timelike and spacelike vectors, respectively. Both current definitions are consistent with the solution of the Thirring model. However, the coupling constant is affected by the current definition. Therefore, in the Thirring model, the coupling constant is determined only when we define the current regularization . It is also noted that these coupling constants are not independent. The coupling constant of the Schwinger definition $`g_\mathrm{S}`$ is given by
$$g_\mathrm{S}=\frac{g_\mathrm{J}}{1g_\mathrm{J}/2\pi },$$
(2)
where $`g_\mathrm{J}`$ is the coupling constant of the Johnson definition. These current ambiguities also appear in the massive Thirring model, which we do not understand yet .
On the other hand, Dell’Antonio, Frishman and Zwanziger analyzed the Thirring model without looking into the structure of the current. They extend the Johnson result . They start with defining the commutation relations of the current, current algebra formulation . There are three parameters and we have two relations among them. Therefore, we can construct the current algebra from the lagrangian and the suitable definition of the current which has one parameter.
In this paper, we present a new current regularization of the Thirring model. We introduce one parameter in the definition. Our formulation is simpler than Klaiber’s one and the new current definition is consistent with other formulations. The Thirring current and field can be written in terms of the free massless bosonic field. Therefore, we can analyze the model exactly.
In this paper, we employ the following notation:
$$x^\pm =x^0\pm x^1,x_\pm =x^{}/2,_\pm =(_0\pm _1)/2,^\pm =2_{}$$
(3)
and gamma matrices are
$$\gamma ^0=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),\gamma ^1=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),\gamma ^5=\gamma ^0\gamma ^1=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).$$
(4)
The anti-symmetric tensor $`ϵ_{\mu \nu }`$ is taken to be $`ϵ_{10}=ϵ^{01}=1`$.
## 2 Thirring model
The Thirring model is (1+1) dimensional field theory with the current-current interaction. The lagrangian of the Thirring model is given by
$$=\overline{\psi }\mathrm{i}_\mu \gamma ^\mu \psi \frac{g}{2}j_\mu j^\mu ,j_\mu =\overline{\psi }\gamma _\mu \psi ,$$
(5)
where $`g`$ is a coupling constant. Then, the equations of motion become
$$\mathrm{i}_+\psi _\mathrm{R}=gj_\mathrm{L}\psi _\mathrm{R},\mathrm{i}_{}\psi _\mathrm{L}=gj_\mathrm{R}\psi _\mathrm{L},$$
(6)
where
$$\psi ^\mathrm{T}=(\psi _\mathrm{R},\psi _\mathrm{L}),j_\mathrm{R}\frac{1}{2}\left(j^0+j^1\right),j_\mathrm{L}\frac{1}{2}\left(j^0j^1\right).$$
(7)
From eq.(6), the current $`j^\mu `$ and its dual (axial) current $`\stackrel{~}{j}^\mu =ϵ^{\mu \nu }j_\nu `$ is conserved,
$$_\mu j^\mu =0,_\mu \stackrel{~}{j}^\mu =0.$$
(8)
The Thirring model is exactly solvable. Thirring constructed the eigenstates while Klaiber found the operator solution. On the other hand, Nakanishi described the quantum operator solution of eq.(6) in terms of the free massless bosonic field $`\phi =\phi _\mathrm{R}(x^{})+\phi _\mathrm{L}(x^+)`$ as
$$\psi (x)=\frac{Z}{\sqrt{2\pi }}\left(\begin{array}{c}\text{:}e^{\mathrm{i}s\phi _\mathrm{R}\mathrm{i}\overline{s}\phi _\mathrm{L}}\text{:}\\ \text{:}e^{\mathrm{i}\overline{s}\phi _\mathrm{R}\mathrm{i}s\phi _\mathrm{L}}\text{:}\end{array}\right),$$
(9)
where $`s,\overline{s}`$ are constant parameters and $`Z`$ is a normalization factor. The free bosonic field satisfies
$$_\mu ^\mu \phi =0$$
(10)
and we can regularize as
$$[\phi _\mathrm{R}^{}(x^{}),\phi _\mathrm{R}^{}(y^{})]=\frac{1}{4\pi }\mathrm{ln}\mathrm{i}(x^{}y^{}\mathrm{i0}),$$
(11)
$$[\phi _\mathrm{L}^{}(x^+),\phi _\mathrm{L}^{}(y^+)]=\frac{1}{4\pi }\mathrm{ln}\mathrm{i}(x^+y^+\mathrm{i0}),$$
(12)
where $`\phi _{\mathrm{R},\mathrm{L}}^{}`$ and $`\phi _{\mathrm{R},\mathrm{L}}^{}`$ are the positive and the negative frequency part respectively. Therefore, we have the Operator Product Expansion (OPE) of the Thirring operator,
$`\psi _\mathrm{R}(x)\psi _\mathrm{R}(y)`$ $`=`$ $`{\displaystyle \frac{|Z|^2}{2\pi }}\mathrm{i}^{\frac{s^2+\overline{s}^2}{4\pi }}(x^{}y^{}\mathrm{i0})^{s^2/4\pi }(x^+y^+\mathrm{i0})^{\overline{s}^2/4\pi }`$ (13)
$`\times \text{:}e^{\mathrm{i}s\phi _\mathrm{R}(x)\mathrm{i}\overline{s}\phi _\mathrm{L}(x)+\mathrm{i}s\phi _\mathrm{R}(y)\mathrm{i}\overline{s}\phi _\mathrm{L}(y)}\text{:},`$
$`\psi _\mathrm{R}^{}(y)\psi _\mathrm{R}(x)`$ $`=`$ $`{\displaystyle \frac{|Z|^2}{2\pi }}\mathrm{i}^{\frac{s^2+\overline{s}^2}{4\pi }}(y^{}x^{}\mathrm{i0})^{s^2/4\pi }(y^+x^+\mathrm{i0})^{\overline{s}^2/4\pi }`$ (14)
$`\times \text{:}e^{\mathrm{i}s\phi _\mathrm{R}(y)+\mathrm{i}\overline{s}\phi _\mathrm{L}(y)+\mathrm{i}s\phi _\mathrm{R}(x)\mathrm{i}\overline{s}\phi _\mathrm{L}(x)}\text{:}`$
and so on. We have the similar relation for $`\psi _\mathrm{L}`$ if $`\overline{s}s`$. For the massless Dirac field case ($`g=0`$), we find $`s=2\sqrt{\pi }`$ and $`\overline{s}=0`$.
To solve the model, we must determine the parameters $`s,\overline{s}`$. The first condition of $`s`$ and $`\overline{s}`$ is given by the aniti-commutativity of $`\psi `$ and we have
$$\frac{s^2\overline{s}^2}{4\pi }=1.$$
(15)
## 3 Current regularization
Next, we insert the operator solution into the field equation eq.(6). To do this, we propose the following current definition,
$`j^\mu (x)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[\overline{\psi }(x+\epsilon )\gamma ^\mu \psi (x)+\overline{\psi }(x+\stackrel{~}{\epsilon })\gamma ^\mu \psi (x)\right]`$ (16)
$`{\displaystyle \frac{\sigma }{2}}\alpha _\nu ^\mu \left[\overline{\psi }(x+\epsilon )\gamma ^\nu \psi (x)\overline{\psi }(x+\stackrel{~}{\epsilon })\gamma ^\nu \psi (x)\right],`$
where $`\epsilon (\stackrel{~}{\epsilon })`$ is an infinitesimal timelike (spacelike) vector and $`\sigma `$ is a parameter and $`\alpha _{\mathrm{\hspace{0.33em}\hspace{0.33em}0}}^0=\alpha _{\mathrm{\hspace{0.33em}\hspace{0.33em}1}}^1=1`$, $`\alpha _{\mathrm{\hspace{0.33em}\hspace{0.33em}1}}^0=\alpha _{\mathrm{\hspace{0.33em}\hspace{0.33em}0}}^1=0`$. Here, we get timelike vector $`\epsilon ^\mu `$ close to zero with $`\epsilon ^10`$ firstly, whereas the spacelike vector $`\stackrel{~}{\epsilon }`$ is done in an opposite way. In our formulation, the current is written by
$$j_\mathrm{R}=\frac{s\sigma \overline{s}}{2\pi }_{}\phi _\mathrm{R},j_\mathrm{L}=\frac{s\sigma \overline{s}}{2\pi }_+\phi _\mathrm{L}.$$
(17)
Therefore, the operator solution (9) is valid if
$$\overline{s}=\frac{g}{2\pi }(s\sigma \overline{s}).$$
(18)
Finally, we have the equations which must be satisfied by the parameters of the solution
$$\frac{s^2\overline{s}^2}{4\pi }=1,\overline{s}=\frac{g}{2\pi }(s\sigma \overline{s}).$$
(19)
First, we consider $`\sigma =0`$ and $`\sigma =1`$ case. The second equation becomes
$$\overline{s}=\frac{g_{\sigma =0}}{2\pi }s,\overline{s}=\frac{g_{\sigma =1}}{2\pi }(s\overline{s}).$$
(20)
They can identify with the map,
$$g_{\sigma =1}=\frac{g_{\sigma =0}}{1g_{\sigma =0}/2\pi }.$$
(21)
This is nothing but the relation between the coupling constant of the Schwinger definition and that of the Johnson definition. This is consistent with Klaiber’s result . Therefore, $`\sigma =1`$ corresponds to Schwinger’s current definition and $`\sigma =0`$ is Johnson’s one in our formulation.
We can also calculate the commutation rules between the current and the spinor field $`\psi `$,
$$[\psi (x^1,t),j^0(y^1,t)]=\frac{(s\sigma \overline{s})(s+\overline{s})}{4\pi }\delta (x^1y^1)\psi (x)$$
(22)
and
$$[\psi (x^1,t),j^1(y^1,t)]=\frac{(s\sigma \overline{s})(s\overline{s})}{4\pi }\delta (x^1y^1)\gamma ^5\psi (x).$$
(23)
## 4 Comparison with other formulations
Dell’Antonio, Frishman and Zwanziger analyzed the Thirring model in a different way. They consider the commutation relations of the current, the Schwinger term. We can identify their result with
$$a=\frac{(s\sigma \overline{s})(s+\overline{s})}{4\pi },\overline{a}=\frac{(s\sigma \overline{s})(s\overline{s})}{4\pi },c=\left(\frac{s\sigma \overline{s}}{2\pi }\right)^2,$$
(24)
where $`a`$, $`\overline{a}`$ and $`c`$ are parameters in their formulation (Note that in , $`ϵ_{\mu \nu }`$ is defined by $`ϵ_{10}=1`$). It is easy to check the consistency condition, eq.(6.1) in ,
$$a\overline{a}=gc.$$
(25)
$`c`$ is written in terms of the coupling constant $`g`$ as
$$1/c=\pi \left[1+\frac{g}{2\pi }(\sigma 1)\right]\left[1+\frac{g}{2\pi }(\sigma +1)\right].$$
(26)
If $`\sigma =0`$, it becomes eq.(6.3) in . Therefore, our result perfectly agrees with Dell’Antonio et al. and the parameter of the current commutation relation is determined by the coupling constant and the parameter $`\sigma `$ appeared in the current definition.
Taguchi, Tanaka and Yamamoto consider the Thirring model with the Tomonaga-Schwinger equation. They consider the deformed hamiltonian and calculate the commutation relations between the current and the spinor field. In this case, we have eq.(22) and (23) in a similar way.
It is well-known that the Thirring model is $`c=1`$ ($`c`$ is the central charge) Conformal Field Theory (CFT) . Klassen and Melzer argued that the Thirring model is equivalent to the fermionic Gaussian CFT. They show the relation between the compactification radius of the fermionic Gaussian CFT and the Thirring coupling constant. We give their result with $`\sigma =0`$ (Johnson current). More generally, the compactification radius $`R`$ is written by
$$R=\frac{1}{\sqrt{\left[1+{\displaystyle \frac{g}{2\pi }}(\sigma 1)\right]\left[1+{\displaystyle \frac{g}{2\pi }}(\sigma +1)\right]}}\left[1\pm \frac{g}{2\pi }\sqrt{1\frac{4\pi }{g}\sigma \sigma ^2}\right].$$
(27)
## 5 Conclusion
We have presented the generalization of the current regularization in the Thirring model. The definition of the current is complicated, but it becomes simple when it is expressed in terms of the free massless bosonic field. The present description is consistent with known results of the Thirring model. The present formulation is simpler than Klaiber’s formulation. Klaiber defines the commutator of the current $`j^\mu `$ and the field $`\psi `$, and makes the anzatz about the current while we employ the operator solution which is given by Nakanishi . The solution is written in terms of the free massless bosonic field, and thus we can easily evaluate various quantities. This is the main difference between Klaiber’s treatment and ours. Further, we obtain the general formula for arbitrary current regularization.
The short distance behavior of the Thirring model is more complicated than the free massless Dirac field. For the Dirac field, the limiting procedures of $`\epsilon 0`$ and $`\stackrel{~}{\epsilon }0`$ are the same. On the other hand, they are different for the Thirring field. This is the consequence of the fact that the Dirac field ($`\overline{s}=0`$) is written in terms of the bosonic field $`\phi _\mathrm{R}`$ and $`\phi _\mathrm{L}`$ separately in contrast with the Thirring field.
In the present description, we introduce a new parameter $`\sigma `$ in our current definition. This current becomes Lorentz covariant limiting operator, $`j^\mu ^\mu \stackrel{~}{\phi }`$. Here, $`\stackrel{~}{\phi }`$ is the dual massless field of $`\phi `$. We obtain the Johnson current for $`\sigma =0`$ and the Schwinger current for $`\sigma =1`$. Note that, for $`\sigma =1`$, the current definition is given by
$$j^0(x)=\overline{\psi }(x+\stackrel{~}{\epsilon })\gamma ^0\psi (x),j^1(x)=\overline{\psi }(x+\epsilon )\gamma ^1\psi (x).$$
(28)
Therefore, in our formulation, $`j^0`$ is defined by limiting from spacelike direction while $`j^1`$ is defined by limiting from timelike direction. This is in contrast with the original Schwinger’s definition which is defined by the spacelike separation only. However, we can show that $`j^0`$ and $`j^1`$ commute with each other for the Thirring case if we adopt the original prescription. On the other hand, in our formulation, the Schwinger term appear as
$$[j_0(x^1,t),j_1(y^1,t)]=\mathrm{i}c\delta ^{}(x^1y^1),$$
(29)
where $`c`$ is given by eq.(26). Accordingly, the current must be defined by limiting from both spacelike and timelike direction in the Thirring model. Note that the parameter $`\sigma `$ determines the current algebra. Unfortunately, we do not understand the physical meaning of the parameter $`\sigma `$ yet.
ACKNOWLEDGMENTS
We would like to thank T. Fujita and M. Hiramoto for helpful discussions and comments.
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# 1 Introduction
## 1 Introduction
In this paper we continue the work presented in Ref. . There we computed two-point functions of chiral operators $`\mathrm{Tr}\mathrm{\Phi }^3`$ in $`𝒩=4`$ $`SU(N)`$ supersymmetric Yang-Mills theory to the order $`g^4`$ in perturbation theory and proved that corrections vanish for all values of $`N`$. In our perturbative approach, the fact that the non renormalization occurs for any value of $`N`$, not only at leading order in the $`N\mathrm{}`$ limit, was strictly connected to the vanishing of the colour combinatoric factor of the nonplanar diagrams. The open question was to see if the same pattern were true also for correlators of operators $`\mathrm{Tr}\mathrm{\Phi }^k`$, with general $`k`$.
We have found that for $`k>3`$ the colour structure of non planar diagrams at order $`g^4`$ does not vanish. It is only a complicated cancellation between $`1/N`$ subleading contributions from planar diagrams and nonplanar ones which allows to obtain a complete all $`N`$ nonrenormalization of the two-point correlators.
The colour structure identities we have used in our calculations can be generalized to arbitrary gauge groups. In so doing we have been able to prove that the nonrenormalization of the correlators up to order $`g^4`$ is valid for any group.
The other issue we focus on in this paper, is related to the presence of contact term contributions . They arise in the computation of the correlators because the theory is affected by ultraviolet divergences. In order to control such infinities one has to choose a regularization scheme: at order $`g^0`$, i.e. at the tree level, the singularity of the two-point function needs to be subtracted and this leads to the introduction of arbitrary contact terms. At higher orders in $`g`$ (we have computed the two-point functions explicitly up to the order $`g^4`$) divergences cancel out and as a consequence, the local contact terms are determined unambiguously. Thus the contact terms in the two-point correlators are of the form
$$[a+f(g^2,N)]^p\delta ^{(4)}(xy)$$
(1.1)
with $`a`$ the arbitrary finite constant of the subtraction and $`p`$ an integer depending on the dimensions of the operators. The function $`f(g^2,N)`$ is in principle exactly computable order by order in $`g^2`$ for any finite $`N`$, such that $`f(g^2=0,N)=0`$.
According to the AdS/CFT prescription one should establish a correspondence between these terms and corresponding ones in the supergravity sector. There ambiguities arise due to the fact that the theory suffers from infrared divergences when approaching the boundary of the AdS space. Within the holographic viewpoint proposed in , the supergravity effective action, evaluated on a solution of the equations of motion with prescribed boundary conditions, becomes the generating functional for the conformal field theory in the large $`N`$ limit. The bulk fields $`\varphi `$ evaluated at the boundary act as source terms for composite operators $`𝒪`$ of the Yang-Mills theory
$$e^{{\scriptscriptstyle d^4x𝒪\varphi _0}}_{\mathrm{CFT}}=e^{S[\varphi _0]}$$
(1.2)
with the boundary action given by
$$S[\varphi _0]=d^4xd^4y\varphi _0(x)\left[\frac{1}{(xy)^{2\mathrm{\Delta }}}+b(^2)^{\mathrm{\Delta }2}\delta ^{(4)}(xy)\right]\varphi _0(y)$$
(1.3)
The local term, which acts as a regulator at short distances in (1.3), has an infrared origin in the 5d AdS supergravity expanded near the boundary . The coefficient $`b`$ is a function of a mass scale parameter identified with the inverse of the IR 5d cutoff .
Now, comparing eq. (1.1) at small coupling with eq. (1.3) one should find that in the large $`N`$ limit, with the ’t Hooft coupling $`g^2N`$ fixed but large
$$a+f(g^2,N)b_f$$
(1.4)
where $`b_f`$ is the arbitrary finite part of the subtraction in (1.3). In the context of our specific calculation we will show how this identification can be realized.
We present this paper as a sequel to the one referred in , since the techniques used here to perform the perturbative calculations of the correlators are the same as the ones used in . Therefore in order to avoid lengthy repetitions in the main text, and at the same time to make this paper self-contained, we have recollected the main formulas and the rules of the game in the Appendices. In the next Section we simply remind the reader which are the quantities we have to focus on. Then in Section 3 we enter directly in medias res: we present the tree–level and the order $`g^2`$ result for the $`<\mathrm{Tr}(\mathrm{\Phi }^1)^k\mathrm{Tr}(\overline{\mathrm{\Phi }}^1)^k>`$ correlators. In Section 4 the order $`g^4`$ contributions are considered. For the relevant diagrams we compute the colour structure factors and the momentum integrals that one obtains after completion of the $`D`$-algebra. The various, complicated contributions give rise in the end to a complete cancellation for any finite $`N`$ . This is in agreement with previous results . In Section 5 we show that the nonrenormalization properties we have found for the $`SU(N)`$ gauge group, are actually valid for general groups. Finally in Section 6 we concentrate on the evaluation of the contact terms and study what they correspond to in the supergravity sector.
## 2 Two-point functions of chiral operators
Our goal is to compute two-point correlators for $`𝒩=4`$ supersymmetric $`SU(N)`$ Yang-Mills theory, perturbatively in $`𝒩=1`$ superspace. The operators under consideration are the chiral primary operators in the $`(0,k,0)`$ representation of the $`SU(4`$) $`R`$–symmetry group. In a $`𝒩=1`$ superspace description of the theory (see Appendix A for details) they are given by $`𝒪=\mathrm{Tr}(\mathrm{\Phi }^{\{i_1}\mathrm{\Phi }^{i_2}\mathrm{}\mathrm{\Phi }^{i_k\}})`$, with flavor indices on the superfields symmetrized and traceless. In fact, in order to simplify matters as in Refs. , , we consider the $`SU(3)`$ highest weight superfield $`\mathrm{\Phi }^1`$ and compute $`<\mathrm{Tr}(\mathrm{\Phi }^1)^k\mathrm{Tr}(\overline{\mathrm{\Phi }}^1)^k>`$. In this way the flavor combinatorics is avoided. At the same time we do not lose in generality since the $`SU(3)`$ transformations, which are invariances of the theory, allow to reconstruct all the other primary chiral correlators from the one above.
The general strategy we have adopted for performing the actual calculation is outlined in Appendix A. Quite generally, at non–coincident points we can write the two-point function as
$$<\mathrm{Tr}(\mathrm{\Phi }^1)^k(z_1)\mathrm{Tr}(\overline{\mathrm{\Phi }}^1)^k(z_2)>=\frac{F(g^2,N)}{(x_1x_2)^{2k}}\delta ^{(4)}(\theta _1\theta _2)$$
(2.1)
where $`z(x,\theta ,\overline{\theta })`$. Away from short distance singularities, the $`x`$-dependence of the result is fixed by the conformal invariance of the theory, and $`F(g^2,N)`$ is the function that we want to determine perturbatively in $`g^2`$. As we have done in Ref. we compute loop integrals in momentum space and use dimensional regularization and minimal subtraction scheme to treat ultraviolet divergences. In $`n`$ dimensions, with $`n=42ϵ`$, the Fourier transform of a power factor $`(x_1x_2)^{2\nu }`$ is given by
$`{\displaystyle \frac{1}{(x^2)^\nu }}=2^{n2\nu }\pi ^{\frac{n}{2}}{\displaystyle \frac{\mathrm{\Gamma }(\frac{n}{2}\nu )}{\mathrm{\Gamma }(\nu )}}{\displaystyle \frac{d^np}{(2\pi )^n}\frac{e^{ipx}}{(p^2)^{\frac{n}{2}\nu }}}`$ $`\nu 0,1,2,\mathrm{}`$ (2.2)
$`\nu {\displaystyle \frac{n}{2}},{\displaystyle \frac{n}{2}}+1,\mathrm{}`$
The main advantage of such an approach is due to the fact that the $`x`$–space structure in (2.1), with a non–vanishing contribution to $`F(g^2,N)`$, can be obtained simply by looking at the contributions that behave like $`1/ϵ`$ from the singular factor $`\mathrm{\Gamma }(\frac{n}{2}\nu )=\mathrm{\Gamma }(\nu +2ϵ)`$, $`\nu 2`$ in (2.2). By analytic continuation one can write the general identity
$$\frac{d^np}{(2\pi )^n}\frac{e^{ipx}}{(p^2)^{2k+\alpha ϵ}}=\frac{2^{2k4}}{\pi ^2}(1)^k(k1)!(k2)!\alpha \frac{ϵ}{(x^2)^{k(\alpha +1)ϵ}}[1+𝒪(ϵ)]$$
(2.3)
Once the UV divergent terms are determined at a given order in $`g`$, one can reconstruct the complete answer using (2.3).
The UV divergent terms have been computed using the method proposed in and various techniques presented in . Infrared divergences have not been considered since the theory we are dealing with is conformally invariant.
Finally we emphasize that finite momentum space contributions to the correlators correspond in $`x`$–space to terms proportional to $`ϵ`$. These are the terms which give rise to contact terms . We will come back to this point in Section 5.
## 3 At tree-level and order $`g^2`$
At tree–level, the 2–point correlation function $`<\mathrm{Tr}(\mathrm{\Phi }^1)^k(z_1)\mathrm{Tr}(\overline{\mathrm{\Phi }}^1)^k(z_2)>`$ is given by the diagram in Fig. 1. The colour structure
$$\mathrm{Tr}(T_{a_1}\mathrm{}T_{a_k})\underset{\sigma }{}\mathrm{Tr}(T_{a_{\sigma (1)}}\mathrm{}T_{a_{\sigma (k)}})$$
(3.1)
which includes all possible permutations $`\sigma `$ in the contractions of the scalar lines, is a $`k`$–degree polynomial in $`N`$. In a double line representation for the colour indices, the $`N`$-leading power term corresponds to the planar double line graph associated to the diagram in Fig. 1, whereas nonplanar graphs give rise to subleading contributions.
Figure 1: tree–level contribution to $`<\mathrm{Tr}(\mathrm{\Phi }_1)^k\mathrm{Tr}(\overline{\mathrm{\Phi }}_1)^k>`$
Now using the result (B.5) for the momentum integrals to leading order in the $`ϵ`$ expansion, we obtain
$$\frac{1}{ϵ}\left[\frac{1}{(4\pi )^2}\right]^{k1}\frac{(1)^k}{[(k1)!]^2}(p^2)^{k2(k1)ϵ}\mathrm{Tr}(T_{a_1}\mathrm{}T_{a_k})\underset{\sigma }{}\mathrm{Tr}(T_{a_{\sigma (1)}}\mathrm{}T_{a_{\sigma (k)}})\delta ^{(4)}(\theta _1\theta _2)$$
(3.2)
In $`x`$–space, using eq. (2.3), the result can be rewritten as
$`<\mathrm{Tr}(\mathrm{\Phi }^1)^k(z_1)\mathrm{Tr}(\overline{\mathrm{\Phi }}^1)^k(z_2)>_0`$
$`=\left({\displaystyle \frac{1}{4\pi ^2}}\right)^k\mathrm{Tr}(T_{a_1}\mathrm{}T_{a_k}){\displaystyle \underset{\sigma }{}}\mathrm{Tr}(T_{a_{\sigma (1)}}\mathrm{}T_{a_{\sigma (k)}}){\displaystyle \frac{1}{(x_1x_2)^{2k}}}\delta ^{(4)}(\theta _1\theta _2)`$ (3.3)
At order $`g^2`$ the only contribution to the correlator is given by the diagram in Fig. 2, with the insertion of a vector line.
Figure 2: $`g^2`$–order contribution to $`<\mathrm{Tr}(\mathrm{\Phi }_1)^k\mathrm{Tr}(\overline{\mathrm{\Phi }}_1)^k>`$
From the two internal vertices $`V_1`$ in (A.6) we obtain the colour structure $`f_{amb}f_{a^{}mb^{}}`$ which contracted with the colour matrices associated to the rest of the diagram gives
$$\mathrm{Tr}(T_{a_1}\mathrm{}T_{a_k})\underset{\sigma }{}\underset{ij}{}\mathrm{Tr}(T_{a_{\sigma (1)}}\mathrm{}[T_{a_{\sigma (i)}},T_m]\mathrm{}[T_{a_{\sigma (j)}},T_m]\mathrm{}T_{a_{\sigma (k)}})$$
(3.4)
Here the sum is over all possible permutations of the external lines and eq. (C.1) has been used. The previous expression can be simplified by noticing that for any set of matrices $`M_j`$, $`j=1,\mathrm{},n`$, and any matrix $`P`$ the following identity holds
$$\mathrm{\Sigma }_{i=1}^n\mathrm{Tr}(M_1\mathrm{}[M_i,P]\mathrm{}M_n)=0$$
(3.5)
Thus (3.4) can be written as
$$\mathrm{Tr}(T_{a_1}\mathrm{}T_{a_k})\underset{\sigma }{}\underset{i=1}{\overset{k}{}}\mathrm{Tr}(T_{a_{\sigma (1)}}\mathrm{}[[T_{a_{\sigma (i)}},T_m],T_m]\mathrm{}T_{a_{\sigma (k)}})$$
(3.6)
which, using the identity (C.6), can be reduced further to the expression in (3.1) up to a factor $`2N`$. Including the various factors from vertices and propagators, we finally have
$$(g^2N)k\mathrm{Tr}(T_{a_1}\mathrm{}T_{a_k})\underset{\sigma }{}\mathrm{Tr}(T_{a_{\sigma (1)}}\mathrm{}T_{a_{\sigma (k)}})$$
(3.7)
The momentum integral associated to the graph after completion of the $`D`$-algebra, is evaluated in (B.6). The final result (here we reinstate a factor $`\frac{1}{(4\pi )^2}`$ for each loop) is
$$12\zeta (3)(g^2N)\left[\frac{1}{(4\pi )^2}\right]^k\frac{(1)^{k1}(k1)}{[(k1)!]^2}(p^2)^{k2kϵ}\mathrm{Tr}(T_{a_1}\mathrm{}T_{a_k})\underset{\sigma }{}\mathrm{Tr}(T_{a_{\sigma (1)}}\mathrm{}T_{a_{\sigma (k)}})$$
(3.8)
In the limit $`ϵ0`$ this expression is finite, therefore it does not contribute to the correlation function at separate points. It gives rise to a contact term which will be discussed in detail in Section 5.
## 4 At order $`g^4`$
The relevant diagrams, drawn in single line representation, are shown in Fig. 3 and 4. They group themselves into planar ($`A`$) and nonplanar ($`B`$) ones. If one were to use a double line representation for the colour indices from $`T_{ij}^a`$, then the single-line planar graphs ($`A`$) would give rise to two distinct types of graphs, double-line planar graphs ($`A_1`$) and nonplanar double-line graphs ($`A_2`$). Now from the $`A_1`$ graphs the colour structure would produce leading $`N`$ contributions, while from the $`A_2`$ graphs subleading contributions would arise. Obviously the single-line type ($`B`$) diagrams could only give structures subleading in $`N`$ since their double-line representation would be necessarily nonplanar.
Figure 3: planar $`g^4`$–order contribution to $`<\mathrm{Tr}(\mathrm{\Phi }_1)^k\mathrm{Tr}(\overline{\mathrm{\Phi }}_1)^k>`$
Figure 4: nonplanar $`g^4`$–order contribution to $`<\mathrm{Tr}(\mathrm{\Phi }_1)^k\mathrm{Tr}(\overline{\mathrm{\Phi }}_1)^k>`$
For the correlators computed in , i.e. for the $`k=3`$ case, the colour factor of the diagrams in the class ($`B`$) turned out to be identically zero. All the graphs in the class ($`A`$) gave rise to the same overall combinatorics, so that the final cancellation occurred as a cancellation among the various terms produced by the momentum integrations.
In the present case, $`k>3`$, we will show that the situation is much more complicated and highly non trivial. All types of diagrams mentioned above have non vanishing colour structures. The diagrams ($`A`$) give contributions proportional to the colour structure (3.1) which, after multiplication by the momentum integrals, sum up to zero. However in addition, from the diagrams in Fig. $`3c`$ and $`3d`$ a new subleading structure arises which is of the same kind as the one from diagrams ($`B`$). Again, a nontrivial cancellation occurs among the diagrams in Fig. $`3c`$, $`3d`$ and ($`B`$) due to the special structure of their momentum integrals.
Now for each of diagram we present the evaluation of the colour factor, while we list the corresponding momentum integrals in Appendix B. The final results are summarized at the end of the Section where the actual cancellation is discussed.
In Fig. $`3a`$ we have the insertion of a two–loop propagator correction
$`2g^4N^2\overline{\mathrm{\Phi }}_a^i(p,\theta )\mathrm{\Phi }_a^i(p,\theta )p^2{\displaystyle \frac{d^nqd^nk}{k^2q^2(kq)^2(kp)^2(pq)^2}}`$
$`=2g^4N^2\overline{\mathrm{\Phi }}_a^i(p,\theta )\mathrm{\Phi }_a^i(p,\theta ){\displaystyle \frac{1}{(p^2)^{2ϵ}}}[6\zeta (3)+𝒪(ϵ)]`$ (4.1)
In this case the colour structure is easy to compute. Inserting all the coefficients from vertices, propagators and combinatorics we obtain
$$12\zeta (3)(g^2N)^2k\mathrm{Tr}(T_{a_1}\mathrm{}T_{a_k})\underset{\sigma }{}\mathrm{Tr}(T_{a_{\sigma (1)}}\mathrm{}T_{a_{\sigma (k)}})$$
(4.2)
The momentum integral emerging after the $`D`$–algebra is given in eq. (B.7).
We now consider the diagram $`3b`$ where the $`𝒪(g^2)`$ effective vertex appears
$`{\displaystyle \frac{g^3}{4}}Nif^{abc}\overline{\mathrm{\Phi }}_a^i(q,\theta )\mathrm{\Phi }_b^i(p,\theta )(4D^\alpha \overline{D}^2D_\alpha `$
$`+(p+q)^{\alpha \dot{\alpha }}[D_\alpha ,\overline{D}_{\dot{\alpha }}])V_c(pq,\theta ){\displaystyle }{\displaystyle \frac{d^nk}{k^2(kp)^2(kq)^2}}`$ (4.3)
Performing the $`D`$–algebra one easily realizes that only the first term in (4.3) gives rise to potentially divergent loop integrals. The second term produces finite contributions which might generate order $`g^4`$ contact terms. We concentrate on the $`D^\alpha \overline{D}^2D_\alpha `$ term, since we will not compute contact contributions to $`g^4`$ order. The colour structure for this diagram is computed following the same procedure as in the case of the diagram in Fig. 2 (see eqs. (3.43.7)). Inserting all the various factors we obtain
$$4(g^2N)^2k\mathrm{Tr}(T_{a_1}\mathrm{}T_{a_k})\underset{\sigma }{}\mathrm{Tr}(T_{a_{\sigma (1)}}\mathrm{}T_{a_{\sigma (k)}})$$
(4.4)
The corresponding momentum integral is given in eq. (B.8).
We now turn to the discussion of the colour structure for the diagram in Fig. $`3c`$. The insertion of the two vector lines gives rise to the structure $`f_{a_1mb_1}f_{a_2mn}f_{npb_2}f_{a_3pb_3}`$. The contraction with scalar lines from the external vertices takes into account all possible permutations. Using eq. (C.1) we can write
$$\mathrm{Tr}(T_{a_1}\mathrm{}T_{a_k})\underset{\sigma }{}\underset{ijl}{}\mathrm{Tr}(T_{a_{\sigma (1)}}\mathrm{}[T_{a_{\sigma (i)}},T_m]\mathrm{}[[T_{a_{\sigma (j)}},T_m],T_n]\mathrm{}[T_{a_{\sigma (l)}},T_n]\mathrm{}T_{a_{\sigma (k)}})$$
(4.5)
This expression can be manipulated by making use of the identity (3.5) and it can be written as the sum of two terms
$`\mathrm{Tr}(T_{a_1}\mathrm{}T_{a_k}){\displaystyle \underset{\sigma }{}}{\displaystyle \underset{jl}{}}\{\mathrm{Tr}(T_{a_{\sigma (1)}}\mathrm{}[[[T_{a_{\sigma (j)}},T_m],T_n],T_m]\mathrm{}[T_{a_{\sigma (l)}},T_n]\mathrm{}T_{a_{\sigma (k)}})`$
$`+\mathrm{Tr}(T_{a_{\sigma (1)}}\mathrm{}[[T_{a_{\sigma (j)}},T_m],T_n]\mathrm{}[[T_{a_{\sigma (l)}},T_n],T_m]\mathrm{}T_{a_{\sigma (k)}})\}`$ (4.6)
The first term can be further reduced by using the identities (C.6), (C.7) and (3.5) again. Inserting the factors from combinatorics, vertices and propagators the final expression for the colour structure can be written as the sum of a term leading in $`N`$ plus a subleading contribution
$`g^4\mathrm{Tr}(T_{a_1}\mathrm{}T_{a_k}){\displaystyle \underset{\sigma }{}}\{2kN^2\mathrm{Tr}(T_{a_{\sigma (1)}}\mathrm{}T_{a_{\sigma (k)}})`$
$`{\displaystyle \underset{jl}{}}\mathrm{Tr}(T_{a_{\sigma (1)}}\mathrm{}[[T_{a_{\sigma (j)}},T_m],T_n]\mathrm{}[[T_{a_{\sigma (l)}},T_n],T_m]\mathrm{}T_{a_{\sigma (k)}})\}`$ (4.7)
For $`k=3`$ the second term vanishes and the above expression reduces to the one computed in Ref. . However for $`k>3`$ the subleading term is nonzero, as shown in Appendix C explicitly for the $`k=4`$ case. For this diagram the loop–integral result can be read in eq. (B.9).
The colour factor for the diagram in Fig. $`3d`$ can be computed by exploiting the previous results. In fact, from the internal vertices $`V_1`$ and $`V_3`$ (see eq. (A.6)) we have the structure $`f_{a_1mb_1}f_{pb_2n}f_{ma_2n}f_{a_3pb_3}`$ which is identical to the one from the diagram $`3c`$. Performing all the contractions and taking into account the coefficients from combinatorics, vertices and propagators we finally obtain
$`{\displaystyle \frac{g^4}{2}}\mathrm{Tr}(T_{a_1}\mathrm{}T_{a_k}){\displaystyle \underset{\sigma }{}}\{2kN^2\mathrm{Tr}(T_{a_{\sigma (1)}}\mathrm{}T_{a_{\sigma (k)}})`$
$`{\displaystyle \underset{jl}{}}\mathrm{Tr}(T_{a_{\sigma (1)}}\mathrm{}[[T_{a_{\sigma (j)}},T_m],T_n]\mathrm{}[[T_{a_{\sigma (l)}},T_n],T_m]\mathrm{}T_{a_{\sigma (k)}})\}`$ (4.8)
The momentum integral for this graph is given in eq. (B.10).
The last potential contributions from diagrams of type $`(A)`$ are shown in Figures $`3e3h`$. The loop integral resulting from the D–algebra on the diagram $`3e`$ given in eq. (B.11), is $`𝒪(ϵ)`$ and then it does not contribute to the correlation function. The diagrams in Figures $`3f`$, $`3g`$ and $`3h`$, exactly as for $`k=3`$, only contribute to finite terms.
We now study the nonplanar diagrams ($`B`$) in Figure 4.
The combinatorial and colour factors, and the loop integrals for the two diagrams $`4a`$ and $`4b`$ are identical. Thus we concentrate on one of them, e.g. $`4b`$. From the internal $`V_1`$ vertices with vector lines contracted as in figure, the colour factor which arises is $`f_{a_1mn}f_{npb_1}f_{a_2pq}f_{qmb_2}`$. When connected with the rest of the diagram it gives
$$\frac{g^4}{2}\underset{\sigma }{}\underset{jl}{}\mathrm{Tr}(T_{a_{\sigma (1)}}\mathrm{}[[T_{a_{\sigma (j)}},T_m],T_n]\mathrm{}[[T_{a_{\sigma (l)}},T_n],T_m]\mathrm{}T_{a_{\sigma (k)}})$$
(4.9)
where combinatorial factors have been included. We note that the previous expression has the same form as the subleading contribution already present in the diagrams $`3c`$ and $`3d`$ (see eqs. (4.7, 4.8)). As previously mentioned and proven in Appendix C, this term is zero for $`k=3`$ but in general nonvanishing when $`k>3`$. The corresponding loop diagram arising after the D–algebra is given in eq. (B.8).
The last nonplanar graph to be considered is the one in Fig. $`4c`$. As in the $`k=3`$ case , its colour coefficient vanishes. A simple way to prove this is to notice that from the internal vertices one obtains
$$f_{a_1mb_1}f_{a_2nb_2}f_{a_3pb_3}f_{mnp}$$
(4.10)
which is antisymmetric under the exchange $`a_1a_2`$ and $`b_1b_2`$. When multiplied by all possible permutations of colour matrices from the external vertices it gives a zero result.
We now collect all the divergent contributions we have computed at order $`g^4`$. For each diagram we list the final result obtained as a product of the colour and combinatorial factors times the results from momentum integrations. We define
$`𝒫_k`$ $`\mathrm{Tr}(T_{a_1}\mathrm{}T_{a_k}){\displaystyle \underset{\sigma }{}}\mathrm{Tr}(T_{a_{\sigma (1)}}\mathrm{}T_{a_{\sigma (k)}})`$ (4.11)
$`𝒬_k`$ $`\mathrm{Tr}(T_{a_1}\mathrm{}T_{a_k}){\displaystyle \underset{\sigma }{}}{\displaystyle \underset{jl}{}}\mathrm{Tr}(T_{a_{\sigma (1)}}\mathrm{}[[T_{a_{\sigma (j)}},T_m],T_n]\mathrm{}[[T_{a_{\sigma (l)}},T_n],T_m]\mathrm{}T_{a_{\sigma (k)}})`$
Factorizing an overall common coefficient
$$\frac{1}{ϵ}(g^2N)^2\left[\frac{1}{(4\pi )^2}\right]^{k+1}\frac{(1)^k(k1)}{[(k1)!]^2(k+1)}(p^2)^{k2(k+1)ϵ}$$
(4.12)
the various contributions are
* Diagram 3a:
$$12k\zeta (3)𝒫_k$$
(4.13)
* Diagram $`3b`$:
$$24k\zeta (3)𝒫_k$$
(4.14)
* Diagram $`3c`$:
$$[12k\zeta (3)40k\zeta (5)]𝒫_k+\frac{1}{N^2}[20\zeta (5)6\zeta (3)]𝒬_k$$
(4.15)
* Diagram $`3d`$:
$$40k\zeta (5)𝒫_k\frac{1}{N^2}20\zeta (5)𝒬_k$$
(4.16)
* Diagrams $`4a+4b`$:
$$\frac{1}{N^2}6\zeta (3)𝒬_k$$
(4.17)
The leading structure $`𝒫_k`$ appears in the diagrams $`3a`$ and $`3b`$ with a coefficient proportional to the $`\zeta (3)`$ Riemann function, whereas in the diagram $`3d`$ it is proportional to $`\zeta (5)`$. The same pattern repeats itself for the non-leading structure $`𝒬_k`$: in the diagrams $`4a`$ and $`4b`$ it is multiplied by $`\zeta (3)`$, whereas in the diagram $`3d`$ it is proportional to $`\zeta (5)`$. Both structures appear in the diagram $`3c`$ and they contribute with such a coefficient to cancel the rest of the divergent terms. In conclusion, only finite contributions survive. As emphasized at the beginning this amounts to say that the correlation function is not renormalized at order $`g^4`$, up to contact terms.
We stress that the complete cancellation of divergent contributions for general $`k`$ occurs for any finite $`N`$.
The results in (4.13-4.17) if restricted to $`k=3`$, reproduce the results discussed in Ref. . In this case the subleading structure $`𝒬_k`$ vanishes, so that the proof of nonrenormalization for the correlator with general $`k`$ cannot be implemented from the $`k=3`$ example.
## 5 General gauge groups
Here we want to show that the nonrenormalization properties of the correlators proven in the two previous Sections for the $`SU(N)`$ super Yang-Mills theory, actually hold for general simple gauge groups. To this end it is sufficient to realize that, using (C.3) and (C.5) which are valid for any group, we can generalize the identities in (C.6) and (C.7) to the following ones
$$[[T_a,T_m],T_m]=k_1T_a$$
(5.1)
and
$$[[[T_a,T_m],T_n],T_m]=\frac{k_1}{2}[T_a,T_n].$$
(5.2)
At the order $`g^2`$ this enables us to replace $`2N`$ with $`k_1`$ in (3.7) and (3.8), thus obtaining a result which depends only on the Casimir of the adjoint representation of the gauge group. In any event to this order the momentum integral gives a finite result which does not affect the correlator at non coincident points.
In the same way one can analyze the general situation at the next perturbative $`g^4`$ order. For the propagator and vertex insertions which appear in the graphs of Fig. $`3a`$, $`3b`$ we use the result as in Ref. , i.e. we set $`2Nk_1`$ in (4.1) and in (4.3). For the rest of the diagrams we substitute (C.6), (C.7) with (5.1), (5.2) respectively. Once this operation is performed consistently everywhere in the various formulas of Section 4 the final, complete cancellation of all the corrections is achieved following exactly the same pattern and the same steps as in the $`SU(N)`$ case.
## 6 Contact terms in the AdS/CFT correspondence
Now we wish to focus our attention on the two-point correlators for $`SU(N)`$ when the two points approach each other. In the limit $`x_1x_2`$ the expression in (2.1) becomes singular and needs to be regulated. Within the dimensional regularization approach we have adopted, this short-distance singularity is signaled by $`1/ϵ`$ poles, according to the general identity
$$\frac{1}{(x^2)^{kkϵ}}\frac{\pi ^2}{2^{2k4}(k1)!(k2)!}\frac{1}{ϵ}(^2)^{k2}\delta ^{(n)}(x)\mathrm{for}ϵ0$$
(6.1)
For the two–point function of the operator $`\mathrm{Tr}(\mathrm{\Phi }^1)^k`$ these short-distance UV divergences manifest themselves already at tree–level (see (3.3)). In order to obtain a well defined function at coincident points we perform a subtraction and define in configuration space
$`<\mathrm{Tr}(\mathrm{\Phi }^1)^k(z_1)\mathrm{Tr}(\overline{\mathrm{\Phi }}^1)^k(z_2)>_{\mathrm{reg}}\delta ^{(4)}(\theta _1\theta _2)\times `$
$`\underset{ϵ0}{lim}{\displaystyle \frac{𝒫_k}{(4\pi )^{2k}}}\left[{\displaystyle \frac{1}{(x^2)^{kkϵ}}}+(1)^k(\mu ^2)^{(1k)ϵ}{\displaystyle \frac{2^{42k}\pi ^2}{[(k1)!]^2}}\left({\displaystyle \frac{1}{ϵ}}+\gamma \right)(^2)^{k2}\delta ^{(n)}(x_1x_2)\right]`$
(6.2)
where $`𝒫_k`$ has been defined in (4.12) and $`\mu `$ is the mass scale of dimensional regularization. The coefficient $`\gamma `$ corresponds to an arbitrary finite subtraction and generates a scheme dependent finite contact term.
Performing explicitly the $`ϵ0`$ limit in (6.2) a residual dependence on the mass scale survives from the (divergent) counter–contact term. Therefore, the subtraction of the infinity at tree–level necessarily introduces a mass scale in the regularized correlation function which breaks conformal invariance . On the other hand, the scheme dependent finite term proportional to $`\gamma `$ is independent of $`\mu `$ and it does not affect conformal invariance.
Now we discuss the appearance of this type of terms in the perturbative computation of the two–point correlator. We have performed our loop calculations in momentum space, with the Fourier transformation given in the basic formula (2.3). Using (2.3) and the general identity in (6.1) we can write
$$\underset{ϵ0}{lim}\frac{d^np}{(2\pi )^n}\frac{e^{ipx}}{(p^2)^{\frac{n}{2}k+\alpha ϵ}}=(1)^{k+1}\alpha (^2)^{k2}\delta ^{(4)}(x)$$
(6.3)
i.e. any finite contribution in momentum space gives rise to finite contact terms in the correlation functions.
In general, the presence of contact terms at any loop order is related to the UV regularization and renormalization procedure. Different subtraction conditions correspond to different finite counterterms which eventually contribute to contact terms in the correlation functions of the theory. In $`n=42ϵ`$ dimensions, one has to choose a particular regularization for evaluating the integrals
$$d^4pf(p)G(ϵ)d^npf(p)$$
(6.4)
where $`G(ϵ)`$ is a regular function<sup>1</sup><sup>1</sup>1A convenient choice is $`G(ϵ)=(4\pi )^ϵ\mathrm{\Gamma }(1ϵ)`$ (G–scheme ) which cancels irrelevant terms proportional to the Euler constant, $`\mathrm{log}4\pi `$ and $`\zeta (2)`$ Riemann function in the $`ϵ`$–expansion of $`\mathrm{\Gamma }`$ functions which appear in the calculation of multiloop integrals. near $`ϵ=0`$, with $`G(0)=1`$. A given prescription has to be used in the computation of both the effective action and physical quantities like correlation functions or scattering amplitudes. By expanding $`G(ϵ)`$ in powers of $`ϵ`$ one can easily realize that in multiloop integrals with only simple pole divergences the coefficients of the divergent terms do not depend on the particular choice of the $`G`$–function, whereas they might depend on the regularization scheme in multiloop integrals with higher order divergences. In any case, a scheme dependence is always present in the finite part of any divergent diagram. It follows that in the evaluation of correlation functions, different choices of the $`G`$–function, i.e. different regularization prescriptions, give rise in general to different finite quantum contact terms. However, if a nonrenormalization theorem holds, then the contact terms become independent of the regularization prescription and they are unambiguously computable.
Indeed let us exemplify for simplicity the case in which at a given loop order the perturbative contribution to a correlation function contains at most a simple pole $`1/ϵ`$ divergence
$$\left(\frac{a}{ϵ}+b+O(ϵ)\right)G(ϵ)$$
(6.5)
By expanding $`G(ϵ)=1+cϵ+\mathrm{}`$, in the limit $`ϵ0`$ we obtain the divergent scheme independent contribution $`a/ϵ`$ and the finite term $`(b+ca)`$ which contains a scheme dependence through the coefficient $`c`$ from $`G(ϵ)`$. However, if there is no perturbative renormalization, i.e. $`a=0`$, the finite term is uniquely determined.
In our specific case, we have shown that up to the $`g^4`$–order the two–point function $`<\mathrm{Tr}(\mathrm{\Phi }^1)^k(z_1)\mathrm{Tr}(\overline{\mathrm{\Phi }}^1)^k(z_2)>`$ is not perturbatively renormalized and we are in the situation where the finite contact terms are uniquely determined. The term at order $`g^2`$ can be easily inferred from (3.8) by using the identity (6.3). In the limit $`ϵ0`$ we find
$$12\zeta (3)(g^2N)\frac{k^21}{[(k1)!]^2}\frac{𝒫_k}{(4\pi )^{2k}}(^2)^{k2}\delta ^{(4)}(x_1x_2)\delta ^{(4)}(\theta _1\theta _2)$$
(6.6)
The same procedure can be applied in order to compute the contact term at order $`g^4`$. One should keep track of all finite contributions from the momentum integrals and then use (6.3) to obtain the result in configuration space. We have not performed the explicit calculation but in general we expect to find a nonvanishing result.
We note that, since these loop contact terms come from finite contributions in the $`ϵ`$–expansion, in the four dimensional limit no dependence on the scale $`\mu `$ survives. Therefore these terms do not affect the conformal invariance properties of the theory. As discussed in Ref. the breaking of conformal invariance can only be ascribed to the UV divergences in the correlation function.
If the non–renormalization theorem holds for any finite $`N`$ at all orders in perturbation theory, one could determine unambiguously the finite contact terms loop by loop, and generate a contribution of the form
$$[a+f(g^2,N)](^2)^{k2}\delta ^{(4)}(x_1x_2)\delta ^{(4)}(\theta _1\theta _2)$$
(6.7)
where the arbitrary tree–level coefficient $`a`$ is given by
$$a\gamma (1)^k\frac{𝒫_k}{(4\pi )^{2k}}\frac{2^{42k}\pi ^2}{[(k1)!]^2}$$
(6.8)
Our result shows that the complete answer for the two–point functions of the theory necessarily contains contact terms, unless we were to choose a non-minimal subtraction at short distances. This would amount to the subtraction of a finite term with a coefficient $`a=f(g^2,N)`$.
The SL(2,Z) invariance of the theory would require $`f(g^2,N)`$ to be a modular function. However, following the arguments in , if loop contact terms were present in the final result they would necessarily break the $`U(1)_Y`$ invariance that the theory inherits from the 5d supergravity. It might be possible to establish a correspondence with $`U(1)`$–breaking terms in the supergravity action .
According to the AdS/CFT conjecture, the generating functional of the regularized correlation functions in the large $`N`$ limit is the semiclassical 5d supergravity action with IR divergences suitably subtracted . At the boundary, i.e. with the IR cut–off removed, it has the form (1.3). In particular, the arbitrary finite contact term in (6.7) corresponds to an arbitrary finite subtraction in 5d
$$S_c[\varphi _0]=b_fd^4x\varphi _0(x)(^2)^{k2}\varphi _0(x)$$
(6.9)
related to the particular IR regularization scheme chosen for the supergravity action. In other words, in the large $`N`$ limit with $`g^2N1`$ we should find
$$[a+f(g^2,N)]b_f$$
(6.10)
If in this limit $`f(g^2,N)`$ is not vanishing, we might conclude that either coupling dependent local terms appear in supergravity, or one is forced to choose a particular subtraction scheme in the Yang–Mills sector.
Finally we notice that perturbative contact terms in the two–point functions give rise in general to coupling dependent contact–type contributions in higher–point correlators, as well as enter the definition of multitrace operators through the point–splitting regularization .
## Acknowledgements
We wish to thank Misha Bershadsky, Dan Freedman and Michael Green for discussions and suggestions.
This work has been supported by the European Commission TMR programme ERBFMRX-CT96-0045, in which S.P. and D.Z. are associated to the University of Torino.
## Appendix A The basic rules for the computation of the two-point correlators
In $`𝒩=1`$ superspace the action of $`𝒩=4`$ supersymmetric Yang-Mills theory can be written in terms of one real vector superfield $`V`$ and three chiral superfields $`\mathrm{\Phi }^i`$ (we follow the notations in )
$`S[J,\overline{J}]`$ $`=`$ $`{\displaystyle d^8z\mathrm{Tr}\left(e^{gV}\overline{\mathrm{\Phi }}_ie^{gV}\mathrm{\Phi }^i\right)}+{\displaystyle \frac{1}{2g^2}}{\displaystyle d^6z\mathrm{Tr}W^\alpha W_\alpha }`$ (A.1)
$`+{\displaystyle \frac{ig}{3!}}\mathrm{Tr}{\displaystyle d^6zϵ_{ijk}\mathrm{\Phi }^i[\mathrm{\Phi }^j,\mathrm{\Phi }^k]}+{\displaystyle \frac{ig}{3!}}\mathrm{Tr}{\displaystyle d^6\overline{z}ϵ_{ijk}\overline{\mathrm{\Phi }}^i[\overline{\mathrm{\Phi }}^j,\overline{\mathrm{\Phi }}^k]}`$
$`+{\displaystyle d^6zJ𝒪}+{\displaystyle d^6\overline{z}\overline{J}\overline{𝒪}}`$
where $`W_\alpha =i\overline{D}^2(e^{gV}D_\alpha e^{gV})`$, and $`V=V^aT^a`$, $`\mathrm{\Phi }_i=\mathrm{\Phi }_i^aT^a`$, $`T^a`$ being $`SU(N)`$ matrices in the fundamental representation. We have added to the classical action source terms for the chiral primary operators generically denoted by $`𝒪`$.
We define the generating functional in Euclidean space
$$W[J,\overline{J}]=𝒟\mathrm{\Phi }𝒟\overline{\mathrm{\Phi }}𝒟Ve^{S[J,\overline{J}]}$$
(A.2)
so that for $`𝒪=\mathrm{Tr}(\mathrm{\Phi }^1)^k`$ the two-point function is given by
$$<\mathrm{Tr}(\mathrm{\Phi }^1)^k(z_1)\mathrm{Tr}(\overline{\mathrm{\Phi }}^1)^k(z_2)>=\frac{\delta ^2W}{\delta J(z_1)\delta \overline{J}(z_2)}|_{J=\overline{J}=0}$$
(A.3)
where $`z(x,\theta ,\overline{\theta })`$. We use perturbation theory to evaluate the contributions to $`W[J,\overline{J}]`$ which are quadratic in the sources, i.e.
$$W[J,\overline{J}]d^4x_1d^4x_2d^4\theta J(x_1,\theta ,\overline{\theta })\frac{F(g^2,N)}{(x_1x_2)^{2k}}\overline{J}(x_2,\theta ,\overline{\theta })$$
(A.4)
The $`x`$-dependence of the result is fixed by the conformal invariance of the theory, and $`F(g^2,N)`$ is the function to be determined.
In order to obtain the result in (A.4) one has to consider all the two-point diagrams from $`W[J,\overline{J}]`$ with $`J`$ and $`\overline{J}`$ on the external legs. First one evaluates all factors coming from combinatorics and colour structures of a given diagram. Then one performs the superspace $`D`$-algebra following standard techniques (see for example ), and reduces the result to a multi-loop integral.
The quantization procedure of the classical action in (A.1) requires the introduction of a gauge fixing (we work in Feynman gauge) and corresponding ghost terms. The ghost superfields only couple to the vector multiplet and are not interesting for our calculation. In momentum space we have the superfield propagators
$$<V^aV^b>=\frac{\delta ^{ab}}{p^2}<\mathrm{\Phi }_i^a\overline{\mathrm{\Phi }}_j^b>=\delta _{ij}\frac{\delta ^{ab}}{p^2}$$
(A.5)
The vertices are read directly from the interaction terms in (A.1), with additional $`\overline{D}^2`$, $`D^2`$ factors for chiral, antichiral lines respectively. The ones that we need are the following
$`V_1=igf_{abc}\delta ^{ij}\overline{\mathrm{\Phi }}_i^aV^b\mathrm{\Phi }_j^cV_2={\displaystyle \frac{i}{2}}gf_{abc}V^a\overline{D}^2D^\alpha V^bD_\alpha V^c`$
$`V_3={\displaystyle \frac{g^2}{2}}\delta ^{ij}f_{adm}f_{bcm}V^aV^b\overline{\mathrm{\Phi }}_i^c\mathrm{\Phi }_j^d`$ (A.6)
$`V_4={\displaystyle \frac{g}{3!}}ϵ^{ijk}f_{abc}\mathrm{\Phi }_i^a\mathrm{\Phi }_j^b\mathrm{\Phi }_k^c\overline{V}_4={\displaystyle \frac{g}{3!}}ϵ^{ijk}f_{abc}\overline{\mathrm{\Phi }}_i^a\overline{\mathrm{\Phi }}_j^b\overline{\mathrm{\Phi }}_k^c`$
All the calculations are performed in $`n`$ dimensions with $`n=42ϵ`$ and in momentum space. We have used the method of uniqueness which is particularly efficient for the computation of massless Feynman integrals of a single variable.
## Appendix B Relevant integrals in momentum space
In this Appendix we list the relevant multiloop integrals which have been used in the course of our calculation.
As described in Section 5, in dimensional regularization $`n=42ϵ`$, one has to choose a particular prescription for the regularized integrals. In our case the integrals have at the most $`1/ϵ`$ divergences so that they do not depend on the particular choice of the $`G`$–function, as far as divergent contributions are concerned. In order to simplify the calculation, we forget about $`(2\pi )`$ factors at intermediate stages and reinsert at the end a $`1/(4\pi )^2`$ factor for each loop. Having this in mind we now list the relevant integrals and their $`ϵ`$ expansion.
The basic integrals from which all our results can be deduced are the following ones:
At one loop
$$I_1=\frac{d^nk}{(k^2)^\alpha [(pk)^2]^\beta }=\frac{\mathrm{\Gamma }(\alpha +\beta \frac{n}{2})}{\mathrm{\Gamma }(\alpha )\mathrm{\Gamma }(\beta )}\frac{\mathrm{\Gamma }(\frac{n}{2}\beta )\mathrm{\Gamma }(\frac{n}{2}\alpha )}{\mathrm{\Gamma }(n\alpha \beta )}\frac{1}{(p^2)^{\alpha +\beta \frac{n}{2}}}$$
(B.1)
At two loops
$$I_2=\frac{d^nkd^nl}{k^2l^2(kl)^2(pk)^2(pl)^2}=\frac{1}{(p^2)^{1+2ϵ}}[6\zeta (3)+𝒪(ϵ)]$$
(B.2)
At four loops
$$I_4=\frac{d^nkd^nld^nrd^ns}{k^2l^2(kl)^2(rk)^2(sl)^2(rs)^2(pr)^2(ps)^2}=\frac{1}{ϵ}\frac{1}{(p^2)^{4ϵ}}[5\zeta (5)+𝒪(ϵ)]$$
(B.3)
and
$`\stackrel{~}{I}_4={\displaystyle \frac{d^nkd^nld^nrd^nsr^2(pl)^2}{k^2l^2(kl)^2(rk)^2(rl)^2(sl)^2(rs)^2(pr)^2(ps)^2}}`$
$`={\displaystyle \frac{1}{ϵ}}(p^2)^{14ϵ}\left[{\displaystyle \frac{5}{2}}\zeta (5){\displaystyle \frac{3}{4}}\zeta (3)+𝒪(ϵ)\right]`$ (B.4)
(the explicit evaluation of the last integral was reported in Ref. ).
From the previous integrals we can derive all the results needed for our calculation.
By repeated use of (B.1) we obtain for the $`g^0`$–order diagram
$$\frac{d^nq_1\mathrm{}d^nq_{k1}}{q_1^2(q_2q_1)^2(q_3q_2)^2\mathrm{}(pq_{k1})^2}=\frac{1}{ϵ}\frac{(1)^k}{[(k1)!]^2}(p^2)^{k2(k1)ϵ}+𝒪(1)$$
(B.5)
For the $`g^2`$–order (finite) diagram we have
$`{\displaystyle d^nq_2\mathrm{}d^nq_{k1}\frac{q_2^2}{(q_3q_2)^2\mathrm{}(pq_{k1})^2}\frac{d^nkd^nl}{k^2l^2(kl)^2(q_2k)^2(q_2l)^2}}`$
$`={\displaystyle \frac{(1)^{k1}(k1)}{[(k1)!]^2k}}12\zeta (3)(p^2)^{k2kϵ}+𝒪(ϵ)`$ (B.6)
In order to obtain this result, one first performs the $`k`$ and $`l`$ integrations with the help of eq. (B.2), then the other integrals are computed with the use of eq. (B.1).
At $`g^4`$–order, the integrals emerging after having performed the $`D`$-algebra are the following ones:
For the graph $`3a`$, by using eq. (B.1) one obtains
$$\frac{d^nq_1\mathrm{}d^nq_{k1}}{(q_1^2)^{1+2ϵ}(q_2q_1)^2(q_3q_2)^2\mathrm{}(pq_{k1})^2}=\frac{1}{ϵ}\frac{(1)^k(k1)}{[(k1)!]^2(k+1)}(p^2)^{k2(k+1)ϵ}+𝒪(1)$$
(B.7)
For the graphs $`3b`$, $`4a`$ and $`4b`$, it is easy to deduce
$`{\displaystyle d^nrd^nq_2\mathrm{}d^nq_{k1}\frac{1}{(q_2r)^2(q_3q_2)^2\mathrm{}(pq_{k1})^2}\frac{d^nkd^nl}{k^2l^2(kl)^2(rk)^2(rl)^2}}`$
$`={\displaystyle \frac{1}{ϵ}}{\displaystyle \frac{(1)^k(k1)}{[(k1)!]^2(k+1)}}6\zeta (3)(p^2)^{k2(k+1)ϵ}+𝒪(1)`$ (B.8)
by first evaluating the two–loop integrals in $`k`$ and $`l`$ with the help of eq. (B.2), and then applying eq. (B.1).
The integral relevant for the graph $`3c`$ is
$`{\displaystyle d^nq_3\mathrm{}d^nq_{k1}\frac{1}{(q_4q_3)^2\mathrm{}(pq_{k1})^2}}`$
$`{\displaystyle \frac{d^nkd^nld^nrd^nsr^2(q_3l)^2}{k^2l^2(kl)^2(rk)^2(rl)^2(sl)^2(rs)^2(q_3r)^2(q_3s)^2}}`$
$`={\displaystyle \frac{1}{ϵ}}{\displaystyle \frac{(1)^k(k1)}{[(k1)!]^2(k+1)}}[6\zeta (3)20\zeta (5)](p^2)^{k2(k+1)ϵ}+𝒪(1)`$ (B.9)
by first evaluating the four–loop integral with momenta $`k,l,r,s`$ with use of eq. (B.4), and then performing the rest of the integrations using eq. (B.1).
Finally, the momentum integral for the graph $`3d`$ is given by
$`{\displaystyle d^nq_3\mathrm{}d^nq_{k1}\frac{q_3^2}{(q_4q_3)^2\mathrm{}(pq_{k1})^2}}`$
$`{\displaystyle \frac{d^nkd^nld^nrd^ns}{k^2l^2(kl)^2(rk)^2(sl)^2(rs)^2(q_3r)^2(q_3s)^2}}`$
$`={\displaystyle \frac{1}{ϵ}}{\displaystyle \frac{(1)^k(k1)}{[(k1)!]^2(k+1)}}40\zeta (5)(p^2)^{k2(k+1)ϵ}+𝒪(1)`$ (B.10)
Here, one first performs the $`k,l,r,s`$ integrations with eq. (B.3), and then uses eq. (B.1).
For the graph $`3e`$ the momentum integral produced after completion of the D–algebra is of order $`ϵ`$. Indeed, it is given by
$`{\displaystyle d^nq_4\mathrm{}d^nq_{k1}\frac{1}{(q_5q_4)^2\mathrm{}(pq_{k1})^2}d^nrr^2(q_4r)^2}`$
$`{\displaystyle \frac{d^nkd^nl}{k^2l^2(kl)^2(rk)^2(rl)^2}\frac{d^nsd^nt}{(sr)^2(tr)^2(st)^2(q_4s)^2(q_4t)^2}}`$
$`=ϵ{\displaystyle \frac{(1)^k(k1)}{[(k1)!]^2(k+1)}}144\zeta (3)^2(p^2)^{k2(k+1)ϵ}+𝒪(ϵ^2)`$ (B.11)
This result can be obtained by using first eq. (B.2) for the $`k,l`$ and $`s,t`$ two-loop integrals, and then eq. (B.1) for the remaining integrations.
## Appendix C Colour structures
In this Appendix we give our conventions and a series of useful identities involving the group generators. Moreover we show that the colour structure (4.9) of the nonplanar graphs $`4a`$ and $`4b`$ is nonvanishing, by evaluating it explicitly in the $`k=4`$ case.
For a general simple Lie algebra we have
$$[T_a,T_b]=if_{abc}T_c$$
(C.1)
where $`T_a`$ are the generators and $`f_{abc}`$ the structure constants. The matrices $`T_a`$’s are normalized as
$$\mathrm{Tr}(T_aT_b)=k_2\delta _{ab}$$
(C.2)
We have also
$$f_{amn}f_{bmn}=k_1\delta _{ab}$$
(C.3)
From the Jacobi identity one obtains
$$f_{abm}f_{cdm}+f_{cbm}f_{dam}+f_{dbm}f_{acm}=0.$$
(C.4)
which in turn allows to write
$$f_{alm}f_{bmn}f_{cnl}=\frac{1}{2}k_1f_{abc}$$
(C.5)
Now we specialize our formulas to the gauge group $`SU(N)`$: we have $`k_1=2k_2N`$ and we choose a normalization in (C.2) such that $`k_2=1`$. The generators $`T_a`$, $`a=1,\mathrm{},N^21`$, in the fundamental representation of $`SU(N)`$ are $`N\times N`$ traceless matrices. For $`SU(N)`$ the relations in (C.3) and (C.5) can be written as
$$[[T_a,T_m],T_m]=2NT_a$$
(C.6)
$$[[[T_a,T_m],T_n],T_m]=N[T_a,T_n].$$
(C.7)
Now we concentrate on the explicit evaluation of some colour factors in the case $`k=4`$. In particular we want to show that the nonplanar colour structure of graphs $`4a`$ and $`4b`$ is nonvanishing, in contradistinction to the case $`k=3`$ (see formula (A.18) in Appendix A of ). The relation which allows to deal with products of $`T_a`$’s is the following
$$T_{ij}^aT_{kl}^a=\left(\delta _{il}\delta _{jk}\frac{1}{N}\delta _{ij}\delta _{kl}\right).$$
(C.8)
From (C.8) one can obtain several useful formulas
$$\mathrm{Tr}(T_aT_bT_cT_d)\mathrm{Tr}(T_aT_bT_cT_d)=\frac{1}{N^2}(N^21)(N^2+3)$$
(C.9)
$$\mathrm{Tr}(T_aT_bT_cT_d)\mathrm{Tr}(T_aT_bT_dT_c)=\frac{1}{N^2}(N^21)(N^23)$$
(C.10)
$$\mathrm{Tr}(T_aT_bT_cT_d)\mathrm{Tr}(T_dT_cT_bT_a)=k_2^4\frac{1}{N^2}(N^21)(N^43N^2+3)$$
(C.11)
the last one being a planar type contribution.
Moreover, by noticing that
$$f_{abc}=i\mathrm{Tr}\left([T_a,T_b]T_c\right)$$
(C.12)
one can compute
$$\mathrm{Tr}(T_{c_1}T_{a_1}T_{c_2}T_{a_2})f_{c_1mb_1}f_{c_2mb_2}=(\delta _{b_1a_1}\delta _{b_2a_2}+\delta _{b_2a_1}\delta _{b_1a_2})$$
(C.13)
$$\mathrm{Tr}(T_{c_1}T_{c_2}T_{a_1}T_{a_2})f_{c_1mb_1}f_{c_2mb_2}=\delta _{b_1b_2}\delta _{a_1a_2}+\mathrm{Tr}(T_{b_1}T_{b_2}T_{a_1}T_{a_2})$$
(C.14)
Consider now the nonplanar colour structure for the graphs $`4a`$ and $`4b`$
$$𝒬_k\mathrm{Tr}(T_{a_1}\mathrm{}T_{a_k})\underset{\sigma }{}\underset{ij}{}\mathrm{Tr}\left(T_{a_{\sigma (1)}}\mathrm{}[[T_{a_{\sigma (i)}},T_m],T_n]\mathrm{}[[T_{a_{\sigma (j)}},T_n],T_m]\mathrm{}T_{a_{\sigma (k)}}\right)$$
(C.15)
As already mentioned $`𝒬_3=0`$. In general however it does not vanish. In the case $`k=4`$ by exploiting the various symmetries of this structure, it can be brought to the more manageable expression
$`𝒬_4`$ $`=`$ $`32[2\mathrm{T}\mathrm{r}(T_{a_1}T_{a_2}T_{b_3}T_{b_4})+\mathrm{Tr}(T_{a_1}T_{b_3}T_{a_2}T_{b_4})]f_{a_1nm}f_{mpc_1}f_{a_2pr}f_{rnc_2}`$ (C.16)
$`\times [\mathrm{Tr}(T_{c_1}T_{c_2}T_{b_3}T_{b_4})+\mathrm{Tr}(T_{c_1}T_{b_3}T_{c_2}T_{b_4})+\mathrm{Tr}(T_{c_1}T_{c_2}T_{b_4}T_{b_3})]`$
By using the Jacobi identity (C.4) for the four $`f`$ structure
$$f_{a_1nm}f_{mpc_1}f_{a_2pr}f_{rnc_2}=Nf_{a_1mc_1}f_{a_2mc_2}+f_{a_1pm}f_{a_2pr}f_{c_1nm}f_{c_2nr}$$
(C.17)
and using equations (C.13, C.14) and (C.9, C.10),it is straightforward to obtain
$$𝒬_4=192(N^21)(2N^23)$$
(C.18)
As claimed the result is nonvanishing.
Now, as a check, we want to verify in the $`k=4`$ case that the colour structure of the diagrams $`3c`$ and $`3d`$ is given by the sum of a term proportional to the tree level structure $`𝒫_k`$ and a term proportional to the nonplanar structure $`𝒬_k`$.
The tree level colour structure
$$𝒫_k\mathrm{Tr}(T_{a_1}\mathrm{}T_{a_k})\underset{\sigma }{}\mathrm{Tr}(T_{a_{\sigma (1)}}\mathrm{}T_{a_{\sigma (k)}})$$
(C.19)
for $`k=4`$ can be reduced to
$$𝒫_4=4\mathrm{Tr}(T_aT_bT_cT_d)[\mathrm{Tr}(T_dT_cT_bT_a)+4\mathrm{Tr}(T_aT_bT_dT_c)+\mathrm{Tr}(T_aT_bT_cT_d)]$$
(C.20)
and then, with (C.9C.11), to
$$𝒫_4=4\frac{1}{N^2}(N^21)(N^46N^2+18).$$
(C.21)
Now consider the structure from the diagrams $`3c`$ and $`3d`$
$`_k\mathrm{Tr}(T_{a_1}\mathrm{}T_{a_k})\times `$
$`{\displaystyle \underset{\sigma }{}}{\displaystyle \underset{ijl}{}}\mathrm{Tr}\left(T_{a_{\sigma (1)}}\mathrm{}[T_{a_{\sigma (i)}},T_m]\mathrm{}[[T_{a_{\sigma (j)}},T_m],T_n]\mathrm{}[T_{a_{\sigma (l)}},T_n]\mathrm{}T_{a_{\sigma (k)}}\right)`$ (C.22)
Setting $`k=4`$ it can be written as
$`_4`$ $`=`$ $`16[\mathrm{Tr}(T_{a_1}T_{a_2}T_{a_3}T_b)+\mathrm{Tr}(T_{a_1}T_bT_{a_2}T_{a_3})+\mathrm{Tr}(T_{a_1}T_{a_2}T_bT_{a_3})+\mathrm{Tr}(T_{a_1}T_{a_3}T_{a_2}T_b)`$ (C.23)
$`+\mathrm{Tr}(T_{a_1}T_bT_{a_3}T_{a_2})+\mathrm{Tr}(T_{a_1}T_{a_3}T_bT_{a_2})]f_{a_1mc_1}f_{a_2rn}f_{nmc_2}f_{a_3rc_3}[\mathrm{Tr}(T_{c_1}T_{c_2}T_{c_3}T_b)`$
$`+\mathrm{Tr}(T_{c_1}T_bT_{c_2}T_{c_3})+\mathrm{Tr}(T_{c_1}T_{c_2}T_bT_{c_3})+\mathrm{Tr}(T_{c_1}T_{c_3}T_{c_2}T_b)+\mathrm{Tr}(T_{c_1}T_bT_{c_3}T_{c_2})`$
$`+\mathrm{Tr}(T_{c_1}T_{c_3}T_bT_{c_2})]`$
Making use of the equations (C.9C.14), it is a lengthy but straightforward calculation to show that
$$_4=32(N^21)(N^418N^2+36).$$
(C.24)
At this point it is immediate, from eqs. (C.18), (C.21) and (C.24), to see that
$$_4=8N^2𝒫_4𝒬_4$$
(C.25)
in accordance with the manipulations leading to eq. (4.7).
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# Coherent scattering of high-energy photon in a medium
## 1 Introduction
The nonlinear effects of QED are due to the interaction of a photon with electron-positron field. These processes are the photon-photon scattering, the coherent photon scattering (the elastic scattering of a photon by the static Coulomb field called often the Delbrück scattering), the photon splitting into two photons, and the coalescence of two photons into photon in the Coulomb field. Among these processes the coherent scattering of photon has been observed and investigated experimentally (see reviews and ) during last decades, the most accurate measurement of the coherent scattering has been performed not long ago in the Budker Institute of Nuclear Physics . The photon splitting into two photons in the Coulomb field was observed for the first time at Budker INP only recently . Observation of photon-photon scattering is still a challenge.
History of the coherent photon scattering study can be found in the mentioned reviews. There is a special interest to the process for heavy elements because contributions of higher orders of $`Z\alpha `$ ($`Z|e|`$ is the charge of nucleus, $`e^2=\alpha =1/137,\mathrm{}=c=1`$) into the amplitude of photon scattering are very important. This means that one needs the theory which is exact with respect to the parameter $`Z\alpha `$. The amplitudes of the coherent photon scattering valid for any $`Z\alpha `$ for high energy photon ($`\omega m`$) and small scattering angle (or small momentum transfer $`\mathrm{\Delta }`$) were calculated in , . The approximate method of summing of the set of Feynman diagrams with an arbitrary number of photons exchanged with the Coulomb source was used. Another representation of these amplitudes (in the Coulomb field) was found in , using the quasiclassical Green function of the Dirac equation in the Coulomb field. This Green function in a spherically symmetrical external field was obtained in where the coherent photon scattering in the screened Coulomb potential was investigated as well. Lately this Green function was calculated for ”localized potential”, and the coherent photon scattering was analyzed using it . Recently the process of the coherent photon scattering was considered in frame of the quasiclassical operator method (see e.g. ) which appears to be very adequate for consideration of this problem.
The coherent photon scattering belong to the class of electromagnetic processes which in high-energy region occurs over rather long distance, known as the formation length. Among other processes there are the bremsstrahlung and the pair creation by a photon. If anything happens to an electron while traveling this distance, the emission can be disrupted. Landau and Pomeranchuk showed that if the formation length of bremsstrahlung becomes comparable to the distance over which a mean angle of the multiple scattering becomes comparable with a characteristic angle of radiation, the bremsstrahlung will be suppressed. Migdal developed a quantitative theory of this phenomenon which is known as the Landau-Pomeranchuk-Migdal (LPM) effect.
A very successful series of experiments (see , ) was performed at SLAC during last years. In these experiments the cross section of bremsstrahlung of soft photons with energy from 200 KeV to 500 MeV from electrons with energy 8 GeV and 25 GeV is measured with an accuracy of the order of a few percent. The LPM was observed and investigated. These experiments were the challenge for the theory since in all the previous papers calculations are performed to logarithmic accuracy which is not enough for description of the new experiment. The contribution of the Coulomb corrections (at least for heavy elements) is larger then experimental errors and these corrections should be taken into account.
Recently authors developed the new approach to the theory of the LPM effect in frame of the quasiclassical operator method. In it the cross section of bremsstrahlung process in the photon energies region where the influence of the LPM effect is very strong was calculated with term $`1/L`$ , where $`L`$ is characteristic logarithm of the problem, and with the Coulomb corrections taken into account. In the photon energy region, where the LPM effect is ”turned off”, the obtained cross section gives the exact Bethe-Heitler cross section (within power accuracy) with the Coulomb corrections. This important feature was absent in the previous calculations. The contribution of an inelastic scattering of a projectile on atomic electrons is also included. We have analyzed (see , ) the soft part of spectrum, including all the accompanying effects: the boundary photon emission, the multiphoton radiation and the influence of the polarization of a medium. Perfect agreement of the theory and SLAC data was achieved in the whole interval of measured photon energies. Very recently we apply this approach to the process of pair creation by a photon in .
In the quasiclassical approximation the amplitude $`M`$ of the coherent photon scattering is described by diagram where the electron-positron pair is created by the initial photon with 4-momentum $`k_1(\omega ,𝐤_1)`$ and then annihilate into the final photon with 4-momentum $`k_2`$. For high energy photon $`\omega m`$ this process occurs over a rather long distance, known as the time of life of the virtual state
$$l_f=\frac{\omega }{2q_c^2},$$
(1.1)
where $`q_cm`$ is the characteristic transverse momentum of the process, the system $`\mathrm{}=c=1`$ is used. When the virtual electron (or positron) is moving in a medium it scatters on atoms. The mean square of momentum transfer to the electron from a medium on the distance $`l_f`$ is
$$q_s^2=4\pi Z^2\alpha ^2n_aLl_f,L(q_c)=\mathrm{ln}\frac{q_c^2}{q_{min}^2},q_{min}^2=a^2+\mathrm{\Delta }^2+\frac{m^4}{\omega ^2},$$
(1.2)
where $`\alpha =e^2=1/137`$, $`Z`$ is the charge of nucleus, $`n_a`$ is the number density of atoms in the medium, $`\mathrm{\Delta }`$ is the photon momentum transfer ($`\mathrm{\Delta }=|𝐤_2𝐤_1|q_c`$), $`a`$ is the screening radius of atom.
The coherent photon scattering amplitude $`M`$ can be obtained from general formulas for probabilities of electromagnetic processes in the frame of the quasiclassical operator method (see e.g. ). It can be estimated as
$$M\frac{\alpha \omega }{2\pi l_fn_a}\frac{q_s^2}{q_c^2}=\frac{\alpha }{\pi n_a}q_s^2.$$
(1.3)
We use the normalization condition Im$`M=\omega \sigma _p`$ for the case $`\mathrm{\Delta }=0`$, where $`\sigma _p`$ is the total cross section of pair creation by a photon.
In the case of small momentum transfer $`q_s\sqrt{q_s^2}m`$ we have in the region of small $`\mathrm{\Delta }m`$
$$q_c^2=m^2,M\frac{2Z^2\alpha ^3\omega }{m^2}\mathrm{ln}\frac{m^2}{a^2+\mathrm{\Delta }^2+\frac{m^4}{\omega ^2}}.$$
(1.4)
At an ultrahigh energy it is possible that $`q_sm`$. In this case the characteristic momentum transfer $`q_c`$ is defined by the momentum transfer $`q_s`$. The self-consistency condition is
$$q_c^2=q_s^2=\frac{2\pi \omega Z^2\alpha ^2n_aL(q_c)}{q_c^2},$$
(1.5)
where $`L(q_c)`$ is defined in (1.2). So using (1.3) one gets for the estimate of the coherent photon scattering amplitude $`M`$ (the influence of the multiple scattering manifests itself at the high photon energies such that $`m^2a/\omega 1`$)
$$M\frac{2Z^2\alpha ^3\omega }{\mathrm{\Delta }_s^2}\sqrt{\mathrm{ln}\frac{\mathrm{\Delta }_s^2}{a^2+\mathrm{\Delta }^2}},\mathrm{\Delta }_s^2=\sqrt{2\pi \omega Z^2\alpha ^2n_a}\mathrm{\Delta }^2.$$
(1.6)
In the present paper we use consider the influence of the multiple scattering on the process of the coherent photon scattering for $`\mathrm{\Delta }q_c`$. The theory of the coherent photon scattering in the Coulomb field in frame of the quasiclassical operator method is stated in Sec.2. We give here an alternative presentation of approach developed in in a more formal way. In Sec.3 we apply the method developed in , to investigation of influence of the multiple scattering on the process of the coherent photon scattering. The general formulas describing this influence were derived and the asymptotic cases of the strong and weak effects are analyzed.
## 2 Coherent scattering of a photon <br>in the Coulomb field
### 2.1 Formulation of approach
The coherent scattering of a photon in the Coulomb field (the photon with momentum $`k_1=(\omega _1,𝐤_1)k_2=(\omega _2,𝐤_2)`$) is represented by the electron loop in the Coulomb field, and $`\omega _1=\omega _2=\omega `$. The corresponding amplitude is
$$T=\frac{2\pi \alpha }{\omega }id^4x_1d^4x_2Tr\left[\widehat{e}_1\mathrm{exp}(ik_1x_1)G(x_1,x_2)\widehat{e}_2^{}\mathrm{exp}(ik_2x_2)G(x_2,x_1)\right],$$
(2.1)
where $`\widehat{e}=e_\mu \gamma ^\mu `$, $`e_\mu `$ is the photon polarization vector, $`G(x_1,x_2)`$ is the electron Green function in the Coulomb field the standard representation of which is
$$G(x_2,x_1)=\left\{\begin{array}{cc}\hfill & i_n\mathrm{\Psi }_n^{(+)}(x_2)\overline{\mathrm{\Psi }}_n^{(+)}(x_1),t_2>t_1\hfill \\ & i_n\mathrm{\Psi }_n^{()}(x_2)\overline{\mathrm{\Psi }}_n^{()}(x_1),t_2<t_1\hfill \end{array}\right\}.$$
(2.2)
Here $`\mathrm{\Psi }_n^{(\pm )}(x_1)`$ are the solution of Dirac equation in the Coulomb field, signs $`(+)`$ and $`()`$ relate respectively to positive and negative frequencies.
As well known, see e.g., Sec.12, in the noncovariant perturbation theory, which we use here, in high-energy region ($`\omega m`$) the diagram contributes into the amplitude $`T`$ where the electron-positron pair is first created by the initial photon with 4-momentum $`k_1`$ and polarization $`e_1`$ in time $`t_1`$ and then annihilate in time $`t_2>t_1`$ into final photon with 4-momentum $`k_2`$ and polarization $`e_2`$, while the contribution of the interval in which $`t_2<t_1`$ is of the order $`m^2/\omega ^2`$. Such contribution is neglected in the scope of our method. Taking this into account and substituting the Green function into the amplitude (2.1) we find
$$T=\frac{2\pi \alpha }{\omega }i\underset{n,m}{}𝑑t_1𝑑t_2\vartheta (t_2t_1)V_{nm}(𝐞_1,𝐤_1,t_1)V_{nm}^{}(𝐞_2,𝐤_2,t_2),$$
(2.3)
where
$$V_{nm}(𝐞,𝐤,t)=d^3r\mathrm{\Psi }_n^{(+)+}(x)𝜶𝐞\mathrm{exp}(ikx)\mathrm{\Psi }_m^{()}(x),$$
(2.4)
It is evident that $`V_{nm}(𝐞,𝐤,t)`$ (2.4) is the matrix element of pair creation by a photon in an external field. In the quasiclassical operator method (see , Sec.3) the wave functions in an external field can be presented in the form
$$\mathrm{\Psi }_n^{(\pm )}(𝐫,t)=<𝐫|\mathrm{\Psi }^{(\pm )}(𝐏,)\mathrm{exp}\left[i(\pm e\phi )t\right]|n>,=\sqrt{m^2+𝐏^2},$$
(2.5)
where $`\phi =\phi (𝐫)`$ is the potential of an atom, $`|n>`$ is the state in the configuration space at the time $`t=0`$, the functions $`\mathrm{\Psi }^{(\pm )}(𝐏,)`$ have form of wave functions for free particles in the momentum space ($`\mathrm{\Psi }^{(\pm )}(𝐩,\epsilon `$). We substitute (2.5) into (2.4) and take into account completeness of states $`|𝐫>`$ ($`{\displaystyle |𝐫><𝐫|d^3r}=I`$). After this we convey consistently in (2.4) the operator $`\mathrm{exp}(i\mathrm{𝐤𝐫})`$ to the right up to $`|m>`$, and then the operator $`\mathrm{exp}\left[i(+e\phi )t\right]`$ from $`<n|`$ to $`\mathrm{exp}\left[i((𝐤𝐏)e\phi )t\right]`$. As a result we obtain the combination of operators
$`V_{nm}(𝐞,𝐤,t)=<n|\mathrm{\Psi }^{(+)+}(𝐏(t))𝜶𝐞\mathrm{\Psi }^{()}(𝐤𝐏(t))`$
$`\times \mathrm{exp}(i(e\phi )t)\mathrm{exp}(i((𝐤𝐏)+e\phi )t)\mathrm{exp}(i\mathrm{𝐤𝐫})|m>`$ (2.6)
the disentanglement of which can be performed in the standard way (see , Eqs.(3.7)-(3.11)):
$`L_p(t)=\mathrm{exp}\left[i(e\phi )t\right]\mathrm{exp}\left[i((𝐤𝐏)+e\phi )t\right]`$
$`\mathrm{exp}(i\omega t)\mathrm{T}\mathrm{exp}\left[i{\displaystyle \underset{0}{\overset{t}{}}}{\displaystyle \frac{kP(t)}{\omega (t)}}𝑑t\right],P=(,𝐏),`$ (2.7)
where $`\mathrm{T}`$ is the operator of the chronological product, and the expression
$$R_p=\overline{\mathrm{\Psi }}^{(+)}(𝐏,)\widehat{e}\mathrm{\Psi }^{()}(𝐤𝐏,(𝐤𝐏))$$
(2.8)
becomes the Heisenberg operator depending on time $`𝐏=𝐏(t),=(t)`$. It can be expressed in terms of two-component spinors. Substituting these results into Eq.(2.4) we find
$`V_{nm}(𝐞,𝐤,t)=n|R_p(t)\mathrm{T}\mathrm{exp}\left[i{\displaystyle \underset{0}{\overset{t}{}}}{\displaystyle \frac{kP(t)}{\omega (t)}}𝑑t\right]\mathrm{exp}(i\mathrm{𝐤𝐫})|m,`$
$`R_p(t)=i\phi _{\overline{s}}^+\left(A(t)i𝝈𝐁(t)\right)\phi _s,A={\displaystyle \frac{1}{2(\omega )}}(𝐞(𝐤\times 𝐏)),`$
$`𝐁={\displaystyle \frac{1}{2(\omega )}}\left[𝐞m\omega +(\mathrm{𝐞𝐏})(𝐤2𝐏)\right],`$ (2.9)
where $`\phi _{\overline{s}},\phi _s`$ are two-component spinors describing polarization of created positron and electron. Here the Coulomb gauge is used.
Substituting (2.9) into (2.3) and performing summation over states $`|n>`$ and $`|m>`$, as well as summation over spin states $`\overline{s},s`$ we have the amplitude of the coherent scattering of a photon in the Coulomb field
$`T={\displaystyle \frac{2\pi \alpha }{\omega }}i{\displaystyle 𝑑t_1𝑑t_2\vartheta (t_2t_1)S_{21}}`$
$`S_{21}=\mathrm{Tr}\left[m_2^+(t_2)\mathrm{exp}[i𝚫𝐫]m_1(t_1)\right],𝚫=𝐤_2𝐤_1,`$ (2.10)
where the operations Tr means sum of the diagonal matrix elements both in the configuration and the spin spaces,
$$m_{1,2}(t)=m(𝐞_{1,2},k_{1,2},t)=\left(A(t)i𝝈𝐁(t)\right)\mathrm{T}\mathrm{exp}\left[i\underset{0}{\overset{t}{}}\frac{k_{1,2}P(t)}{\omega (t)}𝑑t\right],$$
(2.11)
functions $`A(t),𝐁(t)`$ are defined in (2.9). Let us remind that $`𝐏(t)`$ is the momentum operator in the Heisenberg picture:
$$𝐏(t)=\mathrm{exp}(iHt)𝐏\mathrm{exp}(iHt),H=+V,$$
(2.12)
where $`V=e\phi `$
We present the evolution operator as
$$\mathrm{exp}(iHt)=\mathrm{exp}(it)N(t),N(t)=\mathrm{exp}(it)\mathrm{exp}(i(+V)t)$$
(2.13)
Differentiating the last expression over time we find
$$\frac{dN(t)}{dt}=i\mathrm{exp}(it)V(𝐫)\mathrm{exp}(i(+V)t)=iV(𝐫+𝐯t)N(t),𝐯=\frac{𝐏}{}$$
(2.14)
The solution of this differential equation for the initial condition $`N(0)=1`$ is
$$N(t)=\mathrm{T}\mathrm{exp}\left[i\underset{0}{\overset{t}{}}V(𝐫+𝐯t^{})𝑑t^{}\right],$$
(2.15)
If the formation length of photon scattering is much longer than the characteristic length of electron scattering ($`\omega /m^2a`$), one can present the dependence on time of the operator $`𝐏(t)`$ in (2.9) as (see e.g. , Sec.7.1)
$$𝐏(t)=\vartheta (t)𝐏(\mathrm{})+\vartheta (t)𝐏(\mathrm{}),𝐏(\pm \mathrm{})=N^+(\pm \mathrm{})𝐏(t)N(\pm \mathrm{}).$$
(2.16)
Let us introduce states convenient for further calculations
$`|i>=N^+(\mathrm{})|𝐩_i>,|f>=N^+(\mathrm{})|𝐩_f>,`$
$`𝐏|𝐩_i>=𝐩_i|𝐩_i>,𝐏|𝐩_f>=𝐩_f|𝐩_f>.`$ (2.17)
The states $`|i>`$ and $`|f>`$ are the eigenvectors of the operators $`𝐏(\mathrm{})`$ and $`𝐏(\mathrm{})`$ correspondingly
$`𝐏(\mathrm{})|i>=N^+(\mathrm{})𝐏N(\mathrm{})N^+(\mathrm{})|𝐩_i>=𝐩_i|i>`$
$`𝐏(\mathrm{})|f>=N^+(\mathrm{})𝐏N(\mathrm{})N^+(\mathrm{})|𝐩_f>=𝐩_f|f>.`$ (2.18)
The general expression for the amplitude of photon scattering (Eqs.(2.10) and (2.11)) is simplified significantly for small momentum transfers $`\mathrm{\Delta }=|𝚫|`$ ($`\mathrm{\Delta }m,\mathrm{\Delta }/\epsilon 1/\gamma `$). With the accuracy $`\mathrm{\Delta }/m`$ we can put that the operator $`𝐏`$ commutes with $`\mathrm{exp}[i𝚫𝐫]`$ and substitute the vector $`𝐤_2`$ in the expression for $`m_2(t_2)`$ by the vector $`𝐤_1`$. Using (2.16) and the states $`|i>`$ and $`|f>`$ (2.18) for construction of the matrix element we have
$`f|m(𝐏(t))\mathrm{exp}\left[{\displaystyle \frac{i}{2}}𝚫𝐫\right]|i=m_{fi}(𝐞,𝐩(t))f|\mathrm{exp}\left[{\displaystyle \frac{i}{2}}𝚫𝐫\right]|i,`$
$`m_{fi}(𝐞,𝐩(t))=\left[\vartheta (t)m(𝐞,𝐩_i)+\vartheta (t)m(𝐞,𝐩_f)\right],`$ (2.19)
where we describe motion of particles of the created pair as a trajectory in ”the form of an angle”. Bearing in mind that the commutator
$$[r_i,v_j]=\frac{i}{}(\delta _{ij}v_iv_j)$$
(2.20)
we can discard operator T in the expression for operator $`N(t)`$ (2.15) and with relativistic accuracy (i.e. with accuracy up to terms $`m/\omega `$) present the operator $`N(t)`$ as
$$N(t)=\mathrm{exp}\left[i\underset{0}{\overset{t}{}}V(\mathit{\varrho },z+t^{})𝑑t^{}\right],$$
(2.21)
where the axis $`z`$ is directed along the momentum of the initial photon $`𝐤_1`$, $`\mathit{\varrho }`$ is two-dimensional vector transverse to axis z.
Using (2.17) and (2.21) we calculate the matrix element in (2.19)
$`f|\mathrm{exp}\left[{\displaystyle \frac{i}{2}}𝚫𝐫\right]|i={\displaystyle <f|𝐫><𝐫|i>\mathrm{exp}\left[\frac{i}{2}𝚫𝐫\right]d^3r}`$
$`={\displaystyle \frac{i}{2\pi }}f\left(𝐪_{}{\displaystyle \frac{𝚫_{}}{2}}\right)\delta \left(𝐪_{}{\displaystyle \frac{𝚫_{}}{2}}\right),`$ (2.22)
where
$`<f|𝐫>=<𝐩_f|N(\mathrm{})|𝐫>=N(\mathrm{})<𝐩_f|𝐫>,<𝐩_f|𝐫>={\displaystyle \frac{\mathrm{exp}(i𝐩_f𝐫)}{(2\pi )^{3/2}}},`$
$`<𝐫|i>=<𝐫|N^+(\mathrm{})|𝐩_i>=N^+(\mathrm{})<𝐫|𝐩_i>,<𝐫|𝐩_i>={\displaystyle \frac{\mathrm{exp}(i𝐩_i𝐫)}{(2\pi )^{3/2}}},`$
$`f(𝐐)={\displaystyle \frac{1}{2\pi i}}{\displaystyle \mathrm{exp}[i𝐐\mathit{\varrho }+i\chi (\mathit{\varrho })]d^2\varrho },\chi (\mathit{\varrho })={\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}V(\mathit{\varrho },z)𝑑z.`$ (2.23)
We introduced here vectors
$$𝐩_i𝐩,𝐩_f=𝐩+𝐪,$$
(2.24)
the function $`f(𝐐)`$ is the scattering amplitude in the eikonal approximation.
Substituting Eq.(2.19) into $`T_{21}`$ Eq.(2.10) we have
$`S_{21}=\mathrm{Tr}\left[m_2^+(t_2)\mathrm{exp}[i𝚫𝐫]m_1(t_1)\right]={\displaystyle \underset{i,f}{}}\mathrm{Tr}\left[m_{fi}^+(𝐞_2,𝐩(t_2))m_{fi}(𝐞_1,𝐩(t_1))\right]`$
$`\times i|\mathrm{exp}\left[{\displaystyle \frac{i}{2}}𝚫𝐫\right]|ff|\mathrm{exp}\left[{\displaystyle \frac{i}{2}}𝚫𝐫\right]|i`$
$`={\displaystyle \frac{\delta (\mathrm{\Delta }_{})}{(2\pi )^2}}{\displaystyle d^3p\mathrm{Tr}\left[m_{fi}^+(𝐞_2,𝐩(t_2))m_{fi}(𝐞_1,𝐩(t_1))\right]}`$
$`\times {\displaystyle }d𝐪_{}f(𝐪_{}{\displaystyle \frac{𝚫_{}}{2}})f^{}(𝐪_{}+{\displaystyle \frac{𝚫_{}}{2}}),q_{}=0.`$ (2.25)
### 2.2 Helicity amplitudes
In the above analysis we traced the transition to the expressions calculated on the trajectory of particle ”in the form of an angle” in the momentum space. Precisely these trajectories determine the amplitude of the coherent photon scattering for $`\omega /m^2a`$ (see Eq.(2.10)):
$$T_{21}=\frac{i\alpha }{(2\pi )^2}(2\pi \delta (\mathrm{\Delta }_{}))\frac{d^3p}{\omega }𝑑\sigma (𝐪_{},𝚫)R_{21},$$
(2.26)
where
$`d\sigma (𝐪_{},𝚫)=f\left(𝐪_{}{\displaystyle \frac{𝚫_{}}{2}}\right)f^{}\left(𝐪_{}+{\displaystyle \frac{𝚫_{}}{2}}\right)d𝐪_{},`$
$`R_{21}={\displaystyle 𝑑t_1𝑑t_2\vartheta (t_2t_1)\mathrm{Tr}\left[m_{fi}^+(𝐞_2,𝐩(t_2))m_{fi}(𝐞_1,𝐩(t_1))\right]}`$
$`={\displaystyle \frac{1}{2}}{\displaystyle 𝑑t_1𝑑t_2\vartheta (t_2t_1)(𝐞_2,𝐞_1;𝐯_2,𝐯_1)\mathrm{exp}\left[i\frac{\epsilon }{\epsilon ^{}}\underset{t_1}{\overset{t_2}{}}k_1v(t)𝑑t\right]},`$
$`(𝐞_2,𝐞_1;𝐯_2,𝐯_1)={\displaystyle \frac{\omega ^2}{\epsilon ^2}}[(𝐞_2^{}\mathrm{𝐧𝐯}_2)(𝐞_1\mathrm{𝐧𝐯}_1)+{\displaystyle \frac{m^2}{\epsilon ^2}}(𝐞_2^{}𝐞_1)`$
$`+{\displaystyle \frac{(\omega 2\epsilon )^2}{\omega ^2}}(𝐞_2^{}𝐯_2)(𝐞_1𝐯_1)],𝐯_{1,2}=𝐯_{}(t_{1,2})`$
$`v={\displaystyle \frac{p}{\epsilon }}=(1,𝐯),\epsilon ^{}=\omega \epsilon ,𝐧={\displaystyle \frac{𝐤_1}{\omega }},(𝐯_{1,2}𝐧)=0,`$ (2.27)
where $`(\mathrm{𝐞𝐧𝐯})(𝐞(𝐧\times 𝐯))`$.
We calculate now $`R_{21}`$ in this expression for the trajectory in ”the form of an angle” (see (2.19)):
$$𝐩(t)=𝐩\vartheta (t)+(𝐩+𝐪)\vartheta (t).$$
(2.28)
We get for this case
$`m^2R_{21}=(𝐞_2^{}\mathrm{𝐧𝐠})(𝐞_1\mathrm{𝐧𝐠})+{\displaystyle \frac{(\omega 2\epsilon )^2}{\omega ^2}}(𝐞_2^{}𝐠)(𝐞_1𝐠)+(𝐞_2^{}𝐞_1)g_0^2;`$
$`𝐠={\displaystyle \frac{𝐩_{}+𝐪_{}}{m^2+(𝐩_{}+𝐪_{})^2}}{\displaystyle \frac{𝐩_{}}{m^2+𝐩_{}^2}},g_0={\displaystyle \frac{m}{m^2+(𝐩_{}+𝐪_{})^2}}{\displaystyle \frac{m}{m^2+𝐩_{}^2}}.`$ (2.29)
The functions $`𝐠`$ and $`g_0`$ have very important properties:
$$𝐠(𝐪_{}=0)=0,g_0(𝐪_{}=0)=0.$$
(2.30)
In the case of complete screening when the screening radius $`a\omega /m^2`$ the amplitude of photon scattering $`T_{21}`$ (2.26) is imaginary at any $`𝚫`$. Taking into account that at $`𝚫=0`$ this amplitude is connected due to unitary relation with the known probability of pair photoproduction, see e.g. , , we will calculate the difference
$$\delta R_{21}=T_{21}(𝚫)T_{21}(0).$$
The interval $`|𝐪_{}||𝚫|m`$ contributes into this difference and one can expand the functions $`𝐠,g_0`$ as a power series in $`𝐪_{}`$:
$$𝐠\frac{𝐪_{}}{m^2+𝐩_{}^2}\frac{2𝐩_{}(𝐪_{}𝐩_{})}{(m^2+𝐩_{}^2)^2},g_0\frac{2m(𝐪_{}𝐩_{})}{(m^2+𝐩_{}^2)^2}.$$
(2.31)
We introduce the dimensionless variable $`𝐮=𝐩_{}/m`$ and carry on averaging over angles of $`𝐮`$ (the azimuthal angle of the component of electron momentum in the plane which is perpendicular to the direction of initial photon $`𝐧`$) using formulas
$$\overline{u_iu_j}=\frac{u^2}{2}\delta _{ij},\overline{u_iu_ju_ku_l}=\frac{u^4}{8}\left(\delta _{ij}\delta _{kl}+\delta _{ik}\delta _{jl}+\delta _{il}\delta _{jk}\right),$$
(2.32)
where $`\delta _{ij}`$ is the two-dimensional Kronecker delta. Substituting Eq.(2.31) into (2.29) and using (2.32) we get
$`m^4\zeta ^4R_{21}=(𝐞_2^{}\mathrm{𝐧𝐪}_{})(𝐞_1\mathrm{𝐧𝐪}_{})+{\displaystyle \frac{(\epsilon ^{}\epsilon )^2}{\omega ^2}}(𝐞_2^{}𝐪_{})(𝐞_1𝐪_{})`$
$`+𝐪_{}^2(𝐞_2^{}𝐞_1)\left[\left(1{\displaystyle \frac{2\epsilon \epsilon ^{}}{\omega ^2}}\right)(\zeta ^21)+{\displaystyle \frac{4\epsilon \epsilon ^{}}{\omega ^2}}(\zeta 1)\right],\zeta =1+𝐮^2.`$ (2.33)
It is convenient to describe the process of photon scattering in terms of helicity amplitudes. We choose the polarization vectors with helicity $`\lambda `$
$`𝐞_\lambda ={\displaystyle \frac{1}{\sqrt{2}}}\left(𝐞_1+i\lambda 𝐞_2\right),𝐞_1=𝝂={\displaystyle \frac{𝚫}{|𝚫|}},𝐞_2=𝐧\times 𝝂,\lambda =\pm 1,`$
$`𝐞_\lambda 𝐞_\lambda ^{}=1,𝐞_\lambda 𝐞_\lambda ^{}=0,𝐞_\lambda \times 𝐧=i\lambda 𝐞_\lambda ,`$ (2.34)
where $`𝐧`$ is defined in Eq.(2.27). There are two independent helicity amplitudes:
$$M_{++}=M_{},M_+=M_+,$$
where the first subscript is the helicity of the initial photon and the second is the helicity of the final photon. When the initial photons are unpolarized the differential cross section of scattering summed over final photons polarization contains the combination
$$2[|M_{++}|^2+|M_+|^2].$$
(2.35)
Substituting (2.34) into (2.33) we find
$`R_{++}={\displaystyle \frac{𝐪_{}^2}{m^4\zeta ^4}}\left[\left(1{\displaystyle \frac{2\epsilon \epsilon ^{}}{\omega ^2}}\right)(\zeta ^21)+{\displaystyle \frac{4\epsilon \epsilon ^{}}{\omega ^2}}(\zeta 1)\right],`$
$`R_+={\displaystyle \frac{2𝐪_{}^2}{m^4\zeta ^4}}{\displaystyle \frac{\epsilon \epsilon ^{}}{\omega ^2}}\mathrm{cos}2\phi ,\mathrm{cos}\phi ={\displaystyle \frac{𝐪𝝂}{q}}.`$ (2.36)
We define the amplitude of the coherent photon scattering $`M_{21}`$ as
$$T_{21}=\frac{\pi }{\omega }\delta (\mathrm{\Delta }_{})M_{21},$$
(2.37)
then the cross section of the coherent photon scattering has a form
$$d\sigma _s=\frac{1}{2\pi \delta (\mathrm{\Delta }_{})}|T_{21}|^2\frac{d^3k_2}{(2\pi )^3}=\frac{1}{16\pi ^2}|M_{21}|^2d\mathrm{\Omega }_2.$$
(2.38)
Substituting (2.37), (2.33) and (2.36) into Eq.(2.26) we obtain for the helicity amplitudes
$$M_{\lambda \lambda ^{}}=\frac{i\alpha }{2\pi }m^2\underset{0}{\overset{\omega }{}}𝑑\epsilon \underset{1}{\overset{\mathrm{}}{}}𝑑\zeta 𝑑\sigma (𝐪_{},𝚫)R_{\lambda \lambda ^{}},$$
(2.39)
where $`\lambda \lambda ^{}`$ is $`++`$ or $`+`$. For the chosen normalization of the helicity amplitudes the unitary relation with total cross section of pair photoproduction $`\sigma _p`$ is
$$\frac{1}{\omega }\mathrm{Im}M_{++}=\sigma (\gamma e^+e^{})\sigma _p.$$
(2.40)
So we present the helicity amplitude $`M_{++}`$ as
$`M_{++}=\delta M_{++}+i\omega \sigma _p,`$
$`\delta M_{++}={\displaystyle \frac{i\alpha }{2\pi }}m^2{\displaystyle \underset{0}{\overset{\omega }{}}}𝑑\epsilon {\displaystyle \underset{1}{\overset{\mathrm{}}{}}}𝑑\zeta {\displaystyle (d\sigma (𝐪_{},𝚫)d\sigma (𝐪_{},0))R_{++}}.`$ (2.41)
Since the forward scattering amplitude ($`\mathrm{\Delta }=0`$) with helicity flip vanishes we have for amplitude $`M_+`$
$$M_+=\delta M_+=\frac{i\alpha }{2\pi }m^2\underset{0}{\overset{\omega }{}}𝑑\epsilon \underset{1}{\overset{\mathrm{}}{}}𝑑\zeta 𝑑\sigma (𝐪_{},𝚫)R_+.$$
(2.42)
Integrals over $`𝐪_{}`$ in Eqs.(2.41) and (2.42) we denote $`\overline{\delta 𝐪_{}^2}`$ and $`\overline{𝐪_{}^2\mathrm{cos}2\phi }`$. We calculate these integrals using the scattering amplitude $`f(𝐐)`$ in the eikonal approximation Eq.(2.23). We find
$`\overline{\delta 𝐪^2}Z_1={\displaystyle 𝐪_{}^2(d\sigma (𝐪_{},𝚫)d\sigma (𝐪_{},0))}`$
$`={\displaystyle 𝑑𝐪_{}𝐪_{}^2\left[f^{}\left(𝐪_{}+\frac{𝚫_{}}{2}\right)f\left(𝐪_{}\frac{𝚫_{}}{2}\right)f^{}\left(𝐪_{}\right)f\left(𝐪_{}\right)\right]}`$
$`=2\pi {\displaystyle \underset{0}{\overset{\mathrm{}}{}}}(\chi ^{}(\varrho ))^2(J_0(\mathrm{\Delta }\varrho )1)\varrho 𝑑\varrho ,`$
$`\overline{𝐪^2\mathrm{cos}2\phi }Z_2=2\pi {\displaystyle \underset{0}{\overset{\mathrm{}}{}}}(\chi ^{}(\varrho ))^2J_2(\mathrm{\Delta }\varrho )\varrho 𝑑\varrho .`$ (2.43)
For the screened Coulomb potential
$`U(r)={\displaystyle \frac{Z\alpha }{r}}\mathrm{exp}(r/a),\chi ^{}(\varrho )={\displaystyle \frac{2Z\alpha }{a}}K_1\left({\displaystyle \frac{\varrho }{a}}\right),`$
$`Z_1=4\pi Z^2\alpha ^2F_2\left({\displaystyle \frac{\mathrm{\Delta }a}{2}}\right),F_2(z)={\displaystyle \frac{2z^2+1}{z\sqrt{1+z^2}}}\mathrm{ln}\left(z+\sqrt{1+z^2}\right)1,`$
$`Z_2=4\pi Z^2\alpha ^2F_1\left({\displaystyle \frac{\mathrm{\Delta }a}{2}}\right),F_1(z)=1{\displaystyle \frac{1}{z\sqrt{1+z^2}}}\mathrm{ln}\left(z+\sqrt{1+z^2}\right).`$ (2.44)
The functions $`F_1(x)`$ and $`F_2(x)`$ encounter in the radiation theory and in the theory of Landau-Pomeranchuk-Migdal (LPM) effect. In this derivation the following relations have been used
$`{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}xK_1^2(x)J_2(\beta x)={\displaystyle \frac{1}{\beta }}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}(xK_1(x))^2{\displaystyle \frac{d}{dx}}\left({\displaystyle \frac{J_1(\beta x)}{x}}\right)`$
$`={\displaystyle \frac{1}{2}}{\displaystyle \frac{2}{\beta }}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}K_0(x)K_1(x)xJ_1(\beta x)`$
$`={\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{\beta }}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{d}{dx}}(K_0(x))^2xJ_1(\beta x)={\displaystyle \frac{1}{2}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}(K_0(x))^2J_0(\beta x),`$
$`{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}xK_0^2(x)x𝑑x={\displaystyle \frac{1}{2}}.`$ (2.45)
The cross section of $`e^{}e^+`$ pair photoproduction in the case of complete screening ($`a\omega /m^2`$) to within terms $`m/\omega `$ has a form (see e.g. Eq.(19.17) in )
$$\sigma _p=\frac{28Z^2\alpha ^3}{9m^2}\left[\mathrm{ln}(ma)+\frac{1}{2}f(Z\alpha )\frac{1}{42}\right],$$
(2.46)
where
$$f(\xi )=\mathrm{Re}\left[\psi (1+i\xi )\psi (1)\right]=\xi ^2\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n(n^2+\xi ^2)},$$
(2.47)
here $`\psi (x)`$ is the logarithmic derivative of the gamma function. Substituting (2.43)-(2.46) into Eqs.(2.41) and (2.42) we obtain
$$\mathrm{Im}M_{\lambda \lambda ^{}}=\frac{4Z^2\alpha ^3\omega }{m^2}\underset{0}{\overset{1}{}}𝑑x\underset{0}{\overset{1}{}}𝑑y\mu _{\lambda \lambda ^{}}f_{\lambda \lambda ^{}},$$
(2.48)
where
$`\mu _{++}=12x(1x)+4x(1x)y(1y),\mu _+=x(1x)y^2;`$
$`f_{++}={\displaystyle \frac{1}{2}}F_2\left({\displaystyle \frac{\mathrm{\Delta }a}{2}}\right)+\mathrm{ln}(ma)+{\displaystyle \frac{1}{2}}f(Z\alpha ){\displaystyle \frac{1}{42}}`$
$`=\mathrm{ln}(ma){\displaystyle \frac{2s^2+1}{2s\sqrt{1+s^2}}}\mathrm{ln}\left(s+\sqrt{1+s^2}\right)f(Z\alpha )+{\displaystyle \frac{41}{42}},`$
$`f_+=F_1\left(s\right),s={\displaystyle \frac{\mathrm{\Delta }a}{2}}.`$ (2.49)
Here we passed to the variables $`x=\epsilon /\omega ,y=1/\zeta `$.
### 2.3 The coherent photon scattering in different cases
The important property of Eq.(2.48) is that the dependence on the screening radius $`a`$ originates in it from the Born approximation. In this approximation in the case of arbitrary screening the radius $`a`$ enters only in the combination
$$\frac{1}{a^2}+q_{}^2,q_{}=\frac{m^2\omega }{2\epsilon \epsilon ^{}}(1+u^2)=\frac{q_m}{x(1x)y},q_m=\frac{m^2}{2\omega }.$$
(2.50)
Since in the course of derivation of Eq.(2.48) we did not integrate by parts over the variables $`x`$ and $`y`$, we can extend Eq.(2.49) on the case of arbitrary screening making the substitution
$$\frac{1}{a}\sqrt{q_{}^2+a^2}q_{ef},s=\frac{\mathrm{\Delta }}{2q_{ef}}$$
(2.51)
Recently the amplitudes of the coherent photon scattering were derived in for arbitrary interrelation between parameters in the Molière potential.
We calculate the real part of scattering amplitude using the dispersion relations method. We can apply this method directly since in the region of momentum transfer under consideration ($`q_{}m`$) the dependence of amplitude $`T_{21}M_{21}/\omega `$ on $`q_{}`$ has the form (see Eqs.(2.10)-(2.15), (2.21), (2.24))
$$T(q_{})=𝑑t_1𝑑t_2F(t_1,t_2)\vartheta (t_2t_1)\mathrm{exp}[iq_{}(t_2t_1)].$$
(2.52)
Thus the real and imaginary part of the amplitude $`M_{21}/\omega `$ are connected by the Sokhotsky-Plemelj formulas
$$\mathrm{Re}f(z)=\frac{𝒫}{\pi }\underset{\mathrm{}}{\overset{\mathrm{}}{}}\frac{\mathrm{Im}f(z^{})}{z^{}z}𝑑z^{}$$
(2.53)
Using this transformation and Eqs.(2.48), (2.49) we obtain
$$\mathrm{Re}M_{\lambda \lambda ^{}}=\frac{4Z^2\alpha ^3\omega }{m^2}\underset{0}{\overset{1}{}}𝑑x\underset{0}{\overset{1}{}}𝑑y\mu _{\lambda \lambda ^{}}\phi _{\lambda \lambda ^{}},$$
(2.54)
where
$`\phi _{\lambda \lambda ^{}}={\displaystyle \frac{1}{\pi }}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}\left[f_{\lambda \lambda ^{}}(z+q_{})f_{\lambda \lambda ^{}}(|zq_{}|)\right]{\displaystyle \frac{dz}{z}}`$
$`={\displaystyle \frac{1}{\pi }}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}\left[f_{\lambda \lambda ^{}}(q_{}(z+1))f_{\lambda \lambda ^{}}(q_{}|z1|)\right]{\displaystyle \frac{dz}{z}}`$ (2.55)
We consider now different limiting cases depending on interrelation between $`a`$ and $`\omega /m^2`$. In the case of complete screening ($`a\omega /m^2`$) we can use directly Eqs.(2.48), (2.49). In this case the functions $`f_{\lambda \lambda ^{}}`$ are independent of $`x`$ and $`y`$ and the corresponding integrals are
$$\underset{0}{\overset{1}{}}𝑑x\underset{0}{\overset{1}{}}𝑑y\mu _{++}=\frac{7}{9},\underset{0}{\overset{1}{}}𝑑x\underset{0}{\overset{1}{}}𝑑y\mu _+=\frac{1}{18}.$$
(2.56)
For the scattering amplitudes we have
$`\mathrm{Im}M_{++}={\displaystyle \frac{28Z^2\alpha ^3\omega }{9m^2}}f_{++},\mathrm{Im}M_+={\displaystyle \frac{2Z^2\alpha ^3\omega }{9m^2}}f_+,`$
$`\mathrm{Re}M_{++}=0,\mathrm{Re}M_+=0`$ (2.57)
The results obtained are consistent with Eqs.(33) and (36) of where calculation has been done for the Molière potential.
In the case $`a\omega /m^2`$ (the screening radius is very large, or in other words we consider the photon scattering in the Coulomb field) we have, taking into account Eqs.(2.49), (2.51)
$`f_{++}=\mathrm{ln}{\displaystyle \frac{m}{q_{}}}{\displaystyle \frac{2s_c^2+1}{s_c\sqrt{1+s_c^2}}}\mathrm{ln}\left(s_c+\sqrt{1+s_c^2}\right)f(Z\alpha )+{\displaystyle \frac{41}{42}},`$
$`f_+=1{\displaystyle \frac{1}{s_c\sqrt{1+s_c^2}}}\mathrm{ln}\left(s_c+\sqrt{1+s_c^2}\right),s_c={\displaystyle \frac{\mathrm{\Delta }}{2q_{}}}={\displaystyle \frac{\mathrm{\Delta }\omega }{m^2}}x(1x)y.`$ (2.58)
The expressions for amplitudes $`M_{++}`$ and $`M_+`$ in this case are consistent with the results obtained in . At $`\mathrm{\Delta }=0`$ we find known result (see Eq.(8.1) in )
$`f_{++}=\mathrm{ln}\left[{\displaystyle \frac{2\omega }{m}}x(1x)y\right]{\displaystyle \frac{1}{2}}f(Z\alpha )+{\displaystyle \frac{41}{42}},`$
$`\phi _{++}={\displaystyle \frac{1}{\pi }}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}\mathrm{ln}{\displaystyle \frac{z+1}{|z1|}}{\displaystyle \frac{dz}{z}}={\displaystyle \frac{2}{\pi }}{\displaystyle \underset{0}{\overset{1}{}}}\mathrm{ln}{\displaystyle \frac{1+z}{1z}}{\displaystyle \frac{dz}{z}}={\displaystyle \frac{\pi }{2}},`$
$`M_{++}=i{\displaystyle \frac{28Z^2\alpha ^3\omega }{9m^2}}\left[\mathrm{ln}{\displaystyle \frac{2\omega }{m}}f(Z\alpha ){\displaystyle \frac{109}{42}}i{\displaystyle \frac{\pi }{2}}\right],M_+=0.`$ (2.59)
At $`\mathrm{\Delta }=0`$ and arbitrary interrelation between $`a`$ and $`\omega /m^2`$ we get
$`f_{++}=\mathrm{ln}{\displaystyle \frac{m}{q_{ef}}}{\displaystyle \frac{1}{2}}f(Z\alpha )+{\displaystyle \frac{41}{42}},`$
$`\phi _{++}={\displaystyle \frac{1}{2\pi }}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}\mathrm{ln}{\displaystyle \frac{\delta ^2(z+1)^2+x^2(1x)^2y^2}{\delta ^2(z1)^2+x^2(1x)^2y^2}}{\displaystyle \frac{dz}{z}},\delta ={\displaystyle \frac{am^2}{2\omega }}={\displaystyle \frac{m}{2\omega }}{\displaystyle \frac{a}{\lambda _c}}.`$ (2.60)
The photon scattering amplitude in this case for arbitrary value of parameter $`\delta `$ was found recently in .
## 3 Influence of the multiple scattering on the coherent scattering of a photon
### 3.1 Basic equations
When a photon is propagating in a medium it dissociates with probability $`\alpha `$ into an electron-positron pair. The virtual electron and positron interact with a medium and can scatter on atoms. In this scattering the electron and positron interaction with the Coulomb field in the course of the coherent scattering of photon is involved also. There is a direct analogue with the Landau-Pomeranchuk-Migdal (LPM) effect: the influence of the multiple scattering on process of the bremsstrahlung and pair creation by a photon in a medium at high energy , . However there is the difference: in the LPM effect the particles of electron-positron pair created by a photon are on the mass shell while in the process of the coherent scattering of photon this particles are off the mass shell, but in the high energy region (this is the only region where the influence of the multiple scattering is pronounced) the shift from the mass shell is relatively small. To include this scattering into consideration the amplitude $`R_{21}`$ Eq.(2.26) should be averaged over all possible trajectories of electron and positron with the weight function $`d\sigma (𝐪,𝚫)`$. This operation can be performed with the aid of the distribution function averaged over the atomic positions of scatterers in the medium. This procedure was worked out in details in (Sec.2), (Sec.2). The amplitude of the coherent photon scattering for a photon propagating in a medium can be derived in the same way as Eqs.(2.4)-(2.6) of :
$`M(𝐞_1,𝐞_2,𝚫)={\displaystyle \frac{i\alpha }{(2\pi )^2n_a}}{\displaystyle d^3p\mathrm{exp}\left(i\frac{\epsilon }{\epsilon ^{}}\tau \right)𝑑𝐯^{}𝑑𝐫^{}(𝐞_1,𝐞_2;\mathit{\vartheta }^{},\mathit{\vartheta })}`$
$`\times F(𝐫^{},𝐯^{},\tau ;𝐫,𝐯)\mathrm{exp}\left[i{\displaystyle \frac{\epsilon }{\epsilon }}^{}𝐤_1(𝐫^{}𝐫)\right].`$ (3.1)
The distribution function $`F(𝐫^{},𝐯^{},\tau ;𝐫,𝐯)`$ satisfies the kinetic equation
$`{\displaystyle \frac{F(𝐫^{},𝐯^{},\tau )}{\tau }}+𝐯^{}{\displaystyle \frac{F(𝐫^{},𝐯^{},\tau )}{𝐫^{}}}`$
$`=n_a{\displaystyle 𝑑\sigma (𝐯^{},𝐯^{\prime \prime },𝚫)\left[F(𝐫^{},𝐯^{\prime \prime },\tau )F(𝐫^{},𝐯^{},\tau )\right]},`$ (3.2)
and the initial condition
$$F(𝐫^{},𝐯^{},0;𝐫,𝐯)=\delta (𝐫𝐫^{})\delta (𝐯𝐯^{}),$$
here $`n_a`$ is the number density of atoms in a medium and in the case of the screened Coulomb potential
$$d\sigma (𝐯^{},𝐯^{\prime \prime },𝚫)=d\sigma (𝐪,𝚫),𝐪=\epsilon (𝐯^{}𝐯).$$
(3.3)
The combinations entering in Eq.(2.27) are
$`{\displaystyle \underset{t_1}{\overset{t_2}{}}}kv(t)𝑑t=\omega \tau 𝐤_1(𝐫(t_2)𝐫(t_1))\omega \tau 𝐤_1(𝐫^{}𝐫),`$
$`(𝐞_1,𝐞_2;𝐯_2,𝐯_1)(𝐞_1,𝐞_2;\mathit{\vartheta }^{},\mathit{\vartheta })`$ (3.4)
We can proceed with further calculation (integration over $`𝐯^{},𝐫^{}`$) of (3.1) by analogy with the procedure used in Eqs.(2.7)-(2.18) of . As a result we obtain for the differential cross section of the coherent photon scattering with the multiple scattering taken into account
$$\frac{d\sigma }{d\mathrm{\Omega }}=\frac{1}{16\pi ^2}\left|M(𝐞_1,𝐞_2,𝚫)\right|^2,$$
(3.5)
where
$`M(𝐞_1,𝐞_2,𝚫)={\displaystyle \frac{2i\alpha m^2\omega }{n_a}}{\displaystyle }{\displaystyle \frac{d\epsilon }{\epsilon \epsilon ^{}}}{\displaystyle \underset{o}{\overset{\mathrm{}}{}}}dt\mathrm{exp}(it)[(𝐞_2^{}𝐞_1)\phi _0(0,t)`$
$`i(𝐞_1𝐧\mathbf{})(𝐞_2^{}𝐧𝝋(0,t))i{\displaystyle \frac{(\omega 2\epsilon )^2}{\omega ^2}}(𝐞_1\mathbf{})(𝐞_2^{}𝝋(0,t))].`$ (3.6)
In derivation we have changed variables into $`t,\mathit{\varrho }`$ as in Eq.(2.7) of . The function $`\phi _\mu (\mathit{\varrho },t),\phi _\mu =(\phi _0,𝝋)`$ satisfies the equation
$`i{\displaystyle \frac{\phi _\mu }{t}}=\phi _\mu ,=𝐩^2iV(\mathit{\varrho },𝚫),𝐩=i\mathbf{}_\mathit{\varrho },`$
$`V(\mathit{\varrho },𝚫)={\displaystyle \frac{2\epsilon \epsilon ^{}n_a}{\omega m^4}}{\displaystyle \left(1\mathrm{exp}(i𝐪\mathit{\varrho })\right)f\left(𝐪_{}\frac{𝚫_{}}{2m}\right)f^{}\left(𝐪_{}+\frac{𝚫_{}}{2m}\right)d^2q},`$ (3.7)
with the initial conditions
$$\phi _0(\mathit{\varrho },0)=\delta (\mathit{\varrho }),𝝋(\mathit{\varrho },0)=𝐩\delta (\mathit{\varrho }).$$
(3.8)
The potential $`V(\mathit{\varrho },0)V(\mathit{\varrho })`$ was used in the theory of the LPM effect (see (Sec.2)):
$`V(\mathit{\varrho })=Q\mathit{\varrho }^2\left(L_1+\mathrm{ln}{\displaystyle \frac{4}{\mathit{\varrho }^2}}2C\right),Q={\displaystyle \frac{2\pi Z^2\alpha ^2\epsilon \epsilon ^{}n_a}{m^4\omega }},L_1=\mathrm{ln}{\displaystyle \frac{a_{s2}^2}{\lambda _c^2}},`$
$`{\displaystyle \frac{a_{s2}}{\lambda _c}}=183Z^{1/3}\mathrm{e}^f,f=f(Z\alpha )=(Z\alpha )^2{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{k(k^2+(Z\alpha )^2)}},`$ (3.9)
where $`C=0.577216\mathrm{}`$ is Euler’s constant. We can restrict ourselves to calculation of the difference of the potentials only:
$`\mathrm{\Delta }V(\mathit{\varrho },𝚫)=V(\mathit{\varrho },𝚫)V(\mathit{\varrho })`$
$`={\displaystyle \frac{2\epsilon \epsilon ^{}n_a}{\omega m^4}}{\displaystyle \left(1\mathrm{exp}(i𝐪\mathit{\varrho })\right)(d\sigma (𝐪,𝚫)d\sigma (𝐪,0))d^2q},`$ (3.10)
where the cross section $`d\sigma (𝐪,𝚫)`$ is defined in Eq.(2.26). In the above formulas $`\mathit{\varrho }`$ is the two-dimensional space of the impact parameters measured in the Compton wavelength $`\lambda _c`$ which is conjugate to the two-dimensional space of the transverse momentum transfers $`𝐪`$ measured in the electron mass $`m`$.
### 3.2 Amplitudes of coherent photon scattering under influence of the multiple scattering
When parameter $`\mathrm{\Delta }m`$ the main contribution into the integral in (3.10) gives the region $`|𝐪|\mathrm{\Delta }/m1`$. For $`\varrho 1`$ one can expand the integrand in (3.10) in powers of $`𝐪\mathit{\varrho }`$. Using Eq.(2.23) we get
$`\mathrm{\Delta }V(\mathit{\varrho },𝚫)={\displaystyle \frac{\epsilon \epsilon ^{}n_a}{\omega m^4}}{\displaystyle \frac{(𝐱\mathit{\varrho })^2}{x^2}(\chi ^{}(𝐱))^2(\mathrm{exp}(i𝜹𝐱)1)d^2x},`$
$`\chi (𝐱)={\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{Z\alpha }{r}}\mathrm{exp}\left({\displaystyle \frac{r}{a_s}}\right)=2Z\alpha K_0\left(x{\displaystyle \frac{a_s}{\lambda _c}}\right),x=|𝐱|`$
$`{\displaystyle \frac{a_s}{\lambda _c}}=111Z^{1/3}={\displaystyle \frac{a_{s2}}{\lambda _c}}\mathrm{exp}\left(f{\displaystyle \frac{1}{2}}\right),𝜹={\displaystyle \frac{𝚫}{m}},`$ (3.11)
where $`K_0(z)`$ is the modified Bessel function. Integrating $`\mathrm{\Delta }V`$ over azimuthal angle of the vector $`𝐱`$ we obtain
$`\mathrm{\Delta }V(\mathit{\varrho },𝚫)=2Q\mathit{\varrho }^2{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}K_1^2(x)\left(1J_0(\beta x)+J_2(\beta x)\mathrm{cos}2\phi \right)x𝑑x`$
$`=2Q\mathit{\varrho }^2{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}\left[K_0^2(x)\mathrm{cos}2\phi +K_1^2(x)\right]\left(1J_0(\beta x)\right)x𝑑x`$
$`=Q\mathit{\varrho }^2\left[F_1\left({\displaystyle \frac{\beta }{2}}\right)\mathrm{cos}2\phi +F_2\left({\displaystyle \frac{\beta }{2}}\right)\right],\beta =\mathrm{\Delta }a_s`$ (3.12)
where $`\phi `$ is the angle between the vectors $`\mathit{\varrho }`$ and $`𝚫`$, $`J_n(z)`$ is the Bessel function, the functions $`F_1(z)`$ and $`F_2(z)`$ are defined in Eq.(2.44). Some details of calculation are given after Eq.(2.44).
Summing $`\mathrm{\Delta }V(\mathit{\varrho },𝚫)`$ (3.12) and $`V(\mathit{\varrho })`$ (3.9) we get
$$V(\mathit{\varrho },𝚫)=Q\mathit{\varrho }^2\left(L_1+\mathrm{ln}\frac{4}{\mathit{\varrho }^2}2CF_1\left(\frac{\beta }{2}\right)\mathrm{cos}2\phi F_2\left(\frac{\beta }{2}\right)\right)$$
(3.13)
Starting from Eq.(3.6) and advancing as in Sec.II (Eqs.(2.2)-(2.12)) of we find
$`M_{++}={\displaystyle \frac{2i\alpha m^2\omega }{n_a}}{\displaystyle \frac{d\epsilon }{\epsilon \epsilon ^{}}0|s_1\left(G_s^1G_0^1\right)+s_2𝐩\left(G_s^1G_0^1\right)𝐩|0},`$
$`M_+={\displaystyle \frac{2i\alpha m^2\omega }{n_a}}{\displaystyle \frac{d\epsilon }{\epsilon \epsilon ^{}}0|s_3\left(𝐞_{}^{}𝐩\right)\left(G_s^1G_0^1\right)\left(𝐞_+𝐩\right)|0},`$ (3.14)
where
$$s_1=1,s_2=\frac{\epsilon ^2+\epsilon ^2}{\omega ^2},s_3=\frac{2\epsilon \epsilon ^{}}{\omega ^2},G_s=+1,G_0=𝐩^2+1,$$
(3.15)
$``$ is defined in Eq.(3.7). In derivation it is convenient to use properties of $`𝐞_\lambda `$ given in (2.34). Note, that coefficients $`s_1,s_2`$ coincide with the same coefficients in the theory of the LPM effect for pair creation by a photon . As in (Eq.(2.17)) and (Eq.(2.5)) we split the potential $`V(\mathit{\varrho },𝚫)`$ as
$`V(\mathit{\varrho },𝚫)=V_c(\mathit{\varrho },𝚫)+v(\mathit{\varrho },𝚫),V_c(\mathit{\varrho },𝚫)=q_p\mathit{\varrho }^2,q_p=QL_s,`$
$`L_s(\varrho _c)=\mathrm{ln}{\displaystyle \frac{a_{s2}^2}{\lambda _c^2\varrho _c^2}}F_2\left({\displaystyle \frac{\beta }{2}}\right),`$
$`v(\mathit{\varrho },𝚫)={\displaystyle \frac{q_p\mathit{\varrho }^2}{L_s}}\left(2C+\mathrm{ln}{\displaystyle \frac{\mathit{\varrho }^2}{4\varrho _c^2}}+F_1\left({\displaystyle \frac{\beta }{2}}\right)\mathrm{cos}2\phi \right)`$ (3.16)
In accordance with such division of the potential we will write the photon scattering amplitudes in the form
$$M_{++}=M_{++}^c+M_{++}^{(1)},M_+=M_+^c+M_+^{(1)},$$
(3.17)
where $`M_{++}^c`$ is the contribution of the potential $`V_c(\mathit{\varrho },𝚫)`$ and $`M_{++}^{(1)}`$ is the first correction to the scattering amplitude due to the potential $`v(\mathit{\varrho },𝚫)`$ (see , Eqs.(2.28) and (2.33), and Eqs.(2.10) and (2.15)). Since dependence of the potential $`V(\mathit{\varrho },𝚫)`$ on $`\mathit{\varrho }`$ is the same as in and the expression for $`M_{++}`$ (3.14) formally coincides (up to external factor) with Eq.(2.2) in the result for $`M_{++}`$ can be taken from (Eqs.(2.10) and (2.15)) with substitution $`qq_p,L_cL_s`$. The combination of photon polarization vectors entering in the amplitude $`M_+^c`$ (3.14) gives in evaluation the expression $`\left(𝐞_{}^{}\mathbf{}\right)\left(𝐞_+\mathit{\varrho }\right)=0`$. Substituting this combination in calculation of the first correction $`M_+^{(1)}`$ (see Eqs.(2.31) and (2.32) of ) we obtain expression of type $`\left(𝐞_{}^{}\mathit{\varrho }\right)\left(𝐞_+\mathit{\varrho }\right)=\mathit{\varrho }^2\mathrm{exp}(2i\phi )`$. So in integration over the azimuthal angle $`\phi `$ only the term with $`\mathrm{cos}2\phi `$ survives in formula for $`v(\mathit{\varrho },𝚫)`$. As a result we get for the main term
$`M_{++}^c={\displaystyle \frac{\alpha m^2\omega }{2\pi n_a}}{\displaystyle \frac{d\epsilon }{\epsilon \epsilon ^{}}\mathrm{\Phi }_s(\nu _s)},M_+^c=0,`$
$`\mathrm{\Phi }_s(\nu _s)=\nu _s{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑te^{it}\left[s_1\left({\displaystyle \frac{1}{\mathrm{sinh}z}}{\displaystyle \frac{1}{z}}\right)i\nu _ss_2\left({\displaystyle \frac{1}{\mathrm{sinh}^2z}}{\displaystyle \frac{1}{z^2}}\right)\right]`$
$`=s_1\left(\mathrm{ln}p\psi \left(p+{\displaystyle \frac{1}{2}}\right)\right)+s_2\left(\psi (p)\mathrm{ln}p+{\displaystyle \frac{1}{2p}}\right),`$ (3.18)
where $`\nu _s=2\sqrt{iq_s},z=\nu _st,p=i/(2\nu _s),\psi (x)`$ is the logarithmic derivative of the gamma function.
The first corrections to the amplitudes are defined by
$`M_{++}^{(1)}={\displaystyle \frac{\alpha m^2\omega }{4\pi n_aL_s}}{\displaystyle \frac{d\epsilon }{\epsilon \epsilon ^{}}F_s(\nu _s)},M_+^{(1)}={\displaystyle \frac{\alpha m^2\omega }{4\pi n_aL_s}}F_1\left({\displaystyle \frac{\beta }{2}}\right){\displaystyle \frac{d\epsilon }{\epsilon \epsilon ^{}}F_3(\nu _s)},`$
$`F_s(\nu _s)={\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dze^{it}}{\mathrm{sinh}^2z}}\left[s_1f_1(z)2is_2f_2(z)\right],f_2(z)={\displaystyle \frac{\nu _s}{\mathrm{sinh}z}}\left(f_1(z){\displaystyle \frac{g(z)}{2}}\right),`$
$`f_1(z)=\left(\mathrm{ln}\varrho _c^2+\mathrm{ln}{\displaystyle \frac{\nu }{i}}\mathrm{ln}\mathrm{sinh}zC\right)g(z)2\mathrm{cosh}zG(z),`$
$`F_3(\nu _s)=i\nu _ss_3{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dze^{it}}{\mathrm{sinh}^3z}}g(z)={\displaystyle \frac{s_3}{2}}\left[1+{\displaystyle \frac{1}{2p}}p\zeta (2,p)\right]`$
$`g(z)=z\mathrm{cosh}z\mathrm{sinh}z,t=t_1+t_2,z=\nu _st,p={\displaystyle \frac{i}{2\nu _s}},`$
$`G(z)={\displaystyle \underset{0}{\overset{z}{}}}(1y\mathrm{coth}y)𝑑y`$
$`=z{\displaystyle \frac{z^2}{2}}{\displaystyle \frac{\pi ^2}{12}}z\mathrm{ln}\left(1e^{2z}\right)+{\displaystyle \frac{1}{2}}\mathrm{Li}_2\left(e^{2z}\right),`$ (3.19)
here $`\mathrm{Li}_2\left(x\right)`$ is the Euler dilogarithm, $`\zeta (s,a)`$ is the generalized Riemann zeta function. Use of the given representations of functions $`F_3(\nu _s)`$ and $`G(z)`$ simplifies the numerical calculation.
In the region of the weak effect of scattering ($`|\nu _s|1,\varrho _c=1`$) the interval $`z1`$ contributes into the integrals (3.18), (3.19). In this region
$`\mathrm{\Phi }_s(\nu )s_1\left({\displaystyle \frac{\nu ^2}{6}}+{\displaystyle \frac{7\nu ^4}{60}}+{\displaystyle \frac{31\nu ^6}{126}}\right)+s_2\left({\displaystyle \frac{\nu ^2}{3}}+{\displaystyle \frac{2\nu ^4}{15}}+{\displaystyle \frac{16\nu ^6}{63}}\right),`$
$`F_s(\nu ){\displaystyle \frac{\nu ^2}{9}}\left(s_2s_1\right),F_3(\nu ){\displaystyle \frac{\nu ^2}{3}}\left(1+{\displaystyle \frac{4}{5}}\nu ^2+{\displaystyle \frac{16}{7}}\nu ^4\right)s_3.`$ (3.20)
Substituting these expressions into (3.18), (3.19) and integrating over $`\epsilon `$ we obtain
$`M_{++}=M_{++}^c+M_{++}^{(1)}=i{\displaystyle \frac{14Z^2\alpha ^3\omega }{9m^2}}[L_{s1}(1+i{\displaystyle \frac{59\omega }{175\omega _e}}{\displaystyle \frac{L_{s1}}{L_1}}`$
$`{\displaystyle \frac{3312}{2401}}\left({\displaystyle \frac{\omega }{\omega _e}}{\displaystyle \frac{L_{s1}}{L_1}}\right)^2){\displaystyle \frac{1}{21}}],`$
$`M_+=M_+^{(1)}=i{\displaystyle \frac{2Z^2\alpha ^3\omega }{9m^2}}F_1\left({\displaystyle \frac{\beta }{2}}\right)\left(1+i{\displaystyle \frac{16\omega }{25\omega _e}}{\displaystyle \frac{L_{s1}}{L_1}}{\displaystyle \frac{384}{245}}\left({\displaystyle \frac{\omega }{\omega _e}}{\displaystyle \frac{L_{s1}}{L_1}}\right)^2\right),`$ (3.21)
here
$`L_{s1}=\mathrm{ln}{\displaystyle \frac{a_{s2}^2}{\lambda _c^2}}F_2\left({\displaystyle \frac{\beta }{2}}\right),\omega _e={\displaystyle \frac{m}{2\pi Z^2\alpha ^2\lambda ^3n_aL_1}},`$
$`L_{s1}{\displaystyle \frac{1}{21}}=2\left[\mathrm{ln}{\displaystyle \frac{a_s}{\lambda _c}}{\displaystyle \frac{1}{2}}\left(F_2\left({\displaystyle \frac{\mathrm{\Delta }a_s}{2}}\right)+1\right)f(Z\alpha )+{\displaystyle \frac{41}{42}}\right],`$ (3.22)
$`L_1`$ is defined in (3.9). The characteristic energy $`\omega _e`$ encountered in analysis of influence of the multiple scattering on the probability of pair photoproduction , in gold $`\omega _e`$=10.5 TeV. The amplitudes (3.21) coincide with the formulas (2.49) and (2.57) if we neglect the terms $`\omega /\omega _e`$ and $`(\omega /\omega _e)^2`$. The terms $`\omega /\omega _e`$ define the real part of the scattering amplitudes while the terms $`(\omega /\omega _e)^2`$ are the corrections to the imaginary part.
If the parameter $`|\nu _s|>1`$ then the value of $`\varrho _c`$ is defined by the equations (see (3.16))
$`4Q\varrho _c^4L_s(\varrho _c)=1,L_sL_{s3}+{\displaystyle \frac{1}{2}}\mathrm{ln}{\displaystyle \frac{4\epsilon \epsilon ^{}}{\omega ^2}},L_{s3}=L_{s2}+{\displaystyle \frac{1}{2}}\mathrm{ln}L_{s2},`$
$`L_{s2}=L_{s1}+{\displaystyle \frac{1}{2}}\mathrm{ln}{\displaystyle \frac{\omega }{\omega _e}}=2\mathrm{ln}(a_s\mathrm{\Delta }_s)+1F_2\left({\displaystyle \frac{\mathrm{\Delta }a_s}{2}}\right)2f(Z\alpha ),`$
$`\mathrm{\Delta }_s^4=2\pi Z^2\alpha ^2\omega n_a,\nu _s^2=i{\displaystyle \frac{\omega }{\omega _e}}{\displaystyle \frac{4\epsilon \epsilon ^{}}{\omega ^2}}\left(L_{s3}+{\displaystyle \frac{1}{2}}\mathrm{ln}{\displaystyle \frac{4\epsilon \epsilon ^{}}{\omega ^2}}\right).`$ (3.23)
We consider now the region of the strong effect of scattering ($`|\nu _s|1`$). We restrict to the main terms of the decomposition only, then we can substitute the exponential $`\mathrm{exp}(it)`$ for 1 in integrand in Eqs.(3.18),(3.19). Performing the integration we find
$$\mathrm{\Phi }_s(\nu _s)is_2\nu _s,F_s(\nu _s)is_2\nu _s\left(\mathrm{ln}2Ci\frac{\pi }{4}\right),F_3(\nu _s)\frac{i}{2}s_3\nu _s.$$
(3.24)
Note, that the next terms of the decomposition can be obtained using the results of Appendix A of . Substituting these expressions into Eqs.(3.17)-(3.19) we obtain
$`M_{++}(i1){\displaystyle \frac{3\alpha m^2}{4\sqrt{2}n_a}}\sqrt{{\displaystyle \frac{\omega }{\omega _e}}{\displaystyle \frac{L_{s3}}{L_1}}}\left[1{\displaystyle \frac{1}{4L_{s3}}}\left(2C+{\displaystyle \frac{1}{3}}+i{\displaystyle \frac{\pi }{2}}\right)\right]`$
$`=(i1){\displaystyle \frac{3\pi Z^2\alpha ^3\omega }{2\sqrt{2}\mathrm{\Delta }_s^2}}\sqrt{L_{s3}}\left[1{\displaystyle \frac{1}{4L_{s3}}}\left(2C+{\displaystyle \frac{1}{3}}+i{\displaystyle \frac{\pi }{2}}\right)\right],`$
$`M_+(i1){\displaystyle \frac{\pi Z^2\alpha ^3\omega }{8\sqrt{2}\mathrm{\Delta }_s^2}}{\displaystyle \frac{1}{\sqrt{L_{s3}}}}F_1\left({\displaystyle \frac{\beta }{2}}\right).`$ (3.25)
It is seen from these equations that in the case of strong scattering the real and imaginary parts of the amplitudes are equal (if we neglect the term $`1/L_{s3}`$ in $`M_{++}`$). Moreover the amplitudes (3.25) don’t depend on the electron mass $`m`$. In place of it we have the value $`\mathrm{\Delta }_s`$. Notice that the asymptotic expansions (3.25) are valid for $`\mathrm{\Delta }\mathrm{\Delta }_s(\mathrm{\Delta }_s>m)`$. In the interval $`\mathrm{\Delta }_s\mathrm{\Delta }a_s^1`$ we have
$$F_22\mathrm{ln}(\mathrm{\Delta }a_s)1,F_1=1,L_{s2}=2\left[\mathrm{ln}\frac{\mathrm{\Delta }_s}{\mathrm{\Delta }}+1f(Z\alpha )\right].$$
(3.26)
## 4 Conclusion
In this paper we considered the influence of a medium on the process of the coherent photon scattering illustrated in Fig.1 and Fig.2, where Im $`M_{++}`$ and Re $`M_{++}`$ as well as Im $`M_+`$ and Re $`M_+`$ are given as a function of photon energy $`\omega `$ in gold. This influence is due to the multiple scattering of electron and positron of the virtual pair on the formation length of the process (see Eqs.(1.1) and (1.2))
$$l_f=\frac{\omega }{2(q_s^2+m^2+\mathrm{\Delta }^2)}.$$
(4.1)
In the region $`\mathrm{\Delta }^2q_s^2+m^2`$ this formation length is independent of $`\mathrm{\Delta }`$ and it value is coincide practically with the formation length of pair creation by a photon $`l_c`$ . There is some difference connected with the logarithmic dependence of $`\nu _s^2`$ value on $`\mathrm{\Delta }^2`$:
$$|\nu _s^2|=\frac{q_s^2}{m^2}=\frac{4\pi Z^2\alpha ^2}{m^2}n_al_f\underset{q_{min}^2}{\overset{q_{max}^2}{}}\frac{dq^2}{q^2}$$
(4.2)
where $`q_{max}^2=m^2+q_s^2`$ and $`q_{min}^2=\mathrm{\Delta }^2+a^2`$, $`a`$ is the screening radius of atom (1.2). This defines the weak (logarithmic) dependence of Im$`M_{++}`$ on $`\mathrm{\Delta }`$ in the region $`\mathrm{\Delta }<\sqrt{q_s^2+m^2}`$. This can be seen in Fig.1 where the curves 1 and 3 represent behavior of Im $`M_{++}`$ for $`\mathrm{\Delta }`$=0.4435 $`m`$ and $`\mathrm{\Delta }`$=0.0387 $`m`$ respectively. For lower value of $`\mathrm{\Delta }`$ the minimal momentum transfer $`q_{min}`$ (1.2) diminishes thereby the interval of contributing the multiple scattering angles increases, so the multiple scattering affects the photon scattering amplitude at a lower energy (and smaller formation length). Because of this the curve 3 is shifted to the left respect the curve 1. Note that the curve 3 ($`\mathrm{\Delta }^1a_{s2}`$) is very similar to the curve 2 in Fig.2 of which represents the behavior of the probability of pair photoproduction in gold vs photon energy.
The new property of influence of a medium is the appearance of the real part of the coherent photon scattering amplitudes at high energy $`\omega `$. In the region $`\omega \omega _e`$ the value of Re$`M`$ is small accordingly (3.21). In the asymptotic region $`\omega \omega _e`$ we have -Re $`M=`$ Im $`M`$ according to (3.25). This property is seen clearly in Figs.1,2. So the value of -Re $`M`$ is small at low and very high energies of photon. At intermediate energies the value of -Re $`M`$ have the maximum at $`\omega `$220 TeV for $`\mathrm{\Delta }=`$0.4435 $`m`$ and at $`\omega `$80 TeV for $`\mathrm{\Delta }=`$0.0387 $`m`$. In Fig.2 the same curves are shown for amplitude $`M_+`$. These curves are very similar to curves in Fig.1. The curves in both figures are normalized to imaginary part of the corresponding amplitude in the absence of the multiple scattering. The ratio these imaginary parts one can find from Eqs.(3.21), (3.22):
$$r=\frac{\mathrm{Im}M_+}{\mathrm{Im}M_{++}}=\frac{1}{14}F_1\left(\frac{\mathrm{\Delta }a_s}{2}\right)\left[\mathrm{ln}\frac{a_s}{\lambda _c}\frac{1}{2}\left[F_2\left(\frac{\mathrm{\Delta }a_s}{2}\right)+1\right]f(Z\alpha )+1\right]^1.$$
(4.3)
This ratio is $`r=0.04435`$ for $`\mathrm{\Delta }`$=0.4435 $`m`$ and $`r=0.003018`$ for $`\mathrm{\Delta }`$=0.0387 $`m`$.
In numerical calculation of amplitude $`M_{++}`$ we neglect corrections of the order $`1/L_s`$. These corrections are quite small, e.g. for $`\mathrm{\Delta }`$=0.0387 $`m`$ they are of the order of a few percent.
We estimate now the integral (over $`𝚫`$) cross section of the coherent photon scattering at $`\omega \omega _e`$. From previous section (see Eqs.(3.25)-(3.26)) we have for module squared of the amplitude at $`\mathrm{\Delta }<\mathrm{\Delta }_s`$ (see also Introduction)
$$|M|^2\frac{9\pi ^2Z^4\alpha ^6\omega ^2}{\mathrm{\Delta }_s^4}\left(\mathrm{ln}\frac{\mathrm{\Delta }_s}{\mathrm{\Delta }}+1\right),$$
(4.4)
where value $`\mathrm{\Delta }_s`$ is defined in (3.23). It is seen from this formula that in the considered region the differential probability of the coherent photon scattering depends weakly (logarithmically) on medium density $`n_a`$.
$$dW_{ph}(\mathrm{\Delta }<\mathrm{\Delta }_s)=\frac{1}{16\pi ^2}|M|^2n_ad\mathrm{\Omega }=\frac{|M|^2n_a}{16\pi ^2\omega ^2}d𝚫.$$
(4.5)
In the interval $`\mathrm{\Delta }\mathrm{\Delta }_s`$ the influence of the multiple scattering on the coherent scattering process is rather weak because the probability of transfer to a medium of the momentum $`|𝚫|`$ on the formation length $`l_f=\omega /(2\mathrm{\Delta }^2)`$ appears to be small:
$$W_s(\mathrm{\Delta })=\frac{4\pi Z^2\alpha ^2}{\mathrm{\Delta }^2}n_a\frac{\omega }{2\mathrm{\Delta }^2}=\frac{\mathrm{\Delta }_s^4}{\mathrm{\Delta }^4}1.$$
(4.6)
In this interval the amplitudes of the coherent photon scattering behave as $`1/\mathrm{\Delta }^2`$ and it doesn’t contribute to the integral cross section. So the interval $`\mathrm{\Delta }\mathrm{\Delta }_s`$ contributes only. Taking into account Eqs.(4.4),(4.5) we have for estimate of the integral cross section of the coherent photon scattering
$$\sigma _{ph}(\omega \omega _e)\frac{\mathrm{\Delta }_s^2}{16\pi \omega ^2}|M(\mathrm{\Delta }\mathrm{\Delta }_s)|^2\frac{Z^4\alpha ^6}{\mathrm{\Delta }_s^2}=\frac{Z^3\alpha ^5}{\sqrt{2\pi \omega n_a}}.$$
(4.7)
Acknowledgments This work was supported in part by the Russian Fund of Basic Research under Grant 00-02-18007.
Figure captions
* Fig.1 The amplitude $`M_{++}`$ of the coherent photon scattering in gold under influence of the multiple scattering at the different momentum transfer to the photon $`\mathrm{\Delta }`$ in terms of the amplitude Im$`M_{++}`$ (2.57) calculated for the screened Coulomb potential.
+ Curve 1 is Im$`M_{++}`$ for $`\mathrm{\Delta }=`$0.4435 $`m`$.
+ Curve 2 is Re$`M_{++}`$ for $`\mathrm{\Delta }=`$0.4435 $`m`$.
+ Curve 3 is Im$`M_{++}`$ for $`\mathrm{\Delta }=`$0.0387 $`m`$.
+ Curve 4 is Re$`M_{++}`$ for $`\mathrm{\Delta }=`$0.0387 $`m`$.
* Fig.2 The same as in Fig.1 but for the amplitude $`M_+`$ in terms of the amplitude Im$`M_+`$ (2.57) calculated for the screened Coulomb potential.
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# An example concerning Ohtsuki’s invariant and the full 𝑆𝑂(3) quautum invariant The first author is supported partially by the National Science Foundation of P. R. China, the second author is supported partially by NSF grant DMS 9304580
## Abstract
Two lens spaces are given to show that Ohtsuki’s $`\tau `$ for rational homology spheres does not determine Kirby-Melvin’s $`\{\tau _r^{^{}},r\text{odd}3\}`$
By using partial Kirby-Melvin’s quautum $`SO(3)`$ invariants $`\{\tau _r^{^{}}(M),r\text{odd prime}>|H_1(M,Z)|\}`$ \[KM\], Ohtsuki \[O\] defined a topological invariant
$$\tau (M)=\underset{n=0}{\overset{\mathrm{}}{}}\lambda _n(M)(t1)^nQ[[t1]]$$
for rational homology 3-sphere $`M`$. R.Lawrence \[La\] conjectured that
$$\lambda _n(M)\{\begin{array}{cc}Z,\hfill & \text{if }|H_1(M,Z)|=1\hfill \\ Z[\frac{1}{2},\frac{1}{|H_1(M,Z)|}],\hfill & \text{if }|H_1(M,Z)|>1\hfill \end{array}$$
and if $`r`$ is an odd prime which does not divide $`|H_1(M,Z)|`$, then $`\{|H_1(M,Z)|\}_r\tau _r^{^{}}(M)`$ is the $`r`$adic limit of the series
$$\underset{n=0}{\overset{\mathrm{}}{}}\lambda _n^{}(M)h^n$$
where $`\{\}_r`$ stands for the Jacobi symbol, and $`h=e^{\frac{2\pi i}{r}}1`$.
Rozansky \[R\] has proved that this conjecture is true. So $`\tau (M)`$ and $`\{\tau _r^{^{}},r`$ odd prime not dividing $`|H_1(M,Z)|\}`$ determine each other.
A natural question arises: Does $`\tau `$ determine all $`\{\tau _r^{^{}},r\text{odd}3\}`$?
It was proved in \[Li\] that $`\tau _r^{^{}}(M)=\tau _r^{^{}}(M^{})`$ iff $`\xi _r(M,A)=\xi _r(M^{},A)`$ for $`r`$ odd $`3`$, where $`A`$ is any $`r`$th primitive root of unit. So the question is equivalent to: Does $`\tau (M)`$ determine all $`\xi _r(M,e_r)`$? Where, $`e_a`$ stands for $`e^{\frac{2\pi i}{a}}`$.
For lens space $`L(p,q)`$, all $`\xi _r(L(p,q),e_r)`$ has been obtained in \[LL1\] ( explicit formulas for $`\tau _r^{^{}}(L(p,q))`$ were given in \[LL2\]), that is: let $`r3`$ be odd and $`c=(p,r)`$ the common factor, then
$$(1)\xi _r(L(p,q),e_r)=\{\begin{array}{cc}\{p\}_re_r^{12s(q,p)}e_p^{r^{}(q+q^{})}\frac{e_r^{2p^{^{}}}e_r^{2p^{^{}}}}{e_r^2e_r^2},\hfill & \text{if }c=1\hfill \\ & \\ & \\ (1)^{\frac{r1}{2}\frac{c1}{2}}\{p/c\}_{r/c}\{q\}_ce_r^{12s(q,p)}\hfill & \\ e_{pc}^{(r/c)^{^{}}(q+q^{}\eta p^{}p)}e_{rc}^{2\eta (p/c)^{^{}}}\frac{ϵ(c)\sqrt{c}\eta }{e_r^2e_r^2},\hfill & \text{if }c>1\text{ , }cq^{}+\eta \hfill \\ & \\ 0,\hfill & \text{if }c>1\text{ and }c\text{/}q^{}\pm 1\hfill \end{array}$$
where $`\eta =1`$ or $`1`$, $`p^{}p+q^{}q=1`$ with $`0<q^{}<p`$, $`(p/c)^{^{}}p/c+(r/c)^{}r/c=1`$, $`s(p,q)`$ is the Dedekind sum, and
$$ϵ(c)=\{\begin{array}{cc}1,\hfill & \text{if }c1(mod4)\hfill \\ i,\hfill & \text{if }c1(mod4)\hfill \end{array}$$
Since
$$\tau (L(p,q))=t^{3s(q,p)}\frac{t^{\frac{1}{2p}}t^{\frac{1}{2p}}}{t^{\frac{1}{2}}t^{\frac{1}{2}}}$$
(\[O\] and \[LL2\]), $`\tau (L(p_1,q_1))=\tau (L(p_2,q_2))`$ iff $`p_1=p_2`$ and $`s(q_1,p_1)=s(q_2,p_2)`$. The following Theorem answers the question above.
Theorem. $`s(6,25)=s(11,25)`$, while $`\xi _r(L(25,11))=\xi _r(L(25,6))`$ if and only if $`(r,25)5`$.
Proof. We calculate $`s(q,p)`$ by the formula in \[H\]:
$$12s(q,p)=\underset{i=1}{\overset{n}{}}m_i+\frac{q+q^{}}{p}3n$$
if $`\frac{p}{q}=m_n\frac{1}{m_{n1}\mathrm{}{\displaystyle \frac{1}{m_2{\displaystyle \frac{1}{m_1}}}}}`$ with $`m_i2`$.
Now
$$\frac{25}{6}=5\frac{1}{2{\displaystyle \frac{1}{2{\displaystyle \frac{1}{2{\displaystyle \frac{1}{2{\displaystyle \frac{1}{2}}}}}}}}},\frac{25}{11}=3\frac{1}{2{\displaystyle \frac{1}{2{\displaystyle \frac{1}{3{\displaystyle \frac{1}{2}}}}}}}$$
$`11^{}=16,6^{}=21`$, so
$$12s(6,25)=12s(11,25)=3+\frac{27}{25}$$
$`c=(r,25)`$ can be only $`1,5`$ or $`25`$. If $`c=25`$, then $`c\text{/}6^{}\pm 1`$ and $`c\text{/}11^{}\pm 1`$, thus by (1) $`\xi _r(L(25,11),e_r)=\xi _r(L(25,6),e_r)=0`$. If $`c=1`$, it is easy to see from (1) that the two $`\xi _r`$ are equal.
Assume $`c=5`$, then $`c|6^{}1`$ and $`11^{}1`$. Now since $`q^{}q+p^{}p=1`$, we have
$$q+q^{}+(25)^{}\times 25=\{\begin{array}{cc}27+125,\hfill & \text{if }q=6\hfill \\ 27+175,\hfill & \text{if }q=11\hfill \end{array}$$
Therefore, by (1)
$$\frac{\xi _r(L(25,11),e_r)}{\xi _r(L(25,6),e_r)}=e_{125}^{(r/5)^{^{}}\times 50}$$
Since $`(p/c)^{^{}}p/c+(r/c)^{}r/c=1`$, we see that $`(r/5)^{^{}}`$ is prime to $`(p/5)=5`$. This shows that $`e_{125}^{(r/5)^{^{}}\times 50}1`$, and the theorem is proved.
References
\[RT\] Reshetikhin, N.Yu., Turaev, V.G.: Invariants of 3-manifolds via link polynomials and quantum groups, Invent. Math. 103 (1991), 547-597
\[KM\] Kirby, R., Melvin, P.: The 3-manifold invariants of Witten and Reshetikhin-Turaev for $`sl(2,C)`$, Invent. Math. 105 (1991), 473-545
\[Li\] Li, B.H., Relations among Chern-Simons-Witten-Jones invariants , Science in China, series A, 38 (1995), 129-146
\[LL1\] Li, B.H., Li, T.J.: Generelized Gaussian Sums and Chern-Simons- Witten-Jones invariants of Lens spaces, J. Knot theory and its Ramifications, vol.5 No. 2 (1996) 183-224
\[O\] Ohtsuki, T., A polynomial invariant of rational homology 3-spheres, Inv. Math. 123 (1996), 241-257
\[H\] Hickerson, D., Continued fraction and density results, J. Reine. Angew. Math. 290 (1977), 113-116
\[R\] Rozensky, L.: On p-adic properties of the Witten- Reshetikhin-Turaev invariants, math. QA/9806075, 12 Jun.1998
\[La\] Lawrence, R.: Asymptotic Expansions of Witten- Reshetikhin-Turaev invariants for some simple 3-manifolds, J. Math. Phys. 36 (1995) 6106-6129
\[W\] Witten, E: Quatum field theory and the Jones polynomial, Comm.Math. Phys., 121 (1989), 351-399
Author’s address
| Bang-He Li | Tian-Jun Li |
| --- | --- |
| Institute of Systems Science, | School of math. |
| Academia Sinica | IAS |
| Beijnig 100080 | Princeton NJ 08540 |
| P. R. China | U.S.A. |
| Libh@iss06.iss.ac.cn | Tjli@IAS.edu |
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# Influence of the 𝑈(1)_𝐴 Anomaly on the QCD Phase Transition
## Abstract
The $`SU(3)_r\times SU(3)_{\mathrm{}}`$ linear sigma model is used to study the chiral symmetry restoring phase transition of QCD at nonzero temperature. The line of second order phase transitions separating the first order and smooth crossover regions is located in the plane of the strange and nonstrange quark masses. It is found that if the $`U(1)_A`$ symmetry is explicitly broken by the $`U(1)_A`$ anomaly then there is a smooth crossover to the chirally symmetric phase for physical values of the quark masses. However, if the $`U(1)_A`$ anomaly is absent, the region of first order phase transitions is significantly enlarged and it is found that there is a phase transition for physical values of the quark masses provided that the $`\sigma `$ meson mass is at least 600 MeV. In both cases, the region of first order phase transitions in the quark mass plane is enlarged as the mass of the $`\sigma `$ meson is increased.
The ultimate goal of relativistic heavy ion experiments is to probe the phase diagram of Quantum Chromodynamics (QCD). General theoretical considerations indicate that at sufficiently high temperatures there should be a transition from ordinary hadronic matter to a chirally symmetric plasma of quarks and gluons . The order parameter for this phase transition is the quark-antiquark condensate. Results from lattice gauge theory predict the temperature of this transition to be about 150 MeV . The order of the phase transition, however, seems to depend very much on the number of quark flavors and their masses .
Classically, the matter part of the QCD Lagrangian with $`N_f`$ flavors is invariant under the symmetry group $`SU(N_f)_r\times SU(N_f)_{\mathrm{}}\times U(1)_A`$. The axial $`U(1)_A`$ symmetry is broken to $`Z(N_f)_A`$ by a nonvanishing topological susceptibility and the $`SU(N_f)_r\times SU(N_f)_{\mathrm{}}`$ symmetry is spontaneously broken to the diagonal group of vector transformations, $`SU(N_f)_{r+\mathrm{}}=SU(N_f)_V`$, by a nonvanishing expectation value for the quark-antiquark condensate. The $`SU(N_f)_r\times SU(N_f)_{\mathrm{}}\times U(1)_A`$ group is also explicitly broken by the effects of nonzero quark masses. It was shown by Pisarski and Wilczek that for three or more massless flavors, the phase transition for the restoration of the $`SU(N_f)_r\times SU(N_f)_{\mathrm{}}`$ symmetry is first order, while for two massless flavors the phase transition is second order.
The $`U(1)_A`$ symmetry may also be restored, if only partially, since instanton effects are Debye screened at high temperatures . There are now two possibilities: either the $`U(1)_A`$ symmetry is restored at a temperature much greater than the $`SU(N_f)_r\times SU(N_f)_{\mathrm{}}`$ symmetry or the two symmetries are restored at (approximately) the same temperature . Recent lattice gauge theory computations have demonstrated a rapid decrease in the topological susceptibility at $`T_c`$ and random matrix models also indicate that the two symmetries are restored simulaneously . Perhaps more dramatically, it was also shown that the topological susceptibility vanishes at $`T_c`$ in the large-$`N_c`$ limit . On the other hand, the fate of the $`U(1)_A`$ anomaly in nature is not completely clear since instanton liquid model calculations indicate that the topological susceptibility is essentially unchanged at $`T_c`$ . Additionally, other lattice computations which measure the chiral susceptibility find that the $`U(1)_A`$ symmetry restoration is at or below the 15% level .
Unlike the idealized massless quark limit, there are no general theoretical arguments which require that a phase transition exists for massive quarks. Indeed, some lattice simulations indicate that for physical quark masses, no phase transition occurs . The general consensus from lattice computations is that in the plane of light quark masses (see Fig. 1) there is a first order region bounded by a line of second order transitions. Outside this region, there is no phase transition, but rather a crossover characterized by a rapid but smooth and continuous decrease of the quark-antiquark condensate. Given the present difficulties with performing lattice computations with realistic quark masses and a large number of sites, it is useful to complement the present lattice results with effective models that capture some of the relevant dynamics of QCD. Some work for three flavors has been done in this direction . In these works, the $`SU(3)_r\times SU(3)_{\mathrm{}}`$ linear sigma model was used to study the order of the chiral symmetry restoring phase transition as a function of the current quark masses with the ratio of the up-down to strange quark masses held fixed. In Refs. , a loop-expansion is used to compute the effective potential and in Ref. a mean-field analysis of this model was presented. The effects of the restoration of the $`U(1)_A`$ symmetry on the spectrum of hadronic observables in heavy ion collisions was addressed within the context of this model in Ref. .
In this paper, I present results concerning the order of the chiral symmetry restoring phase transition as a function of the current quark masses using the $`SU(3)_r\times SU(3)_{\mathrm{}}`$ linear sigma model without fixing the ratio of the masses. In addition, the effects of the $`U(1)_A`$ anomaly on the order of the QCD phase transition are investigated. Here, the Cornwall-Jackiw-Tomboulis (CJT) formalism is used to derive gap equations for the condensates and the tadpole-resummed scalar and pseudoscalar nonet meson masses at nonzero temperature. The derivation of and the solutions to this set of equations in a variety of limits has been presented elsewhere . The results agree qualitatively with earlier studies on a lattice and with other studies using the $`SU(3)_r\times SU(3)_{\mathrm{}}`$ linear sigma model . In the presence of an explicit $`U(1)_A`$ symmetry breaking term, I find that for physical values of the strange and nonstrange current quark masses, there is no phase transition but rather a smooth crossover. For smaller values of the masses, the phase transition is first order with a line of second order transitions separating the first order and the crossover regions. In the absence of an explicit $`U(1)_A`$ symmetry breaking term, the region of phase transitions is greatly enlarged. In particular, if the $`\sigma `$ meson mass is greater than 600 MeV, then the transition is driven to first order for physical values of the quark masses. In both cases, the region of first order phase transitions is enlarged as the mass of the $`\sigma `$ meson is increased.
The most general renormalizable theory compatible with the flavor symmetries of QCD is the $`SU(3)_r\times SU(3)_{\mathrm{}}`$ linear sigma model. While this model cannot account for the full dynamics of QCD, on the line of second order phase transitions the only relevant dynamics are determined by the symmetries of the theory. So, in the vicinity of this line, the use of the $`SU(3)_r\times SU(3)_{\mathrm{}}`$ linear sigma model is appropriate. Its Lagrangian is given by
$`(\mathrm{\Phi })`$ $`=`$ $`\mathrm{Tr}\left(_\mu \mathrm{\Phi }^{}^\mu \mathrm{\Phi }m^2\mathrm{\Phi }^{}\mathrm{\Phi }\right)\lambda _1\left[\mathrm{Tr}\left(\mathrm{\Phi }^{}\mathrm{\Phi }\right)\right]^2`$ (1)
$``$ $`\lambda _2\mathrm{Tr}\left(\mathrm{\Phi }^{}\mathrm{\Phi }\right)^2+c\left[\mathrm{Det}\left(\mathrm{\Phi }\right)+\mathrm{Det}\left(\mathrm{\Phi }^{}\right)\right]`$ (2)
$`+`$ $`\mathrm{Tr}\left[H(\mathrm{\Phi }+\mathrm{\Phi }^{})\right].`$ (3)
Here, $`\mathrm{\Phi }`$ is a $`U(3)`$ matrix defined by $`\mathrm{\Phi }=T_a(\sigma _a+i\pi _a)`$. The $`T_a=\widehat{\lambda }_a/2`$ are the generators of $`U(3)`$ where $`\widehat{\lambda }_a`$ are the Gell-Mann matrices with $`\lambda _0=\sqrt{2/3}I`$. The $`T_a`$ are normalized such that $`\mathrm{Tr}(T_aT_b)=\delta _{ab}/2`$.
The parameters of the Lagrangian are the bare mass $`m`$, a background matrix field $`H=h_0T_0+h_8T_8`$, a cubic coupling $`c`$ and two quartic couplings, $`\lambda _1`$ and $`\lambda _2`$. The various patterns of symmetry breaking and the parameterizations of the coupling constants for this Lagrangian were studied in and will only be briefly reviewed here. For $`\underset{¯}{H}=0`$, $`c=0`$ and $`m^2>0`$, the Lagrangian has a global $`SU(3)_r\times SU(3)_{\mathrm{}}\times U(1)_A`$ symmetry. The effects of the $`U(1)_A`$ symmetry breaking by a nonvanishing topological susceptibility (i.e. the presence of instantons in the QCD vacuum) are included by setting $`c0`$ which reduces the symmetry to $`SU(3)_r\times SU(3)_{\mathrm{}}`$. For nonzero $`H`$, chiral symmetry is explicitly broken.
I assume that there are nonzero vacuum expectation values for the $`\sigma _0`$ and $`\sigma _8`$ fields which I denote by $`\overline{\sigma }_0`$ and $`\overline{\sigma }_8`$. After shifting these fields by their expectation values and following , the Lagrangian can be rewritten as:
$``$ $`=`$ $`{\displaystyle \frac{1}{2}}[_\mu \sigma _a^\mu \sigma _a+_\mu \pi _a^\mu \pi _a(m_S^2)_{ab}\sigma _a\sigma _b`$ (4)
$``$ $`(m_P^2)_{ab}\pi _a\pi _b]+(𝒢_{abc}{\displaystyle \frac{4}{3}}_{abcd}\overline{\sigma }_d)\sigma _a\sigma _b\sigma _c`$ (5)
$``$ $`3(𝒢_{abc}+{\displaystyle \frac{4}{3}}_{abcd}\overline{\sigma }_d)\pi _a\pi _b\sigma _c2_{abcd}\sigma _a\sigma _b\pi _c\pi _d`$ (6)
$``$ $`{\displaystyle \frac{1}{3}}_{abcd}(\sigma _a\sigma _b\sigma _c\sigma _d+\pi _a\pi _b\pi _c\pi _d)h_a\overline{\sigma }_a,`$ (7)
where
$`𝒢_{abc}`$ $`=`$ $`{\displaystyle \frac{c}{6}}[d_{abc}{\displaystyle \frac{3}{2}}(\delta _{a0}d_{0bc}+\delta _{b0}d_{a0c}+\delta _{c0}d_{ab0})`$ (9)
$`+`$ $`{\displaystyle \frac{9}{2}}d_{000}\delta _{a0}\delta _{b0}\delta _{c0}],`$ (10)
$`_{abcd}`$ $`=`$ $`{\displaystyle \frac{\lambda _1}{4}}\left(\delta _{ab}\delta _{cd}+\delta _{ad}\delta _{bc}+\delta _{ac}\delta _{bd}\right)`$ (11)
$`+`$ $`{\displaystyle \frac{\lambda _2}{8}}\left(d_{abn}d_{ncd}+d_{adn}d_{nbc}+d_{acn}d_{nbd}\right),`$ (12)
$`_{abcd}`$ $`=`$ $`{\displaystyle \frac{\lambda _1}{4}}\delta _{ab}\delta _{cd}+{\displaystyle \frac{\lambda _2}{8}}(d_{abn}d_{ncd}+f_{acn}f_{nbd}`$ (13)
$`+`$ $`f_{bcn}f_{nad}),`$ (14)
$`(m_S^2)_{ab}`$ $`=`$ $`m^2\delta _{ab}6𝒢_{abc}\overline{\sigma }_c+4_{abcd}\overline{\sigma }_c\overline{\sigma }_d,`$ (15)
$`(m_P^2)_{ab}`$ $`=`$ $`m^2\delta _{ab}+6𝒢_{abc}\overline{\sigma }_c+4_{abcd}\overline{\sigma }_c\overline{\sigma }_d.`$ (16)
Here the summation runs over the index $`n`$ only and $`d_{abc}`$ and $`f_{abc}`$ are the symmetric and antisymmetric structure constants, respectively, of $`U(3)`$.
The $`\sigma _a`$ fields are members of the scalar ($`J^\pi =0^+`$) nonet and the $`\pi _a`$ fields are members of the pseudoscalar ($`J^\pi =0^{}`$) nonet. The $`\pi _{1,2,3}`$ are the pions, the $`\pi _{4,5,6,7}`$ are the kaons and the $`\pi _0`$ and the $`\pi _8`$ are admixtures of the $`\eta `$ and the $`\eta ^{}`$ with mixing angle $`\theta _\mathrm{P}`$. The situation with the scalar nonet is not as clear and still somewhat controversial . The $`\sigma _0`$ and the $`\sigma _8`$ are admixtures of the $`\sigma `$ and the $`f_0(1370)`$ with mixing angle $`\theta _\mathrm{S}`$. The $`\sigma _{1,2,3}`$ are identified with the $`a_0(980)`$ and the $`\sigma _{4,5,6,7}`$ with the $`\kappa `$ meson.
The explicit symmetry breaking terms can be determined (see, for instance, Ref. ), to be
$`h_0`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{6}}}(m_\pi ^2f_\pi +2m_K^2f_K)`$ (18)
$`h_8`$ $`=`$ $`{\displaystyle \frac{2}{\sqrt{3}}}(m_\pi ^2f_\pi m_K^2f_K).`$ (19)
The gap equations (Schwinger–Dyson equations) derived from the CJT effective potential in the tadpole-resummed approximation, or Hartree approximation, are found to be
$`(𝒮_{ab}`$ $`(k)`$ $`)^1=k^2+m^2\delta _{ab}6𝒢_{abc}\overline{\sigma }_c+4_{abcd}\overline{\sigma }_c\overline{\sigma }_d`$ (20)
$`+`$ $`4_{abcd}{\displaystyle _k}𝒮_{cd}(k)+4_{abcd}{\displaystyle _k}𝒫_{cd}(k),`$ (21)
$`(𝒫_{ab}`$ $`(k)`$ $`)^1=k^2+m^2\delta _{ab}+6𝒢_{abc}\overline{\sigma }_c+4_{abcd}\overline{\sigma }_c\overline{\sigma }_d`$ (22)
$`+`$ $`4_{abcd}{\displaystyle _k}𝒮_{cd}(k)+4_{abcd}{\displaystyle _k}𝒫_{cd}(k)`$ (23)
$`h_a`$ $`=`$ $`m^2\overline{\sigma }_a3𝒢_{abc}\overline{\sigma }_b\overline{\sigma }_c3𝒢_{abc}{\displaystyle _k}𝒮_{cb}(k)`$ (24)
$`+`$ $`3𝒢_{abc}{\displaystyle _k}𝒫_{cb}(k)+{\displaystyle \frac{4}{3}}_{abcd}\overline{\sigma }_d\overline{\sigma }_b\overline{\sigma }_c`$ (25)
$`+`$ $`4_{abcd}\overline{\sigma }_d{\displaystyle _k}𝒮_{cb}(k)+4_{bcad}\overline{\sigma }_d{\displaystyle _k}𝒫_{cb}(k),`$ (26)
where in the last equation, $`a=0,8`$ and $`𝒮_{ab}(k)`$ ($`𝒫_{ab}(k)`$) are the Green’s functions for the scalar (pseudoscalar) mesons. $`𝒮_{08}(k)`$ and $`𝒫_{08}(k)`$, however, are nonzero on account of the mixing between the singlet and the octet states. All other non–diagonal entries are identically zero. As such, it is necessary to rotate these Green’s functions into the mass eigenbasis since only physical fluctuations can contribute to the masses:
$`U_{ia}^\mathrm{S}𝒮_{ab}(k)U_{jb}^\mathrm{S}`$ $`=`$ $`\stackrel{~}{𝒮}_i(k)\delta _{ij}`$ (28)
$`U_{ia}^\mathrm{P}𝒫_{ab}(k)U_{jb}^\mathrm{P}`$ $`=`$ $`\stackrel{~}{𝒫}_i(k)\delta _{ij},`$ (29)
where $`U_{ia}^\pi =\delta _{ia}`$ for $`i,a0,8`$ and where $`U_{ia}^\pi `$ is given by an $`O(2)`$ rotation by $`\theta _\mathrm{P}`$ in the 0–8 block. The definition for $`U_{ia}^\sigma `$ is similarly given with $`\theta _\mathrm{P}\theta _\mathrm{S}`$. The thermal integral arising from tadpole diagrams is
$`{\displaystyle _k}\stackrel{~}{𝒮}_i(k)={\displaystyle \frac{d^3𝐤}{(2\pi )^3}\frac{1}{ϵ_𝐤[(\stackrel{~}{M}_S^2)_i]}\left(\mathrm{exp}\left\{\frac{ϵ_𝐤[(\stackrel{~}{M}_S^2)_i]}{T}\right\}1\right)^1}`$ (30)
and similarly for the pseudoscalar tadpole integrals, $`_k𝒫_{cd}(k)`$. Here, $`ϵ_𝐤[(\stackrel{~}{M}_S^2)_i]=(𝐤^2+(\stackrel{~}{M}_S^2)_i)^{1/2}`$ is the relativistic energy of the $`i`$th scalar quasiparticle with momentum $`𝐤`$. I have neglected the vacuum contribution arising from the loop integrals. Implementing a systematic renormalization scheme is difficult but possible in this approximation (see ). The results, however, are not significantly altered.
Since in the Hartree approximation, the gap equations do not have an explicit momentum dependence, we can assume that $`(𝒮_{ab}(k))^1=k^2+M_S^2`$ where $`M_S`$ depends on temperature but not momentum, and similarly for $`(𝒫_{ab}(k))^1`$. Equations (26) and (Influence of the $`U(1)_A`$ Anomaly on the QCD Phase Transition) are then fixed point equations and can be numerically solved simultaneously as a function of temperature for $`M_\pi `$, $`M_K`$, $`M_\eta `$, $`M_\eta ^{}`$, $`M_\sigma `$, $`M_\kappa `$, $`M_{f_0}`$, $`M_{a_0}`$, $`\overline{\sigma }_0`$, $`\overline{\sigma }_8`$, $`\theta _\mathrm{P}`$ and $`\theta _\mathrm{S}`$. The numerical solutions for a variety of parameters are given in Ref. .
The condensate and mass gap equations are solved with fixed $`m`$, $`c`$, $`\lambda _1`$ and $`\lambda _2`$, while varying the background fields, $`h_0`$ and $`h_8`$. The determination of the coupling constants is detailed in Ref. . For $`c0`$, the four couplings in the Lagrangian are fitted to yield the physical tree-level masses of the pion, kaon, $`\sigma `$, $`\eta `$ and $`\eta ^{}`$, while for $`c=0`$, the remaining three couplings are determined from the physical tree-level masses of the pion, kaon, $`\sigma `$ and $`\eta `$. The background fields are proportional to the current quark masses: $`m_{\mathrm{up}}=m_{\mathrm{down}}=a(h_0+h_8/\sqrt{2})`$, $`m_{\mathrm{strange}}=b(h_0\sqrt{2}h_8)`$. For simplicity, I assume temperature independent proportionality constants, $`a`$ and $`b`$. Requiring that $`m_\pi =138`$ MeV, $`m_K=496`$ MeV, $`m_{\mathrm{up}}=m_{\mathrm{down}}=10`$ MeV and $`m_{\mathrm{strange}}=150`$ MeV, gives $`a=4.64\times 10^6[\mathrm{MeV}]^2`$ and $`b=2.27\times 10^6[\mathrm{MeV}]^2`$.
To determine the order of the phase transition, I examined the continuity of the order parameters as a function of temperature. For a first order transition, the condensates are multivalued functions of temperature in the vicinity of the phase transition. For a smooth crossover, the condensates are smooth singlevalued functions of temperature and always nonzero. This behavior is demonstrated in Fig. 2. Only the nonstrange condensate is shown since both condensates exhibit qualitatively the same behavior.
The numerical results are plotted in Fig. 3. For $`c0`$, these results agree with those of lattice groups . For $`m_\sigma =1000`$ MeV, the authors of Ref. report that the ratio of the critical current up-down quark mass to the physical up-down quark mass for $`m_{\mathrm{strange}}/m_{\mathrm{up},\mathrm{down}}=32`$ to be $`0.01`$. The corresponding value found in the present work is $`0.20`$. The larger value found here is most likely due to the inclusion of thermal fluctuations from the scalar and pseudoscalar nonets. For $`m_{\mathrm{up}}=m_{\mathrm{down}}=0`$, the transition is first order from zero strange quark mass to some critical strange quark mass. For $`m_\sigma =800`$ MeV, this critical strange quark mass is about 16 MeV, while for $`m_\sigma =900`$ MeV, it is 260 MeV.
For $`c=0`$, however, the results are dramatically different. In particular, the line of second order transitions does not seem to approach the strange quark mass axis. For $`m_\sigma =600`$ MeV, the physical point, $`m_{\mathrm{up}}m_{\mathrm{down}}10`$ MeV and $`m_{\mathrm{strange}}150`$ MeV, is just outside the first order region. For larger values of the $`\sigma `$ meson mass, the physical point is well within the first order region. The results also seem to indicate that for $`c=0`$ there is a first order phase transition for three flavors provided only that one of the flavors is sufficiently heavier than the other two flavors. The departure of the second order phase transition line from the strange quark mass axis was also predicted using arguments from large-$`N_c`$ chiral perturbation theory in Ref. .
At this point, it should be mentioned that the Hartree approximation sometimes predicts a first order transition when the transition is actually second order. For example, renormalization group arguments predict a second order phase transition for the massless limit of the $`O(4)`$ linear sigma model , while the Hartree approximation predicts a first order transition (see, for instance, Ref. ). This is not a problem in the low quark mass region since the transition is expected to be first order. The location of the second order line should not be significantly affected.
Additionally, the cubic and quartic couplings are fixed and temperature independent. The running of the couplings with temperature should be at most logarithmic, while the integrals arising from the tadpole diagrams depend quadratically on the temperature. So, it is reasonable that the running of the couplings does not qualitatively alter these results. On the other hand, if the Coleman-Weinberg mechanism is strongly operative, some portion of the crossover region may actually be driven to first order .
Acknowledgements
I especially want to thank Dirk H. Rischke, Robert D. Pisarski and Jürgen Schaffner–Bielich for many valuable discussions. I also want to thank the Nuclear Theory Group at Brookhaven National Laboratory for their generous support and hospitality while this work was completed. I am grateful to Mark Alford, Tom Blum, Eduardo Fraga, Robert Harlander, Alex Krasnitz, and Ove Scavenius for useful discussions. I am supported by the Director, Office of Energy Research, Division of Nuclear Physics of the Office of High Energy and Nuclear Physics of the U.S. Department of Energy under Contract No. DE-AC02-98CH10886.
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# 𝜌-𝜔 mixing in asymmetric nuclear matter via QCD sum rule approach
## I Introduction
Changes of hadronic properties in hot and dense nuclear medium are an intriguing issue which ties together modern particle and nuclear physics. The interest to these questions has been intensified over the past decade due to the possibility of studying the transition from hadrons to the deconfining phase at heavy ion collisions. In particular, the modification of vector meson properties in nuclear medium has been a subject of a persistent theoretical activity . This was initiated by the idea that in nuclear medium the vector meson masses should drop as a precursor to the chiral symmetry restoration . Several experiments have also been proposed to study the changes of masses, widths and coupling constants of vector resonances in dense (and/or hot) nuclear matter .
The properties of vector resonances in vacuum and the effects of isospin symmetry violation on the mixing of the $`\rho `$, $`\omega `$ resonances in vacuum have been investigated rather carefully by means of QCD sum rules in the past . In the pioneering work of Ref. it was found that the nonzero value for the $`\rho \omega `$ mixing can be linked to the difference of light quark masses, and the possibility of $`m_u=0`$ is seemingly excluded.
Later, the QCD sum rule method was extended to finite temperatures and densities . A number of analyses have found that the masses of $`\rho `$ and $`\omega `$ resonances decrease in nuclear medium<sup>*</sup><sup>*</sup>*See, however, the work of Y. Koike , where opposite behaviour is claimed. A later analysis, based on the relation between the current-nucleon forward scattering amplitude and the scattering length of the vector meson off the nucleon in the static limit, again revealed negative mass shifts in the linear density approximation .. In Refs. finite widths of the vector mesons have been taken into account by hand and by calculation of the $`rho`$-meson self-energy in a chiral model for the spectral function, respectively. For the $`\rho `$-meson channel it was found in Ref. that at nuclear saturation density an increasing width of the $`\rho `$-resonance necessitates an increasing $`\rho `$-meson mass. However, for large values of the width the mass is blurred over a large window of possible values.
While appreciable efforts have been directed to estimate the density dependent modification of the masses and lifetime of the light vector mesons at finite density (and/or temperature), the question of $`\rho \omega `$ mixing at finite densities (and/or temperature) has not received much attention. In fact finite nuclear densities can have a significant impact on this amplitude. The fact that nuclear matter can intrinsically be isospin asymmetric implies that the $`\rho \omega `$ mixing in matter can potentially be larger than the vacuum part of the mixing which is induced by the difference in $`u`$ and $`d`$ quark masses, small in units of characteristic hadronic scales. This idea was suggested first in Ref. where it has been pointed out that the presence of asymmetric nuclear matter has a profound effect on the mixing of the $`\rho `$ and $`\omega `$ resonances. There the mixing angle was determined from the matter induced non-diagonal self energy of the $`\rho ^0`$ resonance by employing an SU(2)<sub>F</sub> symmetric hadronic model. Subsequently such a matter induced mixing has also been analysed on a more elaborate footing in Ref. . Along the same lines the author of Ref. investigated the nucleonic density and temperature dependent $`\rho ^0`$-$`\omega `$ mixing at a fixed asymmetry. Thereby, an enhancement of the modulus of the vacuum mixing amplitude was found due to finite density. In all model descriptions the vacuum part of the mixing serves as an input parameter, to which all the results are normalized by hand in the limit of vanishing density. To this end, it is desirable to obtain an independent analysis of the mixing using finite density QCD sum rules, which hopefully would allow to treat the vacuum part and the density part from the first principles. At this point we already notice one principal problem of finite density QCD sum rules. In the presence of nuclear matter there exist nonscalar condensates which can be related to the twists of different dimension. In general, going from mass dimension $`2n`$ to $`2(n+1)`$ the ratio of contributions $`R_{2k}^{2n}`$ with a nonzero, fixed twist $`2k`$ is
$$R_{2k}^{2n}\frac{A_{2(k+1)}}{A_{2k}}\left(\frac{m_N}{M}\right)^2.$$
(1)
This requires the external momenta to be much larger than $`m_N`$ for the OPE to converge. However, the possibility to link properties of a ground state resonance to nonperturbative effects in the vacuum (the condensates) via the sum rule requires external momenta of $`m_N`$. We will later show that for the contribution of twist operators there is a numerical suppression in the corresponding Wilson coeffcients up to mass dimension six. Since at higher dimensions we have no parametrical smallness the above should limit the applicability of QCD sum rules at finite nucleonic density.
In this paper we study the behavior of the isoscalar-isovector mixed correlator of the two vector currents in order to extract nuclear density effects. The asymptotic behaviour of this correlator at large space-like external momenta can be studied within the perturbative QCD framework, with the power corrections represented by quark, gluon, quark-gluon, etc. condensates. In the presence of finite nucleon density the power correction due to these condensates will change as compared to their vacuum values. Due to the presence of the preferred reference frame, in which the nucleons are at rest, new density-dependent power corrections will appear. In both cases we assume the small density regime and keep only the linear terms in the external nuclear density. As we shall see, this approximation is justified for densities not larger than the nuclear saturation density, which is small in proper “vacuum” units.
The asymptotics of the two-point correlation function, calculated this way, can be related to the “phenomenological part” which includes the contributions of vector resonances, continuum and the screening terms . A success or a failure of the QCD sum rule analysis of vector meson properties would depend on how reliably the contribution of individual resonances ($`\rho `$, $`\omega `$,…) can be separated from the rest of the contributions.
We carefully examine the density-dependent part of the operator product expansion (OPE) and find that the effects of matter-induced mixing due to nucleonic matrix elements of nonscalar and scalar QCD operators in asymmetric nuclear matter follow a certain hierarchy. The asymmetric density-induced effects start dominating vacuum contribution at asymmetries $`\alpha _{pn}(n_pn_n)/(n_p+n_n)0.2`$ and an overall nucleonic density twice the nuclear saturation density $`n_N^0=0.17`$ fm$`{}_{}{}^{3}=`$(111 MeV)<sup>3</sup>. However, the analysis of the phenomenological part of the QCD sum rules shows that the scattering contribution, usually called screening term, turns out to be numerically by far more important. Brought to the OPE-side of the sum rule, the screening term can be regarded as a mass dimension two power correction. Already at intermediate asymmetries $`\alpha _{pn}0.1`$ and saturation density vacuum and matter induced $`\rho `$-$`\omega `$ mixing are of the same sign and comparable in magnitude.
The smallness of the density-dependent pieces in the OPE as compared to the screening term indicates that any conclusion about the density-dependent piece in the $`\rho \omega `$ mixing amplitude will mostly depend on the assumptions made about the spectral density, i.e. what is usually called phenomenological part of the sum rules. This casts strong doubts on the applicability of the finite density QCD sum rules for the extraction of the isovector-isoscalar mixing since the density-dependent “QCD input” is negligibly small. This concern lead us to re-examine the screening terms in the isovector-isovector and isoscalar-isoscalar correlators which were used in previous works to investigate the modification of $`\rho `$ and $`\omega `$ masses in nuclear matter. We have found that all previous analyses have used the same value for the screening terms in the isovector-isovector and isoscalar-isoscalar correlators. This is an unfortunate error because the screening term in the omega channel turns out to be 9 times larger than the value used in Refs. . This changes drammatically all the conclusions about the behaviour of $`m_\omega `$ in nuclear matter, and indicates that $`m_\omega `$ is a growing function of density in the linear density approximation.
## II Isosinglet-isotriplet correlator of the two vector currents at finite densities
We start with the (causal) mixed correlator of isotriplet and isosinglet currents in asymmetric nuclear matter
$$\mathrm{\Pi }_{\mu \nu }id^4x\text{e}^{iqx}\text{T}j_\mu ^T(x)j_\mu ^S(0)_{n_N},$$
(2)
where
$$j_\mu ^T\frac{1}{2}\left(\overline{u}\gamma _\mu u\overline{d}\gamma _\mu d\right),j_\mu ^S\frac{1}{2}\left(\overline{u}\gamma _\mu u+\overline{d}\gamma _\mu d\right).$$
(3)
We choose the same normalization of the two currents, which also means that their couplings to physical $`\rho `$ and $`\omega `$ resonances are approximately equal. In Eq. (2) the Gibbs average $`_{n_N}`$ ($`n_N`$ indicating finite nucleon density) is approximated by a vacuum and one-particle nucleon states . Due to the presence of a singled out rest frame with four-velocity $`u_\mu `$ there are, in general, two independent, current conserving tensor structures (longitudinal and isotropic) into which $`\mathrm{\Pi }_{\mu \nu }`$ can be decomposed. However, in the limit $`\stackrel{}{q}0`$ one of the corresponding invariants $`\mathrm{\Pi }_l,\mathrm{\Pi }_i`$ becomes redundant , and we therefore concentrate on $`\mathrm{\Pi }_l`$ which satisfies the following dispersion relation
$$\mathrm{\Pi }_l(Q_0^2q_0^2)\frac{\mathrm{\Pi }_\mu ^\mu }{3Q_0^2}=\frac{1}{\pi }_0^{\mathrm{}}𝑑s\frac{\text{Im}\mathrm{\Pi }_l}{s+Q_0^2}+\text{subtractions}.$$
(4)
Subtracting the terms attributed to the $`\rho \gamma \omega `$ electromagnetic mixing from the spectral representation and the OPE and appealing to the literature on density dependent OPE’s of $`\rho `$ and $`\omega `$ current correlators we arrive at the following sum rule
$$\mathrm{\Pi }_l^{}(Q_0^2)=\frac{1}{\pi }_0^{\mathrm{}}𝑑s\frac{\text{Im}\stackrel{~}{\mathrm{\Pi }}_l(s)}{s+Q_0^2}+\text{subtractions}.$$
(5)
The asymptotic behaviour of the lhs of Eq. (5) can be calculated by means of the operator product expansion (OPE). The result is given in terms of the perturbative contribution and power corrections, proportional to the condensates, taken in the presence of the external nucleon density. Retaining terms up to the order $`Q_0^6`$, we present the result in the following form:
$`\mathrm{\Pi }_l^{}(Q_0^2)={\displaystyle \frac{\alpha }{16\pi ^3}}{\displaystyle \frac{1}{4}}\mathrm{ln}Q_0^2+{\displaystyle \frac{1}{Q_0^2}}{\displaystyle \frac{3}{2\pi ^2}}{\displaystyle \frac{m_d^2m_u^2}{4}}+`$ (6)
$`{\displaystyle \frac{1}{Q_0^4}}\left[{\displaystyle \frac{m_u}{2}}\overline{u}u_{n_N}{\displaystyle \frac{m_d}{2}}\overline{d}d_{n_N}+{\displaystyle \frac{2}{3}}{\displaystyle \frac{Q_\mu Q_\nu }{Q^2}}𝒮(\overline{u}\gamma _\mu 𝒟_\nu u\overline{d}\gamma _\mu 𝒟_\nu d)_{n_N}\right]`$ (7)
$`{\displaystyle \frac{1}{Q_0^6}}[{\displaystyle \frac{\pi \alpha _s}{2}}(\overline{u}\gamma _\mu \gamma _5\lambda ^au)^2(\overline{d}\gamma _\mu \gamma _5\lambda ^ad)^2+{\displaystyle \frac{2}{9}}[(\overline{u}\gamma _\mu \lambda ^au)^2(\overline{d}\gamma _\mu \lambda ^ad)^2]_{n_N}+`$ (8)
$`2\pi \alpha {\displaystyle \frac{4}{9}}(\overline{u}\gamma _\mu \gamma _5u)^2{\displaystyle \frac{1}{9}}(\overline{d}\gamma _\mu \gamma _5d)^2+{\displaystyle \frac{2}{9}}[{\displaystyle \frac{4}{9}}(\overline{u}\gamma _\mu u)^2{\displaystyle \frac{1}{9}}(\overline{d}\gamma _\mu d)^2]_{n_N}+`$ (9)
$`{\displaystyle \frac{8}{3}}{\displaystyle \frac{Q_\mu Q_\nu Q_\lambda Q_\sigma }{Q^4}}𝒮(\overline{u}\gamma _\mu D_\nu D_\lambda D_\sigma u\overline{d}\gamma _\mu D_\nu D_\lambda D_\sigma d)_{n_N}].`$ (10)
In Eq. (6), symbol $`𝒮`$ denotes the operation of making tensors symmetric and traceless. As usual, (for example ), the averages over mixed operators and twist four contributions have been omitted in Eq. (6). The former can either be reduced to four quark operators by use of the equation of motion (these contributions are already included), or they are suppressed at $`\mu ^21`$ GeV<sup>2</sup> since there the gluon content of the nucleonic wave function is small . The latter has been argued in Ref. to have no substantial effect on the $`\rho `$ and $`\omega `$ mass shifts, and we will therefore omit twist four operators. Further progress in calculating the OPE depends on how accurately we can predict the size of various contributions to $`\mathrm{\Pi }^{}`$. We restrict ourselves to the case of low and medium densities, so that the linear (mean field) approximation is justified, and the density dependent part enters in the final expression multiplied by the matrix elements over the single nucleon states. We further make use of the vacuum saturation hypothesis , which becomes an exact relation in the limit of large number of colors. This hypothesis is known to “work” reasonably well in vacuum. However, its application to the nucleon matrix elements is not fully justified. We use this hypothesis to estimate the order of magnitude of dim 6 contribution, noting that their numerical weight in the final result turns out to be small as the OPE is largely dominated by dim 4 contributions. With all these assumptions, Eq. (6) can be reduced to the following form:
$`\mathrm{\Pi }_l^{}(Q_0^2)={\displaystyle \frac{\alpha }{16\pi ^3}}{\displaystyle \frac{1}{12}}\mathrm{ln}Q_0^2+{\displaystyle \frac{1}{Q_0^2}}\left[{\displaystyle \frac{3}{2\pi ^2}}{\displaystyle \frac{m_d^2m_u^2}{12}}\right]+`$ (11)
$`{\displaystyle \frac{1}{Q_0^4}}\left[{\displaystyle \frac{m_um_d}{2}}\{\overline{q}q_0(\mu ^2)+{\displaystyle \frac{\mathrm{\Sigma }_{\pi N}(\mu ^2)}{2\overline{m}}}n_N\}\overline{m}p|\overline{u}u\overline{d}d|p_{\stackrel{}{k}_p=0}\alpha _{pn}n_N+{\displaystyle \frac{1}{2}}m_p\alpha _{pn}n_NA_2^{ud}(\mu ^2)\right]+`$ (12)
$`{\displaystyle \frac{1}{Q_0^6}}\left[{\displaystyle \frac{112}{81}}\pi [\alpha _s\overline{q}q_0^2(\mu ^2)\{\gamma +{\displaystyle \frac{\alpha }{8\alpha _s(\mu ^2)}}\}+2\overline{q}q_0p|\overline{u}u\overline{d}d|p_{\stackrel{}{k}_p=0}\alpha _{pn}n_N]{\displaystyle \frac{5}{12}}m_p^3\alpha _{pn}n_NA_4^{ud}(\mu ^2)\right].`$ (13)
In Eq. (11) $`n_N`$ ($`n_p`$, $`n_n`$) denotes the total nucleonic (proton, neutron) density; $`\alpha _{pn}`$ the $`p`$$`n`$ asymmetry, defined as $`\alpha _{pn}(n_pn_n)/n_N`$; $`m_p`$ the proton mass; $`\overline{q}q_0`$ the value of the $`u`$-quark condensate; $`\gamma `$ the asymmetry of $`u`$-quark and $`d`$-quark condensate defined as $`\gamma \overline{d}d_0/\overline{u}u_01`$; $`\mathrm{\Sigma }_{\pi N}=(45\pm 7)`$ MeV the nucleon sigma term; and $`\overline{m}`$ is defined as $`\overline{m}1/2(m_u+m_d)`$. The electromagnetic coupling is $`\alpha `$, the strong coupling at scale $`\mu `$ is $`\alpha _s(\mu ^2)`$. In Eq. (11) we have already used the numerical smallness of isospin and chiral symmetry violating parameters as compared to the normal hadronic scale and thus neglected terms proportional to $`m_{u(d)}\gamma `$, $`\gamma ^2`$ and so on. For similiar reasons it is justified to neglect the effects of isopsin violation in nucleon matrix elements and take $`p|\overline{u}u\overline{d}d|p=n|\overline{u}u\overline{d}d|n`$, which leads to the dependence on the asymmetry factor $`\alpha _{pn}`$. The scalar matrix element $`p|\overline{u}u\overline{d}d|p`$ is related to the baryon octet mass splitting, $`(m_\mathrm{\Xi }m_\mathrm{\Sigma })/m_s`$, and it is numerically close to 0.7 .
As for the contribution of symmetric and traceless twist two quark bilinears of dimension four and six, their nucleonic matrix elements are determined by the quark parton distributions $`A_2^{ud}(\mu ^2)`$ and $`A_4^{ud}(\mu ^2)`$ as
$`𝒮\overline{q}\gamma _{\mu _1}D_{\mu _2}q_{N(k)}(\mu ^2)`$ $`=`$ $`iA_2^{ud}(\mu ^2)\left(k_{\mu _1}k_{\mu _2}{\displaystyle \frac{1}{4}}g_{\mu _1\mu _2}k^2\right),`$ (14)
$`𝒮\overline{q}\gamma _{\mu _1}D_{\mu _2}D_{\mu _3}D_{\mu _4}q_{N(k)}(\mu ^2)`$ $`=`$ $`iA_4^{ud}(\mu ^2)\left(k_{\mu _1}k_{\mu _2}k_{\mu _3}k_{\mu _4}\text{traces}\right),`$ (15)
where $`k_\mu `$ denotes the nucleon momentum, $`D_\mu `$ is the gauge covariant derivative. In general, the factor $`A_k^q(\mu ^2)`$ can be obtained from the parton distributions $`Q(x,\mu ^2)`$ and $`\overline{Q}(x,\mu ^2)`$ in the proton as
$$A_k^q(\mu ^2)=2_0^1𝑑xx^{k1}\left(Q(x,\mu ^2)+()^k\overline{Q}(x,\mu ^2)\right).$$
(16)
In Ref. the parton distributions in the nucleon have been fitted to experiment at a resolution scale $`\mu ^2=0.26`$ GeV<sup>2</sup>. Using these distributions and performing the integrations of Eq. (16), we obtain
$$A_2^{ud}(\mu ^2=0.26,\text{GeV}^2)=0.429,A_4^{ud}(\mu ^2=0.26,\text{GeV}^2)=0.097.$$
(17)
To generate the respective values at the scale $`\mu ^21`$ GeV<sup>2</sup> relevant for the sum rule we simply use the conversion factors $`f_2(\mu _2^2,\mu _1^2)\frac{A_2^{u+d}(\mu _2^2)}{A_2^{u+d}(\mu _1^2)}`$ and $`f_4(\mu _2^2,\mu _1^2)\frac{A_4^{u+d}(\mu _2^2)}{A_4^{u+d}(\mu _1^2)}`$. Using $`A_2^{u+d}`$ and $`A_4^{u+d}`$ at 1 GeV from Ref. , we arrive at the following values of matrix elements of interest
$$A_2^{ud}(1\text{GeV}^2)=0.32,A_4^{ud}(1\text{GeV}^2)=0.062.$$
(18)
Keeping this in mind, we will for now proceed to the numerical evaluation of OPE. Performing a Borel transformation of the OPE of Eq. (11) and omitting the numerically strongly suppressed dimension two power correction, we obtain
$`\mathrm{\Pi }_l^{}(M^2)=3.7\times 10^6+{\displaystyle \frac{|\overline{q}q|}{M^4}}\left(2\mathrm{MeV}1.5\mathrm{MeV}{\displaystyle \frac{\alpha _{np}}{0.2}}{\displaystyle \frac{n_N}{n_N^0}}\right)`$ (19)
$`{\displaystyle \frac{|\overline{q}q|}{M^4}}{\displaystyle \frac{0.1\mathrm{GeV}^2}{M^2}}\left(1.4\mathrm{MeV}+{\displaystyle \frac{\alpha _{np}}{0.2}}{\displaystyle \frac{n_N}{n_N^0}}[3.8\mathrm{MeV}2.4\mathrm{MeV}]\right)`$ (20)
In this expression we have used the following set of values:
$`m_u=5`$MeV; $`m_d=9`$MeV; $`\overline{m}=7`$MeV; $`m_p=940`$MeV; $`n_N=n_N^0=(111`$MeV)<sup>3</sup> (the nuclear matter saturation density); $`\mathrm{\Sigma }_{\pi N}=45`$MeV; $`\overline{q}q_0=(225`$MeV)<sup>3</sup>; $`\gamma =10^2`$ ; $`\alpha =1/137`$; and $`\alpha _s(1\text{GeV}^2)=0.5`$ . The quark condensate has been factored out numerically for the sigma-term and the twist contributions.
Several important observations should be made at this point. For $`M1`$ GeV, a pure perturbative contribution is negligibly small as compared to power corrections. The latter are dominated by dimension 4, with the constant term originating from $`m_u`$ and $`m_d`$ mass difference and the $`\alpha _{np}n_N`$-dependent piece given by the $`A_2^{ud}`$ contribution. At the level of dimension 6 we observe three different terms (second line of Eq. (19)): vacuum part, density-dependent scalar condensate and twist contributions. At this dimension the density dependent contribution from scalar condensate and twist tend to cancel each other. This cancellation can be an artefact of chosen parameters and/or of the crude nature of approximations made in estimating the size of the four-quark matrix elements over the nucleon. Nevertheless, $`M1`$ GeV, and the OPE is dominated by dim 4 terms, where at $`\frac{\alpha _{np}}{0.2}\frac{n_N}{n_N^0}1`$ the suppression of dim 6 is about 50%. For higher values of asymmetric density $`\mathrm{\Pi }^{}(M^2=1\mathrm{GeV})`$ changes sign.
What does this behaviour of $`\mathrm{\Pi }^{}(M^2=1\mathrm{GeV})`$ mean in terms of the $`\omega \rho `$ resonance mixing amplitude? To answer this question we should parametrize the spectral function in terms of the resonance contributions and analize the resulting sum rule (5).
Following Refs. , we approximate the imaginary part of the correlator by contributions of $`\rho `$, $`\omega `$, $`\rho ^{}`$, $`\omega ^{}`$ resonances and continuum:
$$\frac{1}{\pi }\text{Im}\stackrel{~}{\mathrm{\Pi }}_l(s,\alpha _{pn},n_N)=\frac{1}{4}\left[f_\rho \delta (sm_\rho ^2)f_\omega \delta (sm_\omega ^2)+f_\rho ^{}\delta (sm_\rho ^{}^2)f_\omega ^{}\delta (sm_\omega ^{}^2)+\frac{\rho _{sc}^{ST}}{8\pi ^2}\delta (s)+\frac{\alpha }{16\pi ^3}\theta (ss_0)\right].$$
(21)
The contribution to the mixing due to the electromagnetic continuum is small . Therefore, we will neglect it in the subsequent consideration. In Eq. (21) $`f_\rho `$ and $`f_\omega `$ refer to the $`\rho `$ and $`\omega `$ residues of the $`\rho `$-$`\omega `$ current propagator, and $`\rho ^{}`$ and $`\omega ^{}`$ symbolize the cumulative effect of higher resonancesIn Ref. the cumulative values $`m_\rho ^{}^2,m_\omega ^{}^2`$ were chosen to be about 1.5 GeV<sup>2</sup>, which is well below the physical masses ($``$ 1.7 GeV) of the resonances $`\rho ^{},\omega ^{}`$. introduced in the original analysis in order to have consistent asymptotic behaviour of $`\mathrm{\Pi }^{}(M^2)`$. Besides the “usual” annihilation continuum above a certain threshold $`s_0`$, Eq. (21) exhibits a scattering term which behaves as a pole at $`s=0`$ (Landau pole) . The corresponding coefficient can be calculated explicitly, and the result in the leading order in Fermi momentum ($`p_f/m_N`$ expansion) is given by
$$\rho _{sc}^{ST}=\frac{2\pi ^2}{m_N}\left[F_p^SF_p^Tn_p+F_n^SF_n^Tn_n\right]=\frac{6\pi ^2}{m_N}\alpha _{pn}n_N.$$
(22)
Here the coefficients $`F_{p(n)}^{S(T)}`$ are defined via nucleon matrix elements of quark vector currents at vanishing momentum transfer:
$`F_p^S=F_n^S=p|\overline{u}\gamma _0u+\overline{d}\gamma _0d|p=3`$ (23)
$`F_p^T=F_n^T=p|\overline{u}\gamma _0u\overline{d}\gamma _0d|p=1.`$ (24)
After Borel transformation the contribution of the Landau screening term is usually carried to the lhs of the sum rule to effectively become a power correction of dimension two in the expansion of $`\mathrm{\Pi }_l^{}(M^2)`$. Defining $`f_{\rho \omega }1/2(f_\rho +f_\omega )`$, $`\overline{m}_r^21/2(m_\rho ^2+m_\omega ^2)`$, $`\mathrm{\Delta }m_r^2m_\omega ^2m_\rho ^2`$, and $`\beta =(f_\omega f_\rho )\overline{m}_r^2/(f_{\rho \omega }\mathrm{\Delta }m_r^2)`$ (the primed quantities are defined analogously), we quote the result of Ref. relating $`f_{\rho \omega }`$ to the measurable quantities $`\overline{m}_r^2`$, $`\mathrm{\Delta }m_r^2`$, $`g_\rho `$, and $`g_\omega `$
$$f_{\rho \omega }\frac{12\overline{m}_r^2}{g_\rho g_\omega }\frac{\delta _{\rho \omega }}{\mathrm{\Delta }m_r^2}\frac{m_r^4}{\mathrm{\Delta }m_r^2}\xi ,$$
(25)
where $`g_\rho `$, $`g_\omega `$ are the respective decay constants, and $`\delta _{\rho \omega }`$ enters the measurable mixing parameter $`\epsilon `$ as follows
$$\epsilon =\frac{\delta _{\rho \omega }}{(m_\omega 1/2i\mathrm{\Gamma }_\omega )^2(m_\rho 1/2i\mathrm{\Gamma }_\rho )^2}.$$
(26)
Thereby, $`\epsilon `$ is defined as
$$\omega =\omega _0+\epsilon \rho _0,\rho =\rho _0\epsilon \omega _0,$$
(27)
and $`\mathrm{\Gamma }`$ denotes the width of the respective resonance. It is fair to remark at this point that the observable combination, $`\epsilon \delta _{\rho \omega }\mathrm{\Gamma }_\rho ^1m_\rho ^1`$ will have an additional dependence on density due to a substantial increase of $`\mathrm{\Gamma }_\rho `$ with $`n_N`$ . Thus, finding the decrease of $`\xi `$ with density would certainly allow to conclude that $`\epsilon `$ is decreasing. The opposite behavior, a rising $`\xi `$, would complicate the prediction of $`\epsilon (n_N)`$.
The final sum rule is given by the following expression:
$`{\displaystyle \frac{1}{4}}\xi {\displaystyle \frac{\overline{m}_r^2}{M^2}}\left({\displaystyle \frac{\overline{m}_r^2}{M^2}}\beta \right)\text{e}^{\overline{m}_r^2/M^2}+(\rho \rho ^{},\omega \omega ^{})=`$ (28)
$`1.110^2\mathrm{GeV}^1\{{\displaystyle \frac{18\mathrm{MeV}}{M^2}}{\displaystyle \frac{\alpha _{np}}{0.2}}{\displaystyle \frac{n_N}{n_N^0}}+{\displaystyle \frac{1}{M^4}}(2\mathrm{MeV}1.5\mathrm{MeV}{\displaystyle \frac{\alpha _{np}}{0.2}}{\displaystyle \frac{n_N}{n_N^0}})`$ (29)
$`{\displaystyle \frac{0.1}{M^6}}(1.4\mathrm{MeV}+{\displaystyle \frac{\alpha _{np}}{0.2}}{\displaystyle \frac{n_N}{n_N^0}}[3.8\mathrm{MeV}2.4\mathrm{MeV}])\},`$ (30)
where all masses and the Borel parameter are taken in units of GeV. It is remarkable that the screening term, brought to the OPE side of this sum rules, completely dominates other density dependent contributions. This shows that in the asymmetric nuclear matter background the influence of the screening term on the $`\omega \rho `$ mixing is by far more important than any changes of the QCD condensates. Moreover, for any realistic $`M^2`$ the screening term becomes comparable to the vacuum contribution to the mixing at $`n_Nn_N^0`$ and asymmetries as low as $`\alpha _{np}0.05`$.
In the limit of vanishing density, relation (28) reduces to the known sum rule for $`\rho \omega `$ mixing. A naive evaluation of this sum rule at $`\rho `$-meson mass, $`M^2=(0.77)^2`$, and at $`n_N=0`$, gives a reasonable agreement with experimentally measured value $`\xi =1.1\times 10^3`$ with $`\beta 0.5`$, advocated in Ref. . Next, we parametrize the linear dependence of $`\xi `$ and $`\beta `$ on the density as follows:
$`\xi =\xi ^{(0)}+\xi ^{(1)}\stackrel{~}{n};\beta =\beta ^{(0)}+\beta ^{(1)}\stackrel{~}{n},`$ (31)
where $`\stackrel{~}{n}`$ denotes $`\alpha _{np}n_N`$ in units of $`0.2n_N^0`$.
The primary reason for the introduction of the $`\rho ^{}\omega ^{}`$ contribution in Ref. was the absence of the $`1/M^2`$ term in the OPE side of the sum rule, so that $`\rho `$ and $`\omega `$ contribution alone would not be consistent with the asymptotic behaviour of $`\mathrm{\Pi }^{}`$. Thus, the role of $`\rho ^{}\omega ^{}`$ is to imitate the cancellation of $`1/M^2`$ terms in contributions of various resonances at large $`M^2`$. For a semiquantitative determination of the linear density dependence of $`\xi `$ and $`\beta `$ we proceed as in Ref. . There the vacuum values of $`\xi `$ and $`\beta `$ were estimated by choosing $`M=m_\rho `$, which strongly suppresses the higher resonances. It should then be legitimate to compare powers of $`M^2`$ in the OPE and the lowest resonance contribution. The result is given by
$$\frac{\xi ^{(1)}}{\xi ^{(0)}}+\frac{\beta ^{(1)}}{\beta ^{(0)}}=2.010^4\frac{4}{\xi ^{(0)}\beta ^{(0)}\overline{m}_r^2}.$$
(32)
Using this relation, we can find $`\beta ^{(1)}`$ and $`\xi ^{(1)}`$ separately, evaluating (28) at $`M^2=0.59`$. The final estimate of $`\xi ^{(1)}`$ reads as
$$\xi ^{(1)}[2.30.8]\times 10^3=1.5\times 10^3,$$
(33)
where 2.3 originates from the screening term and -0.8 comes from the OPE. A similar number can be obtained from the combination of (28) and its first derivative in $`M^2`$. This value of $`\xi ^{(1)}`$ leads to the doubling of mixing amplitude and complete screening at $`n_nn_p\pm 2.5\times 10^2\mathrm{fm}^3`$, respectively.
## III Importance of the screening term for the isoscalar-isoscalar correlator
Having found such an important role of the screening term in the isoscalar-isovector mixed correlator, we would like to return to previous analyses of diagonal correlators (isovector-isovector and isoscalar isoscalar) which were used to extract the behaviour of $`m_\rho `$ and $`m_\omega `$ at finite nucleon density . In all these papers it was found that masses and coupling constants of $`\rho `$ and $`\omega `$ resonances behave similarly in nuclear matter, simply because the OPE sides of the sum rules in both cases are the same after the application of the vacuum saturation hypotheses.
We use the same symmetric normalization of the two currents, Eq. (3). From now on we neglect the asymmetry of the nuclear matter and other isospin breaking effects. Then the sum rules for isovector-isovector and isoscalar-isoscalar correlators in medium take the following symbolic form:
$`{\displaystyle \frac{1}{M^2}}F_\rho ^{}e^{m_\rho ^2/M^2}={\displaystyle \frac{1}{8\pi ^2}}\left(1e^{S_\rho ^{}/M^2}\right){\displaystyle \frac{1}{4}}{\displaystyle \frac{n_N}{m_NM^2}}+{\displaystyle \frac{c_4}{M^4}}+{\displaystyle \frac{c_6}{2M^6}}`$ (34)
$`{\displaystyle \frac{1}{M^2}}F_\omega ^{}e^{m_\omega ^2/M^2}={\displaystyle \frac{1}{8\pi ^2}}\left(1e^{S_\omega ^{}/M^2}\right){\displaystyle \frac{9}{4}}{\displaystyle \frac{n_N}{m_NM^2}}+{\displaystyle \frac{c_4}{M^4}}+{\displaystyle \frac{c_6}{2M^6}},`$ (35)
where $`c_4`$ and $`c_6`$ are the same for both expressions. Obviously, at vanishing nucleon density $`F_\omega F_\rho `$ and $`S_\omega S_\rho `$.
It is remarkable, that the screening terms in Eqs. (34-35) are different by a factor of 9. The enhancement of the screening term in the isoscalar-isoscalar channel is due to
$`{\displaystyle \frac{\rho _{Sc}^{SS}}{\rho _{Sc}^{TT}}}={\displaystyle \frac{F_p^SF_p^S}{F_p^TF_p^T}}=9.`$ (36)
This difference was overlooked in Refs. In Ref. $`\rho _{sc}`$ was taken as a free search parameter and determined from the sum rules at the level consistent with $`\rho _{Sc}^{TT}`$ for both correlators. It casts a strong doubt on the validity of the whole approach, since the actual value of the screening term for $`\omega `$ should be 9 times larger..
The coefficients $`c_2`$ and $`c_4`$ can be computed along the same standard technique (again with considerable degree of uncertainty for $`c_6`$). When plugging these values into the sum rules (34-35), we obtain the following numerical relations:
$`{\displaystyle \frac{1}{M^2}}F_\rho ^{}e^{m_\rho ^2/M^2}={\displaystyle \frac{1}{8\pi ^2}}\left(1e^{S_\rho ^{}/M^2}\right){\displaystyle \frac{3.4\times 10^4}{M^2}}{\displaystyle \frac{n_N}{n_N^0}}+{\displaystyle \frac{10^4}{M^4}}\left[4.1+3.8{\displaystyle \frac{n_N}{n_N^0}}\right]+{\displaystyle \frac{10^4}{2M^6}}\left[2.8+1.2{\displaystyle \frac{n_N}{n_N^0}}\right]`$ (37)
$`{\displaystyle \frac{1}{M^2}}F_\omega ^{}e^{m_\omega ^2/M^2}={\displaystyle \frac{1}{8\pi ^2}}\left(1e^{S_\omega ^{}/M^2}\right){\displaystyle \frac{31\times 10^4}{M^2}}{\displaystyle \frac{n_N}{n_N^0}}+{\displaystyle \frac{10^4}{M^4}}\left[4.1+3.8{\displaystyle \frac{n_N}{n_N^0}}\right]+{\displaystyle \frac{10^4}{2M^6}}\left[2.8+1.2{\displaystyle \frac{n_N}{n_N^0}}\right]`$ (38)
where again all masses and dimensional coupling constants are taken in GeV units. It is remarkable that at $`M1GeV`$ the screening term in the $`\omega `$ sum rule is larger by an order of magnitude than any other density-dependent term from the OPE!
As in the previous case, it is convenient to parameterize the density dependence of masses and coupling constants as follows:
$`m=m^{(0)}\left(1+{\displaystyle \frac{m^{(1)}}{m^{(0)}}}{\displaystyle \frac{n_N}{n_N^0}}\right);F=F^{(0)}\left(1+{\displaystyle \frac{F^{(1)}}{F^{(0)}}}{\displaystyle \frac{n_N}{n_N^0}}\right);S_0=S_0^{(0)}\left(1+{\displaystyle \frac{S_0^{(1)}}{S_0^{(0)}}}{\displaystyle \frac{n_N}{n_N^0}}\right).`$ (39)
Using the sum rules (37) and (38), and the first derivatives of these expressions, we solve for $`m^2`$ as a function of $`S_0^{}`$, and Borel parameter $`M`$. The dependence of the threshold on the the density is obtained by requiring the Borel curves, $`m^{}(M^2,S^{},n)`$, be parallel over the Borel window which we take from 0.6 to 1.2 GeV for different values of densities. The slope of the Borel curve $`m(M^2)`$ in the Borel window at zero density represents a “systematic uncertainty” introduced by sum rules and the requirement of the Borel curves to be parallel at different densities is equivalent to the requirement that this uncertainty does not change while going to finite but small densities. The resulting dependence of $`S_0`$ on the density, $`S_0^{(1)}/S_0^{(0)}=0.2`$ for $`\rho `$ and $`0.1`$ for $`\omega `$, allows us to deduce the following estimates for the linear dependence of masses on the density:
$`{\displaystyle \frac{m_\rho ^{(1)}}{m_\rho ^{(0)}}}0.15,{\displaystyle \frac{m_\omega ^{(1)}}{m_\omega ^{(0)}}}0.05.`$ (40)
Our estimate for $`m_\rho ^{(1)}`$ agrees with the results of previous analyses . The result for $`m_\omega ^{(1)}`$ has the opposite sign and correspond to an $`40`$ MeV increase of $`m_\omega `$ at the nuclear saturation density. This difference could be easily explained by the error in the screening term for $`\omega `$ sum rule in . The disagreement with the results of , where the correct form of the screening terms is used, is harder to explain, and we hypothesize that it could be an artefact of different numerical methods used to extract the density dependence of the resonance masses and thresholds.
## IV Discussion
Apart from the question of (non)convergence of the OPE we would like to point out some concerns about usefulness and validity of the sum rules at finite densities.
1. Vacuum factorization at dim 6. It is unclear what the status of factorization procedure is, especially in the presence of nuclear matter. In principle, one could try to relate four-fermion matrix elements over the nucleon states, which appear in the calculation of the OPE, to some measured processes induced by weak interactions. Indeed, non-leptonic hyperon decays and parity violating pion-nucleon coupling constants could be reduced to similar matrix elements from the four-quark operators. It is unclear, though, whether such an analysis is feasible.
2. The importance of a particular choice of the spectral function. In linear density approximation the analysis of the examples of the $`\rho \omega `$ and $`\omega \omega `$ sum rules suggest that there are large contributions from the respective screening terms. In fact, these contributions dominate all density dependent pieces in the OPE. It means that the “QCD input” in these channels is not important in comparison with the choice of the spectral function at finite density.
3. Is the linear density approximation valid up to $`n_0`$ and beyond? The use of the dilute Fermi gas to model the behaviour of the scattering terms and the QCD condensates has its limitations, and a more realistic description may greatly affect the resulting sum rule. However, it seems unfeasible to calculate QCD operator averages over interacting multi nucleon states which one would have to consider when going beyond the dilute gas approximation. An inclusion of Fermi motion of noninteracting nucleons is practically doable but does not alter the zero momentum result significantly.
In conclusion, we have considered the correlator of isovector and isosinglet vector currents in the presence of the asymmetric nuclear matter in linear density approximation. We see a significant dependence of the OPE on $`n_nn_p`$, which becomes comparable to vacuum contributions at $`n_nn_p0.05n_N^0`$. An attempt to extract the $`\rho `$-$`\omega `$ mixing, using the dispersion relation has shown that this mixing is more affected by the presence of the scattering term than by density-dependent part of the OPE. A similar tendency exists in the isosinglet-isosinglet channel, which is normally used to deduce the dependence of $`m_\omega `$ on density. Hence, in linear density approximation the explicitely density dependent part of the spectral functions (scattering terms) in the $`\rho `$-$`\omega `$ and $`\omega `$-$`\omega `$ channels dominantly drive the density dependence of hadronic parameters. The density dependence of the width of the resonances, which has been neglected here, does not alter this finding. However, it may drastically change the conclusions about the direction of resonance mixing at finite, asymmetric density.
## Acknowledgments
R.H. and M.P. would like to thank V. Eletsky and A. Vainshtein for sharing their scepticism about the relevance of QCD sum rules at finite density. M.P. is grateful to the McGill university nuclear theory group, where this work was started, for the hospitality extended to him during his visit. The work of R.H. was funded by a fellowship of Deutscher Akademischer Austauschdienst (DAAD). The work of A.K.D.M. has been supported by the Natural Sciences and Engineering Research Council of Canada and the Fonds FCAR of the Quebéc Government. The work of M.P. was supported in part by the Department of Energy under Grant No. DE-FG02-94ER40823.
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# References
Teleportation of the Relativistic Quantum Field
R.Laiho, S.N.Molotkov, S.S.Nazin
Wihuri Physical Laboratory,
Department of Physics, University of Turku, 20014 Turku, Finland
Institute of Solid State Physics of Russian Academy of Sciences,
Chernogolovka, Moscow District, 142432 Russia
## Abstract
The process of teleportation of a completely unknown one-particle state of a free relativistic quantum field is considered. In contrast to the non-relativistic quantum mechanics, the teleportation of an unknown state of the quantum field cannot be in principle described in terms of a measurement in a tensor product of two Hilbert spaces to which the unknown state and the state of the EPR-pair belong. The reason is of the existence of a cyclic (vacuum) state common to both the unknown state and the EPR-pair. Due to the common vacuum vector and the microcausality principle (commutation relations for the field operators), the teleportation amplitude contains inevitably contributions which are irrelevant to the teleportation process. Hence in the relativistic theory the teleportation in the sense it is understood in the non-relativistic quantum mechanics proves to be impossible because of the impossibility of the realization of the appropriate measurement as a tensor product of the measurements related to the individual subsystems so that one can only speak of the amplitude of the propagation of the field as a whole.
PACS numbers: 03.67.-a, 03.65.Bz, 42.50Dv
One of the fundamental results of the non-relativistic quantum information theory consists in the possibility of the teleportation of an unknown quantum state by means of a quantum communication channel realized by a non-local entangled state (an EPR-pair ) used together with a classical communication channel .
When the unknown quantum state $`|\psi _s`$ belongs to a finite-dimensional Hilbert space (dim$`_s<\mathrm{}`$), the teleportation can be preformed ideally and with the unit probability employing an EPR-pair with finite energy. On the other hand, if the state space of the teleported system is infinite-dimensional, (dim$`_s=\mathrm{}`$), the ideal teleportation formally requires an EPR-pair with infinite energy . However, the teleported state can be made arbitrarily close to the input unknown state with the probability arbitrarily close to unit by increasing the energy of the EPR-pair and thus making the EPR correlations more and more close the ideal ones.
The non-relativistic quantum mechanics yields only an approximate description of the reality. A more correct and complete description is provided by quantum field theory (since the relativistic quantum mechanics does not allow any sensible physical interpretation, the relativistic theory arises from the very beginning as a quantum field theory). Therefore, it is interesting to consider the possibility of teleportation of a completely unknown state of the relativistic quantum field. In addition, although all the teleportation experiments carried out so far are performed with photons which are essentially relativistic particles, they are always interpreted within the framework of the non-relativistic quantum mechanics.
In the rest of the paper, considering a simple example, we shall show that in the relativistic quantum field theory the teleportation of even a one-particle state of free quantum field cannot be achieved in the sense it is understood in the non-relativistic quantum mechanics. The latter actually follows from the existence of a vacuum (cyclic) state which is common to both the completely unknown state to be teleported and the EPR pair together with the microcausality principle (commutation or anticommutation relations for the field operators).
For convenience we shall briefly remind the teleportation procedure in the non-relativistic case.
Suppose that we are given an unknown quantum state $`|\psi _s_s`$ (dim$`_s<\mathrm{}`$) and a maximally entangled EPR-pair $`|\psi _{EPR}_{12}=_1_2`$ (dim$`_{12}<\mathrm{}`$). The EPR-pair is a composite system consisting of two particles with the state spaces $`_1`$ and $`_2`$. To achieve the teleportation, one performs a joint measurement on the unknown quantum state and one of the particles of the EPR-pair described by an identity resolution in $`=_s_{12}`$ over some measurable outcome space $`\mathrm{\Theta }`$ (which is discrete, $`\mathrm{\Theta }=_i\theta _i`$, for the teleportation of the states belonging to finite-dimensional spaces). The measurement is defined by the identity resolution
$$I_{s12}=\underset{i}{}(\theta _i)=\underset{i}{}I_1_{2s}(\theta _i).$$
(1)
If the measurement yields the $`i`$-th outcome, the subsystem 1 is found in a new state
$$\rho _1^i=\frac{\text{Tr}_{2s}\{(\theta _i)(\rho _s\rho _{EPR})\}}{\text{Tr}_{12s}\{(\theta _i)(\rho _s\rho _{EPR})\}}.$$
(2)
To within a unitary rotation $`U_i`$ which only depends on the measurement outcome $`i`$ and does not depend on the unknown state, the state (2) coincides with the unknown input state:
$$\stackrel{~}{\rho }_1=\rho _s\stackrel{~}{\rho }_1=U_i\rho _1^iU_i^1.$$
(3)
The non-relativistic teleportation procedure substantially employs the fact that the Hilbert state space of each subsystem can be accessed separately.
In the non-relativistic quantum mechanics physically different systems are treated in the same way in the sense that any two systems are formally considered to be identical if their state spaces are identical (isomorphic). Therefore, formally, teleported is the unknown state vector rather than the particle itself. There are no rules that prohibit the superposition of the states belonging to physically different subsystems. In the quantum field theory the situation is completely different.
Let us now turn to the teleportation of an unknown state of a free quantum field. For simplicity we shall first consider the teleportation of a one-particle state of the free scalar quantum field, although the most interesting is perhaps the case of the gauge (photon) field. To avoid unnecessary technical details associated with the indefinite metrics, we shall restrict our analysis to the scalar field. All the remarks concerning the scalar field teleportation are also relevant to photon teleportation (we mean the teleportation of a completely unknown one-photon state when not only the polarization state but also the wave packet shape is unknown).
The states of a relativistic quantum system are described by the rays in the physical Hilbert space $``$ where a unitary representation of the covering Poincaré group is realized . The local quantum field $`\phi (\widehat{x})`$ (here $`\widehat{x}=(t,𝐱)`$ is a point in the Minkowski space-time) is defined as a tensor (if the field has more than one component) operator-valued distribution. To be more precise, corresponding to any function (or a set of functions, if the field is multicomponent) $`f(\widehat{x})𝒥(\widehat{x})`$, where $`𝒥(\widehat{x})`$ is the space of test infinitely differentiable functions decreasing together with all their derivatives at the infinity faster than any polynomial ), is the operator symbolically written as
$$\phi (f)=\underset{j=1}{\overset{r}{}}\phi _j(f_j)=\underset{j=1}{\overset{r}{}}\phi _j(\widehat{x})f_j(\widehat{x})𝑑\widehat{x}.$$
(4)
The operators $`\phi (f)`$ and $`\phi ^{}(f)`$ have a common domain which does not depend on $`f(\widehat{x})`$, is dense in $``$, and is invariant under the action of the field operators, $`\phi (f)\mathrm{\Omega }\mathrm{\Omega }`$ ($`\phi ^{}(f)\mathrm{\Omega }\mathrm{\Omega }`$). For any vectors $`|\varphi ,|\psi \mathrm{\Omega }`$ the quantity $`\varphi |\phi (f)|\psi `$ is a distribution from $`𝒥^{}(\widehat{x})`$ ($`𝒥^{}(\widehat{x})`$ is the space of distributions conjugate to $`𝒥(\widehat{x})`$).
The space $`\mathrm{\Omega }`$ contains a cyclic vector, called the vacuum state, $`|0\mathrm{\Omega }`$ such that the set of all polynomials $`𝒫(\phi ,f)`$ constitute an operator algebra with involution whose action on $`|0\mathrm{\Omega }`$ generates the entire space $`\mathrm{\Omega }`$. The field operator algebra is defined as
$$𝒫(\phi ,f)=f_0+\underset{n=1}{\overset{\mathrm{}}{}}\mathrm{}\phi (\widehat{x}_1)\phi (\widehat{x}_2)\mathrm{}\phi (\widehat{x}_n)f(\widehat{x}_1,\widehat{x}_2,\mathrm{},\widehat{x}_n)𝑑\widehat{x}_1𝑑\widehat{x}_2\mathrm{}𝑑\widehat{x}_n,$$
(5)
$$f(\widehat{x}_1,\widehat{x}_2,\mathrm{},\widehat{x}_n)𝒥(\widehat{x}^n).$$
The fact that the field operators form an algebra implies that any observable can be expressed through the field operators .
The unsmeared field operators $`\phi (\widehat{x})`$ map the regular states from $`\mathrm{\Omega }`$ to the generalized states $`𝒫(\phi (\widehat{x}))\mathrm{\Omega }\mathrm{\Omega }^{}`$ ($`\mathrm{\Omega }^{}`$ is the conjugate space to $`\mathrm{\Omega }`$ consisting of all the linear functionals defined on $`\mathrm{\Omega }`$ and continuous with respect to the scalar product in $``$). The microcausality principle is also postulated; to be more precise, the field operators are assumed to commute (anticommute) if the supports of their corresponding functions $`f(\widehat{x}),g(\widehat{y})`$ are separated by a space-like interval ($`\text{supp}f(\widehat{x})g(\widehat{y})(\widehat{x}\widehat{y})^2<0`$), i.e. for any vector $`|\psi \mathrm{\Omega }`$ we have
$$[\phi (f),\phi (g)]_\pm |\psi =0,(\widehat{x}\widehat{y})^2<0.$$
(6)
The expression (6) is interpreted as the impossibility of any causal relation between the measurements performed in the domains separated by a space-like interval since no interaction can propagate faster than light.
Further, the requirements that the system states are described by the rays in the Hilbert where a Poincaré group representation is realized and the spectrum of the group generators in the momentum representation lies in the front part of the light cone imply that the Lorentz-covariant quantum field can only be realized as an operator valued distribution rather than the field of operators $`\phi (\widehat{x})`$ acting in $``$ .
The interpretation of a quantum field as a field of operators acting in $``$ is only consistent with the trivial two-point function $`0|\phi ^{}(\widehat{x})\phi ^+(\widehat{y})|0=const`$ and results in an obvious violation of the microcausality principle. The smearing function $`f(\widehat{x})`$ can be interpreted (with some reservations) as the amplitude (“shape”) of the one-particle packet.
Let us now construct the EPR-state, the one-particle state of the scalar to be teleported, and the corresponding measurement for the relativistic case emphasizing the differences from the non-relativistic theory.
The state of the EPR-pair is described by the vector $`|\psi _{EPR}\mathrm{\Omega }`$ in the subspace of the two-particle states. The most general form of the relevant vectors is
$$|\psi _{EPR}=𝒫_2(\phi ,)|0=𝑑\widehat{x}_1𝑑\widehat{x}_2(\widehat{x}_1,\widehat{x}_2)\phi ^+(\widehat{x}_1)\phi ^+(\widehat{x}_2)|0,$$
(7)
where
$$\phi ^\pm (\widehat{x})=\frac{1}{(2\pi )^{3/2}}\text{e}^{i\widehat{k}\widehat{x}}\theta (k^0)\delta (\widehat{k}^2m^2)a^\pm (\widehat{k})𝑑\widehat{k}=\frac{1}{(2\pi )^{3/2}}_{V_m^+}\text{e}^{i\widehat{k}\widehat{x}}a^\pm (𝐤)\frac{d𝐤}{\sqrt{2k^0}},$$
$$\widehat{k}\widehat{x}=k^0x^0𝐤𝐱,k^0=\sqrt{𝐤^2+m^2}.$$
The symbol $`V_m^+`$ in the second integral is kept to emphasize the fact that contributing to the integral are only the values at the mass shell inside the front part of the light cone $`k^0>0`$. The ideal EPR correlations correspond to the case where
$$\stackrel{~}{}(\widehat{x}_1,\widehat{x}_2)=\delta (x_1^0x^0)\delta (x_2^0x^0)\delta (𝐱_1𝐱_2)const(𝐱_1+𝐱_2).$$
(8)
However, the function (8) does not belong to the space of test functions $`𝒥(\widehat{x}^2)`$ and should be understood as a limit of functions $`(\widehat{x}_1,\widehat{x}_2)𝒥(\widehat{x}^2)`$, $`(\widehat{x}_1,\widehat{x}_2)\stackrel{~}{}(\widehat{x}_1,\widehat{x}_2)`$. The ideal EPR pair correspond to the generalized state vector $`|\psi _{EPR}\mathrm{\Omega }^{}`$ of the form
$$|\psi _{EPR}=\frac{1}{(2\pi )^3}_{V_m^+}\frac{d𝐤}{2k^0}\text{e}^{2ik^0x^0}a^+(𝐤)a^+(𝐤)|0.$$
(9)
This state is an analogue of the ideal EPR state in the non-relativistic case for the composite system consisting of two particles 1 and 2,
$$|\psi _{EPR}=\frac{1}{(2\pi )^3}𝑑𝐤|𝐤_1|𝐤_2,$$
(10)
where $`|𝐤_{1,2}`$ are the generalized eigenvectors of the momentum operator, and $`|𝐤_{1,2}𝒥^{}(𝐤)`$ ($`𝒥^{}(𝐤)`$ is the distribution space conjugate to $`𝒥(𝐤)`$). Accordingly, in the position representation the state is written as
$$|\psi _{EPR}=\frac{1}{(2\pi )^3}𝑑𝐱|𝐱_1|𝐱_2,$$
(11)
where $`|𝐱_{1,2}`$ are the generalized eigenvectors of the position operator, and $`|𝐱_{1,2}𝒥^{}(𝐱)`$ ($`𝒥^{}(𝐱)`$. The Fourier transform is known to map the space of distributions $`𝒥^{}(𝐤)`$ onto $`𝒥^{}(𝐱)`$.
Qualitatively, at the intuitive level, the state (9) with $`\stackrel{~}{}`$ from (8) can be interpreted as the creation of two particles with $`𝐱_1=𝐱_2`$ at time $`x^0`$ from vacuum, simultaneously at the entire space (because of the presence of a factor $`const(𝐱_1+𝐱_2)`$). The EPR state is essentially non-local.
The one-particle packet state of the quantum field can be written as
$$|\psi _s=\phi ^+(f)|0=𝑑\widehat{x}f(\widehat{x})\phi ^+(\widehat{x})|0=\frac{1}{(2\pi )^{3/2}}_{V_m^+}\frac{d𝐤}{2k^0}f(𝐤)a^+(𝐤)|0,$$
(12)
where $`f(𝐤)`$ is the packet amplitude in the $`𝐤`$-representation. The state is defined by the equivalence class to which the function $`f(\widehat{x})`$ belongs. Different functions $`f(\widehat{x})`$ which have the same values on the mass shell define the same states.
The non-relativistic analogue of the packet is the state
$$|\psi _s=𝑑𝐤f(𝐤)|𝐤_s,$$
(13)
belonging in the non-relativistic case to the Hilbert state space of the particle to be teleported, $`|\psi _s_s`$.
Because of the existence of a common vacuum state in the quantum field theory, the vector corresponding to the system “EPR pair + teleported state” should be written as
$$|\mathrm{\Psi }=\phi ^+(f)𝒫(\phi ^+,)|0=𝑑\widehat{x}_1𝑑\widehat{x}_2𝑑\widehat{x}(\widehat{x}_1,\widehat{x}_2)f(\widehat{x})\phi ^+(\widehat{x}_1)\phi ^+(\widehat{x}_2)\phi ^+(\widehat{x})|0,$$
(14)
The existence of a common vacuum state in the relativistic quantum field theory results in a fundamental difference between the relativistic and non-relativistic cases. In contrast to the non-relativistic case where $`|\psi _s|\psi _{EPR}_s_{12}=_s_1_2`$, the three-particle states of the quantum field $`|\mathrm{\Psi }\mathrm{\Omega }`$. (Of course, a different representation of the state space $`=_n\text{Sym}^n_1`$ as a direct sum of the symmetrized tensor products of the one-particle Fock spaces introduces no changes because of the existence of a common cyclic vacuum vector.)
In addition, in the quantum field theory the states are all essentially non-localizable in the sense that, as it was already established long ago (see e.g. Ref.), it is impossible to construct a state with a compact support in the $`𝐱`$-representation using the normalized functions $`f(𝐤)`$ defined on the mass shell (although the states with the fall off arbitrarily close to the exponential one at the infinity can be constructed) . In some cases one can perhaps approximately assume that the states of a composite system localized to within the exponential tails in distant spatial domains can be regarded as the states defined in the tensor product of the corresponding state spaces which formally have different vacuum states. However, this assumption is certainly wrong if the composite systems in entangled states and their measurements are to be considered. This is exactly the case in the teleportation problem. Moreover, if the state space of a composite system is described as a tensor product of the constituent system state spaces, the microcausality principle (commutation relations) is inevitably violated since the operators acting in different Hilbert spaces (factors in the tensor product) are certainly always commuting independently of their position in the Minkowski space; to be more precise, the operators even do not “know” about each other.
Let us now construct the appropriate measurement. Since in the relativistic case the states are also described by the rays in the Hilbert space (just as in the non-relativistic quantum mechanics), the measurements are also described by the positive operator valued identity resolutions.
In the non-relativistic case the measurement used in the teleportation procedure is described by an identity resolution in $`=_s_1_2`$. Let $`\mathrm{\Theta }`$ be a measurable space of possible outcomes with the measure $`d\theta `$; then
$$I_{s12}=I_sI_1I_2=_\mathrm{\Theta }_{}(d\theta )=I_1_\mathrm{\Theta }_{2s}(d\theta )=I_1I_{2s}.$$
(15)
The measurement (15) is only performed on one of the particles in the EPR-pair and the particle in the unknown state to be teleported while the second particle in the EPR-pair (the factor $`I_1`$) is not involved in the measurement itself. It is important for the teleportation procedure that the identity resolution in the entire state space of the three subsystems can be expressed as a tensor product of the corresponding identity resolutions in $`_1`$ and $`_{s2}=_s_2`$, which implicitly assumes the access to the individual subsystems.
For the relativistic quantum field, the identity resolution in the three-particle states subspace cannot be in any way represented as a tensor product of the appropriate identity resolutions in the one-particle and two-particle subspaces. Such a measurement should only be constructed as a general identity resolution in the entire three-particle states subspace.
It is first instructive to examine the measurement used in the teleportation of a one-particle packet in the non-relativistic case:
$$_{}(d\theta )=I_1_{2s}(d\theta )=I_1|\mathrm{\Phi }_{\mathrm{𝐗𝐏}}_{2s}\text{ }_{2s}\mathrm{\Phi }_{\mathrm{𝐗𝐏}}|\frac{d𝐗d𝐏}{(2\pi )^3},$$
(16)
where the space of possible outcomes is the set $`\mathrm{\Theta }=\{𝐗\times 𝐏𝐑_𝐗\times 𝐑_𝐏\}`$.
Here $`𝐗`$ is the sum of the particle positions $`𝐗=𝐱_2+𝐱_s`$, $`𝐏=𝐩_2𝐩_s`$ is the difference of their momenta, and
$$|\mathrm{\Phi }_{\mathrm{𝐗𝐏}}_{2s}=𝑑𝐤\text{e}^{i\mathrm{𝐤𝐗}}|𝐤_2|𝐤+𝐏_s.$$
(17)
It is easy to check that $`(d𝐗d𝐏)`$ is actually an identity resolution in $`_2_s`$; indeed,
$$I_{2s}=I_2I_s=|\mathrm{\Phi }_{\mathrm{𝐗𝐏}}_{2s}\text{ }_{2s}\mathrm{\Phi }_{\mathrm{𝐗𝐏}}|\frac{d𝐗d𝐏}{(2\pi )^3}=$$
(18)
$$d𝐤_1d𝐤_2(|𝐤_1_2|𝐤_2_s\left)\right({}_{s}{}^{}𝐤_2|{}_{2}{}^{}𝐤_1|).$$
A similar identity resolution for the relativistic quantum field in the subspace of two-particle states is
$$I=_{V_m^+}_{V_m^+}\frac{d𝐤_1}{2k_1^0}\frac{d𝐤_2}{2k_2^0}\left(a^+(𝐤_1)a^+(𝐤_2)|0\right)\left(0|a^{}(𝐤_2)a^{}(𝐤_1)\right).$$
(19)
Let us first write down the analogue of the measurement (17) and then complete it to the identity resolution in the subspace of three-particle states. The corresponding measurement can be represented in the form
$$(d\theta )=$$
(20)
$$\left(𝑑\widehat{\xi }_1𝑑\widehat{\xi }_2\mathrm{\Phi }(\theta ,\widehat{\xi }_1,\widehat{\xi }_2)\phi ^+(\widehat{\xi }_1)\phi ^+(\widehat{\xi }_2)|0\right)\left(𝑑\widehat{\xi }_1^{}𝑑\widehat{\xi }_2^{}\mathrm{\Phi }^{}(\theta ,\widehat{\xi }_1^{},\widehat{\xi }_2^{})0|\phi ^{}(\widehat{\xi }_1^{})\phi ^{}(\widehat{\xi }_2^{})\right)d\theta ,$$
which should give the identity resolution (19), i.e.
$$I=(d\theta )=\left(𝑑\widehat{\xi }_1𝑑\widehat{\xi }_2\phi ^+(\widehat{\xi }_1)\phi ^+(\widehat{\xi }_2)|0\right)\left(𝑑\widehat{\xi }_1^{}𝑑\widehat{\xi }_2^{}0|\phi ^{}(\widehat{\xi }_1^{})\phi ^{}(\widehat{\xi }_2^{})\right)$$
(21)
$$\left(𝑑\theta \mathrm{\Phi }(\theta ,\widehat{\xi }_1,\widehat{\xi }_2)\mathrm{\Phi }^{}(\theta ,\widehat{\xi }_1^{},\widehat{\xi }_2^{})\right),$$
which implies the conditions
$$\mathrm{\Phi }(\theta ,\widehat{\xi }_1,\widehat{\xi }_2)=\delta (\xi _1^0\xi ^0)\delta (\xi _2^0\xi ^0)\mathrm{\Phi }(\theta ,𝝃_1,𝝃_2),$$
(22)
$$𝑑\theta \mathrm{\Phi }(\theta ,𝝃_1,𝝃_2),\mathrm{\Phi }^{}(\theta ,𝝃_1^{},𝝃_2^{})=\delta (𝝃_1𝝃_2)\delta (𝝃_1^{}𝝃_2^{}).$$
It should be emphasized that the time $`\xi ^0`$ is the same for $`\widehat{\xi }_1,\widehat{\xi }_2`$ and $`\widehat{\xi }_1^{},\widehat{\xi }_2^{}`$. We shall not dwell on the interpretation of the measurement (20) and not only that this measurement can be considered as a non-local measurement in the position representation performed at time $`\xi ^0`$.
The conditions (19–22) are satisfied if $`\mathrm{\Phi }`$ is chosen in the from
$$\mathrm{\Phi }(\theta ,𝝃_1,𝝃_2)=\delta (𝝃_1𝝃_2)\text{e}^{i𝐏𝝃_1},\theta =(𝐗,𝐏),$$
(23)
where $`𝐗,𝐏`$ have the same meaning as in the non-relativistic case.
Finally, one obtains
$$(d𝐗d𝐏)=$$
(24)
$$\left(_{V_m^+}\frac{d𝐤}{\sqrt{2k^0(𝐤)}\sqrt{2k^0(𝐤+𝐏)}}\text{e}^{i\mathrm{𝐤𝐏}i(k^0(𝐤)+k^0(𝐤+𝐏))\xi ^0}a^+(𝐤)a^+(𝐤+𝐏)|0\right)$$
$$\left(_{V_m^+}\frac{d𝐤^{}}{\sqrt{2k^0(𝐤^{})}\sqrt{2k^0(𝐤^{}+𝐏)}}\text{e}^{i𝐤^{}𝐏+i(k^0(𝐤^{})+k^0(𝐤^{}+𝐏))\xi ^0}0|a^{}(𝐤^{}+𝐏)𝐚^{}(𝐤^{})\right)\frac{d𝐗d𝐏}{(2\pi )^3}.$$
We shall further need the following identity resolution in the subspace of one-particle states of the quantum field:
$$I_1=_{V_m^+}\frac{d𝐤}{2k^0}\left(a^+(𝐤)|0\right)\left(0|a^{}(𝐤)\right)=𝑑𝐱\left(\phi ^+(\widehat{x})|0\right)\left(0|\phi ^{}(\widehat{x})\right).$$
(25)
The complete measurement in the subspace of the states of the composite system consisting of the EPR-pair and the packet to be teleported is
$$_{}(d\theta )=$$
(26)
$$𝑑𝐱\left(𝑑𝝃_1𝑑𝝃_2\mathrm{\Phi }(\theta ,𝝃_1,𝝃_2)\phi ^+(\widehat{\xi }_1)\phi ^+(\widehat{\xi }_2)\phi ^+(\widehat{x})|0\right)$$
$$\left(𝑑𝝃_1^{}𝑑𝝃_2^{}\mathrm{\Phi }^{}(\theta ,𝝃_1^{},𝝃_2^{})0|\phi ^{}(\widehat{x})\phi ^{}(\widehat{\xi }_1^{})\phi ^{}(\widehat{\xi }_2^{})\right)d\theta =$$
$$𝑑𝐱\left(𝑑𝝃\text{e}^{i𝐏𝝃}\phi ^+(\widehat{\xi })\phi ^+(\widehat{\xi }𝐗)\phi ^+(\widehat{x})|0\right)\left(𝑑𝝃^{}\text{e}^{i𝐏𝝃^{}}0|\phi ^{}(\widehat{\xi }^{})\phi ^{}(\widehat{\xi }^{}𝐗)\phi ^{}(\widehat{x})\right)\frac{d𝐗d𝐏}{(2\pi )^3}=$$
$$𝑑𝐱\left(\phi ^+(\widehat{x})|\mathrm{\Phi }_{\mathrm{𝐗𝐏}}\right)\left(\mathrm{\Phi }_{\mathrm{𝐗𝐏}}|\phi ^{}(\widehat{x})\right)\frac{d𝐗d𝐏}{(2\pi )^3},$$
where
$$|\mathrm{\Phi }_{\mathrm{𝐗𝐏}}=𝑑𝝃\text{e}^{i𝐏𝝃}\phi ^+(\widehat{\xi })\phi ^+(\widehat{\xi }𝐗)|0.$$
(27)
Remember that in the variables $`\widehat{\xi }_1,\widehat{\xi }_2`$, $`\widehat{\xi }_1^{},\widehat{\xi }_2^{}`$, and $`\widehat{\xi },\widehat{\xi }^{}`$ the quantity $`\xi ^0`$ has the same value. For symmetry, we retain the four-dimensional notation for the variables.
It should be noted that because of the existence of the common cyclic (vacuum) vector the identity resolution in the subspace of three-particle states cannot be in any way represented as a tensor product of the corresponding identity resolutions in the subspaces of one- and two-particle states. Unlike the non-relativistic case, the identity resolution (26) for the relativistic quantum field cannot be reduced to the form defined by Eq.(16).
The measurement (26) corresponds to the situation where the observation is only performed on the two particles of three, while the third particle is not involved in the measurement.
The probability of obtaining an outcome in the neighbourhood $`d𝐗d𝐏`$ of the point $`\mathrm{𝐗𝐏}`$ of the space of possible outcomes is given by the standard expression
$$\text{Pr}\{d𝐗d𝐏\}=\text{Tr}\{|\mathrm{\Psi }\mathrm{\Psi }|_{}(d𝐗d𝐏)\}=$$
(28)
$$\left(𝑑𝐱|𝒜(𝐱,\mathrm{𝐗𝐏})|^2\right)d𝐗d𝐏,$$
where the total transition amplitude $`𝒜(𝐱,\mathrm{𝐗𝐏})`$ is defined as
$$𝒜(𝐱,\mathrm{𝐗𝐏})=𝑑\widehat{x}^{}𝑑𝐱_1𝑑𝝃\text{e}^{i𝐏𝝃}f(\widehat{x}^{})$$
(29)
$$0|\phi ^{}(\widehat{\xi })\phi ^{}(\widehat{\xi }𝐗)\phi ^{}(\widehat{x})\phi ^+(\widehat{x}_1)\phi ^+(\widehat{x}_1)\phi ^+(\widehat{x}^{})|0,$$
where two coordinates $`\widehat{x}_1,\widehat{x}_1`$ belong to the EPR-pair, the variables $`\widehat{x}`$ and $`\widehat{x}^{}`$ correspond to the packet with the shape $`f(\widehat{x}^{})`$ being teleported, and, finally, the coordinates $`\widehat{\xi },\widehat{\xi }𝐗`$ refer to the measurement.
It should be noted that for the relativistic quantum field there exists no analogy for the expression (2). The knowledge of measurement (an operator valued measure $`(d\theta )`$) itself is not sufficient to tell what states is the quantum system in after the measurement which gave a particular outcome. To answer this question one should know the instrument (superoperator) generating the indicated operator valued measure. However, the superoperator cannot be uniquely recovered from the given operator valued measure. Fortunately, in the non-relativistic quantum mechanics it is sufficient to know only the measurement itself to completely describe the state of the teleported particle . The latter is explained by the possibility of the representation of the measurement itself (an operator valued measure) as a tensor product of the appropriate identity resolutions in the subspaces of the states of constituent subsystems.
For the relativistic quantum field it is impossible to represent the measurement as a tensor product of the form (1). Therefore, asking what is the state of the teleported system after the measurement for the quantum field is physically meaningless and one can only speak of the transition amplitude of the field as a whole from one state to another.
The vacuum average in Eq. (29) is only determined by the quantum field properties. For a free field the vacuum average is decoupled into the pairwise averages so that only six contributions to the transition amplitude $`𝒜`$ arise.
The quantity $`|𝒜(𝐱,\mathrm{𝐗𝐏})|^2`$ can be interpreted as the probability of detecting the “teleported” particle at a point $`𝐱`$ at time $`x^0`$ if the measurement gave an outcome in the neighbourhood $`(𝐗,𝐏;\mathrm{𝐗𝐏}+d𝐗d𝐏)`$. Similarly, the quantity $`𝒜(𝐱,\mathrm{𝐗𝐏})`$ is the transition amplitude for the packet from the state with the shape $`f(𝐱^{})`$ at time $`x^0^{}`$ to the point $`𝐱`$ at time $`x^0`$. One has the following expression for the amplitude:
$$𝒜(𝐱,\mathrm{𝐗𝐏})=2𝑑\widehat{x}^{}𝑑𝐱_1𝑑𝝃\text{e}^{i𝐏𝝃}f(\widehat{x}^{})$$
(30)
$$\{𝒟_m^+(\widehat{x}_1\widehat{x}^{})[𝒟_m^+(\widehat{x}\widehat{\xi })𝒟_m^+(\widehat{x}_1\widehat{\xi }+𝐗)+𝒟_m^+(\widehat{x}_1\widehat{\xi })𝒟_m^+(\widehat{x}\widehat{\xi }+𝐗)]+$$
$$2𝒟_m^+(\widehat{x}\widehat{x}^{})𝒟_m^+(\widehat{x}_1\widehat{\xi })𝒟_m^+(\widehat{x}_1\widehat{\xi }+𝐗)\},$$
where $`𝒟_m^+(\widehat{x})`$ is the commutator distribution for a free field with mass $`m`$,
$$𝒟_m^\pm (\widehat{x})=\pm \frac{1}{i(2\pi )^{3/2}}\text{e}^{i\widehat{p}\widehat{x}}\theta (\pm p_0)\delta (\widehat{p}^2m^2)𝑑\widehat{p}=$$
(31)
$$\frac{1}{4\pi }\epsilon (x_0)\delta (\widehat{x}^2)\frac{im}{8\pi \sqrt{\widehat{x}^2}}\theta (\widehat{x}^2)\left[N_1(m\sqrt{\widehat{x}^2})i\epsilon (x_0)J_1(m\sqrt{\widehat{x}^2})\right]\pm \frac{im}{4\pi ^2\sqrt{\widehat{x}^2}}\theta (\widehat{x}^2)K_1(m\sqrt{\widehat{x}^2}),$$
$$\epsilon (x_0)\delta (\widehat{x}^2)\frac{\delta (x_0|𝐱|)\delta (x_0+|𝐱|)}{2|𝐱|}.$$
To within the exponential tails, the commutator function is zero beyond the light cone and has a singularity on its surface $`\lambda ^2=(\widehat{x}\widehat{x}^{})^2=0`$; outside the light cone the $`𝒟^\pm (\lambda )`$-function decay exponentially at the Compton length as $`|\lambda |^{3/4}\mathrm{exp}(m\sqrt{|\lambda |})`$ . At a fixed point $`\widehat{x}`$ contributing to the integral are only points $`\widehat{x}^{}`$ lying within the light cone issued from the point $`\widehat{x}`$, which actually follows from the microcausality principle and impossibility of the faster-than-light field propagation.
The amplitude (30) is actually a distribution which should be smeared with a test function to obtain a final result. It should also be noted that the product of any number of positive- or negative-frequency functions $`𝒟^\pm (\widehat{x})`$ (unlike the product of causal functions) is again correctly defined as a distribution from $`𝒥^{}(\widehat{x})`$ since in the momentum representation all these functions have their supports located in the front part of the light cone.
Since we are only interested in the relative probabilities of different processes, we shall directly employ the expression (30) for the amplitude.
Because of the common vacuum vector, it is impossible to arrange a measurement in which there are no contributions to the transition amplitude from the processes which are irrelevent to the teleportation. Formally, the fraction of all these irrelevant processes is 1/2. This circumstance has a fundamental nature and cannot be circumvented by any geometrical tricks in the experiment.
The commutator function $`𝒟_m^+(\widehat{x}\widehat{y})`$ describes the creation of a particle at point $`\widehat{x}`$, its propagation, and destruction at point $`\widehat{y}`$ (for $`y^0>x^0`$)
$$0|\phi ^{}(\widehat{y})\phi ^+(\widehat{x})|0=i𝒟_m^+(\widehat{x}\widehat{y}).$$
(32)
Further, the Lorentz-invariant scalar product
$$(\phi ^{}(f),\phi ^+(g))=0|\phi ^{}(f)\phi ^+(g)|0=𝑑\widehat{x}𝑑\widehat{y}f^{}(\widehat{y})𝒟_m^+(\widehat{x}\widehat{y})g(\widehat{x})=$$
(33)
$$𝑑\widehat{x}𝑑\widehat{y}f^{}(\widehat{y})0|\phi ^{}(\widehat{y})\phi ^+(\widehat{x})|0g(\widehat{x})$$
is interpreted as the amplitude of the packet transition from the state with the “shape” $`g(\widehat{x})`$ to the state with the “shape” $`f(\widehat{y})`$. Since the test functions $`g(\widehat{x})`$ and $`f(\widehat{y})`$ determine the state of the field through their values on the mass shell only, it is convenient to rewrite the amplitude in the form
$$(\phi ^{}(f),\phi ^+(g))=\frac{i}{(2\pi )^{3/2}}𝑑\widehat{p}\theta (p^0)\delta (\widehat{p}^2m^2)f^{}(\widehat{p})g(\widehat{p}),$$
(34)
where $`f(\widehat{p})`$ and $`g(\widehat{p})`$ are the four-dimensional Fourier transforms of the functions $`f(\widehat{x})`$ and $`g(\widehat{x})`$,
$$f(\widehat{p})=\frac{1}{(2\pi )^{3/2}}𝑑\widehat{x}\text{e}^{i\widehat{p}\widehat{x}}f(\widehat{x}).$$
(35)
Integration over the mass shell in Eq. (34) yields
$$(\phi ^{}(f),\phi ^+(g))=\frac{i}{(2\pi )^{3/2}}_{V_m^+}\frac{d𝐩}{2p_0}f^{}(𝐩)\text{e}^{ip^0y^0}g(𝐩)\text{e}^{ip^0x^0}=$$
(36)
$$𝑑𝐱𝑑𝐲f^{}(𝐲)𝒟_m^+(\widehat{x}\widehat{y})g(𝐱),$$
where the commutator function is defined as
$$𝒟_m^+(\widehat{x}\widehat{y})=\frac{i}{(2\pi )^{3/2}}_{V_m^+}\frac{d𝐩}{2p_0}\text{e}^{i[𝐩(𝐱𝐲)ip^0(x^0y^0)]}\widehat{x}=(x^0,𝐱),\widehat{y}=(y^0,𝐲),$$
(37)
and
$$f(𝐱)=\frac{1}{(2\pi )^{3/2}}𝑑𝐩\text{e}^{i\mathrm{𝐩𝐱}}f(𝐩).$$
(38)
The values of the test functions $`f(𝐩)\text{e}^{ip^0y^0}`$ and $`g(𝐩)\text{e}^{ip^0x^0}`$ on the mass shell uniquely determine the state and are interpreted as the packet shape in the momentum representation. In the position representation the quantities $`f(𝐱)`$ and $`g(𝐲)`$ are interpreted as the spatial shape of the packet at times $`x^0`$ and $`y^0`$. Note that the factors $`\text{e}^{ip^0y^0}`$ and $`\text{e}^{ip^0x^0}`$ refer to the packet shape (one could simply write $`\stackrel{~}{f}(𝐩)=f(𝐩)\text{e}^{ip^0y^0}`$ and similarly for $`g(𝐩)`$) and have nothing to do with the dummy integration variable in Eq. (33). This representation is chosen because in this form the Lorentz-invariant scalar product (36) has the meaning of the transition amplitude from the state which at time $`x^0`$ has the spatial shape $`g(𝐱)`$ to the state with the spatial shape $`f(𝐱)`$ by the time $`y^0`$. To within the exponentially decreasing tails at the Compton length outside the light cones, the contributions to this amplitude are only given by the points lying inside the light cones issued from each point $`\widehat{x}=(x^0,𝐱)`$ ($`|𝐱𝐲|^2|x^0y^0|^2<0`$) where the function $`g(𝐱)`$ is different from zero.
Then in a similar way Eq.(30) can be rewritten in the form where the integration is only performed over the spatial coordinates
$$𝒜(𝐱,\mathrm{𝐗𝐏})=2𝑑𝐱^{}𝑑𝐱_1𝑑𝝃\text{e}^{i𝐏𝝃}f(𝐱^{})$$
(39)
$$\{𝒟_m^+(\widehat{x}_1\widehat{x}^{})[𝒟_m^+(\widehat{x}\widehat{\xi })𝒟_m^+(\widehat{x}_1\widehat{\xi }+𝐗)+𝒟_m^+(\widehat{x}_1\widehat{\xi })𝒟_m^+(\widehat{x}\widehat{\xi }+𝐗)]+$$
$$2𝒟_m^+(\widehat{x}\widehat{x}^{})𝒟_m^+(\widehat{x}_1\widehat{\xi })𝒟_m^+(\widehat{x}_1\widehat{\xi }+𝐗)\},$$
where the quantity $`f(𝐱)`$ has the meaning of the spatial shape of the unknown packet to be teleported at time $`x^0`$ ($`\widehat{x}=(x^0,𝐱)`$). The quantity $`x^0`$ appears as a parameter in the arguments of the commutator functions. The rest variables $`x_1^0`$, $`x^0^{}`$, $`\xi ^0`$ also appear in the arguments in $`\widehat{x}_1`$, $`\widehat{x}_1^{^{}}`$, $`\widehat{\xi }^{^{}}`$ as parameters. This form is best suitable for interpretation.
For example, the first term in Eq. (39) yields the amplitude of the process associated with the creation of a non-local EPR-pair state (formally, instantaneously in the entire space, as indicated by the integral over $`𝐱_1`$) at time $`x_1^0`$, propagation of the packet in an unknown state which at time $`x^0`$ has the shape $`f(𝐱)`$ and subsequent joint measurement (also non-local, the integral over $`𝝃`$) at time $`\xi ^0`$ performed on the particle in the unknown state and one of the particles of the EPR-pair. In addition, one of the factors describes the free propagation of the second particle in the EPR-pair, which is not involved in the measurement, to the point $`\widehat{x}=(x^0,𝐱)`$. The second term in Eq. (39) has a similar interpretation. The last two terms describe the processes irrelevant to the teleportation. They describe the contributions to the amplitudes corresponding to the processes where the measurement affects only the two particles of the EPR-pair while the particle whose state is to be teleported propagates freely.
Although the transition to the non-relativistic theory cannot be performed literately, it is still interesting to mention the formal algorithm for this transition: One should omit all the terms in the transition amplitude associated with the particle permutations and replace the commutator distributions $`𝒟_m^{}(\widehat{x})`$ by ordinary $`\delta `$-functions. Note that because of the singularity, this replacement can only be understood symbolically.
The replacement of $`𝒟_m^+(\widehat{x})`$-functions by ordinary $`\delta (𝐱)`$-functions is required because in the non-relativistic case the integration is performed with the Galilei-invariant measure $`d\mu (𝐩)=d𝐩`$, while in the relativistic theory the Lorentz-invariant measure $`d\mu (𝐩)=\theta (p^0)\delta (\widehat{p}^2m^2)d\widehat{p}=d𝐩/2p^0|_{V_m^+}`$ is employed which finally gives
$$𝒟_m^+(\widehat{x})=\frac{1}{(2\pi )^{3/2}}_{V_m^+}\frac{d𝐩}{2p_0}\text{e}^{i\widehat{p}\widehat{x}}\frac{i}{(2\pi )^{3/2}}𝑑𝐩\text{e}^{i\mathrm{𝐩𝐱}}=\delta (𝐱).$$
(40)
The temporal phase factors in the non-relativistic case do not matter because of the absence of any limitations on the propagation speed. Finally, the partial amplitude of the transition from point $`𝐱^{}`$ at time $`x^0^{}`$ to the point $`𝐱`$ at time $`x^0`$ we have
$$𝒜(𝐱^{},𝐱,\mathrm{𝐗𝐏})=$$
(41)
$$2f(𝐱)\text{e}^{i\mathrm{𝐏𝐱}}\delta (𝐱𝐱^{}+𝐗)+2f(𝐱)\text{e}^{i𝐏(𝐱+𝐗)}\delta (𝐱𝐱^{}+𝐗)+4f(𝐱)\delta (𝐗)\delta (𝐏),$$
$$𝒜(𝐱,\mathrm{𝐗𝐏})=𝑑𝐱^{}𝒜(𝐱^{},𝐱,\mathrm{𝐗𝐏}),$$
where the partial amplitude for the transition from point $`𝐱^{}`$ at time $`x^0^{}`$ to the point $`𝐱`$ at time $`x^0`$ “weighted” with the packet shape $`f(𝐱^{})`$ is introduced.
If the contribution of only the first term in Eq. (41) to $`𝒜(𝐱^{},𝐱,\mathrm{𝐗𝐏})`$ is understood literally as the amplitude of the transition to the point $`𝐱`$ under the condition that the measurement gave an outcome in the interval $`(𝐗,𝐏;\mathrm{𝐗𝐏}+d𝐗d𝐏)`$, this amplitude coincides (to within an obvious unitary transformation which is only determined by the measurement outcome, i.e. the value of the pair $`𝐗,𝐏`$) with the amplitude of propagation of the wave packet having the shape $`f(𝐱^{})`$ at the initial moment of time $`x^0^{}`$ to the final point with coordinates $`𝐱`$,$`x^0`$
$$𝒜(𝐱,\mathrm{𝐗𝐏})=f(𝐱^{}𝐗)\text{e}^{i𝐏(𝐱^{}𝐗)}.$$
(42)
For the probabilities (again understood symbolically) of obtaining different measurement outcomes we have
$$\text{Pr}\{d𝐗d𝐏\}=\left(𝑑𝐱|𝒜(𝐱,\mathrm{𝐗𝐏})|^2\right)\frac{d𝐗d𝐏}{(2\pi )^3}=$$
(43)
$$\left(𝑑𝐱|f(𝐱)|^2\right)\frac{d𝐗d𝐏}{(2\pi )^3}=\frac{d𝐗d𝐏}{(2\pi )^3};$$
just as the ideal teleportation requires, the probabilities of obtaining various measurement outcomes do not depend on the unknown state which is to be teleported.
The second term in the amplitude (41) also refers to the teleportation process where one of the particles of the EPR-pair and the particle in the unknown state are exchanged (compared with the teleportation process described by the first term in Eq. (41)).
The last term in Eq. (41) is irrelevant to the teleportation and arises when the measurement only affects the two particles of the EPR-pair (there are two equal contributions because of the exchange of the particles within the EPR-pair). These processes contribute only at the point $`𝐗=0,𝐏=0`$ of the outcome space and their effect can in principle be eliminated by simply discarding the measurements which gave this result.
Nevertheless, the first two terms in Eq. (41) describing the teleportation process have different phase factors which does not allow to correctly modify the transition amplitude by a unitary transformation similar to the non-relativistic case (this would be possible if only the first term were present).
At a first glance, one could simply keep only the measurements which gave the results with $`𝐏=0`$ (when the phase factors are identical). However, in that case the contribution of the “parasitic” processes when the unknown packet is not affected by the measurement and propagates freely becomes essential because of the $`\delta `$-functions ($`\delta (𝐗)\delta (𝐏)`$). Under these conditions it is impossible to distinguish between the teleportation and free propagation contributions to the transition amplitude.
In spite of the fact that the parasitic terms cannot be eliminated, their contribution is only important in the vicinity of the point $`𝐗=0,𝐏=0`$ and has zero measure. Although being rather strange at a first glance, this circumstance has a simple qualitative interpretation related to the fact that each pure state in the infinite-dimensional Hilbert space has in a certain sense zero measure, as it is most simply explained in the non-relativistic theory. The EPR-pair state is written as
$$|\psi _{EPR}=\frac{1}{(2\pi )^3}𝑑𝐤|𝐤_1|𝐤+𝐪_2,$$
(44)
with a fixed $`𝐪`$. The parasitic terms correspond to the measurement performed on the two particles of the EPR-pair which is actually reduced to the projection on the state
$$|\mathrm{\Phi }_{\mathrm{𝐗𝐏}}=\frac{1}{(2\pi )^3}𝑑𝐤^{^{}}\text{e}^{i𝐤^{^{}}𝐗}|𝐤^{^{}}_1|𝐤^{^{}}+𝐏_2,$$
(45)
which in fact is another EPR-pair with a different total momentum (remember that we used the EPR-pair with $`𝐪=0`$). For the projection we have
$$\mathrm{\Phi }_{\mathrm{𝐗𝐏}}|\psi _{EPR}\delta (𝐏)𝑑𝐤\text{e}^{i\mathrm{𝐤𝐗}}=\delta (𝐗)\delta (𝐏),$$
(46)
where the $`\delta `$-functions should be understood as indicated in the above discussion.
The latter means that the measure of an individual EPR-pair with a fixed $`𝐪`$ among the entire set of all EPR-pairs is zero. Therefore, the probability of occurrence of a particular EPR-pair with a specified $`𝐪`$ is zero (the measurement runs over the entire set of EPR-pairs).
Thus, in the relativistic quantum field theory the existence of a common cyclic vector state together with the microcausality principle (commutation relations) make the quantum teleportation impossible in the sense it is understood in the non-relativistic quantum mechanics.
It should be emphasized once again that all the above arguments are only applicable to the teleportation of a completely unknown field state. In that case there is no way to “label” the individual particles involved in the teleportation procedure. However, if the state to be teleported is only partly unknown (e.g. for the case of photon field only the polarization state is unknown while the photon momentum and the total momentum of the EPR-pair are specified beforehand), the available information can be used to construct the “labels” distinguishing the identical particles .
This work was supported by the Russian Foundation for Basic Research (project No 99-02-18127), the project “Physical Principles of the Quantum Computer”, and the program “Advanced Devices and Technologies in Micro- and Nanoelectronics”.
This work was also supported by the Wihuri Foundation, Finland.
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# Gravitinos in non-Ricci-flat backgrounds
## Acknowledgement
The author is indebted to Paul Mansfield for useful conversations.
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# Cumulative gravitational lensing in Newtonian perturbations of Friedman-Robertson-Walker cosmologies
## 1. Introduction
The purpose of this paper is threefold: first to give a general analysis of synchronous gauge Newtonian perturbations of Friedman-Robertson-Walker (FRW) cosmologies; second to obtain a condition which identifies a suitable homogenised FRW cosmology for comparison purposes, and third to obtain the relevant formula for the correction to the apparent magnitude-redshift relation for weak gravitational lensing.
The classic paper of Lifschitz & Khalatnikov (1963) generated much interest in the study of perturbations of FRW cosmologies. Among the influential papers that followed were those of Hawking (1966) and Sachs & Wolfe (1967). A review that covers both classical and quantum aspects is Mukhanov et al. (1992). One of the main topics in the literature is the growth of small perturbations in the early universe and the formation of the galaxies. The present work, by contrast, is concerned with the description of the universe as it is today, with the galaxies providing an essentially time independent perturbation of an FRW background.
Although the use of a gauge invariant approach is frequently advocated, since it eliminates any need to specify and interpret gauge conditions, it is nonetheless more appropriate here to employ a synchronous gauge. This gauge is known to involve a certain ambiguity, as was pointed out by Mukhanov et al. (1992), and has given rise to difficulties of interpretation (Press & Vishniac 1980). It will therefore be necessary in the present work to identify an appropriate supplementary gauge condition which admits a natural interpretation. This will be done in §6.
Gravitational lensing has been considered by many authors and with different techniques (see e.g. Schneider et al. 1992; Bertotti 1966; Dyer & Roeder 1972,1973, 1974; Kantowski 1998). One such technique has been the use of stochastic numerical integration of the lensing equations (Holz & Wald 1998; Dyer & Oattes 1988). The effects of weak gravitational lensing on the cosmic microwave background radiation have been considered by Seljak (1996) and by Sachs & Wolfe (1967). The optical properties of the Swiss cheese model of Einstein & Strauss (1945,1946) have been considered by Dyer & Roeder (1974). Both the Swiss cheese model and the point particle model of Newman & McVittie (1982) will be considered here as applications of a general theory.
Newtonian perturbations of FRW cosmologies have previously been considered by McVittie (1931), Harrison (1967) and Newman & McVittie (1982). A central difficulty is that if one neglects terms involving the spatial curvature constant $`k`$, the total gravitational potential for a uniform distribution of particles will be everywhere infinite. What does not seem to be widely appreciated is that this problem disappears when proper account is taken of either a positive or negative $`k`$. This will be discussed in §8$`b`$.
Much of the analysis will be carried out in the linear approximation in the gravitational coupling constant $`\kappa =8\pi G/c^2`$. Consequently only weak gravitational lensing will be described. In particular, shear is neglected because it has a second order effect on luminosity. Caustics are also neglected. These may both be significant restrictions. There are now well-known examples of astronomical images that show the effects of caustics, although there are as yet no firm estimates of the proportion of images for which caustics are involved. It has been claimed by Ellis et al. (1998) that light rays followed back from our present location in time and space meet a galaxy within a redshift $`0<z<5`$. If this is correct then strong gravitational lensing is likely to be significant for images with $`z5`$.
The question as to whether, in general, weak gravitational lensing gives the same apparent magnitude-redshift relation on average as a best-fit FRW model was considered by Weinberg (1976) who gave two supporting arguments. The first was based on an analysis of double image lensing but was only valid for $`q1`$, i.e. for $`\mathrm{\Omega }1`$. The second, subsequently generalised by Peacock (1986), was essentially based on photon conservation (Schneider et al. 1992, pp.99–100), and was valid for all $`\mathrm{\Omega }`$. Although widely accepted in the literature, these arguments have been criticised by Ellis et al. (1988), the first on the grounds that generic lensing gives rise to three images, not two, and the second at least in part because it effectively assumes the result to be proved. The criticism of Weinberg’s first argument is presumably based on the odd number theorem of Burke (1981) and McKenzie (1985), according to which a transparent gravitational lens can give rise only to an odd number of images. However this theorem is flawed since, as was pointed out by Gottlieb (1994), it attempts to use $`3`$-dimensional topology in a $`4`$-dimensional setting. Nonetheless the criticism of Weinberg’s second argument appears to stand. The present paper therefore gives a new supporting argument that, in the weak field limit, and when averaged over large angular scales, the apparent magnitude-redshift relation does indeed agree with that of a best-fit FRW model.
## 2. General perturbations of FRW cosmologies
The FRW metric, in standard form, is
$$\stackrel{}{𝐠}=\mathrm{𝐝𝐭}^\mathrm{𝟐}+𝐒^\mathrm{𝟐}(𝐭)\stackrel{}{𝐡}$$
(2.1)
where $`\stackrel{}{𝐡}`$ is a $`t`$-independent 3-metric given by
$$\stackrel{}{𝐡}=\left(\mathrm{𝟏}+\frac{\mathrm{𝐤𝐫}^\mathrm{𝟐}}{\mathrm{𝟒}}\right)^\mathrm{𝟐}\underset{𝐢=\mathrm{𝟏}}{\overset{\mathrm{𝟑}}{}}(\mathrm{𝐝𝐱}^𝐢)\text{}^\mathrm{𝟐}.$$
(2.2)
Here and henceforth, indices $`i,j\mathrm{}`$ run from $`1`$ to $`3`$, whilst indices $`a,b\mathrm{}`$ will run from $`0`$ to $`3`$. For the most part it will be possible to describe the level surfaces of $`t`$ with respect to an arbitrary coordinate system $`𝐱=\{x^1,x^2,x^3\}`$. The space-time may then be conveniently described with respect to the coordinate system $`\{x^a:a=0,\mathrm{},3\}`$ with $`x^0=t`$.
Quantities associated with $`\stackrel{}{𝐠}`$ will carry a superscript $`\stackrel{}{\text{}}`$ . Covariant differentiation with respect to $`\stackrel{}{𝐠}`$ will be denoted by $`\stackrel{}{;}`$, e.g. $`\stackrel{}{g}_{ab\stackrel{}{;}c}=0`$. The energy tensor of $`\stackrel{}{𝐠}`$ is
$$\stackrel{}{T}\text{}^{ab}=(\stackrel{}{\rho }+\stackrel{}{p})u^au^b+\stackrel{}{p}\stackrel{}{g}\text{}^{ab}$$
(2.3)
where $`\stackrel{}{g}\text{}^{ab}`$ is the inverse of $`\stackrel{}{g}_{ab}`$, $`u^a:=\stackrel{}{g}\text{}^{ab}t_{\stackrel{}{;}b}`$ is the unit future-directed normal to the level surfaces of $`t`$,
$`\kappa \stackrel{}{\rho }(t)`$ $`={\displaystyle \frac{3k}{S^2(t)}}+{\displaystyle \frac{\dot{S}\text{}^2(t)}{S^2(t)}}\mathrm{\Lambda }`$ (2.4)
$`\kappa \stackrel{}{p}(t)`$ $`={\displaystyle \frac{k}{S^2(t)}}{\displaystyle \frac{\ddot{S}(t)}{S(t)}}+\mathrm{\Lambda }`$ (2.5)
give the density $`\stackrel{}{\rho }(t)`$ and pressure $`\stackrel{}{p}(t)`$ respectively, and $`\mathrm{\Lambda }`$ is the cosmological constant, which is not assumed to vanish. A superscript dot denotes differentiation with respect to $`t`$.
The Hubble and deceleration parameters for $`\stackrel{}{𝐠}`$ are defined by
$`\stackrel{}{H}(t)`$ $`:={\displaystyle \frac{\dot{S}(t)}{S(t)}}`$ (2.6)
$`\stackrel{}{q}(t)`$ $`:={\displaystyle \frac{\ddot{S}(t)S(t)}{\dot{S}\text{}^2(t)}}`$ (2.7)
respectively whilst the dimensionless density parameter is defined by
$$\stackrel{}{\mathrm{\Omega }}(t):=\frac{\kappa \stackrel{}{\rho }(t)}{3\stackrel{}{H}\text{}^2(t)}.$$
(2.8)
Note that by (2.4) and (2.5) one has
$$\frac{d}{dt}\stackrel{}{H}(t)=\stackrel{}{H}\text{}^2(t)(1+\stackrel{}{q}(t))=\frac{k}{S^2(t)}\frac{1}{2}\kappa (\stackrel{}{\rho }+\stackrel{}{p}).$$
(2.9)
Covariant differentiation with respect to $`\stackrel{}{𝐡}`$ in the level surfaces of $`t`$ will be denoted by a subscript $`\stackrel{}{\text{ }\text{ }\text{ }\text{ }\text{ }}`$, e.g. $`\stackrel{}{h}_{ij\stackrel{}{\text{ }\text{ }\text{ }\text{ }\text{ }}k}=0`$. Quantities associated with the 3-metric $`\stackrel{}{𝐡}`$, apart from $`\stackrel{}{𝐡}`$ itself will, in addition to the superscript $`\stackrel{}{\text{}}`$ , carry a prefix $`\text{}^{(3)}`$. Thus the Laplacian associated with $`\stackrel{}{𝐡}`$ will be denoted by $`\text{}^{(3)}\stackrel{}{\mathrm{\Delta }}`$. The Riemann and Ricci tensors and Ricci scalar of $`\stackrel{}{𝐡}`$ are
$`\text{}^{(3)}\stackrel{}{R}\text{}^i\text{}_{jkl}`$ $`=2k\delta ^i\text{}_{[k}\stackrel{}{h}\text{}_{l]j}`$ (2.10)
$`\text{}^{(3)}\stackrel{}{R}_{ij}`$ $`=2k\stackrel{}{h}_{ij}`$ (2.11)
$`\text{}^{(3)}\stackrel{}{R}`$ $`=6k`$ (2.12)
respectively.
The calculations of the remainder of this section pertain to a perturbation of the FRW metric (2.1) of the particular form
$$𝐠=dt^2+S^2(t)𝐡$$
(2.13)
where
$$𝐡=\stackrel{}{𝐡}+\delta 𝐡$$
(2.14)
is a $`3`$-metric intrinsic to the level surfaces of $`t`$, i.e. transverse to $`\stackrel{}{𝐮}`$, and where $`\delta 𝐡`$ may depend upon all four coordinates.
The scale factor $`S(t)`$, the cosmological constant $`\mathrm{\Lambda }`$ and the spatial curvature constant $`k`$ will all be taken to be the same as for the unperturbed metric $`\stackrel{}{𝐠}`$. Note that since $`\stackrel{}{𝐡}`$ is independent of $`t`$ one has $`_th_{ij}=_t\delta h_{ij}`$. Covariant differentiation with respect to $`𝐡`$ in the level surfaces of $`t`$ will be denoted by , e.g. $`h_{ij\text{ }\text{ }\text{ }\text{ }\text{ }k}=0`$. Quantities associated with $`𝐡`$, other then $`𝐡`$ itself, will carry a prefix $`\text{}^{(3)}`$. Thus the Laplacian associated with $`𝐡`$ will be denoted by $`\text{}^{(3)}\mathrm{\Delta }`$ and the Ricci tensor and Ricci scalar of $`𝐡`$ will be denoted by $`\text{}^{(3)}R_{ij}`$ and $`\text{}^{(3)}R`$ respectively. Indices $`i,j\mathrm{}`$ on quantities associated with $`𝐡`$ will be raised and lowered by $`𝐡^1`$ and $`𝐡`$ respectively; indices $`a,b\mathrm{}`$ on quantities associated with $`𝐠`$ will be raised and lowered by $`𝐠^1`$ and $`𝐠`$ respectively. Differences between quantities associated with $`𝐠`$ and $`\stackrel{}{𝐠}`$, or $`𝐡`$ and $`\stackrel{}{𝐡}`$, will be denoted by $`\delta `$, e.g. $`\delta \text{}^{(3)}R_{ij}:=\text{}^{(3)}R_{ij}\text{}^{(3)}\stackrel{}{R}_{ij}`$.
The vector field $`u^a=t^{;a}`$ is irrotational and geodesic. The integral curves of $`𝐮`$ are to be regarded as the world lines of a preferred family of observers, described henceforth as comoving. In particular the galaxies are presumed to be comoving. The coordinate $`t`$ thus carries physical significance.
The second fundamental form of the level surfaces of $`t`$ is given by
$$\chi _{ab}:=u_{a;b}.$$
(2.15)
Clearly $`\chi _{ab}`$ is transverse to $`u^a`$ and so defines a $`3`$-tensor intrinsic to the level surfaces of $`t`$:
$$\chi _{ij}=\mathrm{\Gamma }_{ij}^t=\frac{1}{2}_t(S^2(t)h_{ij})=\dot{S}(t)S(t)h_{ij}+\frac{1}{2}S^2(t)_t\delta h_{ij}.$$
(2.16)
The Hubble and deceleration parameters of $`u^a`$ with respect to $`g_{ab}`$ are defined by
$`H`$ $`:=u^a\text{}_{;a}={\displaystyle \frac{1}{3}}\chi ^a\text{}_a=\stackrel{}{H}(t)+{\displaystyle \frac{1}{6}}h^{ij}_t\delta h_{ij}`$ (2.17)
$`q`$ $`:=_t\left({\displaystyle \frac{1}{H}}\right)1.`$ (2.18)
Note that $`H`$ and $`q`$ are functions of all four coordinates. By (2.16) the Gauss equation for $`\text{}^{(3)}R^i\text{}_{jkl}`$ leads to
$`\text{}^{(3)}R`$ $`=6kS^2(t)\kappa (\stackrel{}{\rho }\stackrel{}{p})+2S^2(t)(\kappa \delta T_{tt}{\displaystyle \frac{\dot{S}(t)}{S(t)}}h^{ij}_t\delta h_{ij})`$
$`+{\displaystyle \frac{1}{2}}h^{i[m}h^{j]n}(_t\delta h_{ij})(_t\delta h_{mn}).`$ (2.19)
From (2.16), (2.17), (2.18) and (2.19) one has
$`\delta \chi _{ij}`$ $`={\displaystyle \frac{1}{2}}_t(S^2(t)\delta h_{ij})`$ (2.20)
$`\delta H`$ $`={\displaystyle \frac{1}{6}}h^{ij}_t\delta h_{ij}`$ (2.21)
$`\delta q`$ $`={\displaystyle \frac{1}{\stackrel{}{H}}}_t\left({\displaystyle \frac{\delta H}{H}}\right)(1+\stackrel{}{q}){\displaystyle \frac{\delta H}{H}}`$ (2.22)
$`\delta \text{}^{(3)}R`$ $`=2S^2(t)(\kappa \delta T_{tt}{\displaystyle \frac{\dot{S}(t)}{S(t)}}h^{ij}_t\delta h_{ij})+{\displaystyle \frac{1}{2}}h^{i[m}h^{j]n}(_t\delta h_{ij})(_t\delta h_{mn}).`$ (2.23)
The components of the Ricci tensor of $`g_{ab}`$ are given by
$`R_{tt}`$ $`=3{\displaystyle \frac{\ddot{S}(t)}{S(t)}}{\displaystyle \frac{1}{2S^2(t)}}_t(S^2(t)h^{ij}_t\delta h_{ij}){\displaystyle \frac{1}{4}}h^{ik}h^{jl}(_t\delta h_{ij})(_t\delta h_{kl})`$ (2.24)
$`R_{it}=R_{ti}`$ $`=h^{jk}(_t\delta h_{k[i})\text{}_{\text{ }\text{ }\text{ }\text{ }\text{ }j]}`$ (2.25)
$`R_{ij}`$ $`=(2k+2\dot{S}\text{}^2(t)+S(t)\ddot{S}(t)+{\displaystyle \frac{1}{2}}\dot{S}(t)S(t)h^{km}_t\delta h_{km})h_{ij}`$
$`+\delta \text{}^{(3)}R_{ij}2k\delta h_{ij}+{\displaystyle \frac{1}{2S(t)}}_t(S^3(t)_t\delta h_{ij})`$
$`+{\displaystyle \frac{1}{4}}S^2(t)h^{km}(_t\delta h_{km})(_t\delta h_{ij}){\displaystyle \frac{1}{2}}S^2(t)h^{km}(\delta h_{ik})(_t\delta h_{jm}),`$ (2.26)
the last of these having been obtained by means of (2.11). One also has
$`R`$ $`=6\left({\displaystyle \frac{k}{S^2(t)}}+{\displaystyle \frac{\dot{S}\text{}^2(t)}{S^2(t)}}+{\displaystyle \frac{\ddot{S}(t)}{S(t)}}\right)+{\displaystyle \frac{1}{S^2(t)}}\delta \text{}^{(3)}R`$
$`+{\displaystyle \frac{1}{S^4(t)}}_t(S^4(t)h^{ij}_t\delta h_{ij})+{\displaystyle \frac{1}{2}}h^{i(j}h^{k)m}(_t\delta h_{km})(_t\delta h_{ij})`$ (2.27)
by means of (2.11) and
$$\delta \text{}^{(3)}R=h^{ij}\text{}^{(3)}R_{ij}\stackrel{}{h}\text{}^{ij}\text{}^{(3)}\stackrel{}{R}_{ij}=h^{ij}(\delta \text{}^{(3)}R_{ij}2k\delta h_{ij}).$$
(2.28)
Hence the energy tensor $`𝐓`$ of $`𝐠`$, as defined by Einstein’s equations
$$\kappa T_{ab}=R_{ab}\frac{1}{2}Rg_{ab}+\mathrm{\Lambda }g_{ab},$$
(2.29)
has components
$`\kappa T_{tt}`$ $`=\kappa \stackrel{}{\rho }+{\displaystyle \frac{1}{2S^2(t)}}\delta \text{}^{(3)}R+{\displaystyle \frac{\dot{S}(t)}{S(t)}}h^{ij}_t\delta h_{ij}+{\displaystyle \frac{1}{4}}(h^{ij}_t\delta h_{ij})^2`$ (2.30)
$`\kappa T_{it}=\kappa T_{ti}`$ $`=(h^{jk}_t\delta h_{k[i})_{\text{ }\text{ }\text{ }\text{ }\text{ }j]}`$ (2.31)
$`\kappa T_{ij}`$ $`=\{S^2(t)\kappa \stackrel{}{p}{\displaystyle \frac{1}{2S(t)}}_t(S^3(t)h^{km}_t\delta h_{km})`$
$`{\displaystyle \frac{1}{4}}S^2(t)h^{p(q}h^{k)m}(_t\delta h_{km})(_t\delta h_{pq})\}h_{ij}`$
$`+\delta \text{}^{(3)}R_{ij}2k\delta h_{ij}{\displaystyle \frac{1}{2}}h_{ij}\delta \text{}^{(3)}R+{\displaystyle \frac{1}{2S(t)}}_t(S^3(t)_t\delta h_{ij})`$
$`+{\displaystyle \frac{1}{4}}S^2(t)h^{km}\{(_t\delta h_{km})(_t\delta h_{ij})2(_t\delta h_{ik})(_t\delta h_{jm})\}`$ (2.32)
The Weyl tensor $`C^a\text{}_{bcd}`$ of $`g_{ab}`$ may be decomposed in the standard form
$$C_{abcd}=8u_{[a}E_{b][d}u_{c]}+2g_{a[c}E_{d]b}2g_{b[c}E_{d]a}2u^e\eta _{eabf}u_{[c}B_{d]}\text{}^f2u^e\eta _{ecdf}u_{[a}B_{b]}\text{}^f$$
(2.33)
where $`\eta _{abcd}`$ is the alternating tensor,
$`E_{ab}`$ $`:=C_{acbd}u^cu^d`$ (2.34)
$`B_{ab}`$ $`:=\text{*}C_{acbd}u^cu^d`$ (2.35)
are the electric and magnetic parts of $`C^a\text{}_{bcd}`$, and
$$\text{*}C_{abcd}:=\frac{1}{2}\eta _{abef}C^{ef}\text{}_{cd}=\frac{1}{2}\eta _{cdef}C^{ef}\text{}_{ab}$$
(2.36)
is the dual of $`C^a\text{}_{bcd}`$.
In order to compute $`E_{ab}`$ one may use the identity
$$2u^bu_{a;[bc]}=R_{dabc}u^bu^d$$
(2.37)
to obtain
$$E_{ij}=\frac{1}{2}R_{ij}\frac{1}{2}S^2(t)(R_{tt}+\frac{1}{3}R)h_{ij}\frac{D}{t}\chi _{ij}\frac{1}{S^2(t)}\chi _{ik}\chi ^k\text{}_j$$
(2.38)
which by (2.16), (2.24), (2.26) and (2.27) gives
$`E_{ij}`$ $`={\displaystyle \frac{1}{4}}S(t)\{_t(S(t)_t\delta h_{ij}){\displaystyle \frac{1}{3}}h_{ij}h^{km}_t(S(t)_t\delta h_{km})\}`$
$`+{\displaystyle \frac{1}{2}}(\delta \text{}^{(3)}R_{ij}2k\delta h_{ij}{\displaystyle \frac{1}{3}}h_{ij}\delta \text{}^{(3)}R)`$
$`+{\displaystyle \frac{1}{8}}S^2(t)\{h^{km}(_t\delta h_{km})(\delta _t\delta h_{ij}){\displaystyle \frac{1}{3}}h^{np}h^{km}(_t\delta h_{km})(_t\delta h_{np})h_{ij}\}.`$ (2.39)
To compute $`B_{ab}`$ one may proceed directly from (2.35) to obtain
$$B_{ab}=u^e\eta _e\text{}^{cd}\text{}_a(\chi _{bd;c}\frac{1}{2}g_{cb}R_{fd}u^f)$$
(2.40)
which by the symmetry of $`B_{ab}`$ gives
$$B_{ab}=u^e\eta _e\text{}^{cd}\text{}_{(a}\chi _{b)d;c}.$$
(2.41)
By means of (2.16) one thus has
$$B_{ij}=\frac{1}{4}S(t)\{\text{}^{(3)}\eta ^{lm}\text{}_i(_t\delta h_{jm})_{\text{ }\text{ }\text{ }\text{ }\text{ }l}+\text{}^{(3)}\eta ^{lm}\text{}_j(_t\delta h_{im})_{\text{ }\text{ }\text{ }\text{ }\text{ }l}\}$$
(2.42)
where
$$\text{}^{(3)}\eta _{ijk}=\frac{1}{S^3(t)}u^a\eta _{aijk}$$
(2.43)
is the alternating tensor for $`h_{ij}`$.
The calculations so far have all been exact. However, in order to facilitate further progress, the perturbation $`\delta 𝐡`$ will be treated as a power series in $`\kappa `$ that vanishes to zeroth order, and the perturbed metric $`𝐠`$ will be studied only to first order in $`\kappa `$. Note that the spatial curvature constant $`k`$ will not be treated as a power series in $`\kappa `$ so, for example, terms with the coefficient $`k\kappa `$ will not be disregarded in first order approximations. This contrasts, in particular, with the work of Holz & Wald (1998).
## 3. Redshift and emission time
Suppose a comoving observer with world line $`\{𝐱=𝐱_1\}`$ in the perturbed space-time makes an observation at time $`t=t_1`$ of a comoving source with redshift $`z_1`$. The image seen by the observer is formed by a congruence of null geodesics centred on a null geodesic $`\gamma `$ connecting a point $`(t_0,𝐱_0)`$ on the source world line $`\{𝐱=𝐱_0\}`$ to the observation point $`(t_1,𝐱_1)`$ on the observer’s world line $`\{𝐱=𝐱_1\}`$. For comparison, consider an observer in the unperturbed space-time who, at the same observation point $`(t_1,𝐱_1)`$, makes an observation of a comoving source with the same redshift $`z_1`$. In this case let $`\stackrel{}{\gamma }`$ be the null geodesic, $`\{𝐱=\stackrel{}{𝐱}\text{}_\mathrm{𝟎}\}`$ the source world line and $`\stackrel{}{t}_0`$ the time of emission. For definiteness let $`\stackrel{}{\gamma }`$ be chosen such that its spatial direction at $`(t_1,𝐱_1)`$ coincides with that of $`\gamma `$.
The immediate problem is to compute the perturbation $`\delta t_0=t_0\stackrel{}{t}\text{}_0`$ in the time of emission. For $`\stackrel{}{𝐠}`$ the emission time $`\stackrel{}{t}_0`$ is determined implicitly by
$$1+z_1=\frac{S(t_1)}{S(\stackrel{}{t}_0)}$$
(3.1)
which for nearby sources gives
$$t_1\stackrel{}{t}_0=\frac{z_1}{\stackrel{}{H}(t_1)}+O(z_1^2)\text{for }z_11.$$
(3.2)
For $`𝐠`$ one may determine $`t_0`$ as follows. Since $`\gamma `$ is a null geodesic one has
$`S^2(t)h_{ij}{\displaystyle \frac{dx^i}{dt}}{\displaystyle \frac{dx^j}{dt}}`$ $`=1`$ (3.3)
$`{\displaystyle \frac{d^2x^a}{d\lambda ^2}}+\mathrm{\Gamma }_{bc}^a{\displaystyle \frac{dx^b}{d\lambda }}{\displaystyle \frac{dx^c}{d\lambda }}`$ $`=0`$ (3.4)
where $`\lambda `$ is an affine parameter along $`\gamma `$ with value $`\lambda _0`$ at $`(t_0,𝐱_0)`$ and value $`\lambda _1`$ at $`(t_1,𝐱_1)`$. By means of (3.3) the $`t`$ component of (3.4) gives
$$\frac{d}{dt}\mathrm{ln}\left(\frac{dt}{d\lambda }\right)+\frac{\dot{S}(t)}{S(t)}+\kappa \nu _\gamma (t)=0$$
(3.5)
where
$$\kappa \nu _\gamma (t):=\frac{1}{2}S^2(t)(_t\delta h_{ij})\frac{dx^i}{dt}\frac{dx^j}{dt}$$
(3.6)
is defined along $`\gamma `$. The redshift as observed at time $`t`$ is given by the standard formula
$$1+z(t)=\frac{\left(\frac{dt}{d\lambda }|_{t_0}\right)}{\left(\frac{dt}{d\lambda }|_t\right)}$$
(3.7)
which combines with (3.5) to give
$$\frac{d}{dt}\mathrm{ln}(1+z(t))=\frac{\dot{S}(t)}{S(t)}+\kappa \nu _\gamma (t).$$
(3.8)
Integration of (3.8) from $`t_0`$ to $`t_1`$ yields
$$1+z_1=\frac{S(t_1)}{S(t_0)}\mathrm{exp}(_{t=t_0}^{t=t_1}\kappa \nu _\gamma (t)𝑑t)$$
(3.9)
which may be regarded as an implicit equation for $`t_0`$ in terms of $`z_1`$ and $`t_1`$. To first order in $`\kappa `$ equation (3.9) gives
$$\delta t_0=\frac{1}{H(\stackrel{}{t}_0)}_{t=\stackrel{}{t}_0}^{t=t_1}\kappa \nu _\stackrel{}{\gamma }(t)𝑑t+O(\kappa ^2)$$
(3.10)
by means of (3.1), with $`\stackrel{}{t}_0`$ is given implicitly in terms of $`z_1`$ and $`t_1`$ by (3.1), and with $`\kappa \nu _\stackrel{}{\gamma }(t)`$ defined by an equation corresponding to (3.6) but for $`\stackrel{}{\gamma }`$ in place of $`\gamma `$: clearly one has $`\kappa \nu _\stackrel{}{\gamma }(t)=\kappa \nu _\gamma (t)+O(\kappa ^2)`$. For nearby sources (3.10) gives
$$\delta t_0=\frac{\kappa \nu _\stackrel{}{\gamma }(t_1)}{\stackrel{}{H}\text{}^2(t_1)}z_1+O(z_1^2)+O(\kappa ^2)\text{for }z_11$$
(3.11)
by means of (3.2).
One may define a spatial proper distance of the source from the observer, with respect to the metric $`𝐠`$, by
$$s_0:=c_{\lambda =\lambda _0}^{\lambda =\lambda _1}(S^2(t)h_{ij}\frac{dx^i}{d\lambda }\frac{dx^j}{d\lambda })^{1/2}𝑑\lambda =c(t_1t_0)$$
(3.12)
where (3.3) has been employed on the right. For nearby sources (3.2) and (3.11) combine to give
$$t_1t_0=\frac{z_1}{\stackrel{}{H}(t_1)}\left(1\frac{\kappa \nu _\stackrel{}{\gamma }(t_1)}{\stackrel{}{H}(t_1)}\right)+O(z_1^2)+O(\kappa ^2)\text{for }z_11$$
(3.13)
so, for such sources, the redshift $`z_1`$ of the source is expressible as a function of its spatial proper distance $`s_0`$ according to
$$z_1(s_0)=\frac{\stackrel{}{H}(t_1)s_0}{c}\left(1+\frac{\kappa \nu _\stackrel{}{\gamma }(t_1)}{\stackrel{}{H}(t_1)}\right)+O(s_0^2)+O(\kappa ^2)\text{for }z_11.$$
(3.14)
Notice that in the case $`\nu _\gamma (t_1)=O(\kappa )`$ equation (3.14) agrees, to first order in $`z_1`$ and $`\kappa `$, with the corresponding formula for the unperturbed FRW space-time.
## 4. Luminosity distance and apparent magnitude
Suppose now that, in the perturbed space-time, the source with world line $`\{𝐱=𝐱_0\}`$ radiates uniformly in all directions with power $`P`$.
Let $`k^a:=\frac{dx^a}{d\lambda }`$ be the tangent to $`\gamma `$ and let $`𝐦`$ be a complex vector at $`(t_0,𝐱_0)`$ satisfying $`𝐦\text{ }𝐦=1`$, $`𝐦𝐤=𝐦𝐦=𝐦𝐮=0`$. In order to maintain these conditions along $`\gamma `$ one requires that $`𝐦`$ be propagated along $`\gamma `$ according to
$$_𝐤𝐦=\left(\frac{𝐦_𝐤𝐮}{𝐮𝐤}\right)𝐤.$$
(4.1)
The congruence of all the null geodesics emanating from $`(t_0,𝐱_0)`$ gives rise to a family of Jacobi fields along $`\gamma `$ with expansion $`\varrho :=k_{a;b}\text{ }m\text{}^am\text{}^b`$ and shear $`\varsigma :=k_{a;b}\text{ }m\text{}^a\text{ }m\text{}^b`$ with respect to the 2-frame $`\{𝐦,\text{ }𝐦\}`$ and the affine parameter $`\lambda `$. Since the congruence is irrotational, the imaginary part of $`\varrho `$ is zero. The standard propagation equations for $`\varrho `$ and $`\varsigma `$ therefore reduce to
$`{\displaystyle \frac{d}{d\lambda }}\varrho `$ $`=\varrho ^2\varsigma \overline{\varsigma }\mathrm{\Theta }`$ (4.2)
$`{\displaystyle \frac{d}{d\lambda }}\varsigma `$ $`=2\varsigma \varrho \mathrm{\Psi }`$ (4.3)
where
$`\mathrm{\Theta }`$ $`:={\displaystyle \frac{1}{2}}R_{ab}k^ak^b={\displaystyle \frac{1}{2}}\kappa T_{ab}k^ak^b`$ (4.4)
$`\mathrm{\Psi }`$ $`:=R_{abcd}k^ak^c\text{ }m\text{}^b\text{ }m\text{}^d=C_{abcd}k^ak^c\text{ }m\text{}^b\text{ }m\text{}^d`$ (4.5)
are the Ricci and Weyl scalars respectively. (The fact that $`𝐦`$ is propagated along $`\gamma `$ according to (4.1), rather than being parallelly propagated as is more usual, has no effect on the equations (4.2) and (4.3).)
Consider a narrow beam of light rays, centred upon $`\gamma `$, from the source point $`(t_0,𝐱_0)`$. Let $`\mathrm{\Delta }A`$ and $`I`$ be the cross-sectional area and apparent luminosity of the beam as determined by comoving observers situated along $`\gamma `$. A standard argument from photon conservation gives that the quantity $`(1+z)^2I\mathrm{\Delta }A`$ is constant along $`\gamma `$. By means of $`\varrho =\frac{1}{2}\frac{d}{d\lambda }\mathrm{ln}\mathrm{\Delta }A`$ one thus obtains
$$\varrho =\frac{d}{d\lambda }\mathrm{ln}(I^{1/2}\frac{dt}{d\lambda })=\frac{d}{d\lambda }\mathrm{ln}(I^{1/2}(1+z)^1).$$
(4.6)
In terms of the quantities
$`J(\lambda )`$ $`:=I^{1/2}(1+z)^1`$ (4.7)
$`\xi (\lambda )`$ $`:=J^2(\lambda )\left({\displaystyle \frac{dt}{d\lambda }}\right)^1\varsigma (\lambda )`$ (4.8)
equations (4.2) and (4.3) become
$`{\displaystyle \frac{d^2}{d\lambda ^2}}J(\lambda )+J^3(\lambda )\xi (\lambda )\overline{\xi }(\lambda )\left({\displaystyle \frac{dt}{d\lambda }}\right)^2+J(\lambda )\mathrm{\Theta }`$ $`=0`$ (4.9)
$`{\displaystyle \frac{d}{d\lambda }}\left({\displaystyle \frac{dt}{d\lambda }}\xi (\lambda )\right)+J^2(\lambda )\mathrm{\Psi }`$ $`=0.`$ (4.10)
Substituting for $`\lambda `$ in terms of $`t`$ in (4.9) and (4.10) by means of (3.5) one thus obtains
$`{\displaystyle \frac{d^2}{dt^2}}J(t)\left({\displaystyle \frac{\dot{S}(t)}{S(t)}}+\kappa \nu _\gamma (t)\right){\displaystyle \frac{d}{dt}}J(t)+\xi (t)\overline{\xi }(t)J^3(t)+\kappa \tau _\gamma (t)J(t)`$ $`=0`$ (4.11)
$`{\displaystyle \frac{d}{dt}}\xi (t)\left({\displaystyle \frac{\dot{S}(t)}{S(t)}}+\kappa \nu _\gamma (t)\right)\xi (t)+J^2(t)\kappa \psi _\gamma (t)`$ $`=0`$ (4.12)
for
$`\kappa \tau _\gamma (t):`$ $`={\displaystyle \frac{1}{2}}\kappa T_{ab}{\displaystyle \frac{dx^a}{dt}}{\displaystyle \frac{dx^b}{dt}}`$ (4.13)
$`\kappa \psi _\gamma (t):`$ $`=C_{abcd}{\displaystyle \frac{dx^a}{dt}}{\displaystyle \frac{dx^c}{dt}}\text{ }m\text{}^b\text{ }m\text{}^d`$ (4.14)
$`=2E_{ij}\text{ }m\text{}^i\text{ }m\text{}^j+{\displaystyle \frac{2}{S^2(t)}}\text{}^{(3)}\eta ^k\text{}_{ij}B_{lk}\text{ }m\text{}^j\text{ }m\text{}^l{\displaystyle \frac{dx^i}{dt}}.`$ (4.15)
The initial conditions for (4.11) and (4.12) are
$$\begin{array}{cc}\hfill \xi (\stackrel{}{t}_0)& =0\hfill \\ \hfill J(\stackrel{}{t}_0)& =0\hfill \\ \hfill \frac{d}{dt}|_{t=\stackrel{}{t}_0}J(t)& =c\sqrt{\frac{4\pi }{P}}\hfill \end{array}\}.$$
(4.16)
To first order in $`\kappa `$, the solutions to (4.11) and (4.12), subject to the initial conditions (4.16), are
$`J(t)`$ $`=\sqrt{{\displaystyle \frac{4\pi }{P}}}{\displaystyle \frac{c}{S(t_0)}}{\displaystyle _{t^{}=t_0}^{t^{}=t}}S(t^{})(1+\kappa \mu _\stackrel{}{\gamma }(t^{}))𝑑t^{}+O(\kappa ^2)`$ (4.17)
$`\xi (t)`$ $`={\displaystyle \frac{4\pi c^2}{P}}{\displaystyle \frac{S(t)}{S(t_0)}}{\displaystyle _{t^{}=t_0}^{t^{}=t}}{\displaystyle \frac{\mathrm{\Sigma }^2(t_0,t^{})}{S(t^{})}}\kappa \psi _\stackrel{}{\gamma }(t^{})𝑑t^{}+O(\kappa ^2)`$ (4.18)
for
$$\mathrm{\Sigma }(t_0,t):=_{t^{}=t_0}^{t^{}=t}S(t^{})𝑑t^{}$$
(4.19)
and
$$\kappa \mu _\stackrel{}{\gamma }(t):=_{t^{}=t_0}^{t^{}=t}(\kappa \nu _\stackrel{}{\gamma }(t^{})\frac{\mathrm{\Sigma }(t_0,t^{})}{S(t^{})}\kappa \tau _\stackrel{}{\gamma }(t^{}))𝑑t^{}$$
(4.20)
where $`\kappa \tau _\stackrel{}{\gamma }(t)`$ and $`\kappa \psi _\stackrel{}{\gamma }(t)`$ are defined by equations analogous to (4.13) and (4.14) for $`\stackrel{}{\gamma }`$ in place of $`\gamma `$. Clearly one has $`\kappa \tau _\stackrel{}{\gamma }(t)=\kappa \tau _\gamma (t)+O(\kappa ^2)`$ and $`\kappa \psi _\stackrel{}{\gamma }(t)=\kappa \psi _\gamma (t)+O(\kappa ^2)`$.
From (4.8), (4.17) and (4.18) one obtains
$$\varsigma (t)=\frac{S(\stackrel{}{t}_0)S(t)}{\mathrm{\Sigma }^2(\stackrel{}{t}_0,t)}\frac{dt}{d\lambda }_{t^{}=\stackrel{}{t}_0}^{t^{}=t}\frac{\mathrm{\Sigma }^2(\stackrel{}{t}_0,t^{})}{S(t^{})}\kappa \psi _\stackrel{}{\gamma }(t^{})𝑑t^{}+O(\kappa ^2)$$
(4.21)
for the shear as a function of $`t`$.
Writing $`t_0=\stackrel{}{t}_0+\delta t_0`$ and setting $`t=t_1`$ in (4.17) one obtains
$$\begin{array}{c}I^{1/2}(z_1)=c\sqrt{\frac{4\pi }{P}}\frac{(1+z_1)^2}{S(t_1)}\{\mathrm{\Sigma }(\stackrel{}{t}_0,t_1)\left(\frac{S^2(\stackrel{}{t}_0)}{\dot{S}(\stackrel{}{t}_0)}\mathrm{\Sigma }(\stackrel{}{t}_0,t_1)\right)_{t=\stackrel{}{t}_0}^{t=t_1}\kappa \nu _\stackrel{}{\gamma }(t)dt\hfill \\ \hfill +_{t=\stackrel{}{t}_0}^{t=t_1}S(t)\kappa \mu _\stackrel{}{\gamma }(t)dt\}+O(\kappa ^2)\end{array}$$
(4.22)
by means of (3.10) and 4.19), where $`\stackrel{}{t}\text{}_0`$ is given in terms of $`t_1`$ and $`z_1`$ by (3.1). From this one obtains
$$\begin{array}{c}I^{1/2}(z_1)=c\sqrt{\frac{4\pi }{P}}\frac{(1+z_1)^2}{S(t_1)}\{\mathrm{\Sigma }(\stackrel{}{t}\text{}_0,t_1)\frac{S^2(\stackrel{}{t}\text{}_0)}{\dot{S}(t_1)}_{t=\stackrel{}{t}\text{}_0}^{t=t_1}\kappa \nu _\stackrel{}{\gamma }(t)dt\hfill \\ \hfill _{t=\stackrel{}{t}\text{}_0}^{t=t_1}\mathrm{\Sigma }(\stackrel{}{t}\text{}_0,t)(\kappa \nu _\stackrel{}{\gamma }(t)\frac{\mathrm{\Sigma }(\stackrel{}{t}\text{}_0,t)}{S(t)}\kappa \tau _\stackrel{}{\gamma }(t))𝑑t\\ \hfill \mathrm{\Sigma }(\stackrel{}{t}\text{}_0,t_1)_{t=\stackrel{}{t}\text{}_0}^{t=t_1}\frac{\mathrm{\Sigma }(\stackrel{}{t}\text{}_0,t)}{S(t)}\kappa \tau _\stackrel{}{\gamma }(t)dt\}+O(\kappa ^2)\end{array}$$
(4.23)
by means of (4.20), (4.19) and an integration by parts. For nearby sources (4.23) gives
$$I^{1/2}(z_1)=\sqrt{\frac{4\pi }{P}}\frac{cz_1}{\stackrel{}{H}(t_1)}\left(1\frac{\kappa \nu _\stackrel{}{\gamma }(t_1)}{\stackrel{}{H}(t_1)}\right)+O(z_1^2)+O(\kappa ^2)\text{for }z_11$$
(4.24)
by means of (3.13). The equation corresponding to (4.23) for the unperturbed metric $`\stackrel{}{𝐠}`$ is
$$\stackrel{}{I}\text{}^{1/2}(z_1)=c\sqrt{\frac{4\pi }{P}}\frac{(1+z_1)^2}{S(t_1)}\mathrm{\Sigma }(\stackrel{}{t}\text{}_0,t_1).$$
(4.25)
From (4.21), the total shear of the image at the observation point $`(t_1,𝐱_1)`$ is given by
$$_{\lambda =\lambda _0}^{\lambda =\lambda _1}\varsigma 𝑑\lambda =_{t=\stackrel{}{t}_0}^{t=t_1}\frac{S(\stackrel{}{t}_0)S(t)}{\mathrm{\Sigma }^2(\stackrel{}{t}_0,t)}\left(_{t^{}=\stackrel{}{t}_0}^{t^{}=t}\frac{\mathrm{\Sigma }^2(\stackrel{}{t}_0,t^{})}{S(t^{})}\kappa \psi _\stackrel{}{\gamma }(t^{})𝑑t^{}\right)𝑑t+O(\kappa ^2).$$
(4.26)
The luminosity distance of the source is defined by
$$d_L:=\sqrt{\frac{P}{4\pi I(z_1)}}$$
(4.27)
which, for nearby sources, combines with (4.24) to give
$$d_L=\frac{cz_1}{\stackrel{}{H}(t_1)}\left(1\frac{\kappa \nu _\stackrel{}{\gamma }(t_1)}{\stackrel{}{H}(t_1)}\right)+O(z_1^2)+O(\kappa ^2)\text{for }z_11.$$
(4.28)
Note that in the case $`\nu _\gamma (t_1)=O(\kappa )`$ this agrees, to first order in $`z_1`$ and $`\kappa `$, with the corresponding formula for the unperturbed FRW space-time. From (4.28) and (3.14) one has
$$d_L=s_0+O(s_0^2)+O(\kappa ^2)$$
(4.29)
which gives the acceptable result that the spatial proper distance and luminosity distance agree to first order in $`s_0`$ and $`\kappa `$, independent of the value of $`\nu _\stackrel{}{\gamma }(t_1)`$. This holds for both the perturbed metric $`𝐠`$ and the unperturbed metric $`\stackrel{}{𝐠}`$.
The apparent magnitude of the source at $`𝐱=𝐱_\mathrm{𝟎}`$ relative to a reference source at $`𝐱=𝐱_{\text{ref}}`$ is defined by
$$m:=m_{\text{ref}}+\frac{5}{2}\mathrm{log}_{10}\frac{I_{\text{ref}}}{I}$$
(4.30)
where $`I`$ and $`I_{\text{ref}}`$ are the apparent luminosities of the respective sources. The reference source is taken to have power $`P_{\text{ref}}`$ and to be at a spatial proper distance $`s_{0,\text{ref}}`$ that is small on a cosmological scale (10 parsecs is conventional). For the unperturbed metric $`\stackrel{}{𝐠}`$, a source of apparent luminosity $`\stackrel{}{I}`$ has an apparent magnitude
$$\stackrel{}{m}=\stackrel{}{m}_{\text{ref}}+\frac{5}{2}\mathrm{log}_{10}\frac{\stackrel{}{I}_{\text{ref}}}{\stackrel{}{I}}$$
(4.31)
relative to a source of apparent luminosity $`\stackrel{}{I}_{\text{ref}}`$. The sources will again be taken to have powers $`P`$ and $`P_{\text{ref}}`$ respectively, with the reference source at the same spatial proper distance $`s_{0,\text{ref}}`$ as for the perturbed metric $`𝐠`$. Let $`m_{\text{ref}}=\stackrel{}{m}_{\text{ref}}`$. The objective sources for $`𝐠`$ and $`\stackrel{}{𝐠}`$ are both taken to be at redshift $`z_1`$. From (4.30) and (4.31) one obtains
$$\delta m(z_1)=\frac{5}{2\mathrm{ln}10}\left\{\underset{s0}{lim}\mathrm{ln}\frac{I_{\text{ref}}(s)}{\stackrel{}{I}_{\text{ref}}(s)}\mathrm{ln}\frac{I(z_1)}{\stackrel{}{I}(z_1)}\right\}$$
(4.32)
in the limit $`s_{0,\text{ref}}0`$. By (4.29), the first term in the braces is just $`\mathrm{ln}(P/P_{\text{ref}})`$ to first order in $`\kappa `$. By (4.23) and (4.25) one thus obtains
$$\begin{array}{c}\delta m(z_1)=\frac{5}{\mathrm{ln}10}\{\frac{S^2(\stackrel{}{t}_0)}{\dot{S}(t_1)\mathrm{\Sigma }(\stackrel{}{t}_0,t_1)}_{t=\stackrel{}{t}_0}^{t=t_1}\kappa \nu _\stackrel{}{\gamma }(t)dt_{t=\stackrel{}{t}_0}^{t=t_1}\frac{\mathrm{\Sigma }(\stackrel{}{t}_0,t)}{S(t)}\kappa \delta \tau _\stackrel{}{\gamma }(t)dt\hfill \\ \hfill \frac{1}{\mathrm{\Sigma }(\stackrel{}{t}_0,t_1)}_{t=\stackrel{}{t}_0}^{t=t_1}\mathrm{\Sigma }(\stackrel{}{t}_0,t)(\kappa \nu _\stackrel{}{\gamma }(t)\frac{\mathrm{\Sigma }(\stackrel{}{t}_0,t)}{S(t)}\kappa \delta \tau _\stackrel{}{\gamma }(t))dt\}+O(\kappa ^2)\end{array}$$
(4.33)
for the correction to the apparent magnitude-redshift relation. Note that this correction depends on the perturbation $`\delta 𝐡`$ only through the functions $`\kappa \nu _\stackrel{}{\gamma }(t)`$ and $`\kappa \delta \tau _\stackrel{}{\gamma }(t)`$.
For nearby sources (4.33) gives
$$\delta m(z_1)=\frac{5}{\mathrm{ln}10}\frac{\kappa \nu _\stackrel{}{\gamma }(t_1)}{\stackrel{}{H}(t_1)}+O(z_1)+O(\kappa ^2)\text{for }z1$$
(4.34)
for the apparent magnitude-redshift correction. This evidently vanishes at $`z_1=0`$ to first order in $`\kappa `$ in the case $`\nu _\gamma (t_1)=O(\kappa )`$.
It is evident from the analysis in the present and preceding sections that the equation
$$\kappa \nu _\gamma (t_1)=0+O(\kappa ^2)$$
(4.35)
may be interpreted as a condition on the perturbation which, as regards local optical properties, ensures that the background FRW metric $`\stackrel{}{𝐠}`$ is a best fit to the perturbed metric $`𝐠`$, in the direction of $`\dot{\gamma }(t_1)`$, at the observation point $`(t_1,𝐱_1)`$. A requirement that (4.35) holds for all null geodesics $`\gamma `$ through $`(t_1,𝐱_1)`$ is, by (3.6), equivalent to the condition
$$_t\delta h_{ij}(t_1,𝐱_1)=0+O(\kappa ^2).$$
(4.36)
In general this represent a physical constraint on the perturbation at $`(t_1,𝐱_1)`$ since there is, in general, no freedom to choose a new time slicing such that the $`3`$-metric $`𝐡`$ is intrinsic to the level surfaces of the new time. However, for the specific class of perturbations to be introduced in §5, there is just such a freedom, at least to first order in $`\kappa `$, which may be exploited to ensure that (4.36) does hold.
## 5. Newtonian perturbations
In order to describe the matter distribution of the cosmos, physical considerations suggest that one seeks a perturbed $`3`$-metric $`𝐡`$ such that the corresponding space-time metric $`𝐠`$ has an energy tensor of the form
$$T^{ab}=(\rho +p)u^au^b+pg^{ab}$$
(5.1)
describing a perfect fluid with velocity $`u^a:=t^{;a}`$, density $`\rho `$ and pressure $`p`$. It is conventional to define a dimensionless density parameter by
$$\mathrm{\Omega }:=\frac{\kappa \rho }{3H^2}.$$
(5.2)
Note that all of $`\rho `$, $`H`$ and $`\mathrm{\Omega }`$ depend on all four coordinates. With regard to the pressure, since the form of (2.13) implies that $`u^a`$ is geodesic, the conservation equation for (5.1) implies $`h_a\text{}^bp_{;b}=0`$ and hence $`h_a\text{}^b(\delta p)_{;b}=0`$ under the assumption that $`\rho +p`$ is nowhere zero. This leads one to consider the particular case $`\delta p=0`$ which describes perturbations arising from the addition or removal of comoving dust.
The full non-linear problem presents formidable difficulties, although one does have by (5.1) and (2.25) that $`\delta h_{ij}`$ satisfies the simple equation
$$h^{jk}(_t\delta h\text{}_{k[i})\text{}_{\text{ }\text{ }\text{ }\text{ }\text{ }j]}=0.$$
(5.3)
In order to make progress, only a linear approximation to a solution will be sought. Specifically the problem is to obtain $`\delta h_{ij}`$, regarded as a power series in $`\kappa `$, vanishing to zeroth order, such that the energy tensor of $`g_{ab}`$ has the form
$$T^{ab}=(\rho +p)u^au^b+pg^{ab}+O(\kappa ).$$
(5.4)
Consider a perturbation of the $`3`$-metric $`\stackrel{}{h}_{ij}`$ of the form
$$\delta h_{ij}=F(t)\kappa \mathrm{\Phi }_{\text{ }\text{ }\text{ }\text{ }\text{ }ij}+G(t)\kappa \mathrm{\Phi }h_{ij}+O(\kappa ^2)$$
(5.5)
for functions $`F(t)`$, $`G(t)`$ and a $`t`$-independent scalar field $`\mathrm{\Phi }(𝐱)`$. The form of (5.5) corresponds to the synchronous gauge perturbations considered by Mukhanov et al. (1992, p.216). It will be convenient, although not necessary, to regard $`F(t)`$ and $`G(t)`$ as power series in $`\kappa `$ even though only the zeroth order terms will be of significance. Note that although $`\delta h_{ij}`$ enters into the right side of (5.5) through the term $`\mathrm{\Phi }_{\text{ }\text{ }\text{ }\text{ }\text{ }ij}`$, this is not significant in the linear approximation since $`\mathrm{\Phi }_{\text{ }\text{ }\text{ }\text{ }\text{ }ij}`$ and $`\mathrm{\Phi }_{\stackrel{}{\text{ }\text{ }\text{ }\text{ }\text{ }}ij}`$ agree to zeroth order in $`\kappa `$.
From first principles one has
$$\delta \text{}^{(3)}R_{ij}2k\delta h_{ik}=\frac{1}{2}G(t)\{\kappa \mathrm{\Phi }_{\text{ }\text{ }\text{ }\text{ }\text{ }ij}+(\kappa \text{}^{(3)}\mathrm{\Delta }\mathrm{\Phi }+4k\kappa \mathrm{\Phi })h_{ij}\}+O(\kappa ^2)$$
(5.6)
which by (2.28) gives
$$\delta \text{}^{(3)}R=2G(t)(\kappa \text{}^{(3)}\mathrm{\Delta }\mathrm{\Phi }+3k\kappa \mathrm{\Phi })+O(\kappa ^2).$$
(5.7)
Substitution of (5.5) into (2.30), (2.31) and (2.32) yields
$`\kappa T_{tt}`$ $`=\kappa \stackrel{}{\rho }+\left({\displaystyle \frac{\dot{S}(t)\dot{F}(t)}{S^2(t)}}{\displaystyle \frac{G(t)}{S^2(t)}}\right)\kappa \text{}^{(3)}\mathrm{\Delta }\mathrm{\Phi }+3\left({\displaystyle \frac{\dot{S}(t)\dot{G}(t)}{S(t)}}{\displaystyle \frac{kG(t)}{S^2(t)}}\right)\kappa \mathrm{\Phi }+O(\kappa ^2)`$ (5.8)
$`\kappa T_{it}`$ $`=\kappa T_{ti}=(k\dot{F}(t)\dot{G}(t))\kappa \mathrm{\Phi }_{\text{ }\text{ }\text{ }\text{ }\text{ }i}+O(\kappa ^2)`$ (5.9)
$`\kappa T_{ij}`$ $`=\{S^2(t)\kappa \stackrel{}{p}+{\displaystyle \frac{1}{2}}(G(t){\displaystyle \frac{1}{S(t)}}(S^3(t)\dot{F}(t))\text{.})\kappa \text{}^{(3)}\mathrm{\Delta }\mathrm{\Phi }`$
$`+(kG(t){\displaystyle \frac{1}{S(t)}}(S^3(t)\dot{G}(t))\text{.})\kappa \mathrm{\Phi }\}h_{ij}`$
$`+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{1}{S(t)}}(S^3(t)\dot{F}(t))\text{.}G(t)\right)\kappa \mathrm{\Phi }_{\text{ }\text{ }\text{ }\text{ }\text{ }ij}+O(\kappa ^2)`$ (5.10)
by means of (5.7), (2.11) and (5.6).
For $`\kappa T_{it}`$ and the trace-free part of $`\kappa T_{ij}`$ to vanish to first order in $`\kappa `$, in accordance with (5.4), it suffices to require that $`F(t)`$ and $`G(t)`$ satisfy
$`G(t)`$ $`={\displaystyle \frac{1}{S(t)}}(S^3(t)\dot{F}(t))\text{.}+O(\kappa )`$ (5.11)
$`\dot{G}(t)`$ $`=k\dot{F}(t)+O(\kappa ).`$ (5.12)
From these one obtains
$$(S(t)(S^2(t)\dot{F}(t))\text{.})\text{.}=0+O(\kappa )$$
(5.13)
by means of (2.9), and hence
$$S(t)(S^2(t)\dot{F}(t))\text{.}=C+O(\kappa )$$
(5.14)
for some constant $`C`$.
If $`C`$ were chosen to vanish to zeroth order in $`\kappa `$ then, by (5.14) and (5.11), both $`(S^2(t)\dot{F}(t))\text{.}`$ and $`G(t)\dot{S}(t)S(t)\dot{F}(t)`$ would vanish to zeroth order in $`\kappa `$ and so, as will be evident from (6.8) in §6, $`𝐠`$ would be an FRW metric to first order in $`\kappa `$ irrespective of the function $`\mathrm{\Phi }(𝐱)`$. To avoid this uninteresting case $`C`$ will be chosen to be non-zero to zeroth order in $`\kappa `$. One may then normalise $`F(t)`$ to give
$$C=1.$$
(5.15)
Equations (5.11), (5.14) and (5.15) give
$`(S^2(t)\dot{F}(t))\text{.}`$ $`={\displaystyle \frac{1}{S(t)}}+O(\kappa )`$ (5.16)
$`G(t)\dot{S}(t)S(t)\dot{F}(t)`$ $`={\displaystyle \frac{1}{S(t)}}+O(\kappa ).`$ (5.17)
These are equivalent to equations (5.11), (5.12) and (5.15) by virtue of (2.9).
By means of (5.16) and (5.17), equations (5.8), (5.9) and (5.10) combine to give that $`T^{ab}`$ has the required form (5.4) for
$`\delta \rho `$ $`={\displaystyle \frac{1}{S^3(t)}}(\text{}^{(3)}\mathrm{\Delta }\mathrm{\Phi }+3k\mathrm{\Phi })+O(\kappa )`$ (5.18)
$`\delta p`$ $`=0+O(\kappa ).`$ (5.19)
From (5.18) one has
$$_t(S^3(t)\delta \rho )=0+O(\kappa )$$
(5.20)
as one would expect since the perturbation has a dust equation of state.
From (5.5), (5.12) and (5.18) one has
$$h^{ij}_t\delta h_{ij}=S^3(t)\dot{F}(t)\kappa \delta \rho +O(\kappa ^2)$$
(5.21)
whereby (2.21), (2.22) and (2.23) give
$`\delta H`$ $`={\displaystyle \frac{1}{6}}S^3(t)\dot{F}(t)\kappa \delta \rho +O(\kappa ^2)`$ (5.22)
$`\delta q`$ $`={\displaystyle \frac{1}{2}}\{S(t)(S^2(t)\dot{F}(t))\text{.}+2S^3(t)\dot{S}(t)\dot{F}(t)(1+\stackrel{}{q}(t))\}\delta \mathrm{\Omega }+O(\kappa ^2)`$ (5.23)
$`\delta \text{}^{(3)}R`$ $`=2S^3(t)G(t)\kappa \delta \rho +O(\kappa ^2)`$ (5.24)
with the help of (5.2) and (5.17).
It is straightforward to check, by means of (2.42) and (2.11), that the perturbed space-time metric g defined by (2.13), (2.14) and (5.5) is silent to first order in $`\kappa `$ in the sense of
$$B_{ij}=0+O(\kappa ^2).$$
(5.25)
The electric part of the Weyl tensor of $`𝐠`$ is, by (2.39), (5.6), (5.7), (5.16) and (5.17), given by
$$E_{ij}=\frac{1}{2S(t)}(\kappa \mathrm{\Phi }_{\text{ }\text{ }\text{ }\text{ }\text{ }ij}\frac{1}{3}h_{ij}\kappa \text{}^{(3)}\mathrm{\Delta }\mathrm{\Phi })+O(\kappa ^2).$$
(5.26)
By (5.4) and the analogue of (4.13) for $`\stackrel{}{\gamma }`$, the function $`\delta \tau _\stackrel{}{\gamma }(t)`$ in (4.33) is given by
$$\delta \tau _\stackrel{}{\gamma }(t)=\frac{1}{2}\delta \rho (\stackrel{}{\gamma }(t))+O(\kappa ).$$
(5.27)
An equation for the function $`\nu _\stackrel{}{\gamma }(t)`$ in (4.33) will be given in §6. By (5.25), (5.26) and the analogue of (4.15) for $`\stackrel{}{\gamma }`$ the function $`\psi _\stackrel{}{\gamma }(t)`$ in equation (4.26) is given by
$$\psi _\stackrel{}{\gamma }(t)=\frac{1}{S(t)}\mathrm{\Phi }_{\stackrel{}{\text{ }\text{ }\text{ }\text{ }\text{ }}ij}\text{ }\stackrel{}{m}\text{}^i\text{ }\stackrel{}{m}\text{}^j+O(\kappa )$$
(5.28)
where $`\stackrel{}{m}\text{}^i`$ is the analogue of $`m^i`$ for $`\stackrel{}{\gamma }`$.
Space-times satisfying the condition $`B_{ab}=0`$ have been termed ‘silent’ by some authors and ‘Newtonian-like’ by others. However it is known that the vanishing of $`B_{ab}`$ can be conserved in time only in specialised cases (Maartens et al. 1998). In any second or higher order study of perturbations of FRW cosmologies one would therefore not expect $`B_{ab}`$ to vanish to any higher than first order in $`\kappa `$.
## 6. The gauge condition
The functions $`F(t)`$ and $`G(t)`$, which enter explicitly into the perturbed metric $`𝐠`$, are not determined uniquely by (5.16) and (5.17) since these equations allow two freely specifiable constants of integration. In alternative terminology, there is a gauge freedom
$`F(t)`$ $`F(t)+A+A_0{\displaystyle _{t^{}=t}^{t^{}=t_1}}{\displaystyle \frac{dt^{}}{S^2(t)}}`$ (6.1)
$`G(t)`$ $`G(t)A_0{\displaystyle \frac{\dot{S}(t)}{S(t)}}`$ (6.2)
where $`A_0`$ and $`A`$ are real constants. In order to understand the meaning of this freedom, consider the metric $`𝐠`$, as given by (2.13), (2.14) and (5.5), expressed with respect to the coordinate system employed in (2.2):
$$\begin{array}{c}𝐠=dt^2+S^2(t)\left(1+\frac{kr^2}{4}\right)^2(1+G(t)\kappa \mathrm{\Phi }(𝐱))\underset{i=1}{\overset{3}{}}(dx^i)^2\hfill \\ \hfill +S^2(t)F(t)\kappa \mathrm{\Phi }(𝐱)_{\text{ }\text{ }\text{ }\text{ }\text{ }ij}dx^idx^j+O(\kappa ^2).\end{array}$$
(6.3)
In terms of new coordinates $`(\stackrel{~}{t},\stackrel{~}{𝐱})`$ defined by
$`\stackrel{~}{t}`$ $`:=t+{\displaystyle \frac{1}{2}}F\text{*}(t)\kappa \mathrm{\Phi }(𝐱)`$ (6.4)
$`\stackrel{~}{𝐱}\text{}^i`$ $`:=x^i+{\displaystyle \frac{1}{2}}F(t)\left(1+{\displaystyle \frac{kr^2}{4}}\right)^2\kappa _i\mathrm{\Phi }(𝐱),`$ (6.5)
for an as yet unspecified function $`F\text{*}(t)`$, the metric $`𝐠`$ assumes the form
$$\begin{array}{c}𝐠=(1\dot{F}\text{*}(\stackrel{~}{t})\kappa \mathrm{\Phi }(\stackrel{~}{𝐱}))d\stackrel{~}{t}\text{}^2+(F\text{*}(\stackrel{~}{t})S^2(\stackrel{~}{t})\dot{F}(\stackrel{~}{t}))\kappa _{\stackrel{~}{𝐱}\text{}^i}\mathrm{\Phi }(\stackrel{~}{𝐱})d\stackrel{~}{x}\text{}^id\stackrel{~}{t}\hfill \\ \hfill +\frac{S^2(\stackrel{~}{t})}{\left(1+\frac{k\stackrel{~}{r}\text{}^2}{4}\right)^2}\left(1+(G(\stackrel{~}{t})\frac{\dot{S}(\stackrel{~}{t})}{S(\stackrel{~}{t})}F\text{*}(\stackrel{~}{t}))\kappa \mathrm{\Phi }(\stackrel{~}{𝐱})\right)\underset{i=1}{\overset{3}{}}(d\stackrel{~}{x}\text{}^i)^2+O(\kappa ^2).\end{array}$$
(6.6)
With the choice
$$F\text{*}(t)=S^2(t)\dot{F}(t)$$
(6.7)
the metric $`𝐠`$ assumes the diagonal form
$$\begin{array}{c}𝐠=(1(S^2(\stackrel{~}{t})\dot{F}(\stackrel{~}{t}))\text{.}\kappa \mathrm{\Phi }(\stackrel{~}{𝐱}))d\stackrel{~}{t}\text{}^2\hfill \\ \hfill +\frac{S^2(\stackrel{~}{t})}{\left(1+\frac{k\stackrel{~}{r}\text{}^2}{4}\right)^2}(1+(G(\stackrel{~}{t})\dot{S}(\stackrel{~}{t})S(\stackrel{~}{t})\dot{F}(\stackrel{~}{t}))\kappa \mathrm{\Phi }(\stackrel{~}{𝐱}))\underset{i=1}{\overset{3}{}}(d\stackrel{~}{x}\text{}^i)^2+O(\kappa ^2).\end{array}$$
(6.8)
No assumptions about the $`t`$-dependence of $`F(t)`$ and $`G(t)`$ have yet been employed. It is evident that if $`(S^2(t)\dot{F}(t))\text{.}`$ and $`G(t)\dot{S}(t)S(t)\dot{F}(t)`$ were both to vanish to zeroth order in $`\kappa `$ then, as quoted in §5, $`𝐠`$ would be an FRW metric to first order in $`\kappa `$.
For $`F(t)`$ and $`G(t)`$ satisfying (5.16) and (5.17), equation (6.8) reduces to
$$𝐠=\left(1\frac{\kappa \mathrm{\Phi }(\stackrel{~}{𝐱})}{S(\stackrel{~}{t})}\right)d\stackrel{~}{t}\text{}^2+\frac{S^2(\stackrel{~}{t})}{(1+\frac{k\stackrel{~}{r}\text{}^2}{4})^2}\left(1+\frac{\kappa \mathrm{\Phi }(\stackrel{~}{𝐱})}{S(\stackrel{~}{t})}\right)\underset{i=1}{\overset{3}{}}(d\stackrel{~}{x}\text{}^i)^2+O(\kappa ^2)$$
(6.9)
which is the metric of Newman & McVittie (1982). Note that the curves $`\{\stackrel{~}{𝐱}=\text{const.}\}`$ have unit tangent
$$\stackrel{~}{u}_a=\left(1+\frac{\kappa \mathrm{\Phi }(\stackrel{~}{𝐱})}{2S(\stackrel{~}{t})}\right)\stackrel{~}{t}_{\stackrel{~}{;}a}+O(\kappa ^2)$$
(6.10)
and so are non-geodesic to first order in $`\kappa `$. They are therefore not the world lines of freely falling observers.
The absence of $`F(t)`$ and $`G(t)`$ in (6.9) shows that the isometry class of $`𝐠`$ is unaffected by the choice of integration constants for (5.16) and (5.17). Indeed the gauge transformation (6.1), (6.2) is induced by the coordinate transformation
$`t`$ $`t+{\displaystyle \frac{A_0}{2}}\kappa \mathrm{\Phi }(𝐱)`$ (6.11)
$`x^i`$ $`x^i{\displaystyle \frac{1}{2}}\left(A+A_0{\displaystyle _{t^{}=t}^{t^{}=t_1}}{\displaystyle \frac{dt^{}}{S^2(t^{})}}\right)\left(1+{\displaystyle \frac{kr^2}{4}}\right)^2\kappa _i\mathrm{\Phi }(𝐱)`$ (6.12)
which preserves the form of the metric (6.3) for $`F(t)`$, $`G(t)`$ transforming according to (6.1) and (6.2). Moreover the transformation (6.11) of $`t`$ preserves the form of (5.5) on the level surfaces of $`t`$, for a given $`\mathrm{\Phi }`$ on the space-time manifold, for $`F(t)`$, $`G(t)`$ transforming according to (6.1) and (6.2). The gauge freedom (6.1), (6.2) is thus to be interpreted as a freedom to choose the time function $`t`$ without violating the form of (5.5). The coordinate freedom (6.11), (6.12) corresponds to that identified by Mukhanov et al. (1992, p.216).
Although the integration constants for (5.16) and (5.17) do not affect the isometry class of $`𝐠`$ to first order in $`\kappa `$, they nonetheless carry physical significance since they help determine the physically significant coordinate $`t`$. In order to make physically appropriate choices of these constants, consider the function $`\nu _\gamma (t)`$ of (3.6) along a light ray $`\gamma `$ through the observation point $`(t_1,𝐱_1)`$. By means of (5.5), (5.12) and the null geodesic equations (3.3) and (3.4) one obtains
$$\nu _\gamma (t)=\dot{F}(t)\{S(t)\frac{d}{dt}\left(S(t)\frac{d}{dt}\mathrm{\Phi }(\gamma (t))\right)+\mathrm{\Phi }(\gamma (t))\}+O(\kappa ).$$
(6.13)
By (5.5) and (5.12) the gauge condition (4.36) becomes
$$\dot{F}(t_1)=0.$$
(6.14)
This holds for a unique value of $`A_0`$ in (6.1), (6.2).
Equations (5.16), (5.17) and (6.14) give
$`\dot{F}(t)`$ $`={\displaystyle \frac{1}{S^2(t)}}{\displaystyle _{t^{}=t}^{t^{}=t_1}}{\displaystyle \frac{dt^{}}{S(t^{})}}+O(\kappa )`$ (6.15)
$`G(t)`$ $`={\displaystyle \frac{1}{S(t)}}{\displaystyle \frac{\dot{S}(t)}{S(t)}}{\displaystyle _{t^{}=t}^{t^{}=t_1}}{\displaystyle \frac{dt^{}}{S(t^{})}}+O(\kappa ).`$ (6.16)
Thus, to zeroth order in $`\kappa `$, $`G(t)`$ is now uniquely specified whilst $`F(t)`$ is determined only up to an arbitrary additive constant. By (6.11) and (6.12) this remaining gauge freedom corresponds to a coordinate transformation of the form
$`t`$ $`t`$ (6.17)
$`x^i`$ $`x^i{\displaystyle \frac{A}{2}}\left(1+{\displaystyle \frac{kr^2}{4}}\right)^2\kappa _i\mathrm{\Phi }(𝐱)`$ (6.18)
where $`A`$ is an arbitrary constant. Since this preserves the level surfaces of $`t`$ it has no significance for the physical properties of the perturbed cosmology. There is therefore no need to specify the remaining integration constant of (5.16) and (5.17).
By (4.34) the gauge fixing condition (6.14) may be interpreted as a necessary and sufficient condition that there is a vanishing correction to the apparent magnitude-redshift relation, in all directions, to zeroth order in the redshift and first order in $`\kappa `$. By (4.28) another interpretation is that there is a vanishing correction to the luminosity distance-redshift relation, in all directions, to zeroth order in the redshift and first order in $`\kappa `$.
An immediate consequence of the gauge fixing condition (6.14) is, by (5.22), that the Hubble expansion $`H`$ of the perturbed cosmology satisfies
$$H(t_1,𝐱)=\stackrel{}{H}(t_1)+O(\kappa ^2)$$
(6.19)
and so is unperturbed to first order in $`\kappa `$ at all points of the surface $`\{t=t_1\}`$. By the use of (5.14) and (6.14) in (5.23) one also has
$$\delta q=\frac{1}{2}\delta \mathrm{\Omega }+O(\kappa ^2).$$
(6.20)
The gauge condition (6.14) thus ensures that the perturbation (5.5) affects the deceleration parameter but not the Hubble parameter at the observation point $`(t_1,𝐱_1)`$.
Subject to the gauge fixing condition (6.14) the apparent magnitude-redshift relation is given by (4.33), with $`\nu _\stackrel{}{\gamma }(t)`$ given by (6.13) and (6.15) and $`\delta \tau _\stackrel{}{\gamma }(t)`$ by (5.27). Substitution for these quantities in (4.33) yields
$`\delta m(z_1)`$ $`=`$
$`{\displaystyle \frac{5}{\mathrm{ln}10}}\{{\displaystyle \frac{S^4(\stackrel{}{t}_0)\dot{F}(\stackrel{}{t}_0)}{2\dot{S}(t_1)\mathrm{\Sigma }(\stackrel{}{t}_0,t_1)}}{\displaystyle \frac{d}{dt}}|_{t=\stackrel{}{t}_0}\kappa \mathrm{\Phi }(\stackrel{}{\gamma }(t))+{\displaystyle \frac{S^2(\stackrel{}{t}_0)}{\dot{S}(t_1)\mathrm{\Sigma }(\stackrel{}{t}_0,t_1)}}{\displaystyle _{t=\stackrel{}{t}_0}^{t=t_1}}{\displaystyle \frac{\dot{S}(t)}{S^2(t)}}\kappa \mathrm{\Phi }(\stackrel{}{\gamma }(t))dt.`$
$`+{\displaystyle \frac{S^2(\stackrel{}{t}_0)}{\dot{S}(t_1)\mathrm{\Sigma }(\stackrel{}{t}_0,t_1)}}\left[\left({\displaystyle \frac{1}{S(t)}}{\displaystyle \frac{1}{2}}G(t)\right)\kappa \mathrm{\Phi }(\stackrel{}{\gamma }(t))\right]_{t=\stackrel{}{t}_0}^{t=t_1}+{\displaystyle \frac{\kappa \mathrm{\Phi }(\stackrel{}{\gamma }(t_1))}{2S(t_1)}}`$
$`{\displaystyle \frac{1}{\mathrm{\Sigma }(\stackrel{}{t}_0,t_1)}}{\displaystyle _{t=\stackrel{}{t}_0}^{t=t_1}}(1{\displaystyle \frac{\dot{S}(t)}{S^2(t)}}\mathrm{\Sigma }(\stackrel{}{t}_0,t))\kappa \mathrm{\Phi }(\stackrel{}{\gamma }(t))𝑑t{\displaystyle \frac{S^3(\stackrel{}{t}_0)\dot{F}(\stackrel{}{t}_0)}{2\mathrm{\Sigma }(\stackrel{}{t}_0,t_1)}}\kappa \mathrm{\Phi }(\stackrel{}{\gamma }(\stackrel{}{t}_0))`$
$`.{\displaystyle _{t=\stackrel{}{t}_0}^{t=t_1}}{\displaystyle \frac{\mathrm{\Sigma }(\stackrel{}{t}_0,t)}{2S(t)}}\kappa \delta \rho (\stackrel{}{\gamma }(t))dt+{\displaystyle \frac{1}{\mathrm{\Sigma }(\stackrel{}{t}_0,t_1)}}{\displaystyle _{t=\stackrel{}{t}_0}^{t=t_1}}{\displaystyle \frac{\mathrm{\Sigma }^2(\stackrel{}{t}_0,t)}{2S(t)}}\kappa \delta \rho (\stackrel{}{\gamma }(t))dt\}+O(\kappa ^2)`$ (6.21)
by successive integration by parts and the use of (5.11), (5.12), (5.16), (5.17) and (6.14). Again $`\stackrel{}{t}_0`$ is given in terms of $`z_1`$ and $`t_1`$ by means of (3.1). Equation (6.21) is the correction to the apparent magnitude redshift relation for Newtonian perturbations of FRW cosmologies subject to the gauge fixing condition (6.14).
## 7. The averaging procedure
Suppose $`\mathrm{\Phi }`$ is smooth on each level surface of $`t`$. Let $`𝒟_t`$ be a $`3`$-domain in one such surface. By definition, the mean mass density perturbation on $`𝒟_t`$ is
$$\delta \rho _{𝒟_t}:=\frac{_{𝒟_t}\delta \rho 𝑑\stackrel{}{v}}{_{𝒟_t}𝑑\stackrel{}{v}}$$
(7.1)
where $`d\stackrel{}{v}`$ is the elemental $`3`$-volume on $`𝒟_t`$ with respect to $`\stackrel{}{𝐡}`$. By means of (7.1) and (5.18) one obtains
$$\delta \rho _{𝒟_t}=\frac{3k}{S^3(t)}\mathrm{\Phi }_{𝒟_t}\frac{_{𝒟_t}\text{}^{(3)}_\stackrel{}{\text{n}}\mathrm{\Phi }d\stackrel{}{v}}{S^3(t)_{𝒟_t}𝑑\stackrel{}{v}}+O(\kappa )$$
(7.2)
by means of the divergence theorem, where $`\stackrel{}{𝐧}`$ is the unit outward pointing normal to $`𝒟_t`$ at $`𝒟_t`$ with respect to $`\stackrel{}{𝐡}`$ in the surface $`\{t=\text{ }\text{const.}\}`$, and where
$$\mathrm{\Phi }_{𝒟_t}:=\frac{_{𝒟_t}\mathrm{\Phi }𝑑\stackrel{}{v}}{_{𝒟_t}𝑑\stackrel{}{v}}$$
(7.3)
is the mean value of $`\mathrm{\Phi }`$ on $`𝒟_t`$. Note that $`\mathrm{\Phi }_{𝒟_t}`$ is independent of $`t`$. It will be assumed that the perturbed matter distribution is sufficiently uniform on the large scale that, in the limit $`𝒟_t\mathrm{}`$, the surface integral term in (7.2) tends to zero, whilst $`\delta \rho _{𝒟_t}`$ and $`\mathrm{\Phi }_{𝒟_t}`$ tend to limits $`\delta \rho _t`$ and $`\mathrm{\Phi }`$ respectively. Under these conditions one obtains
$$\delta \rho _t=\frac{3k}{S^3(t)}\mathrm{\Phi }+O(\kappa ).$$
(7.4)
In terms of the dimensionless density parameter $`\mathrm{\Omega }`$ of (5.2) this gives
$$\delta \mathrm{\Omega }_t=\frac{\kappa \delta \rho _t}{3\stackrel{}{H}\text{}^2(t)}+O(\kappa ^2)=(1+\stackrel{}{q}(t))\frac{\kappa \mathrm{\Phi }}{S(t)}+O(\kappa ^2)$$
(7.5)
by means of (2.9). Equation (7.4) and the right side of (7.5) must be considered invalid for $`k=0`$ since in that case $`\mathrm{\Phi }`$ must be infinite to give a finite $`\delta \rho _t`$.
In order to obtain the mean correction to the apparent magnitude-redshift relation, one replaces the function $`\mathrm{\Phi }(\gamma (t))`$ on the right of (6.21) with its mean value $`\mathrm{\Phi }`$, and $`\frac{d}{dt}|_{t=\stackrel{}{t}_0}\mathrm{\Phi }(\gamma (t))`$ by its expectation value of zero. In practice, since all cosmological sources lie within gravitational potential wells, the $`\frac{d}{dt}|_{t=\stackrel{}{t}_0}\mathrm{\Phi }(\gamma (t))`$ term in (6.21) will always give a positive contribution to the apparent magnitude, but this will be so small as to be negligible for present purposes. By means of (5.20), (7.4) and (2.8), and for $`\stackrel{}{t}_0`$ given in terms of $`t_1`$ and $`z_1`$ by (3.1), one thus obtains
$`\delta m(z_1)`$ $`=`$
$`{\displaystyle \frac{5}{2\mathrm{ln}10}}S^3(t_1)\stackrel{}{H}\text{}^2(t_1)\delta \mathrm{\Omega }_{t_1}\{{\displaystyle \frac{S^2(\stackrel{}{t}_0)}{\dot{S}(t_1)\mathrm{\Sigma }(\stackrel{}{t}_0,t_1)}}{\displaystyle _{t=\stackrel{}{t}_0}^{t=t_1}}\dot{F}(t)dt+3{\displaystyle _{t=\stackrel{}{t}_0}^{t=t_1}}{\displaystyle \frac{\mathrm{\Sigma }(\stackrel{}{t}_0,t)}{S^4(t)}}dt`$
$`{\displaystyle \frac{1}{\mathrm{\Sigma }(\stackrel{}{t}_0,t_1)}}{\displaystyle _{t=\stackrel{}{t}_0}^{t=t_1}}\mathrm{\Sigma }(\stackrel{}{t}_0,t)\dot{F}(t)dt{\displaystyle \frac{3}{\mathrm{\Sigma }(\stackrel{}{t}_0,t_1)}}{\displaystyle _{t=\stackrel{}{t}_0}^{t=t_1}}{\displaystyle \frac{\mathrm{\Sigma }^2(\stackrel{}{t}_0,t)}{S^4(t)}}dt\}+O(\kappa ^2)`$ (7.6)
for the mean correction to the apparent magnitude-redshift relation for Newtonian perturbations of FRW cosmologies. This is the fundamental equation of the paper. It could alternatively have been derived directly from (4.33) with the help of the observation that the first term in the braces on the right of (6.13) gives a zero contribution to $`\delta m(z_1)`$. For nearby sources (7.6) gives
$$\delta m(z_1)=\frac{5}{4\mathrm{ln}10}\delta \mathrm{\Omega }_{t_1}z_1+O(z_1^2)+O(\kappa ^2)$$
(7.7)
by means of (3.13) and (7.5).
In the particular case of a zero mean density perturbation, $`\delta \rho _t=0`$, equation (7.6) clearly gives that there is a zero mean correction to the apparent magnitude-redshift relation to first order in $`\kappa `$. In the case of a non-zero mean density perturbation $`\delta \rho _t0`$, one may compare either with the background metric $`\stackrel{}{𝐠}`$ or with that of the homogenised cosmology determined from $`𝐠`$ by the constant potential $`\mathrm{\Phi }`$. (See §8$`a`$.) In the latter case the mean perturbation of the potential is zero, so the mean density perturbation is also zero. Hence the correction to the apparent magnitude-redshift relation is zero, to first order in $`\kappa `$.
It is worthwhile to consider the case $`\delta \rho _t>0`$ in more detail. Of the four terms in the braces on the right of (7.6), the first three are positive. And though the fourth term is negative, it is evidently dominated by the second. The sum of the terms in the braces is therefore positive. This shows that an object at a given redshift appears brighter than in the reference FRW model. Moreover an object at a given redshift in a high density universe appears brighter than in a low density universe. These effects can be regarded as the gravitational lensing of the universe as a whole.
A fruitful approach to clumped matter perturbations of FRW cosmologies, introduced by Dyer & Roeder (1973), is to assume that of all the matter present, a proportion $`\alpha `$ is uniformly distributed and pressure-free, while the remaining proportion $`1\alpha `$ is gravitationally bound into clumps. In order to compute the lensing of light beams that remain far from all clumps it suffices to take into consideration the effect of only the uniformly distributed matter. In particular, shear can be neglected. The angular diameter distance (Schneider et al. 1992, eq. (3.66)) determined from such clump-avoiding light beams is often called the Dyer-Roeder distance. For light beams of large angular diameter, one expects the cumulative lensing effect to approach that of the homogenised FRW cosmology. The transition between these two regimes has been studied by Linder (1998) and found to occur typically between $`1`$ and $`10`$ arcseconds.
The Dyer-Roeder method may be implemented in the present framework as follows. The background metric $`\stackrel{}{𝐠}`$ is taken to be the metric of the homogenised FRW model. The perturbation $`\delta 𝐡`$ then corresponds to the removal of a proportion $`\alpha `$ of the matter of this model, followed by the addition of the same quantity of clumped matter. The apparent magnitude-redshift relation for narrow, clump-avoiding light beams is then given by (7.6), with $`\delta \mathrm{\Omega }_{t_1}`$ replaced by $`\alpha \mathrm{\Omega }_{t_1}`$. By means of (2.8) this yields
$`\delta m(z_1)`$ $`=`$
$`{\displaystyle \frac{5}{2\mathrm{ln}10}}S^3(t_1)\stackrel{}{H}\text{}^2(t_1)\alpha \delta \mathrm{\Omega }_{t_1}\{{\displaystyle \frac{S^2(\stackrel{}{t}_0)}{\dot{S}(t_1)\mathrm{\Sigma }(\stackrel{}{t}_0,t_1)}}{\displaystyle _{t=\stackrel{}{t}_0}^{t=t_1}}\dot{F}(t)dt+3{\displaystyle _{t=\stackrel{}{t}_0}^{t=t_1}}{\displaystyle \frac{\mathrm{\Sigma }(\stackrel{}{t}_0,t)}{S^4(t)}}dt`$
$`{\displaystyle \frac{1}{\mathrm{\Sigma }(\stackrel{}{t}_0,t_1)}}{\displaystyle _{t=\stackrel{}{t}_0}^{t=t_1}}\mathrm{\Sigma }(\stackrel{}{t}_0,t)\dot{F}(t)dt{\displaystyle \frac{3}{\mathrm{\Sigma }(\stackrel{}{t}_0,t_1)}}{\displaystyle _{t=\stackrel{}{t}_0}^{t=t_1}}{\displaystyle \frac{\mathrm{\Sigma }^2(\stackrel{}{t}_0,t)}{S^4(t)}}dt\}+O(\kappa ^2).`$ (7.8)
Since the expression in the braces is positive for $`z_1>0`$, this correction is positive for $`z_1>0`$. Sources viewed at a given redshift along clump-avoiding light beams thus appear dimmer than for an all-sky average. This is as one would expect since clump-avoiding light beams are less focussed. For larger angular scales, the mean correction is no longer given by (7.8) and should be expected to tend to zero since the mean density perturbation is zero. Thus (7.8) may be interpreted as the mean correction to the apparent magnitude-redshift relation for narrow, clump-avoiding light beams relative to that for wide angle beams.
The standard implementation of the Dyer-Roeder ansatz (e.g. Schneider et al. 1992, p.138 et seq.) gives the Dyer-Roeder distance as a solution to the Dyer-Roeder equation which describes light propagation through a uniformly underdense region of space-time. From the Dyer-Roeder distance one can obtain an apparent magnitude-redshift relation. On the other hand, an apparent magnitude-redshift relation for light propagation through a uniformly underdense region of space-time is also described by (7.8) for an appropriately valued constant $`\mathrm{\Phi }`$ (see 8.8). Nonetheless one cannot expect these two relations to agree unless similar gauge fixing conditions are applied in each case (see (8.3).
## 8. Examples
###### Example 8.1 (Uniform density perturbations).
In the special case
$$\mathrm{\Phi }=\text{const.}$$
(8.1)
the perturbation (5.5) reduces to
$$\delta h_{ij}=\kappa G(t)\mathrm{\Phi }\stackrel{}{h}_{ij}+O(\kappa ^2).$$
(8.2)
To first order in $`\kappa `$ this is equivalent to leaving the $`3`$-metric $`\stackrel{}{𝐡}`$ in (2.1) fixed and perturbing the scale factor $`S(t)`$ according to $`S(t)(1+\frac{1}{2}G(t)\kappa \mathrm{\Phi })S(t)`$. The coordinate freedom (6.11), (6.12) reduces to the freedom to change $`t`$ by an additive constant. One sees directly from (8.2) that the gauge condition (4.36) is satisfied iff $`G(t)`$ satisfies
$$\dot{G}(t_1)=0.$$
(8.3)
From (5.8), (5.9) and (5.10) one has that the perturbation is pressure-free to first order in $`\kappa `$, and so of a form as discussed in §5, if $`G(t)`$ satisfies
$$kG(t)=\frac{1}{S(t)}(S^3(t)\dot{G}(t))\text{.}+O(\kappa ).$$
(8.4)
If $`G(t)`$ were to vanish at $`t=t_1`$ then, by (8.3) and (8.4), $`G(t)`$ would vanish for all $`t`$. For consistency with (6.16) one may choose
$$G(t_1)=\frac{1}{S(t_1)}.$$
(8.5)
By (5.18) and (5.19) one has
$`\delta \rho `$ $`={\displaystyle \frac{3k}{S^3(t)}}\mathrm{\Phi }+O(\kappa )`$ (8.6)
$`\delta p`$ $`=0+O(\kappa ).`$ (8.7)
By (7.5) the perturbation of the dimensionless density parameter is given by
$$\delta \mathrm{\Omega }(t)=\frac{1}{2}(1+q(t))\frac{\kappa \mathrm{\Phi }}{S(t)}+O(\kappa ^2).$$
(8.8)
###### Example 8.2 (Point particle perturbations).
A perturbation by the introduction of a family of comoving point particles with world lines $`\{𝐱=𝐱_\alpha \}`$, $`\alpha =1,2,\mathrm{}`$ and masses $`m_\alpha >0`$ is described by a density perturbation of the form
$$\delta \rho (t,𝐱)=\frac{1}{S^3(t)}\underset{\alpha }{}m_\alpha \text{}^{(3)}\stackrel{}{\delta }_{𝐱_\alpha }(𝐱)$$
(8.9)
where $`\text{}^{(3)}\stackrel{}{\delta }_{𝐱_\alpha }(𝐱)`$ is the Dirac distribution with respect to $`\stackrel{}{𝐡}`$, centred on $`𝐱_\alpha `$, on each level surface of $`t`$. The particles shall represent the galaxies.
In order to satisfy (8.9) and (5.18) one seeks a potential $`\mathrm{\Phi }(𝐱)`$ satisfying
$$(\text{}^{(3)}\stackrel{}{\mathrm{\Delta }}\mathrm{\Phi }(𝐱)+3k\mathrm{\Phi }(𝐱))=\underset{\alpha }{}m_\alpha \text{}^{(3)}\stackrel{}{\delta }_{𝐱_\alpha }(𝐱).$$
(8.10)
Since this equation is linear one may decompose $`\mathrm{\Phi }(𝐱)`$ as a sum
$$\mathrm{\Phi }(𝐱)=\underset{\alpha }{}\mathrm{\Phi }_\alpha (𝐱).$$
(8.11)
In order to consider a typical summand $`\mathrm{\Phi }_\alpha (𝐱)`$ it is convenient to express the $`3`$-metric $`\stackrel{}{𝐡}`$ in the form
$$\stackrel{}{𝐡}=\{\begin{array}{cc}\frac{\mathrm{𝟏}}{𝐤}(𝐝\omega _\alpha ^\mathrm{𝟐}+\mathrm{sin}^\mathrm{𝟐}\omega _\alpha 𝐝𝛀_\alpha ^\mathrm{𝟐})\hfill & \text{if }k>0\hfill \\ \mathrm{𝐝𝐫}_\alpha ^\mathrm{𝟐}+𝐫_\alpha ^\mathrm{𝟐}𝐝𝛀_\alpha ^\mathrm{𝟐}\hfill & \text{if }k=0\hfill \\ \frac{\mathrm{𝟏}}{(𝐤)}(𝐝\omega _\alpha ^\mathrm{𝟐}+\mathrm{sinh}^\mathrm{𝟐}\omega _\alpha 𝐝𝛀_\alpha ^\mathrm{𝟐})\hfill & \text{if }k<0\hfill \end{array}$$
(8.12)
where $`d\mathrm{\Omega }_\alpha ^2`$ is the $`2`$-sphere metric and the radial coordinate $`\omega _\alpha `$ is defined by
$$\omega _\alpha (r_\alpha ):=\{\begin{array}{cc}2\mathrm{tan}^1\left(\frac{\sqrt{k}}{2}r_\alpha \right)\hfill & \text{if }k>0\hfill \\ 2\mathrm{tanh}^1\left(\frac{\sqrt{k}}{2}r_\alpha \right)\hfill & \text{if }k<0\hfill \end{array}$$
(8.13)
with the origin $`\omega _\alpha =0`$ being the point $`𝐱_\alpha `$ of (8.9). The range of $`\omega _\alpha `$ is $`0\omega _\alpha \pi `$ if $`k>0`$ and $`0\omega _\alpha <\mathrm{}`$ if $`k<0`$.
The general radial solution to (8.10), as found by Newman & McVittie (1982), is
$$\mathrm{\Phi }_\alpha (𝐱)=\{\begin{array}{cc}\frac{\sqrt{k}}{4\pi \mathrm{sin}\omega _\alpha }(m_\alpha \mathrm{cos}2\omega _\alpha +C_\alpha \mathrm{sin}2\omega _\alpha )\hfill & \text{if }k>0\hfill \\ \frac{m_\alpha }{4\pi r}+C_\alpha \hfill & \text{if }k=0\hfill \\ \frac{\sqrt{k}}{4\pi \mathrm{sinh}\omega _\alpha }(m_\alpha \mathrm{cosh}2\omega _\alpha +C_\alpha \mathrm{sinh}2\omega _\alpha )\hfill & \text{if }k<0\hfill \end{array}$$
(8.14)
where $`C_\alpha `$ is an arbitrary constant.
The case $`k>0`$ would appear to be the simplest insofar as the level surfaces of $`t`$ are compact, so physical plausibility demands that there are at most finitely many particles present. However it is evident from (8.14) in this case that for each particle there is a complementary particle of equal mass located at the antipodal point in each surface $`\{t=\text{const.}\}`$. This bizarre doppelgänger phenomenon leads one to question whether the $`k>0`$ solution is realistic after all. It is unclear whether the problem is an artifact of the symmetry of the level surfaces of $`t`$ or of the linear approximation. There may even be a deeper issue here concerning the constraint components of the Einstein equations (D’Eath 1976).
For $`k>0`$ an integration of (8.10) over a level surface of $`t`$ with respect to the volume element $`d\stackrel{}{v}`$ associated with $`\stackrel{}{𝐡}`$ yields
$$3k_{t=\text{const.}}\mathrm{\Phi }𝑑\stackrel{}{v}=\underset{\alpha }{}m_\alpha .$$
(8.15)
Setting $`\text{}^{(3)}\stackrel{}{V}:=_{t=\text{const.}}𝑑\stackrel{}{v}`$ one thus obtains
$$\delta \rho _t:=\frac{_\alpha m_\alpha }{S^3(t)\text{}^{(3)}\stackrel{}{V}}=\frac{3k}{S^3(t)}\mathrm{\Phi }$$
(8.16)
for
$$\mathrm{\Phi }=\frac{_{t=\text{const.}}\mathrm{\Phi }𝑑\stackrel{}{v}}{_{t=\text{const.}}𝑑\stackrel{}{v}}.$$
(8.17)
Note that (8.16) agrees with (7.4) even though $`\mathrm{\Phi }`$ is not smooth in the present case.
In the case $`k<0`$ case the surfaces $`\{t=\text{const.}\}`$ have infinite volume, so infinitely many particles (or none) are needed in order to achieve a distribution that is uniform on the large scale. For each $`\alpha `$ one must choose $`C_\alpha =m_\alpha `$ in order that $`\mathrm{\Phi }_\alpha `$ decays to zero at infinity. Each $`\mathrm{\Phi }_\alpha (𝐱)`$ is then integrable on the level surfaces $`\{t=\text{const.}\}`$ and a simple reciprocity argument indicates that for a sufficiently uniform distribution of particles $`\mathrm{\Phi }(𝐱)=_\alpha \mathrm{\Phi }_\alpha (𝐱)`$ should converge everywhere other than on the world lines of the particles. Let $`𝒟_t`$ be a compact $`3`$-domain in a level surface of $`t`$. An integration of (8.10) over $`𝒟_t`$ with respect to $`d\stackrel{}{v}`$ yields
$$3k_{𝒟_t}\mathrm{\Phi }𝑑\stackrel{}{v}=\underset{\alpha }{}m_\alpha +_{𝒟_t}\stackrel{}{}_\stackrel{}{𝐧}\mathrm{\Phi }𝑑\stackrel{}{a}$$
(8.18)
where $`\stackrel{}{𝐧}`$ is the outward unit normal at $`𝒟_t`$, $`d\stackrel{}{a}`$ is the area element on $`𝒟_t`$ induced by $`\stackrel{}{𝐡}`$, and where, in the first term on the right, the sum is carried over all $`\alpha `$ such that $`𝐱_\alpha 𝒟_t`$. Setting $`\stackrel{}{V}(𝒟_t):=_{𝒟_t}𝑑\stackrel{}{v}`$ one obtains
$$\delta \rho _{𝒟_t}=\frac{3k}{S^3(t)}\mathrm{\Phi }_{𝒟_t}\frac{_{𝒟_t}\stackrel{}{}_\stackrel{}{𝐧}\mathrm{\Phi }𝑑\stackrel{}{a}}{S^3(t)\text{}^{(3)}\stackrel{}{V}(𝒟_t)}$$
(8.19)
where
$$\delta \rho _{𝒟_t}:=\frac{_\alpha m_\alpha }{S^3(t)\text{}^{(3)}\stackrel{}{V}(𝒟_t)}$$
(8.20)
is the mean density perturbation on $`𝒟_t`$ and
$$\mathrm{\Phi }_{𝒟_t}:=\frac{_{𝒟_t}\mathrm{\Phi }𝑑\stackrel{}{v}}{\text{}^{(3)}\stackrel{}{V}(𝒟_t)}$$
(8.21)
is the mean value of $`\mathrm{\Phi }`$ on $`𝒟_t`$. Reasonable uniformity conditions on the distribution of the particles should ensure that $`\delta \rho _{𝒟_t}`$ and $`\mathrm{\Phi }_{𝒟_t}`$ tend to limits $`\delta \rho _t`$ and $`\mathrm{\Phi }`$ respectively for arbitrarily large $`𝒟_t`$. Such conditions should also ensure that the second term on the right of (8.19) tends to zero for large $`𝒟_t`$. One then obtains (7.4) for $`k<0`$, again even though $`\mathrm{\Phi }`$ is not smooth.
In the $`k=0`$ case $`\mathrm{\Phi }_\alpha (r_\alpha )`$, as given by (8.14), decays to zero at infinity only if $`C_\alpha =0`$, and then only as $`1/r_\alpha `$. For a uniform distribution of particles the potential $`\mathrm{\Phi }=_\alpha \mathrm{\Phi }_\alpha `$ would be infinite everywhere, so the theory breaks down in this case.
For $`k0`$ the perturbation of the space-time metric corresponding to the introduction of the uniform distribution of comoving point particles is described by (2.13), (5.5), (8.11) and (8.14), and the mean correction to the apparent magnitude-redshift relation, relative to the background metric, is given by (7.6).
###### Example 8.3 (Swiss cheese model).
This model, proposed by Einstein & Strauss (1945,1946), is an exact $`C^1`$ solution to the Einstein equations consisting of a pressure-free FRW model in which spherical regions are replaced by spherical pieces of Schwarzschild geometry. The idea is that each of these spherical ‘holes’ represents the condensation of dust into a star, represented by the singularity at the centre. For present purposes, the singularities will be considered to represent the galaxies and the intervening dust, the ‘cheese’, will represent the intergalactic medium. The model will be considered here in terms of the perturbative formalism of §5.
Let $`\rho (t)`$ be the density of the intergalactic dust. In order to describe a spherically symmetric hole with a central point particle of mass $`m`$ and coordinate radius $`\omega =\widehat{\omega }`$ if $`k0`$, or $`r=\widehat{r}`$ if $`k=0`$, one may seek a radial potential function $`\mathrm{\Phi }(𝐱)`$ which satisfies
$$(\text{}^{(3)}\mathrm{\Delta }\mathrm{\Phi }(𝐱)+3k\mathrm{\Phi }(𝐱))=S^3(t_1)\rho (t_1)+m\text{}^{(3)}\delta _0(𝐱)$$
(8.22)
in the hole and matches in a $`C^2`$ manner to a constant potential $`\mathrm{\Phi }(𝐱)=\widehat{\mathrm{\Phi }}`$ outside the hole. The first term on the right of (8.22) ensures that the hole is a vacuum. The general radial solution to (8.22) has the form
$$\mathrm{\Phi }(𝐱)=\{\begin{array}{cc}\frac{\sqrt{k}m}{4\pi }\frac{\mathrm{cos}2\omega }{\mathrm{sin}\omega }+\frac{C}{8\pi }\frac{\mathrm{sin}2\omega }{\mathrm{sin}\omega }+\frac{S^3\left(t_1\right)\rho \left(t_1\right)}{3k}\hfill & \text{if }k>0\hfill \\ \frac{m}{4\pi r}+\frac{C}{4\pi }+\frac{1}{6}S^3(t_1)\rho (t_1)r^2\hfill & \text{if }k=0\hfill \\ \frac{\sqrt{k}m}{4\pi }\frac{\mathrm{cosh}2\omega }{\mathrm{sinh}2\omega }+\frac{C}{8\pi }\frac{\mathrm{sinh}2\omega }{\mathrm{sinh}\omega }+\frac{S^3\left(t_1\right)\rho \left(t_1\right)}{3k}\hfill & \text{if }k<0\hfill \end{array}$$
(8.23)
within the hole, where $`C`$ is a constant.
In order for the solution (8.23) to join in a $`C^2`$ manner to the constant solution $`\mathrm{\Phi }(𝐱)=\widehat{\mathrm{\Phi }}`$ outside the hole, the radial derivative of $`\mathrm{\Phi }(𝐱)`$ must vanish at the boundary. By means of the divergence theorem, an integration of (8.22) over the hole therefore yields
$$m=S^3(t_1)\rho (t_1)\text{}^{(3)}\stackrel{}{V}_{\text{hole}}3k_{\text{hole}}\mathrm{\Phi }𝑑\stackrel{}{v}$$
(8.24)
for all values of $`k`$, where $`\text{}^{(3)}\stackrel{}{V}_{\text{hole}}:=_{\text{hole}}𝑑\stackrel{}{v}`$ is the volume of the hole with respect to $`\stackrel{}{𝐡}`$. Thus the mean matter density of the hole is precisely $`\rho (t_1)`$ for $`k=0`$ and $`\rho (t_1)+O(\widehat{\omega }\text{}^2)`$ for $`k0`$. Substitution of (8.23) into (8.24) yields
$$C=\{\begin{array}{cc}2\sqrt{k}m\frac{\left(\mathrm{cos}\widehat{\omega }{\scriptscriptstyle \frac{1}{3}}\mathrm{cos}3\widehat{\omega }\right)}{\left(\mathrm{sin}\widehat{\omega }{\scriptscriptstyle \frac{1}{3}}\mathrm{sin}3\widehat{\omega }\right)}\hfill & \text{if }k>0\hfill \\ 2\sqrt{k}m\frac{\left({\scriptscriptstyle \frac{1}{3}}\mathrm{cosh}3\widehat{\omega }\mathrm{cosh}\widehat{\omega }\right)}{\left({\scriptscriptstyle \frac{1}{3}}\mathrm{sinh}3\widehat{\omega }\mathrm{sinh}\widehat{\omega }\right)}\hfill & \text{if }k<0,\hfill \end{array}$$
(8.25)
whilst in the case $`k=0`$ (8.24) gives
$$m=\frac{4\pi }{3}\rho (t_1)S^3(t_1)\widehat{r}\text{}^3\text{ if }k=0\text{ .}$$
(8.26)
For $`k0`$ the continuity of $`\mathrm{\Phi }`$ at the boundary of the hole gives
$$\widehat{\mathrm{\Phi }}=\{\begin{array}{cc}\frac{S^3\left(t_1\right)\rho \left(t_1\right)}{3k}\frac{\sqrt{k}m}{3\pi }\frac{1}{\left(\mathrm{sin}\widehat{\omega }{\scriptscriptstyle \frac{1}{3}}\mathrm{sin}3\widehat{\omega }\right)}\hfill & \text{if }k>0\hfill \\ \frac{S^3\left(t_1\right)\rho \left(t_1\right)}{3k}+\frac{\sqrt{k}m}{3\pi }\frac{1}{\left({\scriptscriptstyle \frac{1}{3}}\mathrm{sinh}3\widehat{\omega }\mathrm{sinh}\widehat{\omega }\right)}\hfill & \text{if }k<0\hfill \end{array}$$
(8.27)
by means of (8.23) and (8.25), whilst in the case $`k=0`$ one obtains
$$\widehat{\mathrm{\Phi }}=\frac{3m}{8\pi \widehat{r}}+\frac{C}{4\pi }\text{if }k=0$$
(8.28)
by means of (8.23) and (8.24).
Suppose now that there are many holes, each labelled by an index $`\alpha `$. Let $`𝒟_t`$ be a compact $`3`$-domain in a level surface of $`t`$ such that $`𝒟_t`$ intersects none of the holes. The mean value of $`\mathrm{\Phi }`$ on $`𝒟_t`$ is
$$\mathrm{\Phi }_{𝒟_t}:=\frac{_{𝒟_t}\mathrm{\Phi }𝑑\stackrel{}{v}}{\text{}^{(3)}\stackrel{}{V}(𝒟_t)}=\widehat{\mathrm{\Phi }}+\frac{_\alpha _{\text{hole}_\alpha }(\mathrm{\Phi }\widehat{\mathrm{\Phi }})𝑑\stackrel{}{v}}{\text{}^{(3)}\stackrel{}{V}(𝒟_t)}$$
(8.29)
where $`\text{}^{(3)}\stackrel{}{V}(𝒟_t):=_{𝒟_t}𝑑\stackrel{}{v}`$ is the volume of $`𝒟_t`$ with respect to $`\stackrel{}{𝐡}`$ and the sum is carried out over all $`\alpha `$ such that the $`\alpha ^{\text{th}}`$ hole is contained in $`𝒟_t`$. By means of (8.24) and (8.27) in the cases $`k0`$, and by means of (8.23), (8.26) and (8.28) in the case $`k=0`$, one has
$$\mathrm{\Phi }_{𝒟_t}=\widehat{\mathrm{\Phi }}+M_{𝒟_t}$$
(8.30)
where
$$M_{𝒟_t}=\{\begin{array}{cc}\frac{1}{\text{}^{\left(3\right)}V\left(𝒟_t\right)}\frac{1}{3}_\alpha \frac{m_\alpha }{k}\{\frac{2\left(\widehat{\omega }_\alpha {\scriptscriptstyle \frac{1}{2}}\mathrm{sin}2\widehat{\omega }_\alpha \right)}{\left(\mathrm{sin}\widehat{\omega }_\alpha {\scriptscriptstyle \frac{1}{3}}\mathrm{sin}3\widehat{\omega }_\alpha \right)}1\}\hfill & \text{if }k>0\hfill \\ \frac{1}{\text{}^{\left(3\right)}V\left(𝒟_t\right)}\frac{1}{10}_\alpha m_\alpha \widehat{r}_\alpha \text{}^2\hfill & \text{if }k=0\hfill \\ \frac{1}{\text{}^{\left(3\right)}V\left(𝒟_t\right)}\frac{1}{3}_\alpha \frac{m_\alpha }{k}\{1\frac{2\left({\scriptscriptstyle \frac{1}{2}}\mathrm{sinh}2\widehat{\omega }_\alpha \widehat{\omega }_\alpha \right)}{\left({\scriptscriptstyle \frac{1}{3}}\mathrm{sinh}3\widehat{\omega }_\alpha \mathrm{sinh}\widehat{\omega }_\alpha \right)}\}\hfill & \text{if }k<0\text{ .}\hfill \end{array}$$
(8.31)
It will be assumed that the distribution of holes is sufficiently uniform that $`M_{𝒟_t}`$ and $`\mathrm{\Phi }_{𝒟_t}`$ tend to limits $`M`$ and $`\mathrm{\Phi }`$ respectively as $`𝒟_t`$ becomes arbitrarily large. In order to obtain $`\mathrm{\Phi }=0`$ one must have
$$\widehat{\mathrm{\Phi }}+M=0.$$
(8.32)
By (7.4) one then has
$$\delta \rho _t=0+O(\kappa ).$$
(8.33)
In the case $`k=0`$ the quantities $`\widehat{m}_\alpha `$ and $`\widehat{r}_\alpha `$ are, for each $`\alpha `$, related by an equation of the form (8.24), whereby one has
$$\frac{M_{𝒟_t}}{S(t_1)}=\left(\frac{3\pi }{4\rho (t_1)}\right)^{2/3}\frac{_\alpha m_\alpha ^{5/3}}{S^3(t_1)\text{}^{(3)}\stackrel{}{V}(𝒟_t)}\text{if }k=0.$$
(8.34)
For $`k0`$ the quantities $`\widehat{m}_\alpha `$ and $`\widehat{r}_\alpha `$ are, for each $`\alpha `$, related by an equation of the form (8.27) which involves the constant $`\widehat{\mathrm{\Phi }}`$. However the holes may be assumed sufficiently small that (8.34) is a valid approximation for $`k0`$. By (8.8) and (8.34), the dimensionless density parameter of the intergalactic matter is then to be perturbed by an amount
$$\delta \mathrm{\Omega }_{\text{cheese}}=(1+q(t_1))\frac{\kappa M}{S(t_1)}\text{if }k0$$
(8.35)
for
$$\frac{M}{S(t_1)}=\left(\frac{3\pi }{4\rho (t_1)}\right)^{2/3}\underset{𝒟_{t_1}\mathrm{}}{lim}\frac{_\alpha m_\alpha ^{5/3}}{\stackrel{}{V}_{t_1}(𝒟_{t_1})},$$
(8.36)
the sum being carried over all $`\alpha `$ for which the $`\alpha ^{\text{th}}`$ hole is contained in $`𝒟_{t_1}`$, with $`\stackrel{}{V}_{t_1}(𝒟_{t_1})=S^3(t_1)\text{}^{(3)}\stackrel{}{V}(𝒟_{t_1})`$ the volume of $`𝒟_{t_1}`$ with respect to the $`3`$-metric $`S^3(t_1)\stackrel{}{𝐡}`$ induced by $`\stackrel{}{𝐠}`$ on $`\{t=t_1\}`$. If one regards the $`k=0`$ case as a limit of the $`k0`$ cases then (8.35) may considered to apply for all $`k`$.
It is the perturbation (8.35) of the density of the intergalactic matter that distinguishes the present approach to the Swiss cheese model from those of other authors, including Dyer & Roeder (1974). Only with such a perturbation will the mean density parameter $`\mathrm{\Omega }_t`$ be unperturbed. And only then, as discussed in §5, will there be a zero mean correction to the apparent magnitude-redshift relation to first order in $`\kappa `$.
## 9. Concluding remarks
The validity of the linear approximation employed in this paper depends upon the smallness of the mean dimensionless density perturbation $`\delta \mathrm{\Omega }_t`$ associated with the matter in the galaxies. Astronomical estimates of this quantity are at present inconclusive, although there is a general consensus that one does have $`\delta \mathrm{\Omega }_t1`$. If this is correct then the theory presented should provide a valid description of the total gravitational field and cumulative weak lensing effects of the galaxies for any given model of the distribution of galactic matter. Depending upon the actual value of $`\delta \mathrm{\Omega }_t`$ and the range of redshifts of interest, there may be a need to carry the theory to second or higher order in $`\kappa `$. This would at least take into account the effect of shear on the apparent brightness of cosmological sources. Caustics however cannot be adequately described by any finite order power series analysis. For this it would be necessary to consider the full non-linear theory.
I express my appreciation to Prof. M.A.H. MacCallum for his interest and encouragement, and to the Department of Mathematics and Applied Mathematics at the University of Cape Town for their hospitality. Funding was provided by the Swedish Natural Science Research Council (NFR).
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# Unravelling Active Galactic Nuclei
## 1 Introduction
The relationship between Seyfert galaxies and radio-loud galaxies has not been explored in detail. However, a complete sample of flat-spectrum radio-loud AGN contains a significant fraction of Seyfert-like sources, including one which would be classified as a narrow-line Seyfert 1 (NLS1).
To understand the physical mechanisms responsible for producing the different characteristics of AGN, a multiwavelength approach is needed. Correlations between properties at different wavelengths can be used to reject and refine physical models of the central regions of AGN. The Parkes Half-Jansky Flat-Spectrum Sample (PHFS) D1 is interesting in this context because it is a radio-selected sample. The selection criteria are as follows:
* Radio-loud: 2.7GHz flux $`>`$ 0.5 Jy.
* Flat-spectrum: $`\alpha _{2.7/5.0}>0.5`$, where $`S_\nu =\nu ^\alpha `$
* Galactic latitude: $`|`$b$`|>20^{}`$
* $`45^{}<`$ Dec(B1950) $`<+10^{}`$
This sample contains 323 sources with a wide range of properties which can be quite different to AGN selected by optical colours. The PHFS AGN have a large range in optical luminosities, with more than 50 objects having absolute magnitudes in the Seyfert luminosity range (M$`{}_{B}{}^{}>`$-23). Another interesting characteristic is the large dispersion in optical colours of this sample compared with optically selected samples such as the Large Bright Quasar Survey W1 ; F1 .
## 2 Spectra of the Parkes Quasars
The mechanism responsible for producing the large dispersion in optical colours, especially the significant number of red sources in the sample, is not at all well established. There have been suggestions of dust reddening W1 , reddening due to the underlying galaxy spectrum M1 and synchrotron reddening W2 . We consider two questions:
1. Are the reddened continuum colours associated with changes in the emission lines of the AGN?
2. Can we identify the reddening mechanisms contributing to each AGN?
We have low resolution optical spectra of an unbiased subsample of the PHFS. This subsample was divided into three equal colour bins based on the B-K colours of the objects: blue (B-K $`<`$ 3), intermediate (3 $`<`$ B-K $`<`$ 4.4) and red (B-K $`>`$ 4.4). A composite spectrum of the objects in each bin was made (see F3 for a description of the technique used). The results are plotted in Figure 1. It should be noted that the continuum has been normalised to the same slope for each bin to compare the emission line properties over the range of colours. The LBQS composite spectrum is also shown as a comparison with an optically selected sample F2 .
There are clear differences between the different composites. The width of H$`\beta `$ line decreases as redness increases and OII(3727Å) is only seen in the red composite. The MgII(2800Å) line is also quite broad in the blue and intermediate composites, and much narrower in the red composite. Further analysis of the spectra in the red bin shows there are two different types of spectra in this bin: (1) low redshift galaxy-type spectra, and (2) higher redshift quasar spectra.
Several very red objects in the Parkes sample seem to appear red owing to a large contribution from a galactic component. The 4000Å break in the spectrum of the galactic light reduces the amount of light in the blue filter giving a B-K colour which is quite red. These objects are at low redshifts and are generally resolved. Masci M1 has fitted an evolved elliptical galaxy SED to the PHFS spectra, showing that these objects have a high proportion of galactic light in their spectra. These objects also have the steepest radio spectrum with values of $`\alpha 0.5`$. The composite spectrum of these resolved sources is fairly typical of a galaxy with no H$`\beta `$ emission.
The second group of red objects in the sample are higher redshift quasar/Seyfert-type objects. They show the strong emission lines characteristic of quasars and have flatter radio spectral indices. Whiting et al. W2 have modelled the reddening by synchrotron emission using broad band SEDs. They show that these objects tend to have a strong component of synchrotron emission which is responsible for their red continuum. These quasars would be missed in a traditional blue-selected quasar sample. The composite of these unresolved sources has a completely different shape, with narrow permitted emission lines at CIV(1549Å), CIII\](1909Å) and MgII(2798Å) superimposed on a fairly smooth continuum. Thus it appears that the red AGN have relatively narrow emission lines. A fuller description of these results will appear in Oshlack et al. O2 .
## 3 A NLS1 in the Parkes Sample
In examining the spectra from the Parkes sample we have identified a NLS1 candidate, PKS 2004-447 (Figure 2). One model for NLS1s suggests they are AGN oriented towards us, so that the doppler broadening of a disk-like Broad Line Region would naturally produce narrower emission lines O1 . An extremely radio-loud NLS1 could provide an independent test of such physical models of NLS1s, and AGN unification in general.
A low resolution spectrum of this source shows an H$`\beta `$ width of $``$2000km s<sup>-1</sup>. There is some indication of FeII emission on the blueward side of H$`\beta `$ and the source has been detected at X-ray energies S1 . This object is extremely radio-loud with a radio to optical flux ratio (R) exceeding 7000. The colour of PKS 2004-447 is very red, in contrast with the previously discovered radio-loud NLS1, RGB J0044+193 which was quite blue S2 . The radio spectrum of PKS 2004-447 is also quite steep, with recent measurements of the (contemporaneous) radio spectral index giving $`\alpha _R=0.67`$, consistent with radio indices of radio-quiet NLS1s (Moran, these proceedings). A higher resolution spectrum of PKS 2004-447 will be obtained, to confirm its identity.
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# Zero-point energies and the multiplicative anomaly.
## I Introduction
Renormalisation and regularisation techniques form a vital part of the physicist’s arsenal when performing calculations in quantum field theory. The technique of $`\zeta `$-function regularisation is a well established method for obtaining finite results in quantum field theoretical calculations and has shown itself to be a very elegant and powerful method, particularly in curved spacetimes and spacetimes with a non-trivial topology.
Calculations involving Feynman path integrals typically involve the determinant of a differential operator. This determinant is an infinite product and has to be regularised in some way. It has been shown recently that results obtained from $`\zeta `$-function regularised determinants can be ambiguous \- the source of this ambiguity is known as a multiplicative anomaly. The multiplicative anomaly is essentially the difference between different zeta regularised factorisations of a determinant. For example, in calculating a determinant one may wish to factorise it in order to make calculations easier. Normally one would write $`det(AB)=det(A)det(B)`$, but in the case of infinite matrices this relation is not always correct after regularisation. The multiplicative anomaly in a $`D`$ dimensional spacetime is defined as,
$$a_D=\mathrm{ln}det(AB)\mathrm{ln}det(A)\mathrm{ln}det(B)$$
(1)
The relevance of the multiplicative anomaly for physics was first brought to light by Elizalde, Vanzo and Zerbini and they showed its connection with the Wodzicki residue. The high temperature limit of the one loop effective action for a charged scalar field with chemical potential was considered by Elizalde, Filippi, Vanzo and Zerbini , and the resulting anomaly was found to depend on the chemical potential. It appeared therefore that an extra term, previously overlooked, might be present in the effective action, a term which could not be removed by renormalisation. (This idea received criticism from Evans and Dowker . Elizalde, Filippi, Vanzo and Zerbini responded to these criticisms in .)
It seemed that there were many different expressions for the effective action, one for each way in which the determinant can be factorised. So, given these varying expressions, each one differing from another by a corresponding multiplicative anomaly, how can we know which one (if any) is correct? The present authors concluded in that the ambiguity associated in choosing different factorisations could only be resolved by making comparisons with calculations performed using canonical methods. In this paper we shall consider two factorisations; the full fourth order operator:
$$\mathrm{\Gamma }_A=\frac{1}{2}\mathrm{ln}detl^4\left(\left(\mathrm{}+m^2\mu ^2\right)^24\mu ^2\frac{^2}{t^2}\right)$$
(2)
and the ‘standard’ factorisation of two second order operators (eg. as used in ):
$$\mathrm{\Gamma }_B=\frac{1}{2}\mathrm{ln}detl^2\left(\mathrm{}+m^2\mu ^2+2i\mu \frac{}{t}\right)\frac{1}{2}\mathrm{ln}detl^2\left(\mathrm{}+m^2\mu ^22i\mu \frac{}{t}\right).$$
(3)
In an exact expression for the effective action could only be calculated for the B factorisation, and this agreed with the standard, well known thermodynamical expression - a sum over the zero-point energies and the thermal contributions of Bose-Einstein sums for particles and anti-particles. The high temperature limits of the effective action for both cases were calculated and it was found that $`\mathrm{\Gamma }_B`$ agreed with the result obtained by Haber and Weldon who did not use path integrals or $`\zeta `$-function regularisation. $`\mathrm{\Gamma }_A`$ differed from $`\mathrm{\Gamma }_B`$ by an amount exactly equal to the multiplicative anomaly calculated in . There does not seem to be an a priori way of determining which factorisation will yield the correct physics; an ‘objective’ comparison with canonical methods needs to be performed. Clearly this is a problem if calculations using $`\zeta `$-function regularisation need to be made on a system where the canonical answer is not known.
Since the high temperature expansions of $`\mathrm{\Gamma }_A`$ and $`\mathrm{\Gamma }_B`$ differ, it might be thought that there would be some discrepancy between their exact expressions also. In Sec. II we shall calculate the exact effective action $`\mathrm{\Gamma }_A`$, and show that it actually gives the correct result, in complete agreement with $`\mathrm{\Gamma }_B`$. Is the multiplicative anomaly therefore just an artefact of the high temperature expansion? In Sec. III we postulate that the multiplicative anomaly arises from the zeta-regularised zero-point energies and is not a thermal phenomenon. It arises in factorisations where the charge operator $`Q`$ has not been normal ordered. The multiplicative anomaly is calculated explicitly as the difference between $`\zeta `$-regularised zero-point energies with and without a chemical potential. We also consider the interacting case. Both calculations give multiplicative anomalies which agree with those calculated in . In Sec. IV we shall draw some conclusions.
## II The exact effective action $`\mathrm{\Gamma }_A`$
We shall use the notation and conventions of , now setting $`e=1`$. We are working with a relativistic, non-interacting, charged scalar field with a chemical potential. The action is
$`S[\varphi ]`$ $`=`$ $`{\displaystyle _0^\beta }dt{\displaystyle _\mathrm{\Sigma }}d\sigma _x({\displaystyle \frac{1}{2}}(\dot{\varphi _1}i\mu \varphi _2)^2+{\displaystyle \frac{1}{2}}(\dot{\varphi _2}+i\mu \varphi _1)^2`$ (5)
$`+{\displaystyle \frac{1}{2}}|\varphi _1|^2+{\displaystyle \frac{1}{2}}|\varphi _2|^2+{\displaystyle \frac{1}{2}}m^2(\varphi _1^2+\varphi _2^2)).`$
We expand about a constant background field with $`\varphi _1=\varphi `$ and $`\varphi _2=0`$. $`\mathrm{\Gamma }_A`$ is defined formally as
$`\mathrm{\Gamma }_A`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{ln}detl^2\left(l^2S_{,ij}[\varphi ]\right)`$ (6)
$`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{ln}detl^4\left(\left(\mathrm{}+m^2\mu ^2\right)^24\mu ^2{\displaystyle \frac{^2}{t^2}}\right).`$ (7)
Using $`\zeta `$-function regularisation we can define:
$$\mathrm{\Gamma }_A=\frac{1}{2}\zeta _A^{}(0)+\frac{1}{2}\zeta _A(0)\mathrm{ln}l^4$$
(8)
with
$$\zeta _A(s)=\underset{n}{}\underset{j=\mathrm{}}{\overset{\mathrm{}}{}}\left[\left(\omega _j^2\mu ^2+E_n^2\right)^2+4\mu ^2\omega _j^2\right]^s$$
(9)
The $`\omega _j`$ are the Matsubara frequencies for scalars, $`\omega _j=2\pi j/\beta `$ and $`E_n^2=\sigma _n+m^2`$. $`\sigma _n`$ are the eigenvalues of $`^2`$ on the spatial part of the manifold. We shall apply the Abel-Plana summation formula,
$$\underset{j=\mathrm{}}{\overset{\mathrm{}}{}}f(j)=_{\mathrm{}}^+\mathrm{}f(x)𝑑x+_{\mathrm{}+iϵ}^{\mathrm{}+iϵ}𝑑z\left(e^{2\pi iz}1\right)^1\left[f(z)+f(z)\right]$$
(10)
to evaluate $`\zeta _A(s)`$. Let us label the two integrals arising on the right hand side of (9) $`Q(s)`$ and $`R(s)`$ respectively after using (10) to perform the sum over $`j`$. Then,
$`Q(s)`$ $`=`$ $`{\displaystyle \underset{n}{}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\left[\left(\left({\displaystyle \frac{2\pi x}{\beta }}\right)^2\mu ^2+E_n^2\right)^2+4\mu ^2\left({\displaystyle \frac{2\pi x}{\beta }}\right)^2\right]^s𝑑x`$ (11)
$`=`$ $`{\displaystyle \underset{n}{}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\left[ax^4+bx^2+c\right]^s𝑑x`$ (12)
where $`a=(2\pi /\beta )^4`$, $`b=2(2\pi /\beta )^2(E_n^2+\mu ^2)`$ and $`c=(E_n^2\mu ^2)^2`$. Making a simple substitution and using the identity
$$\alpha ^s=\frac{1}{\mathrm{\Gamma }(s)}_0^{\mathrm{}}𝑑tt^{s1}e^{\alpha t}$$
(13)
gives us
$$Q(s)=\underset{n}{}\frac{1}{\mathrm{\Gamma }(s)}_0^{\mathrm{}}𝑑te^{ct}t^{s1}_0^{\mathrm{}}𝑑xx^{\frac{1}{2}}e^{(ax^2+bx)t}.$$
(14)
Evaluating the integral over $`x`$ first results in
$$Q(s)=\underset{n}{}\frac{1}{2\mathrm{\Gamma }(s)}\sqrt{\frac{b}{a}}_0^{\mathrm{}}𝑑te^{t\left(c\frac{b^2}{8a}\right)}t^{s1}K_{\frac{1}{4}}\left(\frac{b^2t}{8a}\right)$$
(15)
where $`K_{\frac{1}{4}}`$ is a modified Bessel function. Performing the integral over $`t`$ we get,
$`Q(s)`$ $`=`$ $`{\displaystyle \underset{n}{}}{\displaystyle \frac{1}{2\mathrm{\Gamma }(s)}}\sqrt{{\displaystyle \frac{b}{a}}}{\displaystyle \frac{\sqrt{\pi }}{c^{s+1/4}}}\left({\displaystyle \frac{b^2}{4a}}\right)^{\frac{1}{4}}{\displaystyle \frac{\mathrm{\Gamma }(s+\frac{1}{4})\mathrm{\Gamma }(s\frac{1}{4})}{\mathrm{\Gamma }(s+\frac{1}{2})}}{}_{2}{}^{}F_{1}^{}(s+{\displaystyle \frac{1}{4}},{\displaystyle \frac{3}{4}};s+{\displaystyle \frac{1}{2}};1{\displaystyle \frac{b^2}{4ac}})`$ (16)
$`=`$ $`{\displaystyle \underset{n}{}}\sqrt{{\displaystyle \frac{\pi }{2}}}\left({\displaystyle \frac{\beta }{2\pi }}\right){\displaystyle \frac{(E_n^2+\mu ^2)}{(E_n^2\mu ^2)^{2s+1/2}}}{\displaystyle \frac{\mathrm{\Gamma }(s\frac{1}{4})}{\mathrm{\Gamma }(s)\mathrm{\Gamma }(\frac{1}{4})}}{\displaystyle _0^1}t^{s\frac{3}{4}}(1t)^{\frac{3}{4}}\left(1+{\displaystyle \frac{4\mu ^2E_n^2t}{(E_n^2\mu ^2)^2}}\right)^{\frac{3}{4}}𝑑t`$ (17)
after substituting in the values of $`a`$, $`b`$ and $`c`$. This is analytic at $`s=0`$, and is equal to zero, $`Q(0)=0`$. It is easy to see that,
$`Q^{}(0)`$ $`=`$ $`{\displaystyle \underset{n}{}}\sqrt{{\displaystyle \frac{\pi }{2}}}\left({\displaystyle \frac{\beta }{2\pi }}\right){\displaystyle \frac{(E_n^2+\mu ^2)}{(E_n^2\mu ^2)^{1/2}}}{\displaystyle \frac{\mathrm{\Gamma }(\frac{1}{4})}{\mathrm{\Gamma }(\frac{1}{4})}}\sqrt{2\pi }{\displaystyle \frac{\mathrm{\Gamma }(\frac{1}{4})}{\mathrm{\Gamma }(\frac{3}{4})}}(1+{\displaystyle \frac{4\mu ^2E_n^2}{(E_n^2\mu ^2)^2}})^{\frac{1}{4}}\mathrm{cos}\left[{\displaystyle \frac{1}{2}}\mathrm{arctan}\left({\displaystyle \frac{4\mu ^2E_n^2}{(E_n^2\mu ^2)^2}}\right)^{\frac{1}{2}}\right]`$ (18)
$`=`$ $`2\beta {\displaystyle \underset{n}{}}(E_n^2+\mu ^2)^{\frac{1}{2}}\mathrm{cos}\left[{\displaystyle \frac{1}{2}}\mathrm{arctan}\left({\displaystyle \frac{2\mu E_n}{(E_n^2\mu ^2)}}\right)\right]`$ (19)
$`=`$ $`2\beta {\displaystyle \underset{n}{}}E_n`$ (20)
This last result involves the sum of zero-point energies. Turning now to the contour integral in (10), we can write
$$R(s)=2\underset{n}{}_{\mathrm{}+iϵ}^{\mathrm{}+iϵ}\left(e^{2\pi iz}1\right)^1\left[\left(\left(\frac{2\pi z}{\beta }\right)^2\mu ^2+E_n^2\right)^2+4\mu ^2\left(\frac{2\pi z}{\beta }\right)^2\right]^s𝑑z$$
(21)
since $`f(z)`$ is even. There are poles at all integers on the real axis, and branch points where the expression in square brackets in the above equation is equal to zero, namely at $`z=\pm i(\beta /2\pi )(E_n+\mu )`$,$`\pm i(\beta /2\pi )(E_n\mu )`$. By taking branch cuts between the poles in each half of the complex plane, and closing the contour in the upper half plane, it can be shown that
$$R(s)=2s\underset{n}{}\mathrm{ln}\left[\left(1e^{\beta (E_n\mu )}\right)\left(1e^{\beta (E_n+\mu )}\right)\right]$$
(22)
Hence $`R(0)=0`$ and
$$R^{}(0)=2\underset{n}{}\mathrm{ln}\left[\left(1e^{\beta (E_n\mu )}\right)\left(1e^{\beta (E_n+\mu )}\right)\right]$$
(23)
Therefore $`\zeta _A(0)=Q(0)+R(0)=0`$ and so
$`\mathrm{\Gamma }_A`$ $`=`$ $`{\displaystyle \frac{1}{2}}\zeta _A^{}(0)`$ (24)
$`=`$ $`{\displaystyle \underset{n}{}}\left\{\beta E_n+\mathrm{ln}\left[\left(1e^{\beta (E_n\mu )}\right)\left(1e^{\beta (E_n+\mu )}\right)\right]\right\}`$ (25)
which is the same result as was obtained for $`\mathrm{\Gamma }_B`$ in . A simpler, but perhaps less elegant way of evaluating (9) is presented in the appendix.
So we have a paradox - the exact expressions for $`\mathrm{\Gamma }_A`$ and $`\mathrm{\Gamma }_B`$ agree, while their high temperature expansions differ by the multiplicative anomaly. The resolution of this paradox lies in the fact that the sums over the energy levels (integrals over $`k`$ when $`\sigma _n=k^2`$) have not yet been performed in the exact expressions we are considering. In the high temperature expansions, the $`\zeta `$-functions were expanded in powers of $`\mu `$, the chemical potential, and the sums over the energy levels and Matsubara frequencies were then performed. As we shall see in the next section, the multiplicative anomaly arises from the chemical potential being present in the zero-point energy contributions. This is due to a lack of normal ordering in the charge operator.
## III Normal ordering and the multiplicative anomaly
In this section we shall show that the multiplicative anomaly stems from the zero-point energy contribution to the effective action.
### A Non-interacting model
There are two ways to write down the zero-point energies for the system described by (5); with or without a chemical potential. One can write $`(\beta /2)_n(E_n\pm \mu )`$ for particles and anti-particles ($`\mu `$ and $`+\mu `$ respectively) or one can simply write $`\beta _nE_n`$ (which was derived in our exact expressions for $`\mathrm{\Gamma }_A`$ and $`\mathrm{\Gamma }_B`$). The multiplicative anomaly is the difference between these $`\zeta `$-regularised zero-point energies. Let us define
$`I_\pm `$ $`=`$ $`{\displaystyle \frac{\beta }{2}}{\displaystyle \underset{n}{}}(E_n\pm \mu )`$ (26)
$`J`$ $`=`$ $`\beta {\displaystyle \underset{n}{}}E_n.`$ (27)
Formally of course, there is no difference between the zero-point energies in the two cases: $`I_++I_{}J=0`$. But if we regularise first then the anomaly appears. Define
$`I_\pm (s)`$ $`=`$ $`{\displaystyle \frac{\beta }{2}}{\displaystyle \underset{n}{}}(E_n\pm \mu )^s`$ (28)
$`J(s)`$ $`=`$ $`\beta {\displaystyle \underset{n}{}}E_n^s.`$ (29)
Then, we claim that the multiplicative anomaly is
$$a_4=I_+(1)+I_{}(1)J(1).$$
(30)
Let us calculate $`I_\pm (s)`$ with $`E_n^2=k^2+m^2`$:
$$I_\pm (s)=\frac{\beta }{2}\frac{4\pi V}{(2\pi )^3}_0^{\mathrm{}}𝑑kk^2\left[\left(k^2+m^2\right)^{\frac{1}{2}}\pm \mu \right]^s$$
(31)
We can binomially expand the square bracket in powers of $`\mu `$ up to $`O(\mu ^4)`$. We do not need to consider higher order terms for reasons which will become apparent in due course. After performing the integrals we have,
$`I_\pm (s)`$ $`=`$ $`{\displaystyle \frac{\sqrt{\pi }\beta V}{16\pi ^2}}\{{\displaystyle \frac{\mathrm{\Gamma }(\frac{s}{2}\frac{3}{2})}{\mathrm{\Gamma }(\frac{s}{2})}}\left(m^2\right)^{\frac{3}{2}\frac{s}{2}}s\mu {\displaystyle \frac{\mathrm{\Gamma }(\frac{s}{2}1)}{\mathrm{\Gamma }(\frac{s}{2}+\frac{1}{2})}}\left(m^2\right)^{1\frac{s}{2}}+s(s+1){\displaystyle \frac{\mu ^2}{2}}{\displaystyle \frac{\mathrm{\Gamma }(\frac{s}{2}\frac{1}{2})}{\mathrm{\Gamma }(\frac{s}{2}+1)}}\left(m^2\right)^{\frac{1}{2}\frac{s}{2}}`$ (33)
$`s(s+1)(s+2){\displaystyle \frac{\mu ^3}{6}}{\displaystyle \frac{\mathrm{\Gamma }(\frac{s}{2})}{\mathrm{\Gamma }(\frac{s}{2}+\frac{3}{2})}}\left(m^2\right)^{\frac{s}{2}}+s(s+1)(s+2)(s+3){\displaystyle \frac{\mu ^4}{24}}{\displaystyle \frac{\mathrm{\Gamma }(\frac{s}{2}+\frac{1}{2})}{\mathrm{\Gamma }(\frac{s}{2}+2)}}\left(m^2\right)^{\frac{s}{2}\frac{1}{2}}\}.`$
The terms with odd powers of $`\mu `$ cancel when we write down an expression for $`I_+(s)+I_{}(s)`$. In the $`\mu ^2`$ and $`\mu ^4`$ terms, the $`\mathrm{\Gamma }`$-function in the numerator is divergent at $`s=1`$, but can be analytically continued to cancel away the factor of $`(s+1)`$ multiplying each term. Thus,
$`I_+(s)+I_{}(s)`$ $`=`$ $`{\displaystyle \frac{\sqrt{\pi }\beta V}{8\pi ^2}}\{{\displaystyle \frac{\mathrm{\Gamma }(\frac{s}{2}\frac{3}{2})}{\mathrm{\Gamma }(\frac{s}{2})}}\left(m^2\right)^{\frac{3}{2}\frac{s}{2}}+2s\mu ^2{\displaystyle \frac{\mathrm{\Gamma }(\frac{s}{2}+\frac{3}{2})}{(s1)\mathrm{\Gamma }(\frac{s}{2}+1)}}\left(m^2\right)^{\frac{1}{2}\frac{s}{2}}`$ (35)
$`+s(s+2)(s+3){\displaystyle \frac{\mu ^4}{12}}{\displaystyle \frac{\mathrm{\Gamma }(\frac{s}{2}+\frac{3}{2})}{\mathrm{\Gamma }(\frac{s}{2}+2)}}\left(m^2\right)^{\frac{s}{2}\frac{1}{2}}\}.`$
All even, higher order terms in $`\mu `$ have analytic $`\mathrm{\Gamma }`$-functions in the numerator, and so the $`(s+1)`$ ensures they are all zero at $`s=1`$. This is why we were able to stop expanding at fourth order in $`\mu `$.
Turning now to $`J(s)`$, we see it is exactly the first term of (35):
$`J(s)`$ $`=`$ $`{\displaystyle \frac{4\pi \beta V}{(2\pi )^3}}{\displaystyle _0^{\mathrm{}}}𝑑kk^2\left(k^2+m^2\right)^{\frac{s}{2}}`$ (36)
$`=`$ $`{\displaystyle \frac{\sqrt{\pi }\beta V}{8\pi ^2}}{\displaystyle \frac{\mathrm{\Gamma }(\frac{s}{2}\frac{3}{2})}{\mathrm{\Gamma }(\frac{s}{2})}}\left(m^2\right)^{\frac{3}{2}\frac{s}{2}}`$ (37)
Using (30), it is now a straightforward matter to show that the multiplicative anomaly is
$$a_4=\frac{\beta V}{8\pi ^2}\mu ^2\left(m^2\frac{\mu ^2}{3}\right)$$
(38)
in agreement with .
This calculation sheds some light on why the high temperature expansions of $`\mathrm{\Gamma }_A`$ and $`\mathrm{\Gamma }_B`$ differ, and the exact expressions agree. In the high temperature situation, the integrations over $`k`$ were carried out after the expansions, and for a reason which is not clear to us at the present, the chemical potentials in the zero-point energies of the A-factorisation were not able to cancel. So the zero-point energies were of the form $`(\beta /2)_n(E_n\pm \mu )`$. In the B-factorisation, the energy levels were just $`\beta _nE_n`$ and so there was no anomaly. In the exact expressions for $`\mathrm{\Gamma }_A`$ and $`\mathrm{\Gamma }_B`$, the integrals have not even been performed, and so the $`+\mu `$ and $`\mu `$ simply disappear, leaving no trace of a discrepancy.
Although simply by looking at the A factorisation (7), (9), we cannot say whether or not it will produce a multiplicative anomaly in the high temperature expansion, given that we know it does produce an anomaly, we can say something about why it does. The canonical energy levels for our system are derived from the Hamiltonian operator $`H`$,
$$H=\underset{n}{}E_n(a_n^{}a_n+\frac{1}{2}+b_n^{}b_n+\frac{1}{2})$$
(39)
where $`a_n^{}`$, $`a_n`$ ($`b_n^{}`$, $`b_n`$) are the creation and annihilation operators for particles (anti-particles). For a system of charged fields with a chemical potential, the full Hamiltonian (which is the argument of the exponential in the partition function) is
$$\overline{H}=H\mu :Q:$$
(40)
where $`:Q:`$ is the normal ordered charge operator
$$:Q:=\underset{n}{}(a_n^{}a_nb_n^{}b_n).$$
(41)
Note that we have to normal order by hand. There is no good mathematical reason why we normal order, we just like to have an uncharged vacuum:
$$0|:Q:|0=0.$$
(42)
It now becomes clear why we have two different expressions for the energy levels (26), (27). They correspond to the eigenvalues of $`\overline{H}`$ and $`H`$ respectively in the case when the charge operator $`Q`$ is not normal ordered. So to avoid having an anomaly we need to ensure that both $`H`$ and $`\overline{H}`$ have the same eigenvalues - we need to normal order $`Q`$. This implies that the A factorisation, which gives rise to an anomaly, is not normal ordered. This is a symptom of using Feynman path integrals, indeed Bernard was aware of this in 1974 - the last sentence in section II of his seminal paper reads, ‘…the functional-integral formalism never does normal ordering for us.’
### B The interacting case
The multiplicative anomaly can also be calculated in the interacting case, and can again be seen to be the difference of $`\zeta `$-regularised zero-point energies.
Let us define
$`X_\pm `$ $`=`$ $`{\displaystyle \frac{\beta }{2}}{\displaystyle \underset{n}{}}\left[E_n^2+{\displaystyle \frac{\lambda \varphi ^2}{3}}+\mu ^2\pm \left(4\mu ^2\left(E_n^2+{\displaystyle \frac{\lambda \varphi ^2}{3}}\right)+{\displaystyle \frac{\lambda ^2\varphi ^4}{36}}\right)^{\frac{1}{2}}\right]^{\frac{1}{2}}`$ (43)
$`Y`$ $`=`$ $`{\displaystyle \frac{\beta }{2}}{\displaystyle \underset{n}{}}\left(E_n^2+{\displaystyle \frac{\lambda \varphi ^2}{2}}\right)^{\frac{1}{2}}+{\displaystyle \frac{\beta }{2}}{\displaystyle \underset{n}{}}\left(E_n^2+{\displaystyle \frac{\lambda \varphi ^2}{6}}\right)^{\frac{1}{2}}`$ (44)
in an analogous way to $`I_\pm `$ and $`J`$ in the non-interacting case. (See for example for a full derivation of the energy levels). Then we have
$`X_\pm (s)`$ $`=`$ $`{\displaystyle \frac{\beta }{2}}{\displaystyle \underset{n}{}}\left[E_n^2+{\displaystyle \frac{\lambda \varphi ^2}{3}}+\mu ^2\pm \left(4\mu ^2\left(E_n^2+{\displaystyle \frac{\lambda \varphi ^2}{3}}\right)+{\displaystyle \frac{\lambda ^2\varphi ^4}{36}}\right)^{\frac{1}{2}}\right]^{\frac{s}{2}}`$ (45)
$`Y(s)`$ $`=`$ $`{\displaystyle \frac{\beta }{2}}{\displaystyle \underset{n}{}}\left(E_n^2+{\displaystyle \frac{\lambda \varphi ^2}{2}}\right)^{\frac{s}{2}}+{\displaystyle \frac{\beta }{2}}{\displaystyle \underset{n}{}}\left(E_n^2+{\displaystyle \frac{\lambda \varphi ^2}{6}}\right)^{\frac{s}{2}}`$ (46)
To evaluate $`X_\pm (s)`$, we expand the square root inside the square bracket in powers of $`\lambda `$ up to $`O(\lambda ^2)`$:
$$X_\pm (s)=\frac{\beta }{2}\underset{n}{}\left[E_n^2+\frac{\lambda \varphi ^2}{3}+\mu ^2\pm 2\mu E_n\pm \frac{\mu \lambda \varphi ^2}{3E_n}\pm \frac{\lambda ^2\varphi ^4}{144\mu E_n}\frac{\mu \lambda ^2\varphi ^4}{36E_n^3}\right]^{\frac{s}{2}}.$$
(47)
We note that we can switch from $`X_+(s)`$ to $`X_{}(s)`$ by letting $`\mu \mu `$. Therefore we shall work with $`X_+(s)`$ for simplicity, and let $`\mu \mu `$ to write down $`X_{}(s)`$ at the end of the calculation. $`X_+(s)`$ can be written in the form,
$`X_+(s)`$ $`=`$ $`{\displaystyle \frac{\beta }{2}}{\displaystyle \underset{n}{}}\left[a_0+a_1\lambda +a_2\lambda ^2\right]^{\frac{s}{2}}`$ (48)
$`=`$ $`{\displaystyle \frac{\beta }{2}}{\displaystyle \underset{n}{}}a_0^{\frac{s}{2}}\left[1+a_0^1a_1\lambda +a_0^1a_2\lambda ^2\right]^{\frac{s}{2}}`$ (49)
$`=`$ $`{\displaystyle \frac{\beta }{2}}{\displaystyle \underset{n}{}}\left(a_0^{\frac{s}{2}}{\displaystyle \frac{1}{2}}sa_0^{\frac{s}{2}1}a_1\lambda {\displaystyle \frac{1}{2}}sa_0^{\frac{s}{2}1}a_2\lambda ^2+{\displaystyle \frac{1}{8}}s(s+2)a_0^{\frac{s}{2}2}a_1^2\lambda ^2\right)`$ (50)
where
$`a_0`$ $`=`$ $`(E_n+\mu )^2`$ (51)
$`a_1`$ $`=`$ $`{\displaystyle \frac{\varphi ^2}{3}}\left(1+\mu E_n^1\right)`$ (52)
$`a_2`$ $`=`$ $`{\displaystyle \frac{\varphi ^4}{36}}\left({\displaystyle \frac{E_n^1}{4\mu }}\mu E_n^3\right).`$ (53)
It is then a straightforward (if a little tedious) matter to expand each term of (50) and integrate over $`k`$, as was done in the non-interacting case.
After the dust settles, we find
$`X_+(s)+X_{}(s)`$ $`=`$ $`{\displaystyle \frac{\sqrt{\pi }\beta V}{8\pi ^2}}\{{\displaystyle \frac{\mathrm{\Gamma }(\frac{s}{2}\frac{3}{2})}{\mathrm{\Gamma }(\frac{s}{2})}}\left(m^2\right)^{\frac{3}{2}\frac{s}{2}}+2s\mu ^2{\displaystyle \frac{\mathrm{\Gamma }(\frac{s}{2}+\frac{3}{2})}{(s1)\mathrm{\Gamma }(\frac{s}{2}+1)}}\left(m^2\right)^{\frac{1}{2}\frac{s}{2}}`$ (56)
$`+s(s+2)(s+3){\displaystyle \frac{\mu ^4}{12}}{\displaystyle \frac{\mathrm{\Gamma }(\frac{s}{2}+\frac{3}{2})}{\mathrm{\Gamma }(\frac{s}{2}+2)}}\left(m^2\right)^{\frac{s}{2}\frac{1}{2}}\}{\displaystyle \frac{\sqrt{\pi }\beta V\lambda \varphi ^2}{48\pi ^2}}\{s{\displaystyle \frac{\mathrm{\Gamma }(\frac{s}{2}\frac{1}{2})}{\mathrm{\Gamma }(\frac{s}{2}+1)}}\left(m^2\right)^{\frac{1}{2}\frac{s}{2}}`$
$`+s(s+2)\mu ^2{\displaystyle \frac{\mathrm{\Gamma }(\frac{s}{2}+\frac{3}{2})}{\mathrm{\Gamma }(\frac{s}{2}+2)}}\left(m^2\right)^{\frac{1}{2}\frac{s}{2}}\}+{\displaystyle \frac{5\sqrt{\pi }\beta V\lambda ^2\varphi ^4s(s+2)}{2304\pi ^2}}{\displaystyle \frac{\mathrm{\Gamma }(\frac{s}{2}+\frac{1}{2})}{\mathrm{\Gamma }(\frac{s}{2}+2)}}(m^2)^{\frac{1}{2}\frac{s}{2}}`$
By expanding (46) to order $`\lambda ^2`$, it can easily be shown that
$$Y(s)=\frac{\sqrt{\pi }\beta V}{8\pi ^2}\left\{\frac{\mathrm{\Gamma }(\frac{s}{2}\frac{3}{2})}{\mathrm{\Gamma }(\frac{s}{2})}\left(m^2\right)^{\frac{3}{2}\frac{s}{2}}s\frac{\lambda \varphi ^2}{6}\frac{\mathrm{\Gamma }(\frac{s}{2}\frac{1}{2})}{\mathrm{\Gamma }(\frac{s}{2}+1)}\left(m^2\right)^{\frac{1}{2}\frac{s}{2}}+\frac{5}{288}s(s+2)\lambda ^2\varphi ^4\frac{\mathrm{\Gamma }(\frac{s}{2}+\frac{1}{2})}{\mathrm{\Gamma }(\frac{s}{2}+2)}\left(m^2\right)^{\frac{1}{2}\frac{s}{2}}\right\}$$
(57)
and hence,
$`a_4`$ $`=`$ $`X_+(1)+X_{}(1)Y(1)`$ (58)
$`=`$ $`{\displaystyle \frac{\beta V}{8\pi ^2}}\mu ^2\left(m^2+{\displaystyle \frac{\lambda \varphi ^2}{3}}{\displaystyle \frac{\mu ^2}{3}}\right).`$ (59)
This agrees with the result in except for a term in $`\lambda ^2\varphi ^4`$. This is of no consequence however, since any term proportional to the background field $`\varphi `$ (but not the chemical potential) may be added to the effective action without changing the physics of the system. All such terms can be harmlessly absorbed by renormalisation.
This section has shown that the multiplicative anomaly has its roots in the manipulation of infinite sums - the non-interacting case in particular demonstrates how the anomaly can appear in relatively simple situations. Elizalde showed the existence of the multiplicative anomaly in possibly the simplest of all cases - infinite, diagonal matrices with real numbers . The first worked example in is striking in its similarity to the calculation presented above in the non-interacting model.
## IV Conclusions
We have demonstrated that the multiplicative anomaly originates in the zero-point energies of fields and is a consequence of shifting the energies by a constant amount. When regularisation is performed, these shifts ($`+\mu `$ and $`\mu `$ for example) are unable to cancel and result in a multiplicative anomaly. It seems therefore that in order to perform anomaly-free calculations one must resist the temptation to integrate over the momentum until the very end, after say, a high temperature expansion has been written down.
It should be borne in mind that the functional integral approach to quantum field theory is not as complete as the canonical one. As was mentioned in Sec. III, path integrals completely neglect normal ordering. Coleman discusses the merits of path integrals at length in his Erice lectures and echoes the comments of Bernard; the functional integral approach does not normal order.
From inspection of equations (26) and (27) it is tempting to conclude that $`\zeta `$-function regularisation may be to blame for the multiplicative anomaly (as opposed to functional integration). Certainly it does not seem likely that an anomaly would survive if say, dimensional regularisation were used to calculate the difference between (26) and (27). But it should be remembered that these equations were written down almost naïvely, to demonstrate the source of the anomaly; as was mentioned above, one should wait until the last possible moment before performing the integration over $`k`$, after everything that can cancel has done so. Nevertheless this does not seem very satisfactory. Why is there a difference between (28) and (29)? The mathematical properties of the $`\zeta `$-function are rigorous and well defined - to negatively criticise the whole subject of $`\zeta `$-function regularisation is a step not to be taken lightly. The paper by Elizalde provides some very interesting mathematical examples of multiplicative anomalies derived from infinite matrices, sometimes using nothing more than Riemann’s $`\zeta `$-function.
These mathematical peculiarities aside, it has been clearly demonstrated in this paper that the problem associated with the multiplicative anomaly can be removed by considering a Hamiltonian with a normal ordered charge operator. Normal ordering is in some ways an artificial procedure that physicists perform to make the theory more physical - it is not prescribed by the theory and there is no mathematical reason why it is done. Using canonical techniques it is easy to see how to normal order, unfortunately it is not so obvious in the functional integral approach - equation (2) is not normal ordered, but (3) is. An important step in understanding the multiplicative anomaly would be finding an a priori method of knowing which factorisations lead to anomalies and which do not, without having to compare with canonical calculations.
The problem of the multiplicative anomaly appears to be quite deeply rooted in the mathematics of infinite, divergent series.
###### Acknowledgements.
J. J. M-S. would like to thank Antonino Flachi for inspiring discussions.
## An alternative method of calculating $`\mathrm{\Gamma }_A^{(1)}`$
The method presented here for calculating (30) is simpler than that given in the main text, but has the disadvantage that the $`\zeta `$-function can only be evaluated at $`s=0`$. The method in the main text can be used for an arbitrary value of $`s`$ (although the integral in (17) may have to evaluated numerically).
We can re-write (9) as,
$$\zeta _A(s;a,b)=\left(\frac{2\pi }{\beta }\right)^{4s}\underset{n}{}\underset{j=\mathrm{}}{\overset{\mathrm{}}{}}\left(j^2+a^2\right)^s\left(j^2+b^2\right)^s$$
(60)
where $`a=\beta /2\pi (E_n+\mu )`$ and $`b=\beta /2\pi (E_n\mu )`$. Differentiating with respect to $`(a^2)`$ and $`(b^2)`$:
$`{\displaystyle \frac{}{(a^2)}}\zeta _A(s;a,b)`$ $`=`$ $`sf(s;a,b)`$ (61)
$`{\displaystyle \frac{}{(b^2)}}\zeta _A(s;a,b)`$ $`=`$ $`sg(s;a,b)`$ (62)
where
$`f(s;a,b)`$ $`=`$ $`\left({\displaystyle \frac{2\pi }{\beta }}\right)^{4s}{\displaystyle \underset{n}{}}{\displaystyle \underset{j=\mathrm{}}{\overset{\mathrm{}}{}}}\left(j^2+a^2\right)^{s1}\left(j^2+b^2\right)^s`$ (63)
$`g(s;a,b)`$ $`=`$ $`\left({\displaystyle \frac{2\pi }{\beta }}\right)^{4s}{\displaystyle \underset{n}{}}{\displaystyle \underset{j=\mathrm{}}{\overset{\mathrm{}}{}}}\left(j^2+a^2\right)^s\left(j^2+b^2\right)^{s1}.`$ (64)
Clearly,
$`{\displaystyle \frac{}{(a^2)}}\zeta _A(0;a,b)`$ $`=`$ $`0`$ (65)
$`{\displaystyle \frac{}{(b^2)}}\zeta _A(0;a,b)`$ $`=`$ $`0`$ (66)
and so we can conclude
$$\zeta _A(0;a,b)=𝒞$$
(67)
where $`𝒞`$ is a constant independent of $`a`$ and $`b`$. We can set $`a=b`$ in order to evaluate the left hand side of (67) and determine the constant $`𝒞`$. Using the Plana formula to calculate $`\zeta _A(s;a,a)`$ is straightforward:
$$\underset{j=\mathrm{}}{\overset{\mathrm{}}{}}\left(j^2+a^2\right)^{2s}=_{\mathrm{}}^{\mathrm{}}𝑑j\left(j^2+a^2\right)^{2s}+2_{\mathrm{}+iϵ}^{\mathrm{}+iϵ}\left(e^{2\pi iz}1\right)^1\left(z^2+a^2\right)^{2s}𝑑z$$
(68)
and so
$$\zeta _A(s;a,a)=\left(\frac{2\pi }{\beta }\right)^{4s}\underset{n}{}\left[\sqrt{\pi }\frac{\mathrm{\Gamma }(2s\frac{1}{2})}{\mathrm{\Gamma }(2s)}\left(a^2\right)^{\frac{1}{2}2s}4s\mathrm{ln}\left(1e^{2\pi a}\right)\right]$$
(69)
(See the appendix of for a more thorough evaluation of a similar sum.) Consequently $`\zeta _A(0;a,a)=𝒞=0`$. Next we need to evaluate $`\zeta _A^{}(0;a,b)`$,
$`{\displaystyle \frac{}{(a^2)}}\zeta _A^{}(0;a,b)`$ $`=`$ $`f(0;a,b)`$ (70)
$`{\displaystyle \frac{}{(b^2)}}\zeta _A^{}(0;a,b)`$ $`=`$ $`g(0;a,b)`$ (71)
The functions $`f(s;a,b)`$ and $`g(s;a,b)`$ can easily be calculated at $`s=0`$:
$`{\displaystyle \frac{}{(a^2)}}\zeta _A^{}(0;a,b)`$ $`=`$ $`{\displaystyle \underset{n}{}}{\displaystyle \frac{\pi }{a}}\mathrm{coth}(\pi a)`$ (72)
$`{\displaystyle \frac{}{(b^2)}}\zeta _A^{}(0;a,b)`$ $`=`$ $`{\displaystyle \underset{n}{}}{\displaystyle \frac{\pi }{b}}\mathrm{coth}(\pi b).`$ (73)
Integrating,
$$\zeta _A^{}(0;a,b)=\underset{n}{}\left\{2\mathrm{ln}\left[\mathrm{sinh}(\pi a)\right]2\mathrm{ln}\left[\mathrm{sinh}(\pi b)\right]\right\}+𝒦.$$
(74)
Again there is a constant $`𝒦`$, independent of $`a`$ and $`b`$. We can evaluate it in the same way as before,
$$𝒦=\zeta _A^{}(0;a,a)+4\underset{n}{}\mathrm{ln}\left[\mathrm{sinh}(\pi a)\right].$$
(75)
Using (69) we see that $`𝒦=4\mathrm{ln}2`$. Writing the $`\mathrm{sinh}`$s in (74) as exponentials we arrive at
$$\zeta _A^{}(0;a,b)=\underset{n}{}\left\{2\pi a2\pi b2\mathrm{ln}\left(1e^{2\pi a}\right)2\mathrm{ln}\left(1e^{2\pi b}\right)\right\}.$$
(76)
Substituting in the values of $`a`$ and $`b`$,
$$\zeta _A^{}(0)=\underset{n}{}\left\{\beta (E_n+\mu )\beta (E_n\mu )2\mathrm{ln}\left(1e^{\beta (E_n+\mu )}\right)2\mathrm{ln}\left(1e^{\beta (E_n\mu )}\right)\right\}$$
(77)
and so finally we have,
$`\mathrm{\Gamma }_A^{(1)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\zeta _A^{}(0)`$ (78)
$`=`$ $`{\displaystyle \underset{n}{}}\left\{\beta E_n+\mathrm{ln}\left(1e^{\beta (E_n+\mu )}\right)+\mathrm{ln}\left(1e^{\beta (E_n\mu )}\right)\right\}.`$ (79)
It is interesting to see that the zero-point energies in (77) are written initially with positive and negative chemical potentials.
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# 1 Stationary flow in a 2-dimensional parabolic potential, of which center is at the center of the figure.
Statistical Mechanics for States with Complex Eigenvalues
and
Quasi-stable Semiclassical Systems
T. KOBAYASHI<sup>1</sup> and T. SHIMBORI<sup>2</sup> <sup>1</sup>Department of General Education, Tsukuba College of Technology,
Tsukuba, Ibaraki 305-0005, Japan
<sup>2</sup>Institute of Physics, University of Tsukuba, Tsukuba, Ibaraki 305-8571, Japan
E-mail: <sup>1</sup>kobayash@a.tsukuba-tech.ac.jp and <sup>2</sup>shimbori@het.ph.tsukuba.ac.jp
## Abstract
Statistical mechanics for states with complex eigenvalues, which are described by Gel’fand triplet and represent unstable states like resonances, are discussed on the basis of principle of equal $`a\mathrm{𝑝𝑟𝑖𝑜𝑟𝑖}`$ probability. A new entropy corresponding to the freedom for the imaginary eigenvalues appears in the theory. In equilibriums it induces a new physical observable which can be identified as a common time scale. It is remarkable that in spaces with more than 2 dimensions we find out existence of stable and quasi-stable systems, even though all constituents are unstable. In such systems all constituents are connected by stationary flows which are generally observable and then we can say that they are semiclassical systems. Examples for such semiclassical systems are constructed in parabolic potential barriers. The flexible structure of the systems is also pointed out.
1. Introduction
It is well known that quantum mechanics including complex eigenvalues which represent unstable states like resonances are described by Gel’fand triplet . Since systems including such unstable states are usually unstable, it seems to be not suitable for constructing a theory like statistical mechanics for thermal equilibriums. We can, however, consider interesting situations, that is, the cases where the imaginary parts of energy eigenvalues of the states are so small that there are time enough to realize thermal equilibriums before the decay. In such situations statistical mechanics based on principle of equal $`a\mathrm{𝑝𝑟𝑖𝑜𝑟𝑖}`$ probability can be meaningful and we can observe some decay properties of such systems. It will possibly be a representation for quasi-stable systems having long life-times. Furthermore there is another noticeable feature of the complex eigenvalues such that they are represented by pairs of complex conjugates like $`a\pm ib`$ for $`a,b𝐑`$, that is, since Hamiltonians $`\widehat{H}`$ are real, for any solutions $`\psi `$ satisfying the equation
$$\widehat{H}\psi =(aib)\psi \mathrm{for}a,b𝐑,$$
we always find solutions having complex conjugate eigenvalues such that
$$\widehat{H}\psi ^{}=(a+ib)\psi ^{}.$$
An explicit example was presented for parabolic potential barrier . It is quite interesting that in spaces with more than 2 dimensions real eigenvalue solutions can be obtained by using those pair solutions represented only by Gel’fand triplet and they are in infinite degrees of degeneracy . In general they describe stationary states but they simultaneously have incoming- and outgoing-quantum stationary flows . Actually the real eigenvalue solutions are interpreted as the states which have equal incoming- and outgoing-flows. It is quite natural to have questions; whether systems constructed with such real-eigenvalue states can really be stable and furthermore how they can physically be interpreted in realistic phenomena. In order to study these problems we shall make statistical mechanics for states including complex eigenvalues on the basis of the same principle for the usual statistical mechanics, that is, principle of equal $`a\mathrm{𝑝𝑟𝑖𝑜𝑟𝑖}`$ probability.
In $`\mathrm{\S }`$2 statistical mechanics for complex eigenvalue states is performed. It is shown that the entropy of a system is represented by the sum of the usual entropy and a new one induced from the freedom of imaginary part and a new observable quantity appears in equilibriums, which will be identified as time scale of decay of the system. Explicit examples for stable and quasi-stable systems are presented in the case of parabolic potential barriers in $`\mathrm{\S }`$3, where those quasi-stable systems are connected by stationary flows . Some realistic applications will be remarked in $`\mathrm{\S }`$4.
2. Statistical mechanics for complex-eigenvalue states
Let us construct a statistical mechanics corresponding to microcanonical ensemble for states having complex energies which are generally represented by
$$\epsilon _{ij_i}=ϵ_ii\gamma _{j_i}\mathrm{for}ϵ_i,\gamma _{j_i}𝐑$$
(1)
where $`i,j_i𝐙_+(𝐙_+=\{0,1,2,\mathrm{}\})`$ and the suffix $`i`$ of $`j_i`$ is needed when there is some relation between the real and imaginary energy eigenvalues. We consider a simple case described by a system composed of $`N`$ independent particles being in complex-energy states. In this case the total energy of the $`N`$-particle system is given by the sum of energy eigenvalues of each particles such that
$$=Ei\mathrm{\Gamma },$$
(2)
where
$$E=\underset{i}{}ϵ_i\mathrm{and}\mathrm{\Gamma }=\underset{j_i}{}\gamma _{j_i}.$$
(3)
Here we shall investigate simple cases where the real and imaginary energy eigenvalues are independently determined and then we can take off the suffix $`i`$ from $`j_i`$. Such models will explicitly be presented in the next section. The basic principle is taken as same as that for the usual statistical mechanics, that is, principle of equal $`a\mathrm{𝑝𝑟𝑖𝑜𝑟𝑖}`$ probability. Then we start from counting the number of independent combinations of states for a fixed energy $``$. Since two freedoms concerning to the real and imaginary parts of energies are independent of each other, the number of the combination $`W()`$ is counted by the product of the number $`W^{\mathrm{}}(E)`$ for realizing the real part $`E`$ and that $`W^{\mathrm{}}(\mathrm{\Gamma })`$ for realizing the imaginary part $`\mathrm{\Gamma }`$ such that
$$W()=W^{\mathrm{}}(E)W^{\mathrm{}}(\mathrm{\Gamma }).$$
(4)
Following the procedure of statistical mechanics, we now see that the entropy $`S()=k_\mathrm{B}\mathrm{log}W()`$ of the system is written in terms of the sum of two entropies such that
$$S()=S^{\mathrm{}}(E)+S^{\mathrm{}}(\mathrm{\Gamma }),$$
(5)
where $`S^{\mathrm{}}(E)=k_\mathrm{B}\mathrm{log}W^{\mathrm{}}(E)`$ and $`S^{\mathrm{}}(\mathrm{\Gamma })=k_\mathrm{B}\mathrm{log}W^{\mathrm{}}(\mathrm{\Gamma })`$ are, respectively, the Boltzmann entropy and the new entropy induced from the freedom of the imaginary part.
Let us consider equilibrium between two systems which can each other transfer only energies. The total energy $`=Ei\mathrm{\Gamma }`$ given by the sum of those for two systems $`_\mathrm{I}=E_\mathrm{I}i\mathrm{\Gamma }_\mathrm{I}`$ and $`_{\mathrm{II}}=E_{\mathrm{II}}i\mathrm{\Gamma }_{\mathrm{II}}`$ is fixed. The number of the available combinations is written by the product of those for the two systems as
$$W()=W_\mathrm{I}(_\mathrm{I})W_{\mathrm{II}}(_{\mathrm{II}}),$$
(6)
where $`W_\mathrm{I}(_\mathrm{I})=W_\mathrm{I}^{\mathrm{}}(E_\mathrm{I})W_\mathrm{I}^{\mathrm{}}(\mathrm{\Gamma }_\mathrm{I})`$ and $`W_{\mathrm{II}}(_{\mathrm{II}})=W_{\mathrm{II}}^{\mathrm{}}(E_{\mathrm{II}})W_{\mathrm{II}}^{\mathrm{}}(\mathrm{\Gamma }_{\mathrm{II}})`$. Now we have the entropy expressed as the sum of four terms
$$S()=S_\mathrm{I}^{\mathrm{}}(E_\mathrm{I})+S_\mathrm{I}^{\mathrm{}}(\mathrm{\Gamma }_\mathrm{I})+S_{\mathrm{II}}^{\mathrm{}}(E_{\mathrm{II}})+S_{\mathrm{II}}^{\mathrm{}}(\mathrm{\Gamma }_{\mathrm{II}}),$$
(7)
where $`S_\mathrm{I}^{\mathrm{}}(E_\mathrm{I})=k_\mathrm{B}\mathrm{log}W_\mathrm{I}^{\mathrm{}}(E_\mathrm{I})`$ and so on. In the procedure maximizing the entropy $`S()`$ under the constraints that $`E=E_\mathrm{I}+E_{\mathrm{II}}`$ and $`\mathrm{\Gamma }=\mathrm{\Gamma }_\mathrm{I}+\mathrm{\Gamma }_{\mathrm{II}}`$ are fixed, we obtain two independent relations corresponding to the two constraints such that
$$\frac{S_\mathrm{I}^{\mathrm{}}(E_\mathrm{I})}{E_\mathrm{I}}=\frac{S_{\mathrm{II}}^{\mathrm{}}(E_{\mathrm{II}})}{E_{\mathrm{II}}},$$
(8)
$$\frac{S_\mathrm{I}^{\mathrm{}}(\mathrm{\Gamma }_\mathrm{I})}{\mathrm{\Gamma }_\mathrm{I}}=\frac{S_{\mathrm{II}}^{\mathrm{}}(\mathrm{\Gamma }_{\mathrm{II}})}{\mathrm{\Gamma }_{\mathrm{II}}}.$$
(9)
The first relation leads the usual temperature but the second one produce a new quantity which must be same for the two systems in the equilibriums. The canonical distribution for energy $`_{lm}=E_li\mathrm{\Gamma }_m`$ is written by
$$P(_{lm})=Z^1\mathrm{exp}(\beta ^{\mathrm{}}E_l\beta ^{\mathrm{}}\mathrm{\Gamma }_m),$$
(10)
where $`\beta ^{\mathrm{}}`$ should be chosen as the usual factor $`\beta =(k_\mathrm{B}T)^1`$ of canonical distribution, $`\beta ^{\mathrm{}}`$ denotes the new physical quantity in the equilibriums and the canonical partition function $`Z`$ is given by
$$Z=\underset{l}{}\underset{m}{}\mathrm{exp}(\beta E_l\beta ^{\mathrm{}}\mathrm{\Gamma }_m).$$
Note that we can not derive the canonical distribution (10) by replacing a real energy $`E`$ with a complex energy $``$ in the usual canonical distribution $`P(E)=Z^1\mathrm{exp}(\beta E)`$.
What is the new quantity $`\beta ^{\mathrm{}}`$? In independent particle systems wave functions are written by the product of all constituents such that
$$\mathrm{\Psi }(t,𝐫_1,\mathrm{},𝐫_N|)=\underset{n=1}{\overset{N}{}}\psi (t,𝐫_n|\epsilon _n),$$
(11)
where the wave function for one constituent with $`\epsilon _n=ϵ_ni\gamma _n`$ is generally given by
$$\psi (t,𝐫_n|\epsilon _n)=\mathrm{e}^{i\epsilon _nt/\mathrm{}}\varphi (𝐫_n).$$
(12)
The probability density for $`\mathrm{\Psi }`$ at $`t`$, which is normalized at $`t=0`$, is evaluated as
$`\rho (t,𝐫_1,\mathrm{},𝐫_N|)`$ $`=`$ $`|\mathrm{\Psi }(t,𝐫_1,\mathrm{},𝐫_N|)|^2`$ (13)
$`=`$ $`{\displaystyle \underset{n}{}}\mathrm{e}^{2\gamma _nt/\mathrm{}}|\varphi (𝐫_n)|^2`$
$`=`$ $`\mathrm{e}^{2\mathrm{\Gamma }t/\mathrm{}}{\displaystyle \underset{n}{}}|\varphi (𝐫_n)|^2.`$
We see that all the states with the same total imaginary energy $`\mathrm{\Gamma }`$ have the same time-dependence $`\mathrm{e}^{2\mathrm{\Gamma }t/\mathrm{}}`$. Since the states with complex energy eigenvalues are unstable, the canonical distribution $`P()`$ must depend on time. It is natural that the time-dependence of $`P()`$ is same as that of the probability density, which is determined by the imaginary part $`\mathrm{\Gamma }`$ of the total energy $``$ of the system. We can specify
$$\beta ^{\mathrm{}}=2t/\mathrm{}.$$
(14)
Thus we can introduce a common time scale $`t`$ for two systems being in equilibriums.
Remember that the imaginary parts $`\gamma _j`$ are expressed by pairs of conjugates, that is, $`\pm |\gamma _j|(j𝐙_+)`$ as noted in $`\mathrm{\S }`$1. This fact means that the total imaginary part $`\mathrm{\Gamma }`$ can possibly be in microscopic order (quantum size), even if the total real part $`E`$ is in macroscopic order. In fact such situations can be realized, when most of the constituents are in pairs of $`\pm |\gamma _j|`$ and number of non-pairing particles are very small in comparison with the macroscopic number $`N`$. In special cases $`\mathrm{\Gamma }=0`$ can happen. These systems have real energies and then there is no reason for the systems to be unstable in time. Remembering that the states with positive imaginary parts and those with negative ones, respectively, represent growing- and decaying-resonance states , we see that in such systems two processes, that is, growing- and decaying- resonance processes occur with the same probability. It is also noticeable that such systems with real energies are in infinite degrees of degeneracy, because the combinations of the pairing with zero imaginary energy are infinity. We will definitely see this situation in parabolic-potential-barrier models in the next section.
3. Stable and quasi-stable systems connected by quantum stationary flows
In the previous section the existence of stable many resonance systems have been pointed out. How can we understand such systems? Let us start from studying an explicit example in the 2-dimensinal parabolic potential case discussed in the paper . Before discussing on the imaginary energies, we note that the real energy eigenvalue is always represented by a fixed constant in parabolic potentials and then there is no freedom arising from the real energy eigenvalue. Imaginary eigenvalues $`\gamma _j`$ for the potential are given by the sum of two eigenvalues for two independent 1-dimensional parabolic potentials ($`V(x,y)=\frac{1}{2}m\gamma ^2(x^2+y^2)`$) such that
$$\gamma _{j_1j_2}=j\mathrm{}\gamma ,$$
(15)
where $`j=j_1+j_2`$ with $`j_1=\pm (n_1+\frac{1}{2})`$ and $`j_2=\pm (n_2+\frac{1}{2})`$ for $`n_1,n_2𝐙_+`$. It is interesting that we have $`\gamma _{j_1j_2}=0`$ for the cases with $`j_1=j_2=\pm (n+\frac{1}{2})(n𝐙_+)`$. Note that the zero imaginary eigenvalues are infinitely degenerate. As discussed in ref. 4, such zero imaginary states are well interpreted by the figure in fig. 1, where they are described by the wave function $`(\psi _{n_1n_2}(t,𝐫))`$ with equal incoming- and outgoing-stationary flows which are defined by the probability current
$$𝐣(t,𝐫)=\mathrm{}[\psi _{n_1n_2}(t,𝐫)^{}(i\mathrm{})\psi _{n_1n_2}(t,𝐫)]/m$$
with $`m=`$the mass of particle .
We see that the stationary flow incoming toward the center of the potential is expressed by the positive eigenvalue and that outgoing from the center done by the negative one. (In details, see ref. 4.)
Let us here consider $`N`$ pairs of a particle and a parabolic potential arranged on a line. Since the magnitudes of the two flows are equal at the same distance from the center, we can connect the flow outgoing from a potential to that incoming toward the next potential at the middle point of the two potentials. Such a situation is expressed in fig. 2.
Now we can easily construct 2-dimensional lattices in which all lattice points are connected by stationary flows as figured in fig. 3. It is remarkable that we can have stable systems composed of unstable constituents such as resonances. Note also that every
square made up from four lattice points is connected by a circular flow, that is, the total lattice system seems to be composed of small circular flows trapped in the squares. The fact that these lattices are not connected by wave functions but by stationary flows means that they can have some classical properties, because the flows are basically observable in quantum mechanics. Therefore we can call the lattices semiclassical systems. Note here that length between two neighboring lattice points are arbitrary and then structures of lattices are not rigid but rather flexible.
The following point should be noticed that at the edge of the system incoming- and outgoing-flows can generally not be connected, that is, the system have absorbing points and sending-out points of the stationary flows in environments (heat baths) like breathing of living matters. This fact implies that those systems can generally be unstable and decay at the edge in the 2-dimensions when they are put in unsuitable environments, where the absorbing flows and the sending-out flows are not arranged with an equal weight. Such quasi-stable systems, however, seems to be quite interesting to describe realistic systems slowly decaying from their edges, which can be seen quite often in our daily life. Of course, we can have stable systems on closed surfaces like soccer balls and torus.
In 3-dimensional isotropic parabolic potentials we have no stable states with zero imaginary energy eigenvalue, because energies induced from zero-point energy eigenvalues $`\pm \frac{i}{2}\mathrm{}\gamma `$ do not cancel each other in odd dimensions. Provided that the product of the curvature of potential $`\gamma `$ and the number of constituents $`N`$ is so small that the time scale given by $`\tau =(\frac{1}{2}\gamma N)^1`$ is enough large in comparison with the time scale realizing thermal equilibriums, we can observe metastable systems which decay not only on the edges but in the inside of the systems as well.
Even in odd dimensions we find out two possibility that real eigenvalues appear from the compositions of complex eigenvalues in 3-dimensional spaces. One is the case where three coupling constants ($`\gamma _1,\gamma _2,\gamma _3`$) representing the curvatures in 3-dimensions are different from each other but satisfied by some relation, for instance, $`\gamma _1+\gamma _2=\gamma _3`$ and three eigenvalues ($`\pm i(n_1+\frac{1}{2})\mathrm{}\gamma _1,\pm i(n_2+\frac{1}{2})\mathrm{}\gamma _2,i(n_3+\frac{1}{2})\mathrm{}\gamma _3`$) fulfill the relations $`n_1\gamma _1+n_2\gamma _2=n_3\gamma _3.`$ States satisfying these two relations can have zero imaginary eigenvalues. Note that at least one state specified by $`n_1=n_2=n_3=0`$ satisfies the relations. These relations means that all flows incoming from two directions go out from one direction and vice versa. Another possibility for zero imaginary eigenvalues is seen in the case that potentials is described by 2-dimensional parabolic potential and 1-dimensional harmonic one. Imaginary energy-eigenvalues of total systems can be zero. In these cases we can find very rich and complicated structures in quasi-stable systems. In both cases the decays of those systems can occur only on their surfaces including edges.
4. Remarks
We see that statistical mechanics based on principle of equal a priori probability can be applicable to systems composed of many unstable states. To complete statistical mechanics, we still have a lot of problems to investigate, for instance, how to define and interpret thermal quantities corresponding to free energy, pressure, heat capacity and so forth. Ergodic theorem have also to be discussed in the prsent scheme. We have, however, found out two remarkable results. One is the introduction of a common time scale, which will possibly have some connection with the common time of our universe. The other is the existence of stable and quasi-stable systems which seem to be interesting candidates to describe realistic systems having long life-times. It should be stressed that the structure of those systems can be very much flexible because the distances between any neighboring constituents are arbitrary.
A toy model can represented by systems for gases of which constituents have parabolic potential barriers in each other. Such systems can completely be separable into center of mass motions and relative motions written by the same parabolic potential. The procedure for the statistical mechanics performed in $`\mathrm{\S }`$2 can directly be applicable to this model and reproduce usual temperature from the center of mass motions and time scale from relative motions. In this model growing- and decaying- resonance states, respectively, describe two constituents approaching and leaving each other. Actually 2-dimensional models seems to be applicable to the investigation of problems with respect to surfaces of materials.
It is noticeable that many body systems can possibly be quasi-stable even under such repulsive potentials. Systems composed of many bodies having repulsive forces between any pairs of constituents will be realized electron gases, that is, electron plasma. We have to solve complex eigenvalue problem for the repulsive Coulomb potential. In order to say whether such quasi-stable systems exist or not, it is enough to show whether there exist complex eigenvalue solutions or not. We need not understand any details of the complex eigenvalues, because the existence of complex conjugate pairs is guaranteed in the theory. We shall study this problem in more detail.
Up to now we have discussed on repulsive potentials. As studied in textbooks , it is known that attractive potentials such as shown in fig. 4 have complex eigenvalue solutions.
For instance, electric potentials of atoms for incoming electrons will be approximated by such potentials. Namely there is an attractive force near the center arising from the positive charge of nucleus, which is surrounded by the repulsive force induced from the electrons trapped in the atom. There will be a possibility for interpreting chemical bonds in this scheme.
It is also pointed out that these quasi-stable states are connected by flows which are generally observable in quantum mechanics. This property indicates that those states can have some classical properties. It seems to be very attractive to describe phenomena being in borders between quantum processes and classical ones such as mesoscopic processes in terms of the present scheme. We should notice that the idea of the connection by flows are applicable to real energy eigenvalues states contained in usual Hilbert spaces.
Finally we would like to comment on cases where imaginary eigenvalues are not independent of real ones. In general there may be some relations between real and imaginary energy eigenvalues ($`ϵ_j`$ and $`\gamma _j`$). In such cases the formula (4) for the number of combinations should be written as
$$W()=\underset{ϵ_1}{}\mathrm{}\underset{ϵ_N}{}\delta _{_jϵ_j,E}W^{\mathrm{}}(\mathrm{\Gamma }:ϵ_1,\mathrm{},ϵ_N),$$
(16)
where $`\delta _{_jϵ_j,E}`$ is the Kronecker delta symbol representing the relation $`_jϵ_j=E`$ and $`W^{\mathrm{}}(\mathrm{\Gamma }:ϵ_1,\mathrm{},ϵ_N)`$ stands for the sum over states with respect to the freedom of the imaginary energies when the real energy eigenvalues of all constituents are fixed. It is obvious that the above formula turns to that of (4), when $`W^{\mathrm{}}(\mathrm{\Gamma }:ϵ_1,\mathrm{},ϵ_N)`$ does not depend on the real energy eigenvalues. We can follow the same procedure in deriving the entropy. It is, however, very hard to obtain a general formula for the entropy, because the situations so much depend on the relations between the real and imaginary eigenvalues. We have to investigate very carefully and precisely the relation for each interaction, that is, we have to solve Gel’fand-triplet problems for the interaction. It is, however, stressed that the existence of quasi-stable states composed of complex energy eigenstates is guaranteed because of pairings of complex energy eigenvalues, when states having non-zero imaginary-energies are involved in Gel’fand-triplet solutions. We should also investigate the possibility of entropy transfer between two entropies $`S^{\mathrm{}}`$ and $`S^{\mathrm{}}`$ given in (5).
References A. Bohm and M. Gadella, Dirac Kets, Gamow Vectors and Gel’fand Triplets, Lecture Notes in Physics, Vol. 348 (Springer, Berlin) 1989. T. Shimbori and T. Kobayashi, Il Nuovo Cimento B (to appear, 2000). T. Shimbori, Operator Methods of the Parabolic Potential Barrier, UTHEP-415, quant-ph/9912073. T. Shimbori and T. Kobayashi, Stationary Flows of Parabolic Potential Barrier in Two Dimensions (to appear, 2000). T. Shimbori and T. Kobayashi, “Velocities” in Quantum Mechanics, UTHEP-422, quant-ph/0004086. For instance, see the following textbook. L. D. Landau and E. M. Lifshitz, Quantum Mechanics 3rd ed. (Pergamon, Oxford, 1977).
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# Covariant Schwinger Terms
## 1 Introduction
One essential feature of chiral gauge theories is the violation of gauge invariance when chiral (Weyl) fermions are quantized. This loss of gauge invariance results in a non-invariance of the vacuum functional in an external gauge field under gauge transformations and in a non-conservation (anomalous divergence) of the corresponding gauge current in a Lagrangian (or space-time) formulation (“anomaly” ), or in anomalous contributions to the equal-time commutators of the generators of time-independent gauge transformations in a Hamiltonian formulation (“Schwinger term” or “commutator anomaly”, ). One regularization scheme, where all the anomalous terms are related to functional derivatives of the vacuum functional, leads to the so-called consistent anomalies. These anomalies have to obey certain consistency conditions because of their relation to functional derivatives . One way of determining these consistent anomalies (up to an overall constant) is provided by the descent equations of Stora and Zumino . They provide a simple algebraic scheme – based on some geometrical considerations – for the computation of the consistent chiral anomaly in the space-time formalism and for the corresponding equal-time commutator anomaly, the Schwinger term (as well as for higher cochain terms).
On the other hand, it is possible to choose a gauge-covariant regularization for the gauge current. This covariant current cannot be related to a functional derivative of the vacuum functional (because of the gauge non-invariance of the latter). As a consequence, the covariant current anomaly does not obey the consistency condition. Nevertheless, there exists a covariant counterpart for each consistent cochain in the descent equations. The first derivation of an algebraic computational scheme for covariant cochains appears to be the one by Tsutsui , using the anti-BRST formalism. Further, a covariant version of the descent equations was derived by Kelnhofer in . The covariant cochains resulting from the calculations by Tsutsui on one hand, and by Kelnhofer on the other hand are, in fact, different, as was shown in . Tsutsui’s and Kelnhofer’s formulas predict the same anomaly in space-time, but their Schwinger terms differ by a sign. The higher cochains (with more than 2 ghosts) seem to be unrelated. We shall give an answer to which of the two formulas is correct, in the sense that it is reproduced by a full quantum field theoretic calculation.
The easiest way to do this calculation would be to compute one of the higher covariant anomalies in some quantum field theoretic setting, because these higher covariant anomalies are given by completely different expressions in and . However, although it has been claimed that the higher cochains can have a physical meaning, this is far from understood. It is therefore not sound to use these terms to argue which of the two formulas is correct. Instead, we shall use the sign of the Schwinger term as a referee.
We shall use three methods to determine the correct expression for the covariant Schwinger term. They all have to be used with care since we are after the sign difference between Tsutsui’s and Kelnhofer’s predictions. The first two methods are to apply the quantum field theoretical calculation schemes that have been used by Adam and by Hosono and Seo , respectively.
The third method is the one by Wess , relating the Schwinger term (consistent or covariant) in any even dimensional space-time with the corresponding space-time anomaly. This method was used by Schwiebert for the consistent case and by Kelnhofer in the covariant case. Thus, by using the expression for the covariant anomaly (which everyone agrees on) the covariant Schwinger term can be determined. Again, care has to be taken. For this, we first perform the calculation in the consistent formalism and set conventions so the result agrees with what is predicted by the descent equations. The corresponding covariant computation is then to determine the covariant Schwinger term including the sign. We shall perform these calculation in 1+1 and 3+1 dimensions.
From the explicit calculations we find that our quantum field theoretical methods produce the same expression (i.e. sign) as Kelnhofer’s covariant descent equations. This result is not obvious since Kelnhofer’s approach, as well as Tsutsui’s, seems to be based only on the requirement of covariance. We shall however show that there is in fact a natural interpretation of the Kelnhofer formula, as one would expect.
Our paper is organized as follows. In Section 2 we briefly describe the geometrical setting for the description of anomalies and review the derivation of the consistent and covariant chain terms. In the covariant case, both Tsutsui’s and Kelnhofer’s versions of the chain terms are given and a geometrical description of Kelnhofer’s construction is provided. In Section 3 the consistent and covariant Schwinger terms in 1+1 dimensions are calculated using the methods of and of . Finally, in Section 4 the Schwinger terms in 1+1 as well as in 3+1 dimensions are calculated with the help of the method of Wess .
## 2 Consistent and covariant cochains
We shall start with deriving the consistent chiral anomaly for a non-abelian gauge theory. Consider therefore Weyl fermions $`\psi `$ coupled to an external gauge field $`A𝒜`$. $`𝒜`$ is the affine space of of gauge connections and the gauge group $`G`$ is assumed to be a compact, semi-simple matrix group. We assume that the space-time $`M`$ is a smooth, compact, oriented, even-dimensional and flat Riemannian spin manifold without boundary. The group $`𝒢`$ of gauge transformations consists of diffeomorphisms of a principal bundle $`P\stackrel{G}{}M`$ such that the base remains unchanged. It acts on $`𝒜`$ by pull-back and to make this action free we restrict to gauge transformations that leaves a reference point $`p_0P`$ fixed.
The generating functional is given by
$$\mathrm{exp}(W(A))=_{\psi ,\overline{\psi }}\mathrm{exp}(_M\overline{\psi }/_A^+\psi d^{2n}x),$$
(1)
where $`W`$ is the effective action and $`/_A^+=/_A(1+\gamma _5)/2=\gamma ^\mu (_\mu +A_\mu )(1+\gamma _5)/2`$. We shall use conventions such that $`\gamma ^\mu `$ is hermitean and $`A_\mu `$ is anti-hermitean. It has been argued that a correct interpretation of the generating functional is as a section of the determinant line bundle $`\text{DET}i/_A=\text{det ker}i/_A^+(\text{det coker}i/_A^+)^{}`$. It can be viewed as a functional by comparing with some reference section. Associated with the determinant line bundle is a connection with corresponding curvature
$$F=2\pi i\frac{1}{(n+1)!}\left(\frac{i}{2\pi }\right)^{n+1}_M\text{tr}\left(^{n+1}\right),$$
(2)
. Here, $`=(d+\delta )(A+(d_A^{}d_A)^1d_A^{})+(A+(d_A^{}d_A)^1d_A^{})^2`$ is a curvature of the principal bundle $`P\times 𝒜M\times 𝒜`$ and $`\delta `$ is the exterior differential in $`𝒜`$. The choice of $``$ is motivated by gauge invariance of the determinant line bundle, .
Recall that
$$\text{tr}\left(_2^{n+1}\right)\text{tr}\left(_1^{n+1}\right)=(d+\delta )\omega _{2n+1}(\alpha _2,\alpha _1)$$
(3)
for
$$\omega _{2n+1}(\alpha _2,\alpha _1)=(n+1)_0^1𝑑t\text{tr}\left((\alpha _2\alpha _1)_t^n\right)$$
(4)
and $`_t`$ the curvature of $`(1t)\alpha _1+t\alpha _2`$, holds for any connections $`\alpha _1,\alpha _2`$ with curvatures $`_1,_2`$. Using this in eq. (2) gives the following expression for the connection of the determinant line bundle:
$$2\pi i\frac{1}{(n+1)!}\left(\frac{i}{2\pi }\right)^{n+1}_M\omega _{2n+1}(A+(d_A^{}d_A)^1d_A^{},0).$$
(5)
The (infinitesimal) consistent anomaly is the variation of the effective action under gauge transformations. Thus, it is the negative of the restriction of (5) to gauge directions, i.e. the fibre directions of $`𝒜𝒜/𝒢`$. Along such directions, $`\delta `$ becomes the BRST operator and $`(d_A^{}d_A)^1d_A^{}`$ becomes the ghost $`v`$. Thus, the consistent anomaly is $`c_n_M\omega _{2n+1}(A+v,0)`$ , with $`c_n=\frac{1}{(n+1)!}\left(\frac{i}{2\pi }\right)^n`$. Since $`M`$ is $`2n`$-dimensional, it is only the term with one ghost in the expansion of $`\omega _{2n+1}(A+v,0)`$ that will give a contribution to the anomaly. We let $`\omega _{2n+1k}^k(A+v,0)`$ denote the part of $`\omega _{2n+1}(A+v,0)`$ that contains $`k`$ number of ghosts. Then $`c_n\omega _{2n}^1(A+v,0)`$ is the non-integrated anomaly. It is well-known, and explicitly proven in , that $`c_n_M\omega _{2n1}^2(A+v,0)`$ is the Schwinger term. In this case $`M`$ is to be interpreted as the odd-dimensional physical space at a fixed time. The forms $`\omega _{2n+1k}^k(A+v,0)`$ can be computed by use of eq. (4). In 1+1 and 3+1 dimensions it gives the following result for the consistent anomaly and Schwinger term:
$`c_1\omega _1^2(A+v,0)`$ $`=`$ $`c_1\text{tr}(vdA)`$
$`c_1\omega _2^1(A+v,0)`$ $`=`$ $`c_1\text{tr}\left(v^2A\right)`$
$`c_2\omega _1^4(A+v,0)`$ $`=`$ $`c_2\text{tr}\left(vd\left(AdA+A^3/2\right)\right)`$
$`c_2\omega _2^3(A+v,0)`$ $`=`$ $`c_2\text{tr}\left(\left(v^2A+vAv+Av^2\right)dA+v^2A^3\right)/2.`$ (6)
If the freedom is used to change the forms $`\omega _{2n+1k}^k(A+v,0)`$ by cohomologically trivial terms (i.e. by coboundaries), then these forms can be given by the following compact expressions that were first derived by Zumino in :
$`\omega _{2n+1k}^k(A+v,0)`$ (9)
$``$ $`(n+1)\left(\begin{array}{c}n\\ k\end{array}\right){\displaystyle _0^1}𝑑t(1t)^k\text{str}((dv)^k,A,\left(tdA+t^2A^2\right)^{nk})`$
when $`0kn`$ and
$`\omega _{2n+1k}^k(A+v,0)`$ $``$ $`(1)^{kn1}\left(\begin{array}{c}n\\ kn1\end{array}\right)\left(\left(\begin{array}{c}k\\ kn1\end{array}\right)\right)^1`$ (14)
$`\times \text{str}(v,(v^2)^{kn1},\left(dv\right)^{2nk+1})`$
when $`n+1k2n+1`$. Here str means the symmetrized trace and $``$ means equality up to a coboundary.
Above, we used eq. (3) for $`\alpha _2=A+(d_A^{}d_A)^1d_A^{}`$ and $`\alpha _1=0`$. When $`P`$ is a non-trivial bundle it is no longer possible to let $`\alpha _1`$ be zero. Instead, we let it be some fixed connection $`A_0`$ on $`P`$ (which can be identified with a connection on $`P\times 𝒜`$). By dimensional reasons, this does not change eq. (2). The consistent anomaly and Schwinger term with such a background connection can be computed in a similar way as above, one just uses $`\omega _{2n+1k}^k(A+v,A_0)`$ instead. Since the expressions corresponding to eq. (9) and eq. (14) are long and not particularly illuminating we shall not present them here (parts of it can be found in ). The ideas behind the background connection are completely analogous with the case without a background. For example, they are consistent, but not gauge covariant. To obtain covariance, we choose as a background the field itself. We are then interested in the (non-consistent) terms coming from the expansion of $`\omega _{2n+1k}^k(A+v,A)`$ in various ghost degrees. With use of eq. (4), the following expression was obtained in :
$`\omega _{2n+1k}^k(A+v,A)`$ $`=`$ $`{\displaystyle \underset{j=0}{\overset{[(k1)/2]}{}}}{\displaystyle \frac{n+1}{kj}}\left(\begin{array}{c}nj\\ k2j1\end{array}\right)\left(\begin{array}{c}n\\ j\end{array}\right)\left(\left(\begin{array}{c}k\\ j\end{array}\right)\right)^1`$ (22)
$`\times \text{str}(v,(\delta v)^j,(\delta A)^{k2j1},F^{nk+j+1}),`$
where a negative power on a factor means that the corresponding term is absent in the sum. Recall that $`[(k1)/2]`$ is $`(k1)/2`$ if $`k`$ is odd and $`(k2)/2`$ if $`k`$ is even. The terms
$`c_n\omega _{2n}^1(A+v,A)`$ $`=`$ $`c_n(n+1)\text{tr}(vF^n)`$
$`c_n\omega _{2n1}^2(A+v,A)`$ $`=`$ $`c_n{\displaystyle \frac{n(n+1)}{2}}\text{str}(v,\delta A,F^{n1})`$ (23)
are the non-integrated covariant anomaly and Schwinger term.
Let us summarize the results so far in the case of 1+1 and 3+1 dimensions in Tables 1 and 2, respectively (we use eq. (9) and (14) for the consistent formalism)
| $`n=1`$ | Anomaly | Schwinger term |
| --- | --- | --- |
| | $`c_1_M\text{tr}((dv)A)`$ | $`c_1_M\text{tr}(vdv)`$ |
| | $`\begin{array}{c}c_12_M\text{tr}\left(v\left(dA+A^2\right)\right)\end{array}`$ | $`c_1_M\text{str}(v(dv+2vA))`$ |
Table 1
| $`n=2`$ | Anomaly | Schwinger term |
| --- | --- | --- |
| | $`\begin{array}{c}c_2_M\text{tr}((dv)AdA\\ +\frac{1}{2}(dv)A^3)\end{array}`$ | $`c_2_M\text{tr}\left((dv)^2A\right)`$ |
| | $`\begin{array}{c}c_23_M\text{tr}\\ \left(v\left(dA+A^2\right)^2\right)\end{array}`$ | $`\begin{array}{c}c_2(3)_M\text{str}(v(dv+vA\\ +Av)(dA+A^2))\end{array}`$ |
Table 2
We shall now evaluate these forms on (anti-hermitean) infinitesimal gauge transformations $`X,Y\text{Lie}𝒢`$. Let us do this explicitely for the consistent Schwinger term when $`n=2`$:
$`c_2{\displaystyle _M}\text{tr}\left((dv)^2A\right)(X,Y)=c_2{\displaystyle _M}\text{tr}(_iv_jvA_k)ϵ^{ijk}d^3x(X,Y)`$ (24)
$`=`$ $`c_2{\displaystyle _M}\text{tr}((_iX_jY_iY_jX)A_k)ϵ^{ijk}d^3x.`$
The corresponding evaluation of the other forms for $`n=1`$ and $`n=2`$ is listed in Tables 3 and 4, respectively.
| $`n=1`$ | Anomaly | Schwinger term |
| --- | --- | --- |
| | $`c_1\text{tr}((_\mu X)A_\nu )ϵ^{\mu \nu }`$ | $`2c_1\text{tr}\left(X_xY\right)`$ |
| | $`\begin{array}{c}2c_1\text{tr}(X(_\mu A_\nu \\ +A_\mu A_\nu ))ϵ^{\mu \nu }\end{array}`$ | $`\begin{array}{c}2c_1\text{tr}(X_xY[X,Y]A_x)\end{array}`$ |
Table 3
| $`n=2`$ | Anomaly | Schwinger term |
| --- | --- | --- |
| | $`\begin{array}{c}c_2\text{tr}((_\mu X)(A_\nu _\rho A_\lambda \\ +\frac{1}{2}A_\nu A_\rho A_\lambda ))ϵ^{\mu \nu \rho \lambda }\end{array}`$ | $`\begin{array}{c}c_2\text{tr}((_iX_jY\\ _iY_jX)A_k)ϵ^{ijk}\end{array}`$ |
| Covariant | $`\begin{array}{c}3c_2\text{tr}(X(_\mu A_\nu \\ +A_\mu A_\nu ))(_\rho A_\lambda \\ +A_\rho A_\lambda )ϵ^{\mu \nu \rho \lambda }\end{array}`$ | $`\begin{array}{c}3c_2\text{tr}((X_iYY_iX\\ [X,Y]A_i+XA_iY\\ YA_iX)(_jA_k+A_jA_k))ϵ^{ijk}\end{array}`$ |
Table 4
That the covariant anomaly and Schwinger term can be computed by expansion of $`\omega _{2n+1}(A+v,A)`$ was discovered by Kelnhofer . An alternative computational scheme leading to covariant cochains differing from the ones of Kelnhofer was given by Tsutsui . To review his approach we reconsider eq. (4) for $`\alpha _2=A+v`$ and $`\alpha _1=0`$. We can then view $`\omega _{2n+1}`$ as a function $`\omega _{2n+1}(A+v|)`$ of $`A+v`$ and $`=(d+\delta )(A+v)+(A+v)^2`$:
$$\omega _{2n+1}(A+v|)=(n+1)_0^1𝑑t\text{tr}\left((A+v)_t^n\right),_t=t+(t^2t)(A+v).$$
(25)
The covariance is broken by the operator $`\delta `$ in the expression for $``$. Thus, $`\omega _{2n+1}(A+v|^{})`$, with $`^{}=d(A+v)+(A+v)^2`$, produces covariant terms. This is exactly the same terms as the ones appearing in Tsutsui’s anti-BRST approach . In (see for $`k=2`$) the following formula was given for the terms with a given ghost degree:
$$\omega _{2n+1k}^k(A+v|^{})=\frac{n+1}{k}\text{tr}\left(v(F\delta (A+v))^n\right)_k,$$
(26)
where the index $`k`$ on the right hand side means the part of the expression that has $`k`$ number of ghosts. Comparison with eq. (22) reveals that this formula gives the same covariant anomaly but the covariant Schwinger term differs by a sign. The higher terms seem to be unrelated. This brings us to the question of who is right: Kelnhofer or Tsutsui? The formula of Tsutsui seems to be motivated by nothing else than covariance. Kelnhofer’s formula, on the other hand, seems to appear in a natural way: it is obtained by putting the background field equal to the field under consideration. In the computation of the Schwinger term from determinant line bundles for manifolds with boundary, one extends space to a cylindrical space-time, . On one side of the cylinder one computes the Schwinger term by comparison of a fixed vacuum bundle (with respect to a background connection) on the other side of the cylinder. In this approach it is certainly possible to put the background field equal to the field itself, see for details. This clearly defines a covariant Schwinger term in a natural way, suggesting that Kelnhofer’s approach is the correct one. This geometrical approach would not have been possible with Tsutsui’s result. This explains the importance of the sign of the covariant Schwinger term. In the forthcoming sections we shall demonstrate that indeed Kelnhofer’s result for the covariant Schwinger term is reproduced by quantum field theoretic computations.
## 3 Calculations in 1+1 dimensions
### 3.1 Calculation of Adam
In this section we want to briefly review the calculation of the consistent and covariant Schwinger term that was performed in for the Abelian case (the chiral Schwinger model). The generalization to the non-Abelian case is straight-forward and shall be displayed below, as well. In the Hamiltonian formulation was used (therefore space-time is 1+1 dimensional Minkowski space), and the computation started from the second-quantized chiral fermion field operator in the interaction picture. For fermionic field operators the Dirac vacuum has to be introduced and operator products have to be normal-ordered w.r.t. the Dirac vacuum. For the introduction of the Dirac vacuum the Hilbert space of fermionic states is split into a positive and negative momentum sub-space (for chiral fermions in two dimensions energy equals momentum). For the negative momentum sub-space the role of creation and annihilation operators is then exchanged. At this point there are two possibilities to split. Either one may split w.r.t. eigenvalues of the free momentum operator $`i_{x^1}`$ and perform normal-ordering (denoted by $`N`$) for this Dirac vacuum. A well-known consequence of this normal-ordering is the fact that the current commutators acquire a central extension (Schwinger term). For a fermion of positive chirality (where the current obeys $`J^0=J^1=:J`$), the Schwinger term is
$$[NJ(x^0,x^1),NJ(x^0,y^1)]=\frac{i}{2\pi }\delta ^{}(x^1y^1)$$
(27)
(here the prime denotes derivative w.r.t. the argument). The second possibility is to split w.r.t. eigenvalues of the kinetic momentum operator $`i_{x^1}+eA_1`$. Again, a corresponding Dirac vacuum and normal ordering (denoted by $`\stackrel{~}{N}`$) may be introduced. It turns out that the kinetically normal-ordered current is related to the conventionally normal-ordered current in a simple fashion
$$\stackrel{~}{N}J(x)=NJ(x)+\frac{e}{2\pi }A_1(x)$$
(28)
therefore $`\stackrel{~}{N}J`$ has the same commutator (14) as $`NJ`$. It was proven in that $`NJ`$ is the consistent current operator and $`\stackrel{~}{N}J`$ is the covariant current operator. Now it is very easy to compute the consistent and covariant Gauss law commutators. The Gauss law operators are defined as ($`_{x^1}_1`$)
$$G(x)=_1\frac{\delta }{e\delta A_1(x)}iNJ(x)$$
(29)
$$\stackrel{~}{G}(x)=_1\frac{\delta }{e\delta A_1(x)}i\stackrel{~}{N}J(x)$$
(30)
Here $`A_1(x)`$ is treated as a function of space only and the time variable $`x^0`$ as a parameter, i.e. $`(\delta /\delta A_1(x^0,x^1))A_1(x^0,y^1)=\delta (x^1y^1)`$. The consistent Gauss law commutator is determined by the current commutator (14),
$$[G(x^0,x^1),G(x^0,y^1)]=\frac{i}{2\pi }\delta ^{}(x^1y^1)$$
(31)
whereas for the covariant case the functional derivatives contribute, as well,
$$[\stackrel{~}{G}(x^0,x^1),\stackrel{~}{G}(x^0,y^1)]=\frac{i}{2\pi }\delta ^{}(x^1y^1).$$
(32)
Therefore, the covariant Schwinger term is minus the consistent one, (18). This relative minus sign is precisely as in Table 1. Observe that the covariant current is indeed gauge invariant, $`[\stackrel{~}{G}(x^0,x^1),\stackrel{~}{N}J(x^0,y^1)]=0`$, as it must be. In fact, the relative minus sign between the consistent and covariant Schwinger terms is a consequence of this gauge invariance of $`\stackrel{~}{N}J`$, and therefore independent of all possible conventions.
A generalization of the above results to the nonabelian case is straight forward. The two versions of normal-ordering are defined as in the abelian case, and they lead to the same relation as in (15), up to an additional colour index
$$\stackrel{~}{N}J^a(x)=NJ^a(x)+\frac{e}{2\pi }A_1^a(x).$$
(33)
Further, the current commutator acquires a canonical piece as well,
$$[NJ^a(x^0,x^1),NJ^b(x^0,y^1)]=if^{abc}NJ^c(x^0,x^1)\delta (x^1y^1)\frac{i}{2\pi }\delta ^{ab}\delta ^{}(x^1y^1)$$
(34)
(for the commutator $`[\stackrel{~}{N}J^a(x),\stackrel{~}{N}J^b(y)]`$, the same expression is obtained, again with $`NJ^c`$ on the r.h.s., not $`\stackrel{~}{N}J^c`$, as is obvious from (20)). The generator of time-independent gauge transformations on gauge fields,
$$\delta ^a(x):=(\delta ^{ab}_1+ef^{acb}A_1^c(x))\frac{\delta }{e\delta A_1^b(x)}$$
(35)
obeys the commutation relation
$$[\delta ^a(x^0,x^1),\delta ^b(x^0,y^1)]=f^{abc}\delta ^c(x^0,x^1)\delta (x^1y^1).$$
(36)
The consistent and covariant Gauss law operators are defined as
$$G^a(x)=\delta ^a(x)iNJ^a(x)$$
(37)
and
$$\stackrel{~}{G}^a(x)=\delta ^a(x)i\stackrel{~}{N}J^a(x)$$
(38)
respectively. Their anomalous commutators may be easily computed,
$$[G^a(x^0,x^1),G^b(x^0,y^1)]+f^{abc}G^c(x^0,x^1)\delta (x^1y^1)=\frac{i}{2\pi }\delta ^{ab}\delta ^{}(x^1y^1)$$
(39)
$$[\stackrel{~}{G}^a(x^0,x^1),\stackrel{~}{G}^b(x^0,y^1)]+f^{abc}\stackrel{~}{G}^c(x^0,x^1)\delta (x^1y^1)=$$
$$\frac{i}{2\pi }\delta ^{ab}\delta ^{}(x^1y^1)+\frac{i}{2\pi }f^{abc}A^c(x^0,x^1)\delta (x^1y^1).$$
(40)
As in the abelian case, the anomalous commutators agree with the ones in Table 1, and again this is most easily seen for the relative minus sign of the $`\delta ^{}(x^1y^1)`$ term. This relative sign may be related to the fact that the covariant current has to transform covariantly under a gauge transformation (i.e., the $`\delta ^{}`$ terms must cancel)
$$[\stackrel{~}{G}^a(x^0,x^1),\stackrel{~}{N}J^b(x^0,y^1)]=f^{abc}\stackrel{~}{N}J^c(x^0,x^1)\delta (x^1y^1)$$
(41)
as may be checked easily.
### 3.2 Calculation of Hosono and Seo
In this section we shall use the Hosono and Seo approach for the calculation of the equal-time commutators of the covariant and consistent Gauss law operator. The calculation is performed in Minkowski space, $`g_{\mu \nu }=\mathrm{diag}(1,1)`$, $`\epsilon _{01}=1`$, with the gamma matrices obeying the usual Clifford algebra relation $`\gamma ^\mu \gamma ^\nu +\gamma ^\nu \gamma ^\mu =2g^{\mu \nu }`$, and $`\gamma _5=\gamma ^0\gamma ^1`$. The anti-Hermitian matrices $`t^i`$ are the generators of a non-abelian algebra $`[t^a,t^b]=f^{abc}t^c`$, and we denote $`A_k=A_k^at^a`$.
The Hamiltonian of the chiral fermion interacting with an external gauge potential is
$$\left(A\right)=i𝑑x\left[\overline{\psi }(t,x)\gamma ^1\frac{1+\gamma _5}{2}\left(_1+A_1^a(t,x)t^a\right)\psi (t,x)\right]$$
(42)
where we chose the Weyl gauge $`\left(A^0(t,x)=0\right)`$.
We expand the Fermion field as
$$\psi (t,x)=\underset{n}{}\alpha _n(t)\zeta _n(t,x),$$
(43)
where $`\zeta _n(t,x)`$ are eigenfunctions of the full Hamiltonian (29) with eigenvalues $`E_n(t)`$. In the quantized theory the $`\alpha _n`$ are treated as operators satisfying the canonical anticommutation relation
$$\left\{\alpha _{n,}\alpha _m^+\right\}=\delta _{nm}$$
(44)
and the Dirac vacuum is defined as
$`\alpha _n\left(t\right)|0,A\left(t\right)_S`$ $`=`$ $`0,E_n\left(t\right)>0,`$
$`\alpha _n^+\left(t\right)|0,A\left(t\right)_S`$ $`=`$ $`0,E_n\left(t\right)<0.`$ (45)
Observe that the expansion of the Fermion field operators w.r.t. the eigenfunctions of the full Hamiltonian (29) automatically implies that we use the Schroedinger picture.
Singular operator products are regularized in by an exponential damping of high frequencies. The regularized current reads
$$\left(j^{\mu a}\left(x\right)\right)_{reg}=\underset{n,m}{}\alpha _n^+\left(t\right)\zeta _n^+(t,x)e^{\left(\epsilon /2\right)E_n^2\left(t\right)}\gamma ^0\gamma ^\mu \frac{1+\gamma _5}{2}e^{\left(\epsilon /2\right)E_m^2\left(t\right)}\zeta _m(t,x)\alpha _m\left(t\right)$$
$$=\underset{n,m}{}\alpha _n^+\left(t\right)\zeta _n^+(t,x)e^{\left(\epsilon /2\right)E_n^2\left(t\right)}t^ae_m^{\left(\epsilon /2\right)E_m^2\left(t\right)}\zeta _m(t,x)\alpha _m\left(t\right)$$
(46)
(where $`j^0=j^1`$ was used in the second line, which holds for the chiral current (33)). The current in (33) is regularized covariantly, therefore it will lead to the covariant anomaly and Schwinger term. The consistent current $`J^\mu `$ is obtained by adding the Bardeen–Zumino polynomial $`\mathrm{\Delta }j^\mu `$,
$$J^\mu \left(x\right)=j^\mu \left(x\right)+\mathrm{\Delta }j^\mu \left(x\right),$$
(47)
$$\mathrm{\Delta }j^\mu \left(x\right)=\frac{i}{4\pi }t^a\epsilon ^{\mu \nu }\mathrm{t}r\left(t^aA_\nu \right).$$
(48)
These currents lead to the covariant and consistent anomalies
$$𝒜_{cov}^a\left(x\right)=\left(D^\mu j_\mu \right)^a\left(x\right)=\frac{i}{2\pi }\epsilon _{\mu \nu }\mathrm{t}r\left(t^a\left(^\mu A^\nu +A^\mu A^\nu \right)\right)\left(x\right)$$
(49)
and
$$𝒜_{con}^a\left(x\right)=\left(D^\mu J_\mu \right)^a\left(x\right)=\frac{i}{4\pi }\epsilon _{\mu \nu }\mathrm{t}rt^a^\mu A^\nu \left(x\right).$$
(50)
The covariant ($`\stackrel{~}{G}^a`$) and consistent ($`G^a`$) Gauss law operators read
$$\stackrel{~}{G}^a\left(x\right)=X^a\left(x\right)+j^{0a}\left(x\right)$$
(51)
$$G^a\left(x\right)=X^a\left(x\right)+J^{0a}\left(x\right)$$
(52)
where
$$X^a\left(x\right)=\left(_1\frac{\delta }{\delta A_1^a\left(x\right)}+f^{abc}A_1^b\left(x\right)\frac{\delta }{\delta A_1^c\left(x\right)}\right)$$
(53)
generates time-independent gauge transformations of the external gauge field.
Assuming that the non-canonical parts ($`n.c.`$) of the commutator of the covariant and consistent Gauss laws are $`c`$-numbers it is sufficient to consider their vacuum expectation values (VEVs) only. The calculation in the Hosono and Seo approach is rather lengthy, therefore it is performed in the Appendices A – C. Here we just present the final form of the covariant Schwinger term
$`\stackrel{~}{ST}^{ab}`$ $`=`$ $`[\stackrel{~}{G}^a\left(x\right),\stackrel{~}{G}^b\left(y\right)]_{n.c.}=[j^{0a}\left(x\right),j^{0b}\left(y\right)]_{n.c.}=`$ (54)
$`=`$ $`{\displaystyle \frac{i}{2\pi }}_x\delta (xy)\mathrm{t}rt^at^b+{\displaystyle \frac{i}{2\pi }}\delta (xy)\mathrm{t}rt^a[A_1\left(y\right),t^b],`$
and the consistent one
$$ST^{ab}=[G^a\left(x\right),G^b\left(y\right)]_{n.c.}=\frac{i}{4\pi }\delta \left(xy\right)\mathrm{t}r\left([t^a,A_1]t^b\right).$$
(55)
Comparing the results (41) and (42) with the expressions for the 1+1 dim Schwinger terms in (6) (for the consistent case) and Table 1 (for the covariant case), we find that these terms agree. Therefore, the method of Hosono and Seo reproduces the result of Kelnhofer .
## 4 Method of Wess
In this section we want to review the papers of Schwiebert and Kelnhofer who used the method of Wess for the calculation of the consistent and covariant Schwinger terms (ST), respectively. The central idea of this method is to infer the current commutators from the time derivatives of a (time-ordered) current two-point function, by using the general relation $`_x^0TA(x)B(y)=\delta (x^0y^0)[A(x),B(y)]`$. As the anomaly is a (covariant) derivative of the current VEV (one-point function), and further current insertions are obtained by functional derivatives w.r.t. the external gauge potential $`A_a^\mu `$, the current commutator may be related to a functional derivative of the anomaly.
The authors of and used slightly different conventions. For our purposes it is important to have the same conventions for both the consistent and covariant cases, because we want to determine one relative sign. Therefore we shall repeat the major steps in the calculations of and within our specific set of conventions. We choose anti-Hermitean Lie algebra generators $`\lambda _a`$,
$$[\lambda _a,\lambda _b]=f_{abc}\lambda _c$$
(56)
where $`f_{abc}`$ are the structure constants. Further we choose Euclidean conventions in this section ($`g^{\mu \nu }=\delta ^{\mu \nu }`$), mainly because the path integral computation of both the consistent and covariant anomaly was done in Euclidean space as well (for our conventions see e.g. ). “Space-time” indices (running from 0 to 1 in $`d=2`$ and from 0 to 3 in $`d=4`$) are denoted by Greek letters $`\mu ,\nu ,\mathrm{}`$ and pure space indices are denoted by latin letters $`k,l,m`$. For the Ward operator we choose
$$X_a(x)=(D_x^\mu )_{ab}\frac{\delta }{\delta A_b^\mu (x)}(\delta _{ab}_x^\mu +f_{acb}A_c^\mu (x))\frac{\delta }{\delta A_b^\mu (x)}$$
(57)
$$[X_a(x),X_b(y)]=f_{abc}X_c(x)\delta (xy).$$
(58)
The Euclidean vacuum functional is
$$Z[A]=e^{W[A]}=0|T^{}e^{{\scriptscriptstyle 𝑑x\widehat{J}_a^\mu (x)A_a^\mu (x)}}|0$$
(59)
where $`A_a^\mu `$ is the external gauge potential and $`\widehat{J}_a^\mu `$ is a covariantly regularized current operator, which necessarily depends on $`A_a^\mu `$ for an anomalous gauge theory. Further $`T^{}`$ is the Lorentz covariantized time-ordered product that results from covariant perturbation theory.
### 4.1 Consistent case
For the VEV of the consistent current $`J_a^\mu `$ (one-point function) we have ($`\widehat{J}A𝑑x\widehat{J}_a^\mu (x)A_a^\mu (x)`$)
$$0|T^{}J_a^\mu (x)e^{{\scriptscriptstyle \widehat{J}A}}|0e^W:=\frac{\delta W}{\delta A_a^\mu (x)}$$
$$=0|T^{}(\widehat{J}_a^\mu (x)+𝑑y\frac{\delta \widehat{J}_b^\lambda (y)}{\delta A_a^\mu (x)}A_b^\lambda (y))e^{{\scriptscriptstyle \widehat{J}A}}|0e^W$$
(60)
and for the two-point function we get
$$\frac{\delta ^2W}{\delta A_a^\mu (x)\delta A_b^\nu (y)}=0|T^{}J_a^\mu (x)J_b^\nu (y)e^{{\scriptscriptstyle \widehat{J}A}}|0e^W$$
$$+0|T^{}\frac{\delta J_a^\mu (x)}{\delta A_b^\nu (y)}e^{{\scriptscriptstyle \widehat{J}A}}|0e^W+\frac{\delta W}{\delta A_a^\mu (x)}\frac{\delta W}{\delta A_b^\nu (y)}$$
(61)
$$=:T_{ab}^{\mu \nu }(x,y)+\mathrm{\Theta }_{ab}^{\mu \nu }(y)\delta (xy)+\mathrm{}$$
(62)
where in (49) we have defined abbreviations for the first and second term of (48) and indicated the third (disconnected) term by ellipses. Here it is assumed that $`J_a^\mu `$ depends on $`A_a^\mu `$ only in a local fashion .
Now we should re-express the $`T^{}`$ product by the ordinary $`T`$ product that is defined via $`\theta `$ functions. For the zero- and one-point functions we may simply define the $`T^{}`$ product by the $`T`$ product, because the latter leads to Lorentz-covariant expressions. On the other hand, for the two-point function $`T^{}J(x)J(y)`$ there occurs a difference (seagull term $`\tau _{ab}^{\mu \nu }`$) at coinciding space-time points, and this seagull term is proportional to $`\delta (xy)`$ . Denoting the ordinary $`T`$ product by $`T_{ab}^{\mu \nu }(x,y)`$, we have
$$T_{ab}^{\mu \nu }(x,y)=T_{ab}^{\mu \nu }(x,y)+\tau _{ab}^{\mu \nu }(y)\delta (xy).$$
(63)
For the divergence of $`T_{ab}^{\mu \nu }`$ we get, using the definition of the $`T`$ product,
$$_x^\mu T_{ab}^{\mu \nu }(x,y)=_x^\mu (\theta (x^0y^0)0|(Te^{_{x^0}^{\mathrm{}}\widehat{J}A})J_a^\mu (x)(Te^{_{y^0}^{x^0}\widehat{J}A})J_b^\nu (y)(Te^{_{\mathrm{}}^{y^0}\widehat{J}A})|0$$
$$+((\mu ,a,x)(\nu ,b,y)))e^W$$
$$=\delta (x^0y^0)0|T[J_a^0(x),J_b^\nu (y)]e^{{\scriptscriptstyle \widehat{J}A}}|0e^W+0|T_x^\mu J_a^\mu (x)J_b^\nu (y)e^{{\scriptscriptstyle \widehat{J}A}}|0e^W$$
$$0|T[J_a^0(x),_{z^0=x^0}𝐝z\widehat{J}_c^\lambda (z)A_c^\lambda (z)]J_b^\nu (y)e^{{\scriptscriptstyle \widehat{J}A}}|0e^W$$
(64)
where $`𝐝z`$ is w.r.t. the spacial coordinates only. The term containing $`_x^\mu J_a^\mu (x)`$ does not produce $`\delta `$ functions and may therefore be neglected. Further, $`\widehat{J}`$ in the third term may be replaced by $`J`$ without introducing $`\delta `$ function like contributions. For the commutator we use (in our Euclidean conventions $`J_b^\nu `$ is anti-Hermitean)
$$\delta (x^0y^0)[J_a^0(x),J_b^\nu (y)]=f_{abc}J_c^\nu (y)\delta (xy)+C_{ab}^{0\nu }(y)\delta (xy)+S_{ab}^{0\nu k}(y)_x^k\delta (xy)$$
(65)
Re-inserting this commutator into (51) and omitting disconnected terms we get
$$_x^\mu T_{ab}^{\mu \nu }(x,y)=C_{ab}^{0\nu }(y)\delta (xy)+S_{ab}^{0\nu k}(y)_x^k\delta (xy)$$
$$+f_{abc}\frac{\delta W}{\delta A_c^\nu (y)}\delta (xy)f_{adc}A_d^\lambda (x)T_{cb}^{\lambda \nu }(x,y).$$
(66)
This result has to be related to the functional derivative of the consistent anomaly, where the consistent anomaly itself is defined as
$$𝒜_a(x):=X_a(x)W[A].$$
(67)
Explicitly, we have in 2 and 4 dimensions ($`A_\mu A_a^\mu \lambda _a`$)
$$d=2:𝒜_a(x)=c_1ϵ^{\mu \nu }\mathrm{tr}\lambda _a^\mu A^\nu $$
(68)
$$d=4,𝒜_a(x)=c_2ϵ^{\mu \nu \rho \sigma }\mathrm{tr}\lambda _a^\mu (A^\nu ^\rho A^\sigma +\frac{1}{2}A^\nu A^\rho A^\sigma )$$
(69)
where $`c_1`$ and $`c_2`$ are some constants. From these explicit expressions we may express the functional derivatives of the anomalies as
$$\frac{\delta 𝒜_a(x)}{\delta A_b^\nu (y)}=I_{ab}^{\mu \nu }(x)_x^\mu \delta (xy)+(_x^\mu I_{ab}^{\mu \nu }(x))\delta (xy)=I_{ab}^{\mu \nu }(y)_x^\mu \delta (xy)$$
(70)
where the last equality follows from properties of the $`\delta `$ function. Explicitly we have
$$d=2,I_{ab}^{\mu \nu }(y)=c_1ϵ^{\mu \nu }\mathrm{tr}\lambda _a\lambda _b$$
(71)
$$d=4,I_{ab}^{\mu \nu }(y)=\frac{c_2}{2}ϵ^{\mu \nu \rho \sigma }\mathrm{tr}\left(\{\lambda _a,\lambda _b\}(2^\rho A^\sigma +A^\rho A^\sigma )\lambda _aA^\rho \lambda _bA^\sigma \right).$$
(72)
For later convenience we also note that
$$\frac{\delta 𝒜_b(y)}{\delta A_a^\mu (x)}=I_{ba}^{\nu \mu }(y)_x^\nu \delta (xy)+(_y^\nu I_{ba}^{\nu \mu }(y))\delta (xy).$$
(73)
On the other hand, we may use the definition (54) of the anomaly (and expression (44) for the Ward operator) to relate the functional derivative (57) to the two-point function (49). We get
$$\frac{\delta 𝒜_a(x)}{\delta A_b^\nu (y)}=f_{abc}\delta (xy)\frac{\delta W}{\delta A_c^\nu (x)}$$
$$+(\delta _{ac}_x^\mu +f_{adc}A_d^\mu (x))(T_{cb}^{\mu \nu }(x,y)\mathrm{\Theta }_{cb}^{\mu \nu }(y)\delta (xy))$$
$$=C_{ab}^{0\nu }(y)\delta (xy)+S_{ab}^{0\nu k}(y)_x^k\delta (xy)$$
$$+\sigma _{ab}^{\mu \nu }(y)_x^\mu \delta (xy)+f_{adc}A_d^\mu (y)\sigma _{cb}^{\mu \nu }(y)\delta (xy)$$
(74)
where we introduced
$$\sigma _{ab}^{\mu \nu }(y):=\tau _{ab}^{\mu \nu }(y)\mathrm{\Theta }_{ab}^{\mu \nu }(y).$$
(75)
Comparing the coefficients of $`\delta (xy)`$, $`_x^k\delta (xy)`$ and $`_x^0\delta (xy)`$ in (57) and (61) leads to
$$C_{ab}^{0\nu }(y)+f_{adc}A_d^\mu (y)\sigma _{cb}^{\mu \nu }(y)=0$$
(76)
$$S_{ab}^{0\nu k}(y)+\sigma _{ab}^{k\nu }(y)=I_{ab}^{k\nu }(y)$$
(77)
$$\sigma _{ab}^{0\nu }(y)=I_{ab}^{0\nu }(y).$$
(78)
For a determination of $`S_{ab}^{00k}`$ and $`C_{ab}^{00}`$ we need $`\sigma _{ab}^{k0}`$ about which we have no information yet (here we slightly deviate from the calculation of and follow the arguments of , but the final result will agree with the result of up to the difference in conventions). For this purpose we compute, analogously to (61),
$$\frac{\delta 𝒜_b(y)}{\delta A_a^\mu (x)}=f_{bac}\delta (xy)\frac{\delta W}{\delta A_a^\mu (x)}$$
$$+(\delta _{bc}_y^\nu +f_{bdc}A_d^\nu (y))(T_{ac}^{\mu \nu }(x,y)\mathrm{\Theta }_{ac}^{\mu \nu }(y)\delta (xy))$$
(79)
and use
$$_y^\nu T_{ab}^{\mu \nu }(x,y)=\mathrm{}=C_{ab}^{\mu 0}(y)\delta (xy)S_{ab}^{\mu 0k}(y)_x^k\delta (xy)$$
$$+f_{dbc}A_d^\nu (y)T_{ac}^{\mu \nu }(x,y)f_{abc}\frac{\delta W}{\delta A_c^\mu (y)}\delta (xy)$$
(80)
to arrive at
$$\frac{\delta 𝒜_b(y)}{\delta A_a^\mu (x)}=C_{ab}^{\mu 0}(y)\delta (xy)S_{ab}^{\mu 0k}(y)_x^k\delta (xy)$$
$$+\delta (xy)(\delta _{bc}_y^\nu +f_{bdc}A_d^\nu (y))\sigma _{ac}^{\mu \nu }(y)\sigma _{ab}^{\mu \nu }(y)_x^\nu \delta (xy).$$
(81)
Comparison of coefficients of (60) and (68) leads to
$$C_{ab}^{\mu 0}(y)+_y^\nu \sigma _{ab}^{\mu \nu }(y)+f_{bdc}A_d^\nu (y)\sigma _{ac}^{\mu \nu }(y)=_y^\nu I_{ba}^{\nu \mu }(y)$$
(82)
$$S_{ab}^{\mu 0k}(y)+\sigma _{ab}^{\mu k}(y)=I_{ba}^{k\mu }(y)$$
(83)
$$\sigma _{ab}^{\mu 0}(y)=I_{ba}^{0\mu }(y).$$
(84)
Together with (63)–(65) this may be solved for $`S_{ab}^{00k}`$ and $`C_{ab}^{00}`$
$$S_{ab}^{00k}(y)=I_{ab}^{k0}(y)I_{ba}^{0k}(y)$$
(85)
$$C_{ab}^{00}(y)=f_{adc}A_d^\mu (y)I_{bc}^{0\mu }(y).$$
(86)
In addition we find from (69) and (73) the consistency condition
$$_y^\nu (I_{ab}^{0\nu }(y)I_{ba}^{\nu 0}(y))+A_d^\nu (y)(f_{adc}I_{bc}^{0\nu }(y)+f_{bdc}I_{ac}^{0\nu }(y))=0$$
(87)
which holds for both 2 and 4 dimensions, as may be checked easily. So far we have determined the anomalous $`[J^0,J^0]`$ commutator, see (52), (72) and (73). We still need the commutator of $`J_a^0`$ and the Ward operator $`X_b`$. As $`X_b`$ does not contain fermionic degrees, this commutator is equal to the action of $`X_b`$ on $`J_a^0`$,
$$X_b(y)J_a^0(x)\delta (x^0y^0)[X_b(y),J_a^0(x)].$$
(88)
This commutator may be inferred from the relation
$$X_b(y)\frac{\delta W}{\delta A_a^\mu (x)}=(D_y^\nu )_{bc}\frac{\delta }{\delta A_c^\nu (y)}0|TJ_a^\mu (x)e^{{\scriptscriptstyle \stackrel{~}{J}A}}|0e^W$$
$$=0|T(X_b(y)J_a^\mu (x))e^{{\scriptscriptstyle \stackrel{~}{J}A}}|0e^W+(D_y^\nu )_{bc}0|TJ_a^\mu (x)J_c^\nu (y)e^{{\scriptscriptstyle \stackrel{~}{J}A}}|0e^W.$$
(89)
Here we used the fact that in the one-point function (47) the $`T^{}`$ product is equal to the $`T`$ product. It is important to use the $`T`$ product here, because we want to extract the (Lorentz non-covariant) commutator $`[J_a^0,X_b]`$ directly, without some covariantizing seagulls. Now we assume that the commutator (75) contains no fermionic degrees of freedom, i.e., it may be extracted from the VEV. Using (49) and (50) for the two-point function we find
$$\delta (x^0y^0)[X_b(y),J_a^\mu (x)]+(D_y^\nu )_{bc}T_{ac}^{\mu \nu }(x,y)=(D_y^\nu )_{bc}(T_{ac}^{\mu \nu }+\sigma _{ac}^{\mu \nu }(y)\delta (xy))$$
(90)
or, for $`\mu =0`$ and using (65),
$$\delta (x^0y^0)[J_a^0(x),X_b(y)]=(D_y^k)_{bc}(I_{ac}^{0k}(y)\delta (xy)).$$
(91)
Actually, for the Gauss operator we only need the Ward operator restricted to purely spacial gauge transformations. In addition it is preferable to get rid of the time coordinate altogether. Therefore we define a spacial Ward operator
$$𝐗_a(x):=𝑑x^0(D_x^k)_{ab}\frac{\delta }{\delta A_b^k(x)}$$
(92)
$$[𝐗_a(x),𝐗_b(y)]=f_{abc}𝐗_c(x)\stackrel{}{\delta }(xy)$$
(93)
where $`\stackrel{}{\delta }(xy)`$ is the spacial $`\delta `$ function. The Gauss operator is
$$G_a(x)=J_a^0(x)+𝐗_a(x).$$
(94)
Using (72), (73) and (78) we find for the anomalous part of the commutator (i.e., the Schwinger term)
$$𝒢_{ab}(x,y):=[G_a(x),G_b(y)]f_{abc}G_c(x)\delta (xy)$$
$$=C_{ab}^{00}(y)\stackrel{}{\delta }(xy)+S_{ab}^{00k}(y)_x^k\stackrel{}{\delta }(xy)+[J_a^0(x),𝐗_b(y)]+[𝐗_a(x),J_b^0(y)]$$
$$=(f_{bdc}A_d^k(y)I_{ac}^{0k}(y)+_y^kI_{ab}^{0k}(y))\stackrel{}{\delta }(xy).$$
(95)
Before evaluating this expression explicitly for $`d=2`$ and $`d=4`$, we want to find the analogous result for the covariant case, following .
### 4.2 Covariant case
The VEV of the covariant current $`\stackrel{~}{J}_a^\mu `$ is related to the consistent one by the Bardeen–Zumino polynomial $`\mathrm{\Lambda }_a^\mu `$,
$$0|T^{}\stackrel{~}{J}_a^\mu (x)e^{{\scriptscriptstyle \widehat{J}A}}|0e^W=0|T^{}J_a^\mu (x)e^{{\scriptscriptstyle \widehat{J}A}}|0e^W+\mathrm{\Lambda }_a^\mu (x).$$
(96)
This leads to the covariant anomaly $`\stackrel{~}{𝒜}_a(x)`$,
$$\stackrel{~}{𝒜}_a(x)=(D_x^\mu )_{ab}0|T^{}\stackrel{~}{J}_b^\mu (x)e^{{\scriptscriptstyle \widehat{J}A}}|0e^W=𝒜_a(x)(D_x^\mu )_{ab}\mathrm{\Lambda }_b^\mu (x).$$
(97)
Explicitly the covariant anomalies are
$$d=2,\stackrel{~}{𝒜}_a(x)=2c_1ϵ^{\mu \nu }\mathrm{tr}\lambda _a(^\mu A^\nu +A^\mu A^\nu )$$
(98)
$$d=4,\stackrel{~}{𝒜}_a(x)=3c_2ϵ^{\mu \nu \rho \sigma }\mathrm{tr}\lambda _a(^\mu A^\nu +A^\mu A^\nu )(^\rho A^\sigma +A^\rho A^\sigma )$$
(99)
where the constants $`c_1`$, $`c_2`$ are the same as in the consistent case, see (55) and (56). The two-point functions are defined analogously to (48)–(50) as
$$\frac{\delta }{\delta A_b^\nu (y)}0|T^{}\stackrel{~}{J}_a^\mu (x)e^{{\scriptscriptstyle \widehat{J}A}}|0e^W=0|T^{}\stackrel{~}{J}_a^\mu (x)\stackrel{~}{J}_b^\nu (y)e^{{\scriptscriptstyle \widehat{J}A}}|0e^W$$
$$+0|T^{}\frac{\delta \stackrel{~}{J}_a^\mu (x)}{\delta A_b^\nu (y)}e^{{\scriptscriptstyle \widehat{J}A}}|0e^W+\mathrm{}$$
$$=:\stackrel{~}{T}_{ab}^{\mu \nu }(x,y)+\stackrel{~}{\mathrm{\Theta }}_{ab}^{\mu \nu }(y)\delta (xy)+\mathrm{}$$
(100)
$$=:\stackrel{~}{T}_{ab}^{\mu \nu }(x,y)\stackrel{~}{\sigma }_{ab}^{\mu \nu }(y)\delta (xy)+\mathrm{}$$
(101)
where the ellipses denote disconnected terms and all definitions are analogous to the consistent case. Further, the computation of $`_x^\mu \stackrel{~}{T}_{ab}^{\mu \nu }(x,y)`$ is completely analogous to the consistent case, see (51)–(53). Parametrizing the covariant commutator in an analogous way,
$$\delta (x^0y^0)[\stackrel{~}{J}_a^0(x),\stackrel{~}{J}_b^\nu (y)]=f_{abc}\stackrel{~}{J}_c^\nu (x)\delta (xy)+\stackrel{~}{C}_{ab}^{0\nu }(y)\delta (xy)+\stackrel{~}{S}_{ab}^{0\nu k}(y)_x^k\delta (xy)$$
(102)
leads to a result analogous to (53),
$$_x^\mu \stackrel{~}{T}_{ab}^{\mu \nu }(x,y)=\stackrel{~}{C}_{ab}^{0\nu }(y)\delta (xy)+\stackrel{~}{S}_{ab}^{0\nu k}(y)_x^k\delta (xy)$$
$$+f_{abc}\delta (xy)0|T\stackrel{~}{J}_c^\nu (y)e^{{\scriptscriptstyle \widehat{J}A}}|0e^Wf_{adc}A_d^\lambda (x)\stackrel{~}{T}_{cb}^{\lambda \nu }(x,y).$$
(103)
Again, this should be related to the functional derivative of the (covariant) anomaly. We express this functional derivative as
$$\frac{\delta \stackrel{~}{𝒜}_a(x)}{\delta A_b^\nu (y)}=\stackrel{~}{I}_{ab}^{\mu \nu }(y)_x^\mu \delta (xy)+\stackrel{~}{B}_{ab}^\nu (y)\delta (xy)$$
(104)
(we do not display the explicit expressions for $`\stackrel{~}{I}`$ and $`\stackrel{~}{B}`$ for $`d=2`$ or $`d=4`$, because we do not need them in the sequel).
On the other hand, using the definition of $`\stackrel{~}{𝒜}_a`$, relating its functional derivative to the two-point function (88) and inserting (90) for $`^\mu \stackrel{~}{T}_{ab}^{\mu \nu }`$ leads to
$$\frac{\delta \stackrel{~}{𝒜}_a(x)}{\delta A_b^\nu (y)}=\stackrel{~}{C}_{ab}^{0\nu }(y)\delta (xy)+S_{ab}^{0\nu k}(y)_x^k\delta (xy)$$
$$+\stackrel{~}{\sigma }_{ab}^{\mu \nu }(y)_x^\mu \delta (xy)+f_{adc}A_d^\mu (y)\stackrel{~}{\sigma }_{cb}^{\mu \nu }(y)\delta (xy)$$
(105)
and therefore to the equations
$$\stackrel{~}{C}_{ab}^{0\nu }(y)=\stackrel{~}{B}_{ab}^\nu (y)f_{adc}A_d^\mu (y)\stackrel{~}{\sigma }_{cb}^{\mu \nu }(y)$$
(106)
$$\stackrel{~}{S}_{ab}^{0\nu k}(y)=\stackrel{~}{I}_{ab}^{k\nu }(y)\stackrel{~}{\sigma }_{ab}^{k\nu }(y)$$
(107)
$$\stackrel{~}{\sigma }_{ab}^{0\nu }(y)=\stackrel{~}{I}_{ab}^{0\nu }(y).$$
(108)
Again we miss information on $`\stackrel{~}{\sigma }_{ab}^{k0}(y)`$, which we may infer from $`(\delta \stackrel{~}{𝒜}_b(y)/\delta A_a^\mu (x))`$. We find
$$\frac{\delta \stackrel{~}{𝒜}_b(y)}{\delta A_a^\mu (x)}=f_{bac}\delta (xy)0|T^{}\stackrel{~}{J}_c^\mu (y)e^{{\scriptscriptstyle \widehat{J}A}}|0e^W$$
$$(D_y^\nu )_{bc}\frac{\delta }{\delta A_a^\mu (x)}0|T^{}\stackrel{~}{J}_c^\nu (y)e^{{\scriptscriptstyle \widehat{J}A}}|0e^W$$
$$=f_{bac}\delta (xy)0|T^{}\stackrel{~}{J}_c^\mu (y)e^{{\scriptscriptstyle \widehat{J}A}}|0e^W$$
$$(D_y^\nu )_{bc}\left(\frac{\delta }{\delta A_c^\nu (y)}0|T^{}\stackrel{~}{J}_a^\mu (x)e^{{\scriptscriptstyle \widehat{J}A}}|0e^W_{ab}^{\mu \nu }(x,y)\right)$$
(109)
$$_{ab}^{\mu \nu }(x,y):=\frac{\delta \mathrm{\Lambda }_a^\mu (x)}{\delta A_b^\nu (y)}\frac{\delta \mathrm{\Lambda }_b^\nu (y)}{\delta A_a^\mu (x)}$$
(110)
where we used relation (83) between consistent and covariant current VEV and the commutativity of functional derivatives (see ). Computing $`_y^\nu \stackrel{~}{T}_{ab}^{\mu \nu }(x,y)`$ as in the consistent case yields
$$\frac{\delta \stackrel{~}{𝒜}_b(y)}{\delta A_a^\mu (x)}=\stackrel{~}{C}_{ab}^{\mu 0}(y)\delta (xy)\stackrel{~}{S}_{ab}^{\mu 0k}(y)_x^k\delta (xy)$$
$$+(D_y^\nu )_{bc}_{ac}^{\mu \nu }(x,y)+\delta (xy)(D_y^\nu )_{bc}\stackrel{~}{\sigma }_{ac}^{\mu \nu }(y)\stackrel{~}{\sigma }_{ab}^{\mu \nu }(y)_x^\nu \delta (xy).$$
(111)
However, as a consequence of the gauge covariance of the covariant current it holds that
$$\frac{\delta \stackrel{~}{𝒜}_b(y)}{\delta A_a^\mu (x)}(D_y^\nu )_{bc}_{ac}^{\mu \nu }(x,y)$$
(112)
as may be checked explicitly . Therefore, the coefficients in (98) are not directly related to the anomaly and have to obey
$$\stackrel{~}{C}_{ab}^{\mu 0}(y)=(D_y^\nu )_{bc}\stackrel{~}{\sigma }_{ac}^{\mu \nu }(y)$$
(113)
$$\stackrel{~}{S}_{ab}^{\mu 0k}(y)=\stackrel{~}{\sigma }_{ab}^{\mu k}(y)$$
(114)
$$\stackrel{~}{\sigma }_{ab}^{\mu 0}(y)=0$$
(115)
and we find
$$\stackrel{~}{S}_{ab}^{00k}(y)=\stackrel{~}{I}_{ab}^{k0}(y)$$
(116)
$$\stackrel{~}{C}_{ab}^{00}(y)=\stackrel{~}{B}_{ab}^0(y)$$
(117)
and the consistency condition
$$\stackrel{~}{B}_{ab}^0(y)=(D_y^\nu )_{bc}\stackrel{~}{I}_{ac}^{0\nu }(y)$$
(118)
which holds indeed, as may be checked by explicit computation . For the anomalous part of the current commutator this leads to
$$\delta (x^0y^0)[\stackrel{~}{J}_a^0(x),\stackrel{~}{J}_b^0(y)]f_{abc}\stackrel{~}{J}_c^0(y)\delta (xy)=$$
$$=\stackrel{~}{C}_{ab}^{00}(y)\delta (xy)+\stackrel{~}{S}_{ab}^{00k}(y)_x^k\delta (xy)\frac{\delta \stackrel{~}{𝒜}_a(x)}{\delta A_b^0(y)}.$$
(119)
Again, we have to calculate the $`[X_b,\stackrel{~}{J}_a^0]`$ commutators as in the consistent case. However, the result is simply that each such term in the Gauss operator commutator produces a contribution that is equal to minus the above expression (106), see . Therefore we find for the covariant Gauss operator
$$\stackrel{~}{G}_a(x)=\stackrel{~}{J}_a^0(x)+𝐗_a(x)$$
(120)
the Schwinger term
$$\stackrel{~}{𝒢}_{ab}(x,y):=[\stackrel{~}{G}_a(x),\stackrel{~}{G}_b(y)]f_{abc}\stackrel{~}{G}_c(x)\stackrel{}{\delta }(xy)=𝑑y^0\frac{\delta \stackrel{~}{𝒜}_a(x)}{\delta A_b^0(y)}$$
(121)
where the $`dy^0`$ integration just serves to get rid of the unwanted $`y^0`$ dependence (this just kills a $`\delta (x^0y^0)`$, because there is no time derivative in the above expression (108)).
### 4.3 Explicit evaluation
Now we are in a position to explicitly compute the Schwinger terms both for $`d=2`$ and $`d=4`$. Starting with the $`d=2`$ case, we find from (58) and (82) for the consistent ST
$$𝒢_{ab}(x,y)=c_1ϵ^{0k}f_{bdc}A_d^k(y)\mathrm{tr}\lambda _a\lambda _c\delta ^{(1)}(xy)$$
$$=c_1ϵ^{0k}\mathrm{tr}\lambda _a[\lambda _b,A^k]\delta ^{(1)}(xy)$$
(122)
($`\delta ^{(1)}(xy):=\delta (x^1y^1)`$) and for the covariant ST, using (85) and (108)
$$\stackrel{~}{𝒢}_{ab}(x,y)=2c_1ϵ^{0k}\mathrm{tr}\lambda _a(\lambda _b_y^k\delta ^{(1)}(xy)+[\lambda _b,A^k]\delta ^{(1)}(xy))$$
(123)
where here and in the following functions always depend on $`y`$ when the coordinate argument is not written down explicitly. In order to compare with the expressions of Section 2, we omit $`ϵ^{0k}`$ and multiply by $`(1/2)dy^kv_a(x)v_b(y)`$, in the indicated order. Here $`dy^k`$ is a one-form, $`v_a(x)`$ is a ghost, and all these objects anti-commute, e.g., $`dy^kv_a(x)=v_a(x)dy^k`$. We find ($`A(y):=A^k(y)dy^k`$, $`v(x):=v_a(x)\lambda _a`$)
$$𝒢(x,y)=\frac{c_1}{2}\mathrm{tr}(v(x)v(y)+v(y)v(x))A(y)\delta ^{(1)}(xy)$$
(124)
or, after integrating w.r.t. $`𝑑x^1𝑑y^1`$
$$𝒢=c_1𝑑y\mathrm{tr}v^2A.$$
(125)
In the same fashion, we get for $`\stackrel{~}{𝒢}_{ab}(x,y)`$ ($`d_y:=dy^k_y^k`$)
$$\stackrel{~}{𝒢}(x,y)=c_1\mathrm{tr}v(x)\left((d_y\delta ^{(1)}(xy))v(y)+\delta ^{(1)}(xy)(A(y)v(y)+v(y)A(y))\right)$$
(126)
and (where a partial integration has to be performed)
$$\stackrel{~}{𝒢}=c_1𝑑y\mathrm{tr}v(dv+Av+vA)=c_1𝑑y\mathrm{tr}vDv.$$
(127)
Comparing with eq. (6) (for the consistent case) and Table 1 (for the covariant case), we find that the relative sign of $`𝒢`$ and $`\stackrel{~}{𝒢}`$ is in precise agreement.
For the case $`d=4`$ we find from (59) and (82)
$$𝒢_{ab}(x,y)=\frac{c_2}{2}ϵ^{0klm}\mathrm{tr}([\lambda _a,\lambda _b](^kA^lA^m+A^k^lA^m+A^kA^lA^m)$$
$$+(\lambda _bA^k\lambda _a\lambda _aA^k\lambda _b)^lA^m)\delta ^{(3)}(xy)$$
(128)
(each derivative acts only on its immediate right hand neighbour), or, after omitting $`ϵ^{0lkm}`$ and multiplying by $`(1/2)dy^kdy^ldy^mv_a(x)v_b(y)`$
$$𝒢(x,y)=\frac{c_2}{4}\mathrm{tr}((v(x)v(y)+v(y)v(x))(dAA+AdA+A^3)$$
$$+(v(y)Av(x)+v(x)Av(y))dA)\delta ^{(3)}(xy)$$
(129)
and upon integration $`d^3xd^3y`$
$$𝒢=\frac{c_2}{2}𝑑y\mathrm{tr}(v^2(dAA+AdA+A^3)+vAvdA).$$
(130)
For the covariant ST we find from (86) and (108)
$$\stackrel{~}{𝒢}_{ab}(x,y)=3c_2ϵ^{0klm}\mathrm{tr}\lambda _a((\lambda _b_y^k\delta ^{(3)}(xy)+(\lambda _bA^kA^k\lambda _b)\delta ^{(3)}(xy))(^lA^m+A^lA^m)$$
$$+(^lA^m+A^lA^m)(\lambda _b_y^k\delta ^{(3)}(xy)+(\lambda _bA^kA^k\lambda _b)\delta ^{(3)}(xy)))$$
(131)
and
$$\stackrel{~}{𝒢}(x,y)=\frac{3c_2}{2}\mathrm{tr}(v(x))(((d_y\delta ^{(3)}(xy))v(y)(vA+Av)\delta ^{(3)}(xy))(dA+A^2)$$
$$+(dA+A^2)((d_y\delta ^{(3)}(xy))v(y)(vA+Av)\delta ^{(3)}(xy)))$$
(132)
$$\stackrel{~}{𝒢}=\frac{3c_2}{2}𝑑y\mathrm{tr}v((dv+Av+vA)(dA+A^2)+(dA+A^2)(dv+Av+vA))$$
$$=\frac{3c_2}{2}𝑑y\mathrm{tr}v(DvF+FDv).$$
(133)
Again, the relative sign of consistent and covariant ST precisely agrees with the one in eq. (6) (consistent case) and Table 1 (covariant case).
## Acknowledgement
This work started at the Erwin Schrödinger Institute (ESI) program ”Quantization, generalized BRS cohomology and anomalies” and we would like to thank the ESI for kind hospitality. Further, we thank R.A. Bertlmann for helpful discussions. One of us (T.S.) is grateful to R.A. Bertlmann for the possibility to spend some time in the creative atmosphere of the ESI.
## Appendix
## Appendix A The Schwinger terms of the Gauss Laws
We start with the covariant case and we consider only the non-canonical part of the commutator:
$`[\stackrel{~}{G}^a\left(x\right),\stackrel{~}{G}^b\left(y\right)]_{n.c.}`$ $`=`$ $`[X^a\left(x\right)+j^{0a}\left(x\right),X^b\left(y\right)+j^{0b}\left(y\right)]_{n.c.}=`$ (134)
$`=`$ $`[X^a\left(x\right),X^b\left(y\right)]_{n.c.}+[X^a\left(x\right),j^{0b}\left(y\right)]_{n.c.}+`$
$`+[j^{0a}\left(x\right),X^b\left(y\right)]_{n.c.}+[j^{0a}\left(x\right),j^{0b}\left(y\right)]_{n.c.}.`$
The gauge field is an external field, therefore the commutator
$$[X^a\left(x\right),X^b\left(y\right)]_{n.c.}$$
(135)
is zero. For the VEV of the commutator
$$[j^{0a}\left(x\right),j^{0b}\left(y\right)]_{n.c.}$$
(136)
we get, after some manipulations,
$`[j^{0a}\left(x\right),j^{0b}\left(y\right)]_{n.c.}`$ $`=`$ $`\mathrm{t}re^{\left(\epsilon /2\right)\mathrm{\Delta }_y}P_{}(t,y,x)t^a\left[e^{\left(\epsilon /2\right)\mathrm{\Delta }_x}1\right]\delta \left(xy\right)t^b`$ (137)
$`(x,ay,b),`$
where $`P_{}(t,x,y)`$ denotes the projection operator (see Appendix C)
$$P_{}(t,x,y)=\underset{E_n<0}{}\zeta _n^+(t,x)\zeta _n(t,x).$$
(138)
Then (137) gives
$`\mathrm{t}re^{\epsilon \mathrm{\Delta }_y}P_{}^{\left(0\right)}(y,x)t^a[e^{\epsilon \mathrm{\Delta }_x}1]\delta (xy)t^b(x,ay,b)=`$
$`=`$ $`\alpha \mathrm{t}r{\displaystyle }dE\theta (E)e^{\epsilon \mathrm{\Delta }_y}e^{iE(xy)}t^a\times `$
$`\times {\displaystyle }dq[e^{\epsilon \mathrm{\Delta }_x}1]e^{iq(xy)}t^b(x,ay,b)=`$
$`=`$ $`\alpha \mathrm{t}r{\displaystyle 𝑑E𝑑q\theta \left(E\right)e^{\epsilon \mathrm{\Delta }_y}e^{iE(xy)}t^a\left[e^{\epsilon \mathrm{\Delta }_x}1\right]e^{iq(xy)}t^b}`$
$`(x,ay,b)=`$
$`=`$ $`\alpha \mathrm{t}r{\displaystyle }dEdq\theta (E)e^{\epsilon E^2}e^{iE(xy)}(12i\epsilon EA\left(y\right))t^a\times `$
$`\times [(1+2i\epsilon qA\left(x\right))e^{\epsilon q^2}1]e^{iq(xy)}t^b(x,ay,b)=`$
$`=`$ $`\alpha {\displaystyle }dEdq\theta (E)e^{\epsilon E^2}e^{i\left(E+q\right)(xy)}[e^{\epsilon q^2}1]\mathrm{t}rt^at^b(x,ay,b)`$
$`i\alpha {\displaystyle 𝑑E𝑑q\theta \left(E\right)e^{\epsilon E^2}\left[e^{\epsilon q^2}1\right]e^{i\left(E+q\right)(xy)}2\epsilon E\mathrm{t}rA\left(y\right)t^at^b}+`$
$`+(x,ay,b)+`$
$`+i\alpha {\displaystyle }dEdq\theta (E)e^{\epsilon E^2}2\epsilon qe^{\epsilon q^2}\mathrm{t}rt^aA\left(x\right)t^b(x,ay,b)=`$
$`=`$ $`\alpha {\displaystyle }d\xi e^{i\xi (xy)}{\displaystyle }dE\theta (E)e^{\epsilon E^2}[e^{\epsilon \left(\xi E\right)^2}1]\mathrm{t}rt^at^b(x,ay,b)`$
$`i\alpha {\displaystyle 𝑑\xi e^{i\xi (xy)}𝑑E\theta \left(E\right)2\epsilon Ee^{\epsilon E^2}\left[e^{\epsilon \left(\xi E\right)^2}1\right]\mathrm{t}rA\left(y\right)t^at^b}+`$
$`+(x,ay,b)+`$
$`+i\alpha {\displaystyle 𝑑\xi e^{i\xi (xy)}𝑑E\theta \left(E\right)2\epsilon qe^{\epsilon E^2}e^{\epsilon \left(\xi E\right)^2}\mathrm{t}rt^aA\left(x\right)t^b}`$
$`(x,ay,b)=`$
$`=`$ $`\alpha {\displaystyle 𝑑\xi e^{i\xi (xy)}𝑑Ee^{\epsilon E^2}\left[e^{\epsilon \left(\xi E\right)^2}1\right]\left(\theta \left(E\right)\theta \left(E\right)\right)\mathrm{t}rt^at^b}`$
$`i\alpha \{[{\displaystyle }d\xi e^{i\xi (xy)}{\displaystyle }dE\theta (E)2\epsilon Ee^{\epsilon E^2}[e^{\epsilon \left(\xi E\right)^2}1]`$
$`{\displaystyle }d\xi e^{i\xi (xy)}{\displaystyle }dE\theta \left(E\right)2\epsilon (\xi E)e^{\epsilon E^2}e^{\epsilon \left(\xi E\right)^2}]\mathrm{t}rA(y)t^at^b`$
$`.(x,ay,b)\}\stackrel{\epsilon 0}{}`$
$``$ $`\alpha {\displaystyle 𝑑\xi e^{i\xi (xy)}𝑑Ee^{\epsilon E^2}\left[e^{\epsilon E^2}\left(1\epsilon \xi ^2+2\epsilon E\xi \right)1\right]\epsilon (E)\mathrm{t}rt^at^b}`$
$`i\alpha \{{\displaystyle }d\xi e^{i\xi (xy)}\times `$
$`\begin{array}{ccccc}\times [dE2\epsilon Ee^{\epsilon E^2}e^{\epsilon \left(\xi E\right)^2}(\theta (E)+\theta \left(E\right))\hfill & & & & \backslash \stackrel{\epsilon 0}{}0\hfill \\ 𝑑E\theta \left(E\right)2\epsilon Ee^{\epsilon E^2}\hfill & & & & \backslash \stackrel{\epsilon 0}{}1\hfill \\ dE\theta \left(E\right)2\epsilon \xi e^{\epsilon E^2}e^{\epsilon \left(\xi E\right)^2}]\times \hfill & & & & \backslash \stackrel{\epsilon 0}{}0\hfill \end{array}`$
$`\times \mathrm{t}rA\left(y\right)t^at^b(x,ay,b).\}`$
$``$ $`\alpha {\displaystyle 𝑑\xi \xi e^{i\xi (xy)}𝑑Ee^{2\epsilon E^2}2\epsilon E\epsilon \left(E\right)\mathrm{t}rt^at^b}`$
$`i\alpha \left\{{\displaystyle 𝑑\xi e^{i\xi (xy)}\mathrm{t}rA\left(y\right)t^at^b}{\displaystyle 𝑑\xi e^{i\xi (xy)}\mathrm{t}rA\left(x\right)t^bt^a}\right\}\stackrel{\epsilon 0}{=}`$
$`=`$ $`\alpha {\displaystyle 𝑑\xi e^{i\xi (xy)}\xi \mathrm{t}rt^at^b}+2\pi i\alpha \delta \left(xy\right)\mathrm{t}rt^a[A\left(y\right),t^b]=`$
$`=`$ $`{\displaystyle \frac{i}{2\pi }}_x\delta (xy)\mathrm{t}rt^at^b+{\displaystyle \frac{i}{2\pi }}\delta (xy)\mathrm{t}rt^a[A\left(y\right),t^b],`$ (143)
where $`\alpha =1/\left(2\pi \right)^2`$. Then
$`[j^{0a}\left(x\right),j^{0b}\left(y\right)]_{n.c.}`$ $`=`$ $`{\displaystyle \frac{i}{2\pi }}_x\delta (xy)\mathrm{t}rt^at^b+{\displaystyle \frac{i}{2\pi }}\delta (xy)\mathrm{t}rt^a[A\left(y\right),t^b]`$
and the covariant Schwinger term of the commutator of the full Gauss law operators has the form
$`\stackrel{~}{ST}^{ab}`$ $`=`$ $`[\stackrel{~}{G}^a\left(x\right),\stackrel{~}{G}^b\left(y\right)]_{n.c.}=[j^{0a}\left(x\right),j^{0b}\left(y\right)]_{n.c.}=`$ (145)
$`=`$ $`{\displaystyle \frac{i}{2\pi }}_x\delta (xy)\mathrm{t}rt^at^b+{\displaystyle \frac{i}{2\pi }}\delta (xy)\mathrm{t}rt^a[A_1\left(y\right),t^b],`$
where we used the result for the cross-term
$$[X^a\left(x\right),j^{0b}\left(y\right)]=\frac{i}{2\pi }\left(_1^x\delta ^{ac}+f^{aec}A_1^e\left(x\right)\right)\delta \left(xy\right)\mathrm{t}rt^bt^c$$
(146)
obtained in Appendix B.
For the commutator of the consistent Gauss laws we get
$`[G^a\left(x\right),G^b\left(y\right)]=`$ (147)
$`=`$ $`[\stackrel{~}{G}^a\left(x\right),\stackrel{~}{G}^b\left(y\right)]+[\stackrel{~}{G}^a\left(x\right),\mathrm{\Delta }j^{0b}\left(y\right)]+[\mathrm{\Delta }j^{0a}\left(x\right),\stackrel{~}{G}^b\left(y\right)]=`$
$`=`$ $`f^{abc}\stackrel{~}{G}^c\left(x\right)\delta \left(xy\right)=`$
$`=`$ $`f^{abc}G^c\left(x\right)\delta \left(xy\right)f^{abc}\mathrm{\Delta }j^{0c}\left(y\right)\delta \left(xy\right),`$
where we used the equality
$$\stackrel{~}{ST}^{ab}+[\stackrel{~}{G}^a\left(x\right),\mathrm{\Delta }j^{0b}\left(y\right)]+[\mathrm{\Delta }j^{0a}\left(x\right),\stackrel{~}{G}^b\left(y\right)]=0.$$
(148)
which results from
$`[\stackrel{~}{G}^a\left(x\right),\mathrm{\Delta }j^{0b}\left(y\right)]`$ $`=`$ $`[X^a\left(x\right),\mathrm{\Delta }j^{0b}\left(y\right)]=X^a\left(x\right)\mathrm{\Delta }j^{0b}\left(y\right)=`$ (149)
$`=`$ $`{\displaystyle \frac{i}{4\pi }}\epsilon ^{01}\left(\delta ^{ac}^\mu +f^{aec}A^{\mu e}\left(x\right)\right){\displaystyle \frac{\delta }{\delta A^{\mu c}\left(x\right)}}\mathrm{t}r\left(t^bA_1\left(y\right)\right)=`$
$`=`$ $`{\displaystyle \frac{i}{4\pi }}\left(\delta ^{ac}_x^1+f^{aec}A^{1e}\left(x\right)\right)\delta \left(xy\right)\mathrm{t}rt^bt^c.`$
Therefore
$`ST^{ab}`$ $`=`$ $`f^{abc}\mathrm{\Delta }j^{0c}\left(y\right)\delta \left(xy\right)=`$ (150)
$`=`$ $`{\displaystyle \frac{i}{4\pi }}\epsilon ^{0\nu }\mathrm{t}r\left(f^{abc}t^cA_\nu \right)\delta \left(xy\right)=`$
$`=`$ $`{\displaystyle \frac{i}{4\pi }}\delta \left(xy\right)\mathrm{t}r\left([t^a,A_1]t^b\right).`$
## Appendix B The cross-term
For the VEV of the cross-term
$$[X^a\left(x\right),j^{0b}\left(y\right)]$$
(151)
we get
$`[X^a\left(x\right),j^{0b}\left(y\right)]`$ $`=`$ $`\left(_x^1\delta ^{ac}+f^{aec}A^{1e}\left(x\right)\right){\displaystyle \frac{\delta }{\delta A^{1c}\left(x\right)}}j^{0b}\left(y\right)=`$ (152)
$`=`$ $`\left(_x^1\delta ^{ac}+f^{aec}A^{1e}\left(x\right)\right){\displaystyle \frac{\delta }{\delta A^{1c}\left(x\right)}}j^{0b}\left(y\right).`$
Because
$`\left({\displaystyle \frac{\delta }{\delta A^{1c}\left(x\right)}}e^{\left(\epsilon /2\right)\mathrm{\Delta }_y}\right)P_{}(y,z)\stackrel{\epsilon 0}{}`$ (153)
$``$ $`\left({\displaystyle \frac{\delta }{\delta A^{1c}\left(x\right)}}e^{\left(\epsilon /2\right)\mathrm{\Delta }_y}\right)P_{}^{\left(0\right)}(y,z)=`$
$`=`$ $`{\displaystyle \frac{1}{2\pi }}\left({\displaystyle \frac{\delta }{\delta A^{1c}\left(x\right)}}e^{\left(\epsilon /2\right)\mathrm{\Delta }_y}\right){\displaystyle 𝑑E\theta \left(E\right)e^{iE\left(yz\right)}}\stackrel{\epsilon 0}{=}`$
$`=`$ $`{\displaystyle \frac{i}{2\pi }}{\displaystyle 𝑑E\theta \left(E\right)\epsilon Ee^{iE\left(yz\right)}e^{\left(\epsilon /2\right)p^2}\delta \left(xy\right)t^c}`$
and
$`{\displaystyle \frac{i}{2\pi }}{\displaystyle 𝑑E\theta \left(E\right)\epsilon Ee^{iE\left(yz\right)}e^{\left(\epsilon /2\right)E^2}e^{\left(\epsilon /2\right)\underset{z}{\overset{}{\mathrm{\Delta }}}}}`$ (154)
$`=`$ $`{\displaystyle \frac{i}{2\pi }}{\displaystyle 𝑑E\theta \left(E\right)\epsilon Ee^{iE\left(yz\right)}e^{\epsilon E^2}},`$
we obtain
$`{\displaystyle \frac{\delta }{\delta A^{1c}\left(x\right)}}j^{0b}\left(y\right)=`$ (155)
$`=`$ $`\underset{zy}{lim}\mathrm{t}rt^b[\left({\displaystyle \frac{\delta }{\delta A^{1c}\left(x\right)}}e^{\left(\epsilon /2\right)\mathrm{\Delta }_y}\right)P_{}(y,z)e^{\left(\epsilon /2\right)\underset{z}{\overset{}{\mathrm{\Delta }}}}+`$
$`+e^{\left(\epsilon /2\right)\mathrm{\Delta }_y}P_{}(y,z)\left({\displaystyle \frac{\delta }{\delta A^{1c}\left(x\right)}}e^{\left(\epsilon /2\right)\underset{z}{\overset{}{\mathrm{\Delta }}}}\right)]=`$
$`=`$ $`{\displaystyle \frac{i}{\pi }}{\displaystyle \theta \left(E\right)\epsilon Ee^{\epsilon E^2}𝑑E\delta \left(xy\right)\mathrm{t}rt^bt^c}.`$
The integral is
$$\theta \left(E\right)\epsilon Ee^{\epsilon E^2}𝑑E=\frac{1}{2}$$
(156)
and therefore
$$\frac{\delta }{\delta A^{1c}\left(x\right)}j^{0b}\left(y\right)=\frac{i}{2\pi }\delta \left(xy\right)\mathrm{t}rt^bt^c.$$
(157)
So, for the commutator
$$[X^a\left(x\right),j^{0b}\left(y\right)]$$
we finally get
$$[X^a\left(x\right),j^{0b}\left(y\right)]=\frac{i}{2\pi }\left(_1^x\delta ^{ac}+f^{aec}A_1^e\left(x\right)\right)\delta \left(xy\right)\mathrm{t}rt^bt^c.$$
(158)
## Appendix C The projection operator
For our purposes we expand the projection operator (138)
$`P_{}(t,x,y)`$ $`=`$ $`x\left|{\displaystyle _C_{}}{\displaystyle \frac{dE}{2\pi i}}{\displaystyle \frac{1}{EH\left(t\right)}}\right|y=`$ (159)
$`=`$ $`x\left|{\displaystyle _C_{}}{\displaystyle \frac{dE}{2\pi i}}{\displaystyle \frac{1}{EH_0V\left(t\right)}}\right|y=`$
$`=`$ $`x\left|{\displaystyle _C_{}}{\displaystyle \frac{dE}{2\pi i}}{\displaystyle \frac{1}{EH_0}}\right|y+`$
$`+x\left|{\displaystyle _C_{}}{\displaystyle \frac{dE}{2\pi i}}{\displaystyle \frac{1}{EH_0}}V\left(t\right){\displaystyle \frac{1}{EH_0}}\right|y+\mathrm{}=`$
$`=`$ $`P_{}^{\left(0\right)}(x,y)+P_{}^{\left(1\right)}(t,x,y)+\mathrm{},`$
where $`C_{}`$ is a contour surrounding the negative real axis in the complex $`E`$ plane.
For the calculation of the commutators it is sufficient to consider only the first term of (159)
$`P_{}^{\left(0\right)}(x,y)`$ $`=`$ $`{\displaystyle \underset{E_n<0}{}}\zeta _n\left(x\right)\zeta _n^+\left(y\right)=`$ (160)
$`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle 𝑑E\theta \left(E\right)e^{iE(xy)}}.`$ (161)
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# Damage of cellular material under simultaneous application of pressure and pulsed electric field
## I Introduction
Mechanical expression (hydraulic pressing) is widely used in the processes of solid-liquid separation for extraction of fruit juices and vegetable oils, dewatering of fibrous materials, etc. (Schwartzberg, 1983). Efficiency of this process can be increased by raw material plasmolysis, cellular damage or permeabilization prior to its expression. Different methods are traditionally used to increase the degree of raw material plasmolysis: heating, osmotic drying or freezing dehydration, alkaline breakage, enzymatic treatment, etc. (Rao & Lund, 1986; Aguilera & Stanley, 1999; Tsuruta, Ishimoto & Masuoka, 1998; Ponant, Foissac & Esnault, 1988; Jones, 1988; Barbosa-Cánovas & Vega-Mercado, 1996). Earlier on, the method of electric field treatment (both d.c. and a.c.) was also proposed for cellular material plasmolysis (known as electro-plasmolysis). The methods of electro-plasmolysis were shown to be good for juice yield intensification and for improving the product quality in juice production (Scheglov, Koval, Fuser, Zargarian, Srimbov, Belik et al., 1988; McLellan, Kime & Lind, 1991, Bazhal & Vorobiev, 2000), processing of vegetable and plant raw materials (Papchenko, Bologa, Berzoi, 1988; Grishko, Kozin, Chebanu, 1991), food stuffs processing (Miyahara, 1985), winemaking (Kalmykova, 1993), and sugar production (Gulyi, Lebovka, Mank, Kupchik, Bazhal, Matvienko et al., 1994; Jemai, 1997). But all these electric field applications are usually restricted by the high and uncontrolled increase in cellular tissue temperature and product quality deterioration because of electrode material electrolytic reactions, etc.
Recently, a variety of new high and moderate pulsed electric field (PEF) applications were successfully demonstrated for liquid and solid foods (Barbosa-Cánovas, Pothakamuri, Palou & Swanson, 1998; Wouters & Smelt, 1997; Knorr, Geulen, Grahl & Sitzmann, 1994; Knorr & Angersbach, 1998, Barsotti & Cheftel, 1998). The PEF application provides a possibility of fine regulation of electric power input and may result in effective permeabilization of cellular membranes (Zimmermann, 1975; Chang et al., 1992; Weaver & Chizmadzhev, 1996) without significant temperature elevation (Barbosa-Cánovas et al., 1998).
One of emerging and promising method is the combined PEF and pressure application, which demonstrates significant yield intensification for juice extracted from apples and beets and clarification of the extracted juice (Vorobiev, Bazhal & Lebovka, 2000; Gulyi et al., 1994). But the major problem arising from simultaneous application of mechanical expression and PEF treatment is the choice of optimal modes of treatment. The mechanism of solid/liquid expression from cellular materials is rather complex and may include many different phases of consolidation process (Lanoiselle, Vorobiev, Bouvier & Piar, 1996). The electric breakdown of a cellular system can influence consolidation phases and change drastically the expression curves. Unfortunately, up to now there are no accepted mechanism of electric breakdown in the cellular systems and reliable criteria for choosing optimal parameters of electric field treatment (Lebovka et al., 1996, Lebovka, Bazhal & Vorobiev, 2000). Another problem is the poor reproducibility of the experimental data, which is typical for objects of biological origin.
The properties of cellular materials influence significantly the the electrotreatment efficiency. The electrometry can be used for characterization of changes in tissue properties under the influence of external factors (electric field, pressure). This is a simple method, as far as the electrical conductivity, $`\sigma `$, reflects a degree of a water saturated tissue permeability (Sahimi, 1995). But the general dependence between the structure of a cellular material and $`\sigma `$ may be rather complex, because the conductivity of a biological tissue may be influenced by a number of processes, such as resealing of membranes in cells (Heinz, Angersbach & Knorr, 1999), diffusional redistribution of moisture inside the samples, etc.
The objective of this study is the optimization of moisture expression from biological raw materials under simultaneous pressing and PEF treatment.
The liquid expression from fine-cut cellular tissue after PEF application for different modes and durations of precompression has been investigated. Apple was used as the example of cellular material.
The useful method of data treatment, which allows to reduce data scattering caused by differences in quality of the samples is described. Discussion of the consolidation kinetics before and after PEF treatment is also given.
## II Preliminary remarks
The reasons for the simultaneous application of mechanical expression and PEF treatment are as follows. The excessive quantity of extraparticle liquid and absence of contacts between solid particles increases electrical energy losses. So the effectiveness of the PEF treatment is restricted by the uniform and tight packing of raw material between electrodes and previous removing of extraparticle air and excessive liquid (from cells destroyed by cutting). The method of raw material compact formation is its pre-consolidation. Moreover, effectiveness of the PEF treatment with respect to the water-saturated cellular materials is restricted by the low values of moisture content, $`W`$. The PEF treatment is ineffective for a water-saturated system (which is the case for the fine-cut apple raw material) because of electric breakage and uncontrolled increase of current flow through the system. The initial steps of consolidation remove an excess liquid from the extracellular volume. Therefore, we can expect increase of the PEF treatment efficiency after pre-consolidation of the raw material.
## III Materials and methods
### A Preparation of apple slices
Freshly harvested apples of Golden Delicious variety were selected for investigation and stored at $`4^{}`$C until required. The moisture content of apples $`W`$ was within 80-85%. The fine-cut apple pieces (3-5 mm diameter) were prepared from an apple pap using rasp.
### B Experimental setup and instrumentations
Figure 1 is a schematic representation of the experimental set-up. All experiments were carried out using laboratory filter-press cell equipped with an electrical treatment system. The polypropylene frame had a cylindrical cavity compartment (20 mm thick, 56 mm in diameter). The cavity compartment of frame was initially filled up with apple slices and then was tightly closed from both sides by the steel plates. One of the plates, covered by a filter cloth, was used as a stationary electrode. The other plate was attached with an elastic rubber diaphragm. A mobile wire gauze electrode was installed between the diaphragm and the layer of apple slices. Pressure was applied to the layer of apple slices through the elastic diaphragm using the hydraulic pressure controller GDS ’Standard’ (GDS Instruments Ltd, UK) with water as a pressure fluid. The pressure controller provided a constant pressure from 1 to 30 bars. The yield of liquid was controlled by balance PT610 (Sartorius AG, Germany).
A high voltage pulse generator, $`1500`$V-$`15`$A (Service Electronique UTC, France) provided the monopolar pulses of rectangular shape and allowed pulse duration $`t_i`$ varied within the interval of $`101000`$ $`\mu `$s (to precision $`\pm 2`$ $`\mu `$s), pulse repetition time $`\mathrm{\Delta }t`$ within the interval of $`1100`$ ms (to precision of $`\pm 0.1`$ ms) and number of pulses $`n`$ within the interval of $`1100000`$. The conductivities were measured by contacting electrode method with an LCR Meter HP $`4284`$A (Hewlett Packard, 38 mm guarded/guard Electrode-A HP 16451B) for thin apple slice samples at the frequency of $`100`$ Hz and with Conductimetre HI$`8820`$N (Hanna Instruments, Portugal) for the apple juice samples at a frequency of $`1000`$ Hz (these frequencies were selected as optimal in order to remove the influence of polarising effect on electrodes and inside the samples). Pulse protocols and all the output data (current, voltage, impedance, pressure, juice yield and temperature) were controlled using a data logger and special software HPVEE v.4.01 (Hewlett-Packard) adapted by Service Electronique UTC, France).
High resolution scanning electron microscopy (SEM) images were obtained using the instrument XL30 ESEM-FEG (Philips, V=15 kV, P=3.5 Torr). The ”WET” chamber mode allowing observation of hydrated apple specimens in their natural state was applied.
The optical absorbance of an expressed liquid was measured with Photocolorimeter CO75 (WPA Ltd, UK) at the wavelength 520 nm. The characteristic absorption spectra were determined with respect to distilled water. Transmittance of an expressed liquid was calculated as a ratio of filtered and nonfiltered liquid absorptions. The liquid was filtrated using a Whatman 2V filter paper.
### C Methods
All experiments were done using electric field voltages $`U`$ from $`200`$ to $`1500`$ V, pulse duration $`t_i=100`$ s, pulse repetition time $`\mathrm{\Delta }t=10`$ ms, number of pulses $`N=50`$, constant pressure $`P=3`$ bars and total time of mechanical expression $`t`$ up to $`10^4`$ s. Pressure value of 3 bars was accepted as the most efficient for exhibition of the effect of simultaneous pressing and PEF treatment (Vorobiev et al., 2000). The experiments were repeated, at least, five times.
## IV Results and discussion
### A Phases of consolidation
Figure 2 presents typical experimental curves of liquid yield $`Y_r`$, instantaneous flow rate $`v_r`$, pressure $`P_r`$ and specific electrical conductivity $`\sigma _r`$ vs. time $`t`$. For convenience of presentation, here all properties are reduced to their maximal values, e.g., $`Y_r=Y/Y_{max}`$, etc., and flow rate is determined as $`v=dY/dt`$.
Initially we observe a rather rapid increase of liquid yield $`Y`$, and decrease of electrical conductivity, $`\sigma `$. This behaviour corresponds to the layer pre-compaction (at a constant velocity of elastic diaphragm displacement) and expulsion of extraparticle air-liquid mixture. The maximum of instantaneous flow rate $`v`$ is observed approximately at $`t=t_{v_{max}}=5060`$ s. In the absence of PEF treatment ($`U=0`$, or $`E=0`$ ) the curve of $`\sigma (t)`$ temporarily stabilises in time interval $`300<t<1000`$ s and by this moment the pressure $`P`$ reaches its maximal value of 3 bars. This behaviour correspond to the end of pre-compaction period. Then, at $`t>1000`$, the $`\sigma (t)`$ curve ($`E=0`$) slightly rises, which can be explained by mechanical rupture of residual cells, by the material deterioration and by the effects of the biological activity of microorganisms.
Typical SEM micrographs of apple tissue structure before and after pressing are presented in Figs. 3(a) and (b), respectively. The mean size of undamaged cells is of order $`100200`$ $`\mu `$m. We see, that after pressing some of the cells are destroyed, but there exist also intact cells. So, for a given mode of treatment ($`P=3`$ bars, $`t=600`$ s) the cellular structure is not completely disrupted after pressing and there exists some isolated cells, which remain intact during pressing period. These cells can be damaged or partially permeabilized by another methods, for example, by PEF, thermal or another mode of treatment.
In analysing the results presented in Figs. 2, 3(a,b) we can discern the following phases of press-cake layer consolidation process:
* Phase I. Initial compaction of press-cake and expulsion of air-liquid mixture, or pre-consolidation period ($`0<t2t_{v_{max}}100120`$). During this phase the velocity of elastic diaphragm displacement is constant and a maximum of liquid flow rate is observed.
* Phase II. Mechanical rupture of cells and expulsion of liquid from ruptured cells ($`100120t300400`$). The decrease of both liquid flow rate and velocity of elastic diaphragm displacement and acceleration of the pressure increase are observed.
* Phase III. Final consolidation of press-cake at constant pressure, packing of press-cake and retardation of liquid flow rate ($`100120t`$). The liquid flowing from intracellular, extracellular and extraparticle volumes is expressed from the press-cake. At the beginning of this phase a minimum value of the specific electrical conductivity is observed. Moisture occupies all channels and so the press-cake is said to be in an impregnated state.
These phases are shown schematically at the top of Fig. 2.
### B Methods of data analysis for combined pressing and PEF treatment
Significant changes in kinetics of moisture expression and press-cake consolidation can be observed after PEF treatment. Figure 4(a) presents some examples of experimental curves of liquid yield $`Y(\%)`$ versus time $`t`$. It can be seen that $`Y(t)`$ curves rise significantly after PEF application (applied in this case at $`t=600`$ s ) as the result of damage or partial permeabilization of intact cells and subsequent expression of liquid.
But here, the main problem is in poor reproducibility of the experimental $`Y(t)`$ data. The measured curves of $`Y(t)`$ can deviate substantially because of differences in initial humidity of samples. This difficulty can be overcome by consideration of the normalized or reduced liquid yield. This normalization procedure was executed in two steps. We began with the first normalized form of liquid yield defined as (see Fig. 4(b)):
$$Y_E^I=Y_E(t)/Y_E(t_{max}),$$
(1)
where $`t_{max}`$ is the maximal time of pressing (here we use the value $`t_{max}=5400`$ s), $`Y_E`$ values correspond to the values of $`Y`$ at different electric field strengths $`E`$.
We assume that all liquid yield curves should be equal in the time range of $`t<t_p`$ for identical conditions of pressing. For experiments with PEF application $`(E0)`$, we take the curve $`Y_{E=0}^I(t)`$ for $`E=0`$ as a reference. Then we calculate the area under this curve for the time period $`t<t_p`$ and renormalize all the curves $`Y_E^I(t)`$ so as to obtain the same values of
$$S_E=\underset{0}{\overset{t_p}{}}Y_E^I(t)𝑑t,$$
(2)
for all curves.
Normalization coefficient is given by
$$k_E=S_{E=0}/S_E.$$
(3)
The second normalized form of liquid yield curve is defined as follows (see Fig. 4(c)):
$$Y_E^{II}(t)=k_EY_E^I(t).$$
(4)
The degree of intensification caused by PEF treatment, $`Y^{}`$ can be determined as the following ratio of the second normalized forms for pressing with and without PEF treatment:
$$Y^{}=Y_E^{II}(t_{max})/Y_{E=0}^{II}(t_{max}).$$
(5)
Here the values of $`Y_{E,E=0}^{II}(t_{max})`$ were determined in the point of the maximal time of pressing $`t=t_{max}`$.
Another interesting property of the expressing is the mean characteristic time, which characterizes durability of the process. The theory developed for mechanical expression of cellular materials (Lanoiselle et al., 1996) describes expression process with the set of characteristic times for the different expression phases. These are very valuable characteristics of the expression processes but, in practice, it is very difficult to find them proceeding from the expression curves. The main problem here is that we can never determine the exact value of limiting expression quantity $`Y_{\mathrm{}}`$ at $`t\mathrm{}.`$ For vegetable stuff we are faced with the problem of high enzymatic destruction at continuous expression. So we always stop the expression process at rather long, but finite time (in our case we choose $`t_{max}=5400`$ s) and the obtained values of $`Y(t_{max})`$ are of course less then the actual values of $`Y_{\mathrm{}}`$. More general approach implies evaluation of the first Moment of the function $`F(t)=Y(t_{max})Y(t)`$:
$$M^1=\frac{\underset{0}{\overset{t_{\mathrm{max}}}{}}F(t)t𝑑t}{\underset{0}{\overset{t_{\mathrm{max}}}{}}F(t)𝑑t}$$
(6)
It is easy to show, that for the simple exponential function $`F(t)=\mathrm{exp}(t/\tau )`$ and $`t_{\mathrm{max}}\mathrm{}`$ the first Moment is equal to the characteristic time, $`M^1=\tau .`$ For the finite but large values of $`t_{\mathrm{max}},`$ we have $`M^1\tau (1\frac{\left(t_{\mathrm{max}}/\tau \right)^2}{2e^{t_{\mathrm{max}}/\tau }})`$, and at $`t_{\mathrm{max}}/\tau =5`$ the first Moment equals to $`M^10.92\tau `$. So, in our case the value of $`M^1`$ may serve for approximate estimation of the effective characteristic time constant of the whole expression process.
It is useful to use this approach for crude estimation of characteristic time or durability of expression process after PEF intervention. In such a case we can treat $`M^1`$ as a mean characteristic time of liquid expulsion processes reflecting the summarized effect of all the mechanisms in a system. We can define the coefficient of PEF-enhanced durability as
$$\tau ^{}=M_E^1/M_{E=0}^1.$$
(7)
This coefficient shows the degree of the durability increase after the PEF treatment.
### C Influence of field strength $`E`$ and time $`t_p`$ of PEF treatment
#### 1 Liquid yield kinetics
Figure 5 presents the curves of excess normalized liquid yield $`\mathrm{\Delta }Y=Y_E^{II}Y_{E=0}^{II}`$ versus time $`t`$ at different times of PEF application, $`t_p`$, and different values of electric field strength, $`E`$. We have applied PEF treatment in different characteristic moments:
* $`t_p=0`$ s (before pressing);
* $`t_p=20`$ s (initial phase of pressing);
* $`t_p=120`$ s (second consolidation phase and specific conductivity of a tissue is low);
* $`t_p=330`$ s (press-cake pressure is achieved a constant value);
* $`t_p=600`$ s and $`t_p=5400`$ s (final consolidation phase).
The range of voltage values used corresponds to conditions of steady PEF-application regime without any disruption of electrical treatment caused by overflow of acceptable maximal current value. For the given pulses protocol, the steady electrical treatment regime was observed for the voltages not exceeding the following maximal values: $`U_{max}=600`$ V for $`t_p=0`$ s, $`U_{max}=1000`$ V for $`t_p=600`$ s, and $`U_{max}=1500`$ V for $`t_p=5400`$ s.
As can be seen from Fig. 5 the form of the liquid yield $`Y_E^{II}`$ curves in the time interval of $`t<300`$ s practically does not depend on the time of PEF input $`t_p`$ and the electric field strength values $`E`$ (that corresponds to small values of deviations $`Y_{E0}^{II}Y_{E=0}^{II}`$). Only at longer time intervals $`t>300`$ s in the phase III of consolidation, we can observe a behaviour reflecting the mode of PEF treatment.
At short times of PEF treatment, $`t_p100200`$ s, the excess liquid yield increases rapidly with $`E`$ increase as compared with untreated sample (see Fig. 5(a-c)). But ineffectiveness of PEF treatment at short $`t_p`$ is less in comparison with the cases of later PEF application. It can be explained by the influence of excess quantities of air and extraparticle liquid in the cake pores. The sample is highly water-saturated at earlier period of pre-consolidation and PEF application during this period may cause dielectric breakage and uncontrolled increase of current flow through the system. At such conditions, the intensification degree of PEF treatment $`Y^{}`$ is rather low (see black square for $`t_p=0`$ in Fig. 6(a)).
When we apply the PEF treatment later on, for example, at the beginning of the phase III of consolidation, at $`t300400`$ s a liquid excess yield seems to be less dependent on $`E`$ (Fig. 5(d)). Pressure achieves its constant value at this time and most of cells that could be destroyed mechanically at a given pressure (3 bars in our experiments) are already disrupted and most of liquid is evacuated from these cells. The residual isolated intact cells are connected with electrodes through the network of channels containing conductive moisture. At such conditions, the transmembrane potential on intact cells should be high enough, even at low values of $`E`$, and, therefore, the effective electropermeabilization of cell can be attained even at minimal values of electric field intensity ($`E=170`$ V cm<sup>-1</sup> in our experiments).
As we can see from comparison of typical micrographs of apple tissue structure received for pressing with and without electrical treatment (Figs. 3(b) and (c)), PEF application at $`t=330`$ s ($`E=500`$ V cm<sup>-1</sup>) causes almost complete destruction of the material. The similar pictures are observed in the wide interval of $`E250500`$ V cm<sup>-1</sup> and micrographs allow us to identify only certain quantity of single isolated intact cells.
Dependencies of the coefficient of PEF enhanced durability $`\tau ^{}`$ versus electric field strength, $`E`$ (Fig. 6(b)) substantiate conclusions set forth above. We see that the value of $`\tau ^{}`$ depends considerably on the electric field strength $`E`$ only at small time of PEF input ($`t_p<100200`$ s). The best liquid excess yield as compared with untreated sample may be obtained at the lowest applied field $`E`$ when the PEF is applied at an instant when the pressure in the system reaches a constant value ($`t_p=300400`$ s).
#### 2 Specific conductivity $`\sigma `$, flow rate $`v`$ and pressure $`P`$ kinetics
Figure 7 presents the experimental curves of (a) a specific electrical conductivity, $`\sigma `$, and (b) instantaneous flow rate, $`v`$, versus time $`t`$ for the compressed layer of apple slices. The PEF treatment was applied at $`t=t_p=120`$ s at different external field strengths, $`E`$. As can be seen from Fig. 7(a), the $`\sigma (t)`$ values begin to rise abruptly after PEF treatment, and this behaviour becomes more pronounced with increasing $`E`$. Such rise of $`\sigma (t)`$ values corresponds to the combined effect of mechanical rupture of cells and their electrical permeabilization.
We present schematically the model of PEF treatment with and without pressing in Fig. 8. In the absence of mechanical pressure the effect of the hidden (or passive) electrical breakdown can be rather important. The electrical conductivity of the whole system depends not only on the destruction degree of individual cells, but also on the character of their connectivity and the presence of continuously conducting channels.
At first, after PEF treatment the real degree of electrical breakdown is hidden and does not affect the conductivity of the whole system. PEF treatment permeabilizates cell membranes and intensifies diffusion processes. The time of electric conductivity build-up in a cellular material after PEF treatment can be estimated from the time constant of diffusion processes: $`\tau _Dd^2/6D1`$ s, where $`d100`$ is a mean cell dimension, and $`D10^9`$ m<sup>2</sup> s<sup>-1</sup> is a diffusion coefficient of an endocellular fluid. In general, this effect can exhibit a wide distribution of time constants $`\tau _D`$ because of differences in cell dimensions, diffusivity of intracellular solutions, degree of cell permeabilization, etc. Moreover, it can be masked by another related phenomena of resealing process. Heinz et al. (1999) observed that the insulating properties of the cell membranes can be recovered within a few seconds after the small power PEF treatment and it results in decreasing of a cellular system conductivity.
Simultaneous effect of pressing and PEF treatment($`P0`$ and $`E0`$) can cause the primary changes. First of all, mechanical and electrical stresses can be coupled to cause membrane breakdown in cells (Akinlaja & Frederick, 1998). Then, the external pressure enhances flowing of the fluid from destroyed cells to extracellular and extraparticle channels. All these decrease the retardation time of electric conductivity build-up after PEF treatment. Under simultaneous PEF treatment and material compression we can also eliminate or diminish the effect of a hidden electric breakdown and to depress the cells resealing processes. So, in a general case, we should observe in final state $`\sigma _{P0}>\sigma _{P=0}`$. This conclusion is confirmed by the data on $`\sigma (t)`$ kinetics at $`P=0`$ presented in Fig. 7(a) by the dashed line. In this experiment we have dropped the pressure after PEF treatment. This results in considerable decrease of $`\sigma (t)`$ values as compared with the curve obtained at $`P0`$ (Fig. 8).
A liquid flow rate $`v`$ also depends upon the degree of cellular system destruction as a result of PEF treatment. But the behaviour of $`v(t)`$ after PEF application does not change drastically (Fig. 7(b)). We can explain such behaviour by the fact of substantial decrease and termination of liquid expulsion from the press-cake by that time. The most pronounced peak is observed only at the highest electric field strength ($`E=400`$ V cm<sup>-1</sup>, Fig. 7(b)), and behaviour of the $`v(t)`$ curve reveals only increase of its long time tails with increase of $`E`$ .
Figure 9 presents the experimental curves of (a) reduced pressure values $`P^{}=P/P_{max}`$ and (b) pressure difference $`P_0P_E`$ versus time $`t`$ for different durations of pre-compression stage before the PEF application. It can be seen that PEF treatment diminishes tissue rigidity which corresponds to the decrease of an effective pressure in a system. Electrical treatment during the pressing period up to $`330`$ s delays pressure increase and accelerates the liquid yield. PEF application ($`t_p>330`$ s) can decrease steady-state pressure abruptly and initiate a rise of a liquid flow rate. The peaks of pressure decrease become sharper with pre-compression time increase. It corresponds to the more rapid destruction of cellular material pre-compressed during a longer time. But then, at the final stage, the mechanical properties of a press-cake reflate and pressure increases again up to the maximal level.
### D Optical properties of expressed liquid
The expressed liquid (apple juice) absorption (or coloration) and transmittance change in the course of a simple pressing process, due to filtration properties of the cellular material in the course of time. In the final phase of consolidation (phase III, see Fig. 2), the liquid coloration reduces considerable as compared with coloration of an initial portion of moisture. As we have demonstrated before (Vorobiev, Bazhal & Lebovka, 2000), the simultaneous pressing and electric field treatment result in considerable reduction of coloration of those differential portions of liquid which were obtained after PEF application. Studying the temporal dependencies of optical properties in differential liquid for different modes of pressing and PEF application is per se of great interest.
Here we will discuss only the optical properties of the cumulative expressed liquid, which is obtained as a result of combined pressing and PEF treatment at the final stage of the process. Figure 10 presents dependencies of absorption and transmittance versus electric strength $`E`$ for extracted liquid at different values of $`t_p`$. On the one hand, the PEF treatment decreases the liquid coloration, as can be seen from absorption curves in Fig.10. This is a positive factor of PEF treatment. We can explain this phenomena by improvement of the tissue filtration properties during pressing. Moreover, filtration properties of the PEF treated press-cake also get improved with $`t_p`$ increase because of increase of the pressed tissue consolidation. On the other hand, transmittance of expressed liquid reduces with $`t_p`$ decrease and increase of $`E`$. This is an undesirable phenomena. It can be explained by the influence of electrical treatment on the press-cake filtration properties. It is known that PEF application causes many defects of the tissue. Application of the PEF treatment increases yield of a liquid with high contents of suspended particles.
That’s why it is so important to choose a proper instant for PEF application allowing to obtain the cumulative liquid with low coloration and high transmittance. The PEF application at a moment when pressure in the system achieves a constant value is consistent with requirement of the best quality of a juice.
## V Conclusion
Investigations of the moisture expression from fine-cut apple raw material under simultaneous mechanical expression and PEF treatment were done. All experiments were performed using both laboratory filter-press cell and high voltage pulse generator which provided monopolar pulses of rectangular shape. The PEF treatment was applied to materials that were expressed at time $`t=t_p`$. Then the yield of liquid was analysed in comparison with that of untreated material. The experimental results were obtained for electric field strength $`E`$ varying from $`200`$ to $`1500`$ V cm<sup>-1</sup>, pulse duration $`t_i`$ = 100 s, pulse repetition time $`\mathrm{\Delta }t=10`$ ms, number of pulses $`N=50`$, constant pressure $`P=3`$ bars and maximal time of mechanical expression $`t_{max}=5.4\times 10^310^4`$ s.
The summary of results is as follows:
(1) The data obtained allows us to conclude that all kinetics curves ( $`\sigma (t)`$, $`Y(t)`$, $`v(t)`$ and $`P(t)`$) clearly reflect three main consolidation phases in cellular material.
(2) The combination of pressing and PEF treatment gives the most optimum results and permits to enhance significantly the liquid yield in comparison with samples untreated by PEF. The PEF treatment application permits to intensify pressing process whenever the PEF is applied. But best liquid excess yield results at lowest value of applied field $`E`$ may be obtained when PEF is applied after some pre-compression period. Such pressure pretreatment before PEF application is necessary for structuring uniformity of the press-cake, removing excess moisture and decreasing the electrical conductivity of cellular material. In our study, the pre-compression period duration of $`300400`$ s was optimal and for which the pressure in the system reaches a constant value. The PEF application in this moment of time results in the best quality of the expressed liquid (apple juice), which is confirmed by its coloration and transmittance.
(3) The simultaneous pressure and PEF treatment application reveals the passive form of the electrical damage. Electrical damage under a low field without pressure application develops very slow. The pressure provokes damage of defected cells, enhances diffusion migration of moisture and depresses cells resealing processes.
(4) The proposed unified approach for liquid yield data analysis allows to reduce the data scattering caused by the differences in the quality of samples.
## Acknowledgements
The authors would like to thank the “Pole Regional Genie des Procedes“ (Picardie, France) and the Society CHOQUENET for providing the financial support. Authors also thank Dr. N. S. Pivovarova and Dr. A. B. Jemai for help with the preparation of the manuscript.
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# 1 Introduction
## 1 Introduction
The study of the physics of $`B_c`$-meson<sup>1</sup><sup>1</sup>1For review see . has already a long story. The collaborative efforts of many scientists around the world have led to remarkable predictions, describing spectroscopy , decays and production mechanisms of this object. The theoretical work done was very helpful in the experimental search and discovery of this meson by CDF collaboration . As it was many times earlier with other particles, the observation of $`B_c`$-meson in nature have not diminished the interest of physicists to this object. People wonder how much we can learn from this meson about the Standard Model (SM) of particle interactions. In this paper we explore a potential of $`B_c`$ \- meson in extracting the parameters of SM lagragian, namely, the heavy quark masses.
What concerns its spectroscopic properties, this meson stands among the families of charmonium $`\overline{c}c`$ and bottominum $`\overline{b}b`$ : two heavy quarks move nonrelativistically, since the confinement scale, determining the presence of light degrees of freedom (sea of gluons and quarks), is suppressed with respect to the heavy quark masses $`m_Q`$ as well as the Coulomb-like exchanges result in transfers about $`\alpha _sm_Q^2`$, which is again much less than the heavy quark mass. It is precisely the nonrelativistic nature of heavy quark dynamics in the $`B_c`$-meson, that offers us a possibility to use the NRQCD framework and gain more inside on QCD dynamics of the constituent heavy quarks. Recently, within this framework the NLO and NNLO NRQCD sum rules were derived and analyzed for the case of $`\mathrm{{\rm Y}}`$-family . The result of this analysis was a precise determination of pole, running, 1S, potential subtracted and kinetic $`b`$-quark masses. Here we employ the same NRQCD sum rule method to determine the numerical values for different definitions of $`c`$-quark masses together with couplings and masses of lowest lying $`B_c`$-meson resonances. The present work is similar to the analysis performed for $`\mathrm{{\rm Y}}`$-family with differences in particular analytical expressions for theoretical moments of two-point correlator.
The paper is organized as follows. In section 2 we remind the reader the QCD sum rule framework for the two-point sum rule. Section 3 introduces the NRQCD framework and we comment on its connection with QCD one. Section 4 explains the strategy used for calculation of NRQCD two-point correlation function. In section 5 we calculate the perturbative theoretical expressions for moments of correlation function. Section 5 deals with the corrections to the correlator, coming from the gluon condensate operator. In section 7 we present our results on the different $`c`$-quark mass definitions, $`B_c`$-meson mass and coupling constant. The detail discussion on optimization methods is also given. And finally, section 8 contains our summary.
## 2 Two-point sum rules
In the QCD sum rules approach meson bound states are described by local interpolating currents of the form $`J=\overline{q}_1\mathrm{\Gamma }q_2`$, where $`\mathrm{\Gamma }`$ is a suitable Dirac structure to account for meson quantum numbers. Thus, the studied object in QCD sum rules is a coupling of meson under consideration to corresponding current. In the case of $`B_c`$ meson this coupling is defined by the following equation
$$0|\overline{b}i\gamma _5c|B_c=\frac{f_{B_c}M_{B_c}^2}{m_b+m_c}$$
(1)
The estimates of $`B_c`$-meson structure constant in QCD sum rule framework were already performed in a number of papers . The distinctive feature of the present work is a complete next-to-leading order analysis, containing correct treatment of Coulomb corrections. For discussion of importance of such corrections see, for example .
In QCD sum rules framework, the $`B_c`$-meson structure constant is naturally obtained from two-point correlator of the following form
$$\mathrm{\Pi }(q^2)=id^4xe^{iqx}0|J_{B_c}(x)J_{B_c}(0)^+|0,$$
(2)
where $`J_{B_c}=\overline{b}i\gamma _5c`$.
The left hand side of Eq.(2) can be computed in QCD, for $`|q^2|`$ much larger than $`\mathrm{\Lambda }_{QCD}^2`$ (or , what alternatively, for $`|q^2(m_c+m_b)^2|`$ much larger than $`\mathrm{\Lambda }_{QCD}^2`$), with the use of short-distance Operator Product Expansion (OPE) for correlation function under consideration.
$$\mathrm{\Pi }(q^2)_{QCD}=\mathrm{\Pi }_{pert}(q^2)+C_{G^2}(q^2)\frac{\alpha _s}{\pi }G^2+\mathrm{},$$
(3)
where $`C_{G^2}(q^2)`$ is a Wilson coefficient of gluon condensate operator and dots in right hand side of Eq.(3) present contributions of operators with higher dimension $`(d>4)`$.
The connection to physical spectrum of $`B_c`$-meson can be obtained by writing the following dispersion relation
$$\mathrm{\Pi }(q^2)=\frac{1}{\pi }\frac{\rho (s)_{had}ds}{sq^2}+\text{subtractions},$$
(4)
where
$`\rho (s)_{had}=\pi {\displaystyle \frac{f_{B_c}^2M_{B_c}^4}{(m_b+m_c)^2}}\delta (sM_{B_c}^2)+\rho (s)_{pert}\theta (ss_{thr}),`$ (5)
Here we have taken into account only lowest lying $`B_c`$-meson state, $`s_{thr}`$ is a continuum threshold and $`\rho (s)_{pert}`$ is connected to $`\mathrm{\Pi }_{pert}(q^2)`$ via the following dispersion relation
$$\mathrm{\Pi }_{pert}(q^2)=\frac{1}{\pi }\frac{\rho (s)_{pert}ds}{sq^2}+\text{subtractions}$$
(6)
In the numerical analysis we however explore a different anzaz, on which we will comment in the section with numerical results.
There are several schemes of QCD sum rules, which can be used for the extraction of quantities, you are interesting in. The most popular among them are momentum and Borel schemes. In this paper we will employ the first one and the studied object will be the momentum of two-point correlation function, given by the following expression
$`P_n=n!\left({\displaystyle \frac{d}{dq^2}}\right)^n\mathrm{\Pi }(q^2)|_{q^2=0}`$ (7)
Thus far we have discussed QCD sum rule framework for determination of $`B_c`$-meson structure constant. As a primary goal of this paper is to perform a consistent analysis of the same quantity in NRQCD approximation, in the next section we define NRQCD sum rule framework and comment on connection of the latter with QCD sum rule analysis.
## 3 NRQCD approximation
In this section we set up a consistent NRQCD framework, in which the two-point correlation function $`\mathrm{\Pi }(q^2)`$ can be determined in a systematic manner at next-to-leading order. Our presentation in this and the next sections closely follow that of , so we advice the reader to read that paper for more detail.
NRQCD is an effective field theory of QCD designed to handle nonrelativistic heavy-quark-antiquark systems to in principle arbitrary precision. Considering all quarks of the first and second generation as massless the NRQCD Lagrangian reads
$`_{\text{NRQCD}}={\displaystyle \frac{1}{2}}\text{Tr}G^{\mu \nu }G_{\mu \nu }+{\displaystyle \underset{q=u,d,s,c}{}}\overline{q}i/Dq`$ (8)
$`+\psi ^{}\left[iD_t+a_1{\displaystyle \frac{𝑫^2}{2M_t}}+a_2{\displaystyle \frac{𝑫^4}{8M_t^3}}\right]\psi +\mathrm{}`$
$`+\psi ^{}\left[{\displaystyle \frac{a_3g}{2M_t}}𝝈𝑩+{\displaystyle \frac{a_4g}{8M_t^2}}(𝑫𝑬𝑬𝑫)+{\displaystyle \frac{a_5g}{8M_t^2}}i𝝈(𝑫\times 𝑬𝑬\times 𝑫)\right]\psi +\mathrm{}`$
$`+\chi \chi ^{}\text{ bilinear terms and higher dimensional operators}.`$
The gluons and massless quarks are described by the conventional relativistic Lagrangian, where $`G_{\mu \nu }`$ is the gluon field strength tensor, $`q`$ the Dirac spinor of a massless quark and $`D_\mu `$ the gauge covariant derivative. For convenience, all color indices in Eq. (8) and throughout this work are suppressed. The nonrelativistic $`c`$ and $`\overline{b}`$ quarks are described by the Pauli spinors $`\psi `$ and $`\chi `$, respectively. $`D_t`$ and $`𝑫`$ are the time and space components of the gauge covariant derivative $`D`$ and $`E^i=G^{0i}`$ and $`B^i=\frac{1}{2}ϵ^{ijk}G^{jk}`$ the electric and magnetic components of the gluon field strength tensor (in Coulomb gauge). The straightforward $`\chi ^{}\chi `$ bilinear terms are omitted and can be readily obtained. The short-distance coefficients $`a_1,\mathrm{},a_5`$ are normalized to one at the Born level. The actual form of the higher order contributions to the short-distance coefficients $`a_1,\mathrm{},a_5`$ is irrelevant for this work, as we will later use the “direct matching” procedure at the level of the final result for the correlation function.
Now let us discuss our correlation function in NRQCD approximation. Eq.(2) can be rewritten as
$$\mathrm{\Pi }(q^2)=i0|T\stackrel{~}{J}_{B_c}(q)\stackrel{~}{J}_{B_c}(q)^+|0,$$
(9)
where $`\stackrel{~}{J}_{B_c}(q)=(\overline{b}i\gamma _5c)(q)`$. Expressing Dirac fields for $`\overline{b}`$ and $`c`$ \- quarks in terms of Pauli spinors $`\chi `$ and $`\psi `$
$`u_c(\text{q})`$ $`=`$ $`\sqrt{{\displaystyle \frac{E_c+m_c}{2E_c}}}\left(\begin{array}{c}\psi \\ \frac{\text{q}𝝈}{E_c+m_c}\psi \end{array}\right),`$ (12)
$`v_b(\text{q})`$ $`=`$ $`\sqrt{{\displaystyle \frac{E_b+m_b}{2E_b}}}\left(\begin{array}{c}\frac{(\text{q})𝝈}{E_b+m_b}\chi \\ \chi \end{array}\right),`$ (15)
we have<sup>2</sup><sup>2</sup>2Here and later in the paper, except stated otherwise, by $`m_c`$ and $`m_b`$ we mean the pole heavy quark masses.
$`\stackrel{~}{J}_{B_c}(q)`$ $``$ $`i\left((\chi ^{}\psi )(q){\displaystyle \frac{1}{8}}\left({\displaystyle \frac{m_bm_c}{m_bm_c}}\right)^2((𝐃\chi )^{}𝐃\psi )(q)+\mathrm{}\right)`$ (16)
$`\stackrel{~}{J}_{B_c}(q)^{}`$ $``$ $`i\left((\psi ^{}\chi )(q){\displaystyle \frac{1}{8}}\left({\displaystyle \frac{m_bm_c}{m_bm_c}}\right)^2((𝐃\psi )^{}𝐃\chi )(q)+\mathrm{}\right)`$ (17)
Inserting these expansions into Eq. (9) we obtain
$`i\mathrm{\Pi }(q^2)`$ $`=`$ $`C_1(\mu _{\mathrm{hard}},\mu _{\mathrm{fact}})𝒜_1(E,\mu _{\mathrm{soft}},\mu _{\mathrm{fact}})`$ (18)
$`{\displaystyle \frac{1}{4}}\left({\displaystyle \frac{m_bm_c}{m_bm_c}}\right)^2C_2(\mu _{\mathrm{hard}},\mu _{\mathrm{fact}})𝒜_2(E,\mu _{\mathrm{soft}},\mu _{\mathrm{fact}})+\mathrm{},`$
where
$`𝒜_1`$ $`=`$ $`0|(\chi ^{}\psi )(\psi ^{}\chi )|0,`$ (19)
$`𝒜_2`$ $`=`$ $`{\displaystyle \frac{1}{2}}0|(\chi ^{}\psi )((𝐃\psi )^{}𝐃\chi )+\text{h.c.}|0,`$ (20)
The right-hand side of Eq. (18) just represents an application of the factorization formalism proposed in . The second term in this expression is suppressed by $`v^2`$, i.e. of next-to-next-to-leading order and thus of no relevance to us in present analysis. The nonrelativistic current correlator $`𝒜_1`$ contains the resummation of the singular Coulomb terms. It incorporates all the long-distance, dynamics governed by soft scales like the relative three momentum $`m_{red}v`$ or the binding energy of the $`c\overline{b}`$ system $`m_{red}v^2`$. The constant $`C_1`$ (it is normalized to one at the Born level), on the other hand, describes short-distance effects involving hard scales of the order of heavy quark mass. It is represented only by a simple power series in $`\alpha _s`$ and does not contain any resummations in $`\alpha _s`$. We would also like to note, that $`C_1`$ is independent of $`q^2`$. In Eq. (18) we have also indicated the dependence of the NRQCD correlators and the short-distance coefficients on the various renormalization scales: The factorization scale $`\mu _{\mathrm{fact}}`$ essentially represents the boundary between hard and soft momenta. The dependence on the factorization scale becomes explicit because of ultraviolet (UV) divergences contained in NRQCD. Because, as in any effective field theory, this boundary is not defined unambiguously, both the correlators and the short-distance coefficients in general depend on $`\mu _{\mathrm{fact}}`$. The soft scale $`\mu _{\mathrm{soft}}`$ and the hard scale $`\mu _{\mathrm{hard}}`$, on the other hand, are inherent to the correlators and the short-distance constants, respectively, governing their perturbative expansion. If we would have all orders in $`\alpha _s`$ and $`v`$ at hand, the dependence of correlation function $`\mathrm{\Pi }(q^2)`$ on variations of each the three scales would vanish exactly. Unfortunately, we only perform the calculation up to next-to-leading order in $`\alpha _s`$ and $`v`$ which leads to a residual dependence<sup>3</sup><sup>3</sup>3Here we would like to note, that the object studied $`P_n`$ at NLO does not depend on $`\mu _{fact}`$ on the three scales $`\mu _{\mathrm{fact}}`$, $`\mu _{\mathrm{soft}}`$ and $`\mu _{\mathrm{hard}}`$.
## 4 Calculation of NRQCD correlator
To calculate the NRQCD correlator $`𝒜_1`$ we use methods originally developed for QED bound state calculations in the framework of NRQED and apply them to $`B_c`$-meson bound state description in the kinematic regime close to the threshold. At next-to-leading order quarks, composing $`B_c`$-meson experience only instantaneous interactions, given by the following potentials
$`V_c^{(0)}(\stackrel{}{r})`$ $`=`$ $`{\displaystyle \frac{C_F\alpha _s}{r}},`$ (21)
$`V_c^{(1)}(\stackrel{}{r})`$ $`=`$ $`V_c^{(0)}\left({\displaystyle \frac{\alpha _s}{4\pi }}\right)\left[2\beta _0\mathrm{ln}(\stackrel{~}{\mu }r)+a_1\right],\stackrel{~}{\mu }e^{\gamma _E}\mu _{\mathrm{soft}},`$ (22)
where
$`\beta _0`$ $`=`$ $`{\displaystyle \frac{11}{3}}C_A{\displaystyle \frac{4}{3}}Tn_l,`$
$`a_1`$ $`=`$ $`{\displaystyle \frac{31}{9}}C_A{\displaystyle \frac{20}{9}}Tn_l,`$ (23)
$`n_l`$ $`=`$ $`3.`$
Here $`r|\stackrel{}{r}|`$, $`C_F=\frac{4}{3}`$, $`C_A=3`$, $`T=\frac{1}{2}`$, $`\alpha _s\alpha _s(\mu _{\mathrm{soft}})`$ and $`\gamma _E`$ is the Euler-Mascheroni constant.
Thus, we can conclude that the problem of $`B_c`$-meson description close to threshold can be treated as a pure quantum two-body problem, so that we can use the well known analytic solutions of the nonrelativistic Coulomb problem for positronium and Rayleigh-Schrödinger time-independent perturbation theory (TIPT) to determine the corrections caused by all higher order interactions and effects.
The calulational procedure for two-point NRQCD correlation function may be devided in the following two steps
Step 1: Solution of the Schrödinger equation. – The Green function of the next-to leading Schrödinger equation is calculated incorporating the potentials displayed above. The correlator $`𝒜_1`$ is directly related to the zero-distance Green function of the Schrödinger equation.
Step 2: Matching calculation. – The short-distance constant $`C_1`$ is determined at $`𝒪(\alpha _s)`$ by matching QCD current $`J_{B_c}`$ to corresponding NRQCD one.
### 4.1 Solution of the Schrödinger equation
The nonrelativistic correlator $`𝒜_1`$ is calculated by determining the Green function of the Schrödinger equation ($`E\sqrt{q^2}(m_b+mc)`$)
$$\left(\frac{\stackrel{}{}^2}{2m_{red}}+\left[V_c^{(0)}(\stackrel{}{r})+V_c^{(1)}(\stackrel{}{r})\right]E\right)G(\stackrel{}{r},\stackrel{}{r}^{},E)=\delta ^{(3)}(\stackrel{}{r}\stackrel{}{r}^{})$$
(24)
The relation between the correlator $`𝒜_1`$ and Green function reads
$`𝒜_1`$ $`=`$ $`6\left[\underset{|\stackrel{}{r}|,|\stackrel{}{r}^{}|0}{lim}G(\stackrel{}{r},\stackrel{}{r}^{},E)\right].`$ (25)
Eq. (25) can be quickly derived from the facts that $`G(\stackrel{}{r},\stackrel{}{r}^{},\stackrel{~}{E})`$ describes the propagation of $`\overline{b}`$ and $`c`$ quark pair, which is produced and annihilated at relative distances $`|\stackrel{}{r}|`$ and $`|\stackrel{}{r}^{}|`$, respectively, and that the same quark pair is produced and annihilated through the $`J_{B_c}`$ current at zero distances. Therefore $`𝒜_1`$ must be proportional to $`lim_{|\stackrel{}{r}|,|\stackrel{}{r}^{}|0}G(\stackrel{}{r},\stackrel{}{r}^{},E)`$. The correct proportionality constant can then be determined by matching the result of full QCD for perturbative spectral density , presented in Appendix A, to the imaginary part of the Green function of the Coulomb nonrelativistic Schrödinger equation. We would like to emphasize that the zero-distance Green function on the right hand side of Eqs. (25) contains UV divergences which have to regularized. In the actual calculations we impose the explicit short-distance cutoff $`\mu _{\mathrm{fact}}`$. As mentioned before, this is the reason why the correlators and the short-distance constants depend explicitly on the (factorization) scale $`\mu _{\mathrm{fact}}`$. In this work we solve equation (24) perturbatively by starting from well known Green function $`G_c^{(0)}`$ of the nonrelativistic Coulomb problem
$$\left(\frac{^2}{m_{red}}V_c^{(0)}(\stackrel{}{r})E\right)G_c(\stackrel{}{r},\stackrel{}{r}^{},E)=\delta ^{(3)}(\stackrel{}{r}\stackrel{}{r}^{})$$
(26)
and incorporate all the higher order terms using TIPT.
The most general form of the Coulomb Green function reads ($`r|\stackrel{}{r}|`$, $`r^{}|\stackrel{}{r}^{}|`$)
$`G_c^{(0)}(\stackrel{}{r},\stackrel{}{r}^{},E)={\displaystyle \frac{m_{red}}{4\pi \mathrm{\Gamma }(1+i\rho )\mathrm{\Gamma }(1i\rho )}}{\displaystyle \underset{0}{\overset{1}{}}}dt{\displaystyle \underset{1}{\overset{\mathrm{}}{}}}ds\left[s(1t)\right]^{i\rho }\left[t(s1)\right]^{i\rho }\times `$ (27)
$`\times {\displaystyle \frac{^2}{ts}}\left[{\displaystyle \frac{ts}{|s\stackrel{}{r}t\stackrel{}{r}^{}|}}\mathrm{exp}\left\{ip(|\stackrel{}{r}^{}|(1t)+|\stackrel{}{r}|(s1)+|s\stackrel{}{r}t\stackrel{}{r}^{}|)\right\}\right],r^{}<r,`$
where
$$pm_{red}v=\sqrt{m_{red}(E+iϵ)},\rho \frac{C_Fa_s}{2v}$$
(28)
and $`\mathrm{\Gamma }`$ is the gamma function. The case $`r<r^{}`$ is obtained by interchanging $`r`$ and $`r^{}`$. $`G_c^{(0)}(\stackrel{}{r},\stackrel{}{r}^{},E)`$ represents the analytical expression for the sum of ladder diagrams with Coulomb exchanges. The analytic form of the Coulomb Green function shown in Eq. (27) has been taken from Ref. . Fortunately we do not need the Coulomb Green function in its most general form but only its S-wave component with one of the relative quark distances set to zero.<sup>4</sup><sup>4</sup>4In the section, discussing the nonperturbative corrections, coming from gluon condensate operator, one more representation of Coulomb Green function will be introduced.
$`G_c^{(0)}(0,r,E)`$ $`=`$ $`G_c^{(0)}(0,\stackrel{}{r},E)=i{\displaystyle \frac{m_{red}p}{2\pi }}e^{ipr}{\displaystyle \underset{1}{\overset{\mathrm{}}{}}}𝑑te^{2iprt}\left({\displaystyle \frac{1+t}{t}}\right)^{i\rho }`$ (29)
$`=`$ $`i{\displaystyle \frac{m_{red}p}{2\pi }}e^{ipr}\mathrm{\Gamma }(1i\rho )U(1i\rho ,2,2ipr)`$
where $`U(a,b,z)`$ is a confluent hypergeometric function . It is an important fact that $`G_c^{(0)}(0,\stackrel{}{r},E)`$ diverges for the limit $`r0`$ because it contains power ($`1/r`$) and logarithmic ($`\mathrm{ln}r`$) divergences . As we have explained before these ultraviolet (UV) divergences are regularized by imposing the small distance cutoff $`\mu _{\mathrm{fact}}`$. The regularized form of $`lim_{r0}G_c(0,\stackrel{}{r},E)`$ reads
$$G_c^{(0),reg}(0,0,E)=\frac{m_{red}^2}{4\pi }\left\{ivC_Fa_s\left[\mathrm{ln}(i\frac{m_{red}v}{\mu _{\mathrm{fact}}})+\gamma _\text{E}+\mathrm{\Psi }\left(1i\frac{C_Fa_s}{2v}\right)\right]\right\},$$
(30)
where the superscript “reg” indicates the cutoff regularization and $`\mathrm{\Psi }(z)=d\mathrm{ln}\mathrm{\Gamma }(z)/dz`$ is the digamma function. For the regularization we use the convention where all power divergences $`\mu _{\mathrm{fact}}`$ are freely dropped and only logarithmic divergences $`\mathrm{ln}(\mu _{\mathrm{fact}}/m_{red})`$ are kept. Further, we define $`\mu _{\mathrm{fact}}`$ such that in the expression between the brackets all constants except the Euler-Mascheroni constant $`\gamma _\text{E}`$ are absorbed. The results for any other regularization scheme which suppressed power divergences (like the $`\overline{\text{MS}}`$ scheme) can be obtained by redefinition of the factorization scale. For convenience we suppress the superscript “reg” from now in this work.
The Coulomb Green function contains $`c\overline{b}`$ bound state poles at the energies $`\sqrt{s}_n=m_b+m_cC_F^2a_s^2m_{red}/4n^2`$ ($`n=1,2,\mathrm{}\mathrm{}`$). These poles come from the digamma function in Eq. (30) and correspond to the nonrelativistic positronium state poles known from QED . They are located entirely below the threshold point $`\sqrt{s}_{\text{thr}}=m_b+m_c`$. This can be seen explicitly from the expression for imaginary part of Coulomb Green function $`G_c^{(0)}(0,0,E)`$
$`\text{Im}\left[G_c^{(0)}(0,0,E)\right]`$ $`=`$ $`4\pi m_{red}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}|\mathrm{\Psi }_n(0)|^2\delta (ss_n)+`$ (31)
$`\mathrm{\Theta }(E){\displaystyle \frac{1}{4\pi }}m_{red}^2{\displaystyle \frac{C_Fa_s\pi }{1\mathrm{exp}(\frac{C_Fa_s\pi }{v})}},`$
where $`|\mathrm{\Psi }_n(0)|^2=(m_{red}C_Fa_s)^3/8\pi n^3`$ is the modulus squared of the LO nonrelativistic bound state wave functions for the radial quantum number $`n`$. The continuum contribution on the right-hand side of Eq. (31) is sometimes called “Sommerfeld factor” or “Fermi factor” in the literature. And the second term from the first line in Eq. (31) describes the resonance contributions. And finally, the corrections to the zero-distance Coulomb Green function calculated below lead to higher order contributions to the bound state energy levels poles and the continuum.
Let us now come to the determination of the corrections to the zero-distance Coulomb Green function coming from the remaining term in the Schrödinger equation (24). At next-to-leading order only the one-loop contributions to the Coulomb potential, $`V_c^{(1)}`$ have to considered. Using first order TIPT in configuration space representation the NLO corrections to $`G_c^{(0)}(0,0,E)`$ reads
$`G_c^{(1)}(0,0,E)`$ $`=`$ $`{\displaystyle d^3\stackrel{}{r}G_c^{(0)}(0,r,E)V_c^{(1)}(\stackrel{}{r})G_c^{(0)}(r,0,E)}.`$ (32)
We will not calculate explicitly NLO correction to the Coulomb Green function here as the goal of this paper is to calculate the theoretical expressions for moments. The later can be most conveniently calculated by dispersion integration, using the following representation for the theoretical moments
$`P_n^{th}={\displaystyle \frac{6}{\pi }}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{ds}{s^n}}\text{Im}\{C_1(\mu _{hard},\mu _{soft})G_c(0,0,E)\},`$ (33)
where $`E=\sqrt{s}m_bm_c`$.
### 4.2 Determination of the short distance coefficients
The short-distance coefficient $`C_1`$ and $`C_2`$ can be determined by matching perturbative calculations of the matrix elements in full QCD and NRQCD. A convenient choice for matching is the matrix element between the vacuum and the state $`|c\overline{b}`$ consisting of a c and a $`\overline{b}`$ on their perturbative mass shells with nonrelativistic four-momenta $`p`$ and $`p^{}`$ in the center of momentum frame: $`𝐩+𝐩^{}=0`$. The matching condition is
$$0|\overline{b}\gamma ^0\gamma _5c|c\overline{b}|_{\text{QCD}}=C_10|\chi _b^{}\psi _c|c\overline{b}|_{\text{NRQCD}}+C_20|(𝐃\chi _b)^{}𝐃\psi _c|c\overline{b}|_{\text{NRQCD}}+\mathrm{},$$
(34)
To determine the short distance coefficients to order $`\alpha _s`$, we must calculate the matrix elements on both sides of (34) to order $`\alpha _s`$. It is sufficient to calculate the order-$`\alpha _s`$ correction only for the coefficient $`C_1`$, since as we already said the contribution proportional to $`C_2`$ is suppressed by $`v^2`$. The coefficient $`C_1`$ can be isolated by taking the limit $`𝐩0`$, in which case the matrix element of $`(𝐃\chi _b)^{}𝐃\psi _c`$ vanishes. Such calculations were done in , where the following result was obtained for $`C_1`$:
$$C_1=\mathrm{\hspace{0.33em}1}+\frac{\alpha _s(m_{\mathrm{red}})}{\pi }\left[\frac{m_bm_c}{m_b+m_c}\mathrm{log}\frac{m_b}{m_c}2\right],$$
(35)
where $`m_{\mathrm{red}}`$ is the scale of the running coupling constant. The same result gives the direct matching of NRQCD correlator with QCD one, taking into account factor 2 for $`\alpha _s`$ correction.
## 5 Dispersion integration
In this section we will discuss issues related to dispersion integration in expressions for NRQCD moments. In general, the integration of spectral density over complete covariant form of integration measure $`\frac{ds}{s^{n+1}}`$ is quite cumbersome However, in NRQCD approximation, as we will see soon, this task significantly simplifies. Let us make a change of variables $`E=\sqrt{s}m_bm_c`$ and consider a limit $`Em_b+m_c`$. In this case the integration measure takes the form
$`{\displaystyle \frac{ds}{s^{n+1}}}`$ $`=`$ $`{\displaystyle \frac{1}{(m_b+m_c)^{2n}}}{\displaystyle \frac{2dE}{m_b+m_c}}\mathrm{exp}\{(2n+1)\mathrm{ln}(1+{\displaystyle \frac{E}{m_b+m_c}})\}`$
$`{\displaystyle \frac{1}{(m_b+m_c)^{2n}}}{\displaystyle \frac{2dE}{m_b+m_c}}\mathrm{exp}\{{\displaystyle \frac{2En}{m_b+m_c}}\}+O({\displaystyle \frac{2E}{m_b+m_c}})`$
The dispersion integration for NRQCD moments in this limit is
$$P_n^{th}=\frac{1}{(m_b+m_c)^{2n}}_{E_{bind}}\frac{2dE}{m_b+m_c}\mathrm{exp}\{\frac{2En}{m_b+m_c}\}R_{NLO}^{th}(E),$$
(37)
where $`E_{bind}=\frac{m_{red}C_F^2\alpha _s^2}{2}`$ is the negative binding energy of the lowest lying resonance. This integration is performed most efficiently by deforming the path of integration into negative complex energy plane, such that the part of the integration path parallel to imaginary axis is far away from bound state poles. That is
$$P_n^{th}=\frac{1}{(m_b+m_c)^{2n}}\frac{1}{2i}_{\gamma i\mathrm{}}^{\gamma +i\mathrm{}}\frac{2dE}{m_b+m_c}\mathrm{exp}\{\frac{2En}{m_b+m_c}\}C_1A_1(E),$$
(38)
where $`\gamma E_{bind}`$. Note, that in the above equation the real part of correlator $`A_1`$ is also present, which is needed for integration over the new path. After performing the second change of variables $`E\stackrel{~}{E}`$
$$P_n^{th}=\frac{\pi }{(m_b+m_c)^{2n}}\frac{1}{2\pi i}_{\gamma i\mathrm{}}^{\gamma +i\mathrm{}}\frac{2d\stackrel{~}{E}}{m_b+m_c}\mathrm{exp}\{\frac{2\stackrel{~}{E}n}{m_b+m_c}\}C_1A_1(\stackrel{~}{E})$$
(39)
we see, that the problem of evaluation of NRQCD moments can be related to the inverse Laplace transform of integrand expression, for with there are a lot of tables in literature. The procedure of taking inverse Laplace transform can be further simplified by noting that integration path is far away from bound state energies and hence the integrand can be safely expanded in $`\alpha _s`$.
Following the steps, described above, and using relations for inverse Laplace transform from Appendix B the NRQCD moments in leading order of NRQCD expansion have the form
$$[P_n^{th}]^{LO}=\frac{3(m_{red})^{3/2}(m_b+m_c)^{1/2}}{2\sqrt{\pi }(m_b+m_c)^{2n}n^{3/2}}\mathrm{\Phi }^0(\gamma ),$$
(40)
where
$$\mathrm{\Phi }^0(\gamma )=1+2\sqrt{(}\pi )\gamma +\frac{2\pi ^2}{3}\gamma ^2+4\sqrt{\pi }_{p=1}^{\mathrm{}}(\frac{\gamma }{p})^3\mathrm{exp}\{\left(\frac{\gamma }{p}\right)^2\}[1+\text{erf}\left(\frac{\gamma }{p}\right)]$$
(41)
and
$$\gamma \frac{C_F\alpha _sm_{red}^{1/2}n^{1/2}}{(m_b+m_c)^{1/2}}$$
(42)
The calculation of NLO order correction to moments, coming from correction to potential (22) is a bit more involved. With the help of presentation (29) for Coulomb Green function we have
$`[P_n^{th}]^{NLO}`$ $`=`$ $`{\displaystyle \frac{6}{(m_b+m_c)^{2n}}}{\displaystyle \frac{1}{2i}}{\displaystyle _{\gamma i\mathrm{}}^{\gamma +i\mathrm{}}}{\displaystyle \frac{2d\stackrel{~}{E}}{m_b+m_c}}\mathrm{exp}\{{\displaystyle \frac{2\stackrel{~}{E}n}{m_b+m_c}}\}\times `$ (43)
$`\left\{4\pi {\displaystyle _0^{\mathrm{}}}r^2dr\left({\displaystyle \frac{m_{red}k}{\pi }}\right)^2{\displaystyle _0^{\mathrm{}}}dt{\displaystyle _0^{\mathrm{}}}due^{2kr(t+u+1)}\left({\displaystyle \frac{(1+t)(1+u)}{tu}}\right)^\lambda \right\}\times `$
$`C_F\left({\displaystyle \frac{\alpha _s^2}{4\pi }}\right)\left\{2\beta _0({\displaystyle \frac{1}{r}}\mathrm{log}(\stackrel{~}{\mu }r))+{\displaystyle \frac{a_1}{r}}\right\},`$
where $`\lambda =\frac{C_Fm_{red}\alpha _s}{k}`$ and $`k=2m_{red}\stackrel{~}{E}`$. The integration over $`r`$ can be easily performed explicitly, while the situation with integrations over $`t`$ and $`u`$ is far more complicated. However, as we noted above the integrand expression in the case under consideration can be expanded in series over $`\lambda `$<sup>5</sup><sup>5</sup>5It’s the same as the expansion in $`\alpha _s`$ . Then the result of integration over all terms in $`\lambda `$ expansion, except $`\lambda ^0`$ can easily expressed in terms of functions $`w_p^1`$ and $`w_p^0`$, introduced in <sup>6</sup><sup>6</sup>6The expressions for them can be found in Appendix B. As for the term with $`\lambda ^0`$, which contain a manifestly divergent integral, it can be easily calculated in terms of $`K(\tau )=0|\mathrm{exp}(H\tau )|0`$, using the free evolution function $`K_0(\tau )=0|\mathrm{exp}(H_0\tau )|0`$. Here $`H=H_0\frac{C_F\alpha _s}{r}`$, $`H_0=\frac{\stackrel{}{}^2}{2m_{red}}`$ and $`\tau `$ is an Euclidean time. The evolution function $`K(\tau )`$ can be related to the $`n`$’th NRQCD moment by the following relation
$`P_n^{th}={\displaystyle \frac{6}{(m_b+m_c)^{2n}}}{\displaystyle \frac{2\pi }{m_b+m_c}}K({\displaystyle \frac{2n}{m_b+m_c}})`$ (44)
Using this relation and explicit expression for free quark propagation function:
$$𝐱|\mathrm{exp}(H_0\tau )|𝐲=\left(\frac{m_{red}}{2\pi \tau }\right)^{3/2}\mathrm{exp}\left(\frac{m_{red}}{2\tau }(𝐱𝐲)^\mathrm{𝟐}\right),$$
(45)
we can now easily evaluate the contribution of $`\lambda ^0`$ term for NRQCD moments. The combined result for NLO order correction to NRQCD moments reads
$`[P_n^{th}]^{NLO}={\displaystyle \frac{3(m_{red})^{3/2}(m_b+m_c)^{1/2}}{2\sqrt{\pi }(m_b+m_c)^{2n}n^{3/2}}}\mathrm{\Phi }^1(\gamma ),`$ (46)
where
$`\mathrm{\Phi }^1(\gamma )={\displaystyle \frac{2\beta _0\alpha _s}{\sqrt{\pi }}}\gamma \left\{{\displaystyle \frac{1}{2}}\mathrm{log}({\displaystyle \frac{\mu _1e^{\gamma _E}\sqrt{n}}{2\sqrt{m_{red}(m_b+m_c)}}})+{\displaystyle \underset{p=1}{\overset{\mathrm{}}{}}}\gamma ^p[w_p^1+w_p^0(\mathrm{log}\left({\displaystyle \frac{2m_{red}(m_b+m_c)}{\mu _1\sqrt{n}}}\right)+{\displaystyle \frac{1}{2}}\mathrm{\Psi }\left({\displaystyle \frac{p}{2}}\right))]\right\}`$ (47)
Here $`\mu _1=\mu _{soft}\mathrm{exp}\left(\frac{a_1}{2\beta _0}\right)`$. The full perturbative result for NRQCD moments with account of hard gluon corrections reads
$$[P_n^{th}]^{pert}=(1+\frac{\alpha _s(m_{\mathrm{red}})}{\pi }\left[\frac{m_bm_c}{m_b+m_c}\mathrm{log}\frac{m_b}{m_c}2\right])\{[P_n^{th}]^{LO}+[P_n^{th}]^{NLO}\}$$
(48)
## 6 Nonperturbative corrections
In this subsection we would like to consider corrections, given by gluon condesate operator. The calculation is very similar to the one done previously by Voloshin and Leutwyler for the case of equal quark masses. For the determination of corrections to Coulomb Green function, coming from gluon condensate operator, we will exploit the fact, that the size of vacuum fluctuations of gluon field is much larger then the size of $`B_c`$-meson state. So, we can perform a multipole expansion for interaction of $`B_c`$-meson with gluon condensate, whose first term is
$$H_{\mathrm{int}}=\frac{1}{2}g\xi ^a\stackrel{}{r}\stackrel{}{E}^a,$$
(49)
where $`\stackrel{}{r}=\stackrel{}{x}_c\stackrel{}{x}_{\overline{b}}`$, $`g^2=4\pi \alpha _s`$, $`\xi ^a=t_1^at_2^a`$. By emploing colour and Lorenz invariance of the vacuum state one can relate the average over the vacuum state of two chromoelectric fields to the manifestly Lorenz-invariant value
$$0|E_i^aE_k^b|0=\frac{1}{96}\delta _{ik}\delta ^{ab}0|G_{\mu \nu }^cG_{\mu \nu }^c|0$$
(50)
Thus, the Coulomb Green function with account for gluon condesate corrections has the following form
$`G(\stackrel{}{x},\stackrel{}{y},E)`$ $`=`$ $`G_{(0)}{\displaystyle \frac{1}{18}}0|\pi \alpha _sG_{\mu \nu }^aG_{\mu \nu }^a|0\times `$
$`{\displaystyle d^3\stackrel{}{r}d^3\stackrel{}{r}^{}(\stackrel{}{r}\stackrel{}{r}^{})G_{(0)}(\stackrel{}{x},\stackrel{}{r},E)G_{(8)}(\stackrel{}{r},\stackrel{}{r}^{},E)G_{(0)}(\stackrel{}{r}^{},\stackrel{}{y},E)},`$
where $`G_{(0)}`$ and $`G_{(8)}`$ are Coulomb Green functions in singlet and octet colour states correspondingly. They are defined as solutions to the following equation
$$\left(\frac{p^2}{2m_{red}}+V_{(0,8)}(|\stackrel{}{x}|)+\frac{k^2}{2m_{red}}\right)G_{(0,8)}(\stackrel{}{x},\stackrel{}{y},\frac{k^2}{2m_{red}})=\delta (\stackrel{}{x}\stackrel{}{y}),$$
(52)
where
$$V_0(r)=\frac{\alpha ^{(0)}}{r}=\frac{4}{3}\frac{\alpha _s}{r},V_8(r)=\frac{\alpha ^{(8)}}{r}=\frac{2}{3}\frac{\alpha _s}{r}.$$
(53)
The expressions to NRQCD moments, coming form gluon condensate, can be calculated with the use of wave decomposition of NRQCD Green function
$$G_c^{(0,8)}(\stackrel{}{x},\stackrel{}{y};E)=\underset{l=0}{\overset{\mathrm{}}{}}(2l+1)G_c^{(0,8)l}(x,y;E)P_l((\stackrel{}{x}\stackrel{}{y})/xy),$$
(54)
where
$`G_c^{(0,8)l}(x,y;k^2/2m_{red})`$ $`=`$ $`{\displaystyle \frac{m_{red}k}{\pi }}(2kx)^l(2ky)^le^{k(x+y)}{\displaystyle \underset{s=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{L_s^{2l+1}(2kx)L_s^{2l+1}(2ky)s!}{(s+l+1m_{red}\alpha ^{(0,8)}/k)(s+l+1)!}},`$
where $`x=|\stackrel{}{x}|`$, $`y=|\stackrel{}{y}|`$. $`P_l`$ and $`L_s^p`$ are Legendre and Laguerre polynomials. The result for gluon condensate correction to Coulomb Green function is<sup>7</sup><sup>7</sup>7Here we assume that one should perform the dispertion integration in order to obtain the gluon condensate correction to NRQCD moments:
$$[P_n^{th}]^{G^2}=\frac{\sqrt{\pi }(m_{red})^{1/2}(m_b+m_c)^{1/2}n^{3/2}}{12(m_b+m_c)^{2n+3}}\chi (\gamma )0|\alpha _sG_{\mu \nu }^aG_{\mu \nu }^a|0$$
(56)
The function $`\chi (\gamma )`$ is given by the following equation
$`\chi (\gamma )`$ $`=`$ $`4\sqrt{\pi }\{{\displaystyle \underset{p=1}{\overset{\mathrm{}}{}}}(p+1)(p+2)(p+3)\{{\displaystyle \frac{1}{9p+16}}[\varphi _2({\displaystyle \frac{\gamma }{p}})\varphi _1({\displaystyle \frac{\gamma }{p}})(1+{\displaystyle \frac{p}{12}}(25{\displaystyle \frac{6}{9p+16}}))]`$
$`+{\displaystyle \frac{16}{9p+17}}[\varphi _2({\displaystyle \frac{\gamma }{p+1}})\varphi _1({\displaystyle \frac{\gamma }{p+1}})(1+{\displaystyle \frac{p+1}{6}}(5{\displaystyle \frac{3}{9n+17}}))]+`$
$`{\displaystyle \frac{4}{p+2}}[\varphi _2({\displaystyle \frac{\gamma }{p+2}}){\displaystyle \frac{17}{18}}\varphi _1({\displaystyle \frac{\gamma }{p+2}})]+{\displaystyle \frac{16}{9p+19}}[\varphi _2({\displaystyle \frac{\gamma }{p+3}})\varphi _1({\displaystyle \frac{\gamma }{p+3}})(1`$
$`{\displaystyle \frac{p+3}{6}}(5+{\displaystyle \frac{3}{9p+19}}))]+{\displaystyle \frac{1}{9p+20}}[\varphi _2({\displaystyle \frac{\gamma }{p+4}})\varphi _1({\displaystyle \frac{\gamma }{p+4}})(1{\displaystyle \frac{p+4}{12}}(25+`$
$`{\displaystyle \frac{6}{9p+20}}))]\}+{\displaystyle \frac{4}{81}}{\displaystyle \underset{p=2}{\overset{\mathrm{}}{}}}{\displaystyle \frac{p^21}{(81p^21)(81p^24)}}\varphi _1({\displaystyle \frac{\gamma }{8p}})\},`$
where
$`\varphi _1(x)`$ $`=`$ $`x^3[e^{x^2}(1+\text{erf}(x))1]({\displaystyle \frac{2}{\sqrt{\pi }}})x^2x^1,`$ (58)
$`\varphi _2(x)`$ $`=`$ $`[e^{x^2}(1+\text{erf}(x))1]x^1\varphi _1(x).`$ (59)
The above expression in the limit of equal quark masses coincides can be checked with the one derived previously by M.B.Voloshin . However, this expression is quite complicated and in the range $`\gamma 1.5`$ it is far more convenient to use an approximated formulum $`\chi (\gamma )=e^{0.80\gamma }\mathrm{\Phi }^0(\gamma )`$.
## 7 Numerical results
In this section we will discuss our numerical estimates for various $`c`$-quark mass definitions, the $`B_c`$-meson mass and coupling constant. Let us first discuss our anzaz for experimental spectral density, the need for which is dictated by our present lack of experimental measurements on higher $`B_c`$-meson states. The experience in potential models and quasi-local sum rules allow us to write
$`M_n`$ $`=`$ $`M_1+2T\mathrm{ln}n,M_1=6.3\mathrm{GeV},T=0.415\mathrm{GeV},`$
$`{\displaystyle \frac{f_n^2}{f_l^2}}`$ $`=`$ $`{\displaystyle \frac{1}{n}}{\displaystyle \frac{M_l}{M_n}},f_1(0^{})=0.4\mathrm{GeV},`$ (60)
where $`M_n`$ is the mass of $`nS`$ $`B_c`$-meson state, and $`f_n`$ is its coupling. So that, the experimental spectral density has the form
$`\rho (s)_{hard}=\pi {\displaystyle \underset{n=1}{\overset{3}{}}}{\displaystyle \frac{f_n^2M_n^4}{(m_b+m_c)^2}}\delta (sM_n^2)+\rho (s)_{pert}\theta (ss_{thr})`$ (61)
Now with the knowledge of experimental spectral density and thus of experimental moments we can compare them with theoretical ones. To begin with let’s see what values of the pole $`c`$-quark mass will give us the required values of experimental moments. For the value of $`b`$-quark pole mass $`M_b^{pole}`$, needed for calculation of theoretical moments<sup>8</sup><sup>8</sup>8The precise values of other parameters, that is gluon condensate and continuum threshold are not so important, as the correlation function depends on them weakly., we have taken the result of NNLO analysis for $`\mathrm{{\rm Y}}`$ family, so that we have $`M_b^{pole}=4.8\pm 0.06`$ GeV. Varying the soft scale in $`\mu _{soft}`$ the range $`1.11.2`$ GeV and fixing<sup>9</sup><sup>9</sup>9Note, that at NLO, due to vanishing of anomalous dimension of the current under consideration, the dependence of two-point correlator from $`\mu _{hard}`$ enters only through $`\alpha _s(\mu _{hard})`$. So, the particular value of $`\mu _{hard}`$ is not very important. the hard scale $`\mu _{hard}`$ to $`2.`$ GeV, from the stability of the ratios of experimental to theoretical moments we have that $`M_c^{pole}=1.72.1`$ GeV. The range of pole $`c`$-quark mass obtained is quite large and to reduce it we need some extra condition to satisfy. Again, from the spectroscopy of $`B_c`$-meson family it is known that the characteristic value of strong coupling constant for heavy quark dynamics inside $`B_c`$-meson is in the range $`\alpha _s=0.430.48`$. This extra requirement heavily constrains the value of $`c`$-quark pole mass, so that $`M_c^{pole}=2.03\pm 0.06`$ GeV. At Fig. 1 we have shown the ratios of experimental to theoretical moments as functions of momentum number at fixed value of $`\mu _{soft}=1.1`$ GeV and different values of $`c`$-quark pole mass.
The performed analysis shows that to satisfy chosen criteria we need very low values of $`\mu _{soft}1.1`$ GeV or lower. On the other hand, as was shown in , to match perturbative heavy quark potential with full QCD potential at characteristic distances of quarks bound inside heavy-heavy mesons, we should have the soft scale $`\mu _{soft}`$ in the region $`1.32.0`$ GeV<sup>10</sup><sup>10</sup>10Such choice of soft scale is also desirable in order to have reliable perturbative predictions for theoretical moments. To achieve this goal we modify the original TIPT for theoretical moments by shifting the leading order Coulomb potential by a part of constant term in NLO potential, that is
$`V_{modified}^{(0)}={\displaystyle \frac{C_F\alpha _v}{r}}={\displaystyle \frac{C_F\alpha _s}{r}}(1+t\left({\displaystyle \frac{\alpha _s}{4\pi }}\right))`$ (62)
The resulting theoretical moments will now depend not only from hard and soft scales but also from the value of the shift parameter $`t`$. To get rid of the latter dependence, while keeping in mind our goal, we perform optimization in parameter $`t`$, require the NLO corrections to theoretical moments at given momentum number $`n`$ do not exceed $`1\%`$<sup>11</sup><sup>11</sup>11Note, that here we will need already $`\alpha _v(\mu _{soft})=0.430.48`$. As a result we get $`M_c^{pole}=1.96\pm 0.05`$ GeV. At Fig. 1 we show the ratios of experimental to theoretical moments as functions of momentum number at fixed value of $`\mu _{soft}=1.3`$ GeV and different values of $`c`$-quark pole mass. We see, that within the error bars our results for the $`c`$-quark pole mass from these two estimates agree with each other as well as with the results of other estimates .
Having extracted the value of $`c`$-quark pole mass from $`B_c`$-sum rules, an analogous analysis can be performed for other definitions of heavy quark masses. Below we give the numerical values of these masses obtained in the same way and comment on their relation with each other. First, let us consider the running $`c`$-quark mass related to the pole one by the following relation
$`{\displaystyle \frac{M^{pole}}{\overline{m}(\overline{m})}}`$ $`=`$ $`1+1.333\left({\displaystyle \frac{\alpha _s(\overline{m})}{\pi }}\right)+[13.441.041n_l]\left({\displaystyle \frac{\alpha _s(\overline{m})}{\pi }}\right)^2`$ (63)
$`+[190.126.7n_l+0.653n_l^2]\left({\displaystyle \frac{\alpha _s(\overline{m})}{\pi }}\right)^3`$
Our estimates for the $`\overline{m}(\overline{m})`$ mass are
$$\overline{m}_c(\overline{m}_c)=1.40\pm 0.07\text{G}eV$$
(64)
However, not pole not running quark masses are qood ones when one works in the region of energies close to threshold. As we have seen in the analysis performed above the two-point correlation function has a strong dependence on renormalization scales as well as correlation between the pole mass and strong coupling constant (the situation is similar for running quark masses), so to reduce this dependence somewhat it’s would be more appropriate to explore quark masses, whose use can eliminate as much as possible all such dependencies and correlations among parameters. For this reason here we give estimates of $`1S`$ and potential subtracted $`c`$-quark masses. The former can be related to the pole mass $`M^{pole}`$ with the help of following formulae
$`M^{1S}`$ $`=`$ $`M^{pole}[1\mathrm{}^{LO}\mathrm{}^{LO}\delta ^1\mathrm{}^{LO}\delta ^2]`$ (65)
where
$`\mathrm{}^{LO}`$ $`=`$ $`{\displaystyle \frac{C_F^2\alpha _s^2}{8}},`$ (66)
$`\delta ^1`$ $`=`$ $`\left({\displaystyle \frac{\alpha _s}{\pi }}\right)[\beta _0(L+1)+{\displaystyle \frac{a_1}{2}}],`$ (67)
$`\delta ^2`$ $`=`$ $`\left({\displaystyle \frac{\alpha _s}{\pi }}\right)^2[\beta _0^2({\displaystyle \frac{3}{4}}L^2+L+{\displaystyle \frac{\zeta _3}{2}}+{\displaystyle \frac{\pi ^2}{24}}+{\displaystyle \frac{1}{4}})+\beta _0{\displaystyle \frac{a_1}{2}}({\displaystyle \frac{3}{2}}L+1)+{\displaystyle \frac{\beta _1}{4}}(L+1)`$ (68)
$`+{\displaystyle \frac{a_1^2}{16}}+{\displaystyle \frac{a_2}{8}}+(C_A{\displaystyle \frac{C_F}{48}})C_F\pi ^2],`$
$`L`$ $``$ $`\mathrm{ln}({\displaystyle \frac{\mu }{C_F\alpha _sM^{pole}}}),`$ (69)
and
$`a_2`$ $`=`$ $`\left({\displaystyle \frac{4343}{162}}+4\pi ^2{\displaystyle \frac{\pi ^4}{4}}+{\displaystyle \frac{22}{3}}\zeta _3\right)C_A^2\left({\displaystyle \frac{1798}{81}}+{\displaystyle \frac{56}{3}}\zeta _3\right)C_ATn_l`$ (70)
$`\left({\displaystyle \frac{55}{3}}16\zeta _3\right)C_FTn_l+\left({\displaystyle \frac{20}{9}}Tn_l\right)^2`$
With the use of these formulae we have the following estimate for $`1S`$ $`c`$-quark mass at optimized value of $`\mu _{soft}=1.3`$ GeV<sup>12</sup><sup>12</sup>12This value is taken from the previous analysis for $`c`$-quark pole mass performed in modified TIPT, as the shift of leading order Coulomb potential implicitly includes part of NNLO corrections. However, as we do not have up to date the complete NNLO order analysis this result should be taken with care.
$$M^{1S}=1.49\pm 0.07\text{GeV}$$
(71)
Having done this estimate it is instructive to compare it with the estimates of NLO analysis done in 1S scheme. To obtain the expressions for theoretical moments in this scheme we make a following substitution
$`{\displaystyle \frac{1}{(m_b+m_c)^{2n}}}`$
$`{\displaystyle \frac{1}{(M_b^{1S}+M_c^{1S})^{2n}}}\mathrm{exp}^{2n\mathrm{}^{LO}(\alpha _s)}\{12n{\displaystyle \frac{M_b^{1S}}{M_b^{1S}+M_c^{1S}}}\mathrm{}^{LO}(\alpha _s)\delta ^1(M_b^{1S},\alpha _s,\mu _{soft})`$
$`2n{\displaystyle \frac{M_c^{1S}}{M_b^{1S}+M_c^{1S}}}\mathrm{}^{LO}(\alpha _s)\delta ^1(M_c^{1S},\alpha _s,\mu _{soft})\}`$ (72)
In all other places the heavy quark pole masses should be changed to $`1S`$ masses. Varying the scales in the same ranges as we did for pole $`c`$-quark mass we obtain, that the sum rules are stable themselves, but the ratio of experimental to theoretical moments is far from unity for reasonable values of $`1S`$ $`c`$-quark mass<sup>13</sup><sup>13</sup>13Recall that this mass by definition differs from a half of $`J/\mathrm{\Psi }`$ mass on a small value given by nonperturbative corrections.. To solve this problem we again consider the modified TIPT for theoretical moments. It is important, when writing expressions for theoretical moments, to make for consistency a similar shift in NLO order relation between pole and $`1S`$ masses. Performing calculations in modified $`1S`$ scheme we obtain the following estimate for $`1S`$ $`c`$-quark mass: $`M_c^{1S}=1.52\pm 0.05`$ GeV, which, within the errors, is in agreement with previous estimate for $`1S`$ $`c`$-quark mass. Fig. 3 shows us the ratios of experimental to theoretical moments as functions of momentum number at fixed values of $`\mu _{soft}=1.3`$ GeV and $`1S`$ $`c`$-quark mass.
To finish the discussion of $`1S`$ $`c`$-quark mass we note, that even this mass was defined as $`M_c^{pole}`$ minus one half of perturbative Coulomb energy in $`J/\mathrm{\Psi }`$-meson, it can be equally well applied for $`B_c`$-meson as well as other systems, containing $`c`$-quark. The arguments here are the analog of Upsilon expansion for charmed quarks and the fact, that the Coulomb energy for static quarks does not depend on their flavors.
The estimate for the value of potential subtracted $`c`$-quark mass was obtained from the relation of the latter to the running quark mass
$`m_{PS}(\mu _f)`$ $`=`$ $`\overline{m}(\overline{m})\{1+{\displaystyle \frac{4\alpha _s(\overline{m})}{3\pi }}[1{\displaystyle \frac{\mu _f}{\overline{m}(\overline{m})}}]+\left({\displaystyle \frac{\alpha _s(\overline{m})}{\pi }}\right)^2[13.441.04n_l`$ (73)
$`{\displaystyle \frac{\mu _f}{3\overline{m}(\overline{m})}}(a_1+4\pi \beta _0[\mathrm{ln}{\displaystyle \frac{\mu _f^2}{(M^{pole})^2}}2])]\}`$
Thus, at $`\mu _f=1.5`$ GeV we have
$$m_{PS}(1.5\text{GeV})=1.42\pm 0.07\text{GeV}$$
(74)
Having completed the estimates of different $`c`$-quark masses, we can now perform estimates of the mass and coupling constant for the ground $`B_c`$-meson state, fixing the heavy quark masses at their central values.
The results for mass and coupling of $`B_c`$-meson in momentum scheme of NRQCD sum rules can be easily seen from Fig. 4 and 5.
Recently, the estimate of $`B_c`$ mass was obtained in the perturbative potential approach by fitting the masses of $`J/\mathrm{\Psi }`$ and $`\mathrm{{\rm Y}}`$ in order to get a good covergency in $`\alpha _s`$ and extract the heavy quark masses . The perturbative mass $`m_{B_c}=6.323\pm 0.007`$ GeV is very close to the our estimate in the framework of QCD sum rules. This fact indicates a small nonperturbative correction, which was limited by $`60`$ MeV in .
## 8 Conclusion
In this paper we have presented complete NLO analysis for $`B_c`$-meson two-point NRQCD sum rules. Analitycal results for perturbative spectral density and gluon condensate contribution with account for summed Coulomb corrections are derived and analyzed. A detail numerical analysis as well as discussion on the determination of various $`c`$-quark masses together with couplings and masses of lowest lying $`B_c`$-meson resonances from the mentioned sum rules are provided. The analysis shows that to reduce the uncertainties of calculations it is mandatory to have complete NNLO expressions for theoretical moments, which we plan to accomplish in nearest future.
The author is grateful to Prof. V.V.Kiselev for stimulating discussions and comments. I especially thank my wife for strong moral support and help in doing physics.
This work was in part supported by the Russian Foundation fro Basic Research, grants 99-02-16558 and 00-15-96645, by International Center of Fundamental Physics in Moscow, International Science Foundation and INTAS-RFBR-95I1300.
## Appendix A
In this Appendix we present the two-loop result for perturbative spectral density of $`B_c`$ meson :
$`\rho _{\mathrm{pert}}(t)`$ $`=`$ $`{\displaystyle \frac{3}{8\pi t}}\overline{q}^4v\{1+{\displaystyle \frac{4\alpha _s}{3\pi }}\{{\displaystyle \frac{3}{8}}(7v^2)`$
$`+{\displaystyle \underset{i=b,c}{}}[(v+v^1)(L_2(\alpha _1\alpha _2))L_2(\alpha _i)\mathrm{log}\alpha _1\mathrm{log}\beta _i`$
$`A_i\mathrm{log}\alpha _i+B_i\mathrm{log}\beta _i]\}+O(\alpha _s^2)\}`$
where
$$L_2(x)=_0^x\frac{dy}{y}\mathrm{log}(1y)$$
(76)
and
$`A_i`$ $`=`$ $`{\displaystyle \frac{3}{2}}{\displaystyle \frac{3m_i+m_j}{m_i+m_j}}{\displaystyle \frac{19+2v^2+3v^4}{32v}}{\displaystyle \frac{m_i(m_im_j)}{\overline{q}^2v(1+v)}}\left(1+v+{\displaystyle \frac{2v}{1+\alpha _i}}\right);`$
$`B_i`$ $`=`$ $`2+2{\displaystyle \frac{m_i^2m_j^2}{\overline{q}^2v}};`$ (77)
$`\alpha _i`$ $`=`$ $`{\displaystyle \frac{m_i}{m_j}}{\displaystyle \frac{1v}{1+v}};\beta _i=\sqrt{1+\alpha _i}{\displaystyle \frac{(1+v)^2}{4v}}`$
$`\overline{q}^2`$ $`=`$ $`t(m_bm_c)^2;v=\sqrt{14{\displaystyle \frac{m_bm_c}{\overline{q}^2}}}`$ (78)
## Appendix B
In this Appendix we have collected some formulae, needed for calculation of NRQCD moments in next to leading order.
$`{\displaystyle \frac{1}{2\pi i}}{\displaystyle _{\gamma i\mathrm{}}^{\gamma +i\mathrm{}}}{\displaystyle \frac{1}{x^\nu }}e^{xt}𝑑x`$ $`=`$ $`{\displaystyle \frac{t^{\nu 1}}{\mathrm{\Gamma }(\nu )}},`$ (79)
$`{\displaystyle \frac{1}{2\pi i}}{\displaystyle _{\gamma i\mathrm{}}^{\gamma +i\mathrm{}}}{\displaystyle \frac{\mathrm{ln}x}{x^\nu }}e^{xt}𝑑x`$ $`=`$ $`{\displaystyle \frac{t^{\nu 1}}{\mathrm{\Gamma }(\nu )}}[\mathrm{\Psi }(\nu )\mathrm{ln}t],`$ (80)
$`w_p^0`$ $`=`$ $`{\displaystyle \frac{1}{p!\mathrm{\Gamma }(\frac{p}{2})}}{\displaystyle _0^{\mathrm{}}}𝑑t{\displaystyle _0^{\mathrm{}}}𝑑u{\displaystyle \frac{1}{(1+t+u)^2}}\mathrm{ln}^p\left({\displaystyle \frac{(1+t)(1+u)}{tu}}\right)={\displaystyle \frac{(p+1)\zeta _{p+1}}{\mathrm{\Gamma }(\frac{p}{2})}},`$ (81)
$`w_p^1`$ $`=`$ $`{\displaystyle \frac{1}{p!\mathrm{\Gamma }(\frac{p}{2})}}{\displaystyle _0^{\mathrm{}}}𝑑t{\displaystyle _0^{\mathrm{}}}𝑑u{\displaystyle \frac{1\mathrm{ln}(1+t+u)}{(1+t+u)^2}}\mathrm{ln}^p\left({\displaystyle \frac{(1+t)(1+u)}{tu}}\right)`$
$`=`$ $`\left\{{\displaystyle \frac{(1+p)}{\mathrm{\Gamma }(\frac{p}{2})}}\left[\gamma _E\zeta _{p+1}+{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{\Psi }(2+m)}{(1+m)^{p+1}}}\right]+{\displaystyle \frac{2}{\mathrm{\Gamma }(\frac{p}{2})}}{\displaystyle \underset{l=0}{\overset{p1}{}}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}(1)^{pl}{\displaystyle \frac{(1+l)\mathrm{\Psi }^{(pl)}(2+m)}{(pl)!(1+m)^{1+l}}}\right\},`$
where $`\zeta _p`$ is the Rieman zeta function for argument $`p`$.
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# Generations of Higgs Bosons in Supersymmetric Models
## I Introduction
One of the great mysteries of the standard model is the number of fermion generations. Nothing in the structure of the standard model gives any clue as to what this number should be, and the same is true of the most popular extensions of the standard model $``$ grand unified theories and the MSSM. The only information we have about the number of fermion generations is phenomenological. It is interesting that this situation is the same for the Higgs bosons, i.e., there is no clue as to the number of generations of Higgs bosons ( in supersymmetric models each generation consisting of two doublets of opposite hypercharge). One might *a priori* expect the number of Higgs generations to be equal to the number of fermion generations. This is because in some supersymmetric grand unified theories, such as $`E_6`$, the Higgs bosons and fermions belong to the same representations. Just as in the fermionic sector, the number of Higgs generations can, at present, only be determined phenomenologically.
Of course, the number of Higgs generations detected experimentally at present is zero. However, a strong clue comes from the absence of tree-level flavor-changing neutral currents (FCNC). In the standard model with a single Higgs doublet, the Yukawa coupling matrices are proportional to the fermionic mass matrices. Diagonalizing the latter thus automatically diagonalizes the former, eliminating FCNC. If one or more Higgs doublets are added to the standard model, diagonalizing the mass matrices does not in general diagonalize the Yukawa coupling matrices. In minimal supersymmetric extensions of the standard model, there are two Higgs doublets of opposite hypercharge, one couples to the $`Q=2/3`$ quarks while the other couples to the $`Q=1/3`$ quarks, thereby eliminating FCNC. Again, with an extra pair of doublets added to these minimal models, FCNC naturally emerge.
Since we are interested in the possibility of having additional generations of Higgs bosons, it is important to solve the problem of FCNC . There are two approachesAlthough generally discussed in the context of two-generation models, these approaches apply more generally.. The first eliminates them completely by imposing a discrete symmetry. The second approach does not eliminate FCNC, but makes them sufficiently small by assuming that the flavor-changing neutral couplings of quarks $`q_i`$ and $`q_j`$ are proportional to the geometric mean of the Yukawa couplings of the same quarks. In this case, the most dangerous FCNC, which involve the down and strange quarks, are suppressed by the small down and strange Yukawa couplings<sup>§</sup><sup>§</sup>§Another possible solution to the FCNC may be that the additional Higgses are too heavy, and thus the effective low energy theory is the usual MSSM..
In Ref. , the first approach (applying a discrete symmetry to remove FCNC) was applied to the possibility of adding extra generations of Higgs bosons to the MSSM. They considered three generations of Higgs bosons: six doublets, three of each hypercharge. It was shown that if one assumes that some symmetry (discrete, continuous, global or local) suppresses FCNC, then the extra generations, which they called “pseudoHiggs bosons” do not acquire vacuum expectation values, do not mix with the “standard” doublets and do not couple to fermions. The lightest of these extra scalars is stable, and much of Ref. was devoted to the possibility that this scalar could constitute the dark matter.
In this article, we analyze the MSSM with extra generations of Higgs doublets. Initially, we assume that there is a symmetry which eliminates tree-level FCNC (at this point, we do not worry about the precise nature of the symmetry). In Section II, we present the model and show that the lightest neutral scalars are degenerate in mass, and that the lightest charged scalar is only a few GeV heavier. It is shown that standard searches for heavy leptons(charginos) with neutrinos(neutralinos) a few GeV lighter will be sensitive to these particles, In Section III we relax the assumption that tree-level FCNC are eliminated, and consider a simple flavor symmetry, based on U(2), and show how the resulting tree-level FCNC can be calculated in terms of fermion masses.
## II The Model
Let us first summarize the model discussed in Ref. . We consider the supersymmetric standard model with three generations of Higgs doublets. The most general superpotential is given by
$`W=\mu _{ij}H_i\overline{H}_j+f_{ijk}Q_iU_jH_k+g_{ijk}Q_iD_j\overline{H}_k+h_{ijk}L_iE_j\overline{H}_k,`$ (1)
where the lowercase Latin indices run over 1,2,3. It is assumed that some symmetry eliminates the tree-level FCNC in the quark and lepton sectors. Given this assumption, a basis can be chosen in which only one generation, conventionally chosen to be $`H_3`$ and $`\overline{H}_3`$, couples to quarks and leptons. This means that $`H_3`$ and $`\overline{H}_3`$ must have different quantum numbers under that symmetry than the other $`H_{1,2}`$ and $`\overline{H}_{1,2}`$, in order that only the former couple to fermions.
The most general soft SUSY-breaking terms involving only the Higgs fields are
$`W_{soft}m_{H_{ij}}^2H_i^{}H_j+m_{\overline{H}_{ij}}^2\overline{H}_i^{}\overline{H}_jB\mu _{ij}(H_i\overline{H}_j+h.c.).`$ (2)
Since $`H_3`$ and $`\overline{H}_3`$ have different quantum numbers than the other $`H_{1,2}`$ and $`\overline{H}_{1,2}`$ under the symmetry, then the quadratic terms involving only one of them vanish, i.e. $`m_{H_{i3}}^2=m_{\overline{H}_{i3}}^2=m_{H_{3i}}^2=m_{\overline{H}_{3i}}^2=\mu _{i3}=\mu _{3i}=0`$ for $`i=1,2`$. Thus there are no quadratic terms mixing the third generation of Higgs fields with the other two. In addition, Ref. shows that equality of scalar masses at the unification scale automatically implies that the first and second generation fields do not get vacuum expectation values. Note that we have four new scalar fields, four new pseudoscalar and four new pairs of charged scalars. In the following, we specialize to the case in which there is only a single generation of extra Higgs fields, denoted by $`H_X`$ and $`\overline{H}_X`$. This is primarily for simplicity$``$including the additional generations does not affect our results. In fact, if the two additional generations have different quantum numbers under the symmetry that eliminates FCNC, then the generations will decouple and this specialization is completely general. Even if the coupling between the two generations exists, however, the results we present below are completely unaffected$``$the results would then simply apply to the lightest of the scalar fields.
So, assuming only one extra generation of Higgs fields, the scalar potential can be written as
$`V`$ $`=`$ $`m_X^2|H_X|^2+\overline{m}_X^2|\overline{H}_X|^2+m_1^2|H|^2+m_2^2|\overline{H}|^2+(\mu _XH_X\overline{H}_X+\mu H\overline{H}+h.c.)`$ (3)
$`+`$ $`{\displaystyle \frac{g^2}{8}}{\displaystyle \underset{a}{}}\left|H_X^{}\tau _aH_X+\overline{H}_X^{}\tau _a\overline{H}_X+H^{}\tau _aH+\overline{H}^{}\tau _a\overline{H}\right|^2`$ (4)
$`+`$ $`{\displaystyle \frac{g^{}_{}{}^{}2}{8}}\left||H|^2+|H_X|^2|\overline{H}|^2|\overline{H}_X|^2\right|^2,`$ (5)
where $`H`$ and $`\overline{H}`$ are the standard MSSM doublets, $`\tau _a`$ are the Pauli matrices, and where $`\mu `$ and $`\mu _X`$ are arbitrary parameters of dimension $`(\mathrm{GeV})^2`$. The full Lagrangian has a symmetry under which the $`H_X`$ fields change sign (it is actually a global U(1) symmetry), and thus the lightest of these fields is stable.
The mass matrices of the scalars can be calculated. The mass matrix for the additional neutral scalar fields is given by
$`M_S^2=\left(\begin{array}{cc}m_X^2\frac{1}{2}M_Z^2\mathrm{cos}2\beta & \mu _X\\ \mu _X& \overline{m}_X^2+\frac{1}{2}M_Z^2\mathrm{cos}2\beta \end{array}\right),`$ (8)
and the matrix for the additional neutral pseudoscalars is given by
$`M_P^2=\left(\begin{array}{cc}m_X^2\frac{1}{2}M_Z^2\mathrm{cos}2\beta & \mu _X\\ \mu _X& \overline{m}_X^2+\frac{1}{2}M_Z^2\mathrm{cos}2\beta \end{array}\right).`$ (11)
The mass matrix for the charged scalars is
$`M_+^2=\left(\begin{array}{cc}m_X^2+\frac{1}{2}M_Z^2\mathrm{cos}2\overline{\beta }& \mu _X\\ \mu _X& \overline{m}_X^2\frac{1}{2}M_Z^2\mathrm{cos}2\overline{\beta }\end{array}\right),`$ (14)
where $`\mathrm{cos}2\overline{\beta }\mathrm{cos}2\beta \mathrm{cos}2\theta _W`$. Note that the mass matrices for the scalar and pseudoscalar are *identical* except for the sign of the even-odd elements (this is true even in the case of many additional generations all coupled together). As a result, the secular equation is identical for both mass matrices, and so the eigenvalues are the same. The charged scalar mass matrix is identical to the neutral scalar mass matrix with $`M_Z^2M_Z^2\mathrm{cos}2\theta _W`$. This results in a *larger* mass for the lightest charged scalar, but only slightly larger$``$we will see that a few GeV is a typical size.
None of this is new. It was discussed in much more detail in Ref. . However, there was an error in that paper (that also appeared in an earlier version of this paper). They argued that the $`SP`$ mass degeneracy would be split by radiative corrections involving a loop with a $`W`$ boson, noting a difference between the magnitude of the couplings of the $`\varphi _S`$ to the $`W`$ and charged scalar and the magnitude of the couplings of the $`\varphi _P`$. However, this difference was due to a sign error in a Feynman rule; the couplings are of the same magnitude. Thus the degeneracy is not split by radiative corrections. One can see this by noting that the $`U(1)`$ global symmetry which ensures the stability of the lightest scalar has the $`H_X`$ with charge +1, the $`\overline{H}_X`$ with charge -1 and the rest with charge zero. This automatically gives the mass degeneracy, and this global symmetry is unbroken. In Ref. , there was no discussion of the mass difference between the charged and neutral scalars, nor the decay modes. We now look at this mass difference, and discuss the signatures.
The mass matrices Eq. (8),(11), and (14) have four unknown parameters: $`m_X^2`$, $`\overline{m}_X^2`$, $`\mu _X`$, and $`\mathrm{tan}\beta `$. However, as noted above, the beta functions for $`m_X^2`$ and $`\overline{m}_X^2`$ are identical, and equality of the masses at a high scale implies that they are equal at all scales. We thus set them equal to each other (our results are not significantly affected by relaxing this assumption). We also see that the results are very insensitive to $`\mathrm{tan}\beta `$. Thus, there are effectively two parameters. These two parameters give the masses of the neutral scalars, and the two charged scalar pairs, as well as all of the mixing angles. The coupling to the vector bosons are thus determined in terms of these parameters, and are given in Ref. (but note that, as discussed in the previous paragraph, the Feynman rule involving $`\varphi _P,\varphi _+`$ and the $`W`$ has a sign error in front of the $`\theta _+`$ term).
We now turn to the mass difference between the lightest charged Higgs and $`\varphi _S`$. At tree level, the splitting can be easily found from the mass matrices and is given by
$$M_+^2M_S^2=\frac{1}{2}\left(\sqrt{M_Z^4\mathrm{cos}^22\beta +4\mu _X^2}\sqrt{M_Z^4\mathrm{cos}^22\overline{\beta }+4\mu _X^2}\right)$$
(15)
Numerically, for $`\mathrm{tan}\beta =2`$, the mass splitting is $`\frac{150\mathrm{GeV}}{M_S}`$ times $`(2.9,0.7,0.2)`$ GeV for $`\mu _X^{1/2}=(40,100,200)`$ GeV. For $`\mathrm{tan}\beta =10`$, the splitting for the same $`\mu _X`$ values is $`\frac{150\mathrm{GeV}}{M_S}`$ times $`(5.5,1.9,0.6)`$ GeV. Thus, we see that the mass splitting varies from a hundred MeV to a few GeV over most of parameter space.
Radiative corrections to this splitting can be calculated. They arise from corrections to the $`\varphi _+`$ propagator due to loops involving the $`Z`$ and $`\gamma `$ and corrections to the $`\varphi _S`$ propagator due to $`Z`$ loops (the corrections due to the $`W`$ are identical in the limit that the masses of the $`\varphi _+`$ and $`\varphi _S`$ are the same, and thus will be very small). We have calculated these corrections and found that they are seldom greater than $`250`$ MeV. However, loops due to the supersymmetric partners can also be significant, and these will depend on gaugino/Higgsino masses and mixing parameters. Precise predictions can therefore not be made. Nonetheless, this will not affect the conclusion that the mass splitting tends to be $`12`$ GeV, but can be as low as a hundred MeV and as high as a few GeV.
From an experimental point of view, the $`\varphi ^+`$ looks like a heavy, charged lepton with an associated neutrino which is very slightly lighter. Searches for heavy leptons with nearly degenerate neutrinos will be sensitive to this particle. (The fact that the $`\varphi `$’s are scalars will appear in the angular distribution in the production rate, as we will see). Recently, DELPHI has conducted a search for charginos which are nearly mass-degenerate with the lightest neutralino. The signatures here will be very similar. They show that the standard chargino searches will cover mass splittings down to about $`3`$ GeV, and that initial state radiation can extend this reach. As we can see, this will cover some, but not most, of parameter space. Searches for “stable” leptons will also be sensitive if the mass splitting is less than $`150`$ MeV (above this, pion decays dominate and the lifetime becomes substantially shorter). In between, one must look for the vertex, as with any intermediate range heavy lepton.
The fact that the model has so few parameters allows a precise determination of the production cross-section, unlike the chargino case, which depends on various mixing angles, scalar neutrino masses, etc. The production cross-section for $`\varphi ^+\varphi ^{}`$ at an $`e^+e^{}`$ collider is given by
$$\sigma =\frac{\pi \alpha ^2}{3s}\left(1\frac{4M_+^2}{s}\right)^{3/2}\left[1+\frac{A}{(1\frac{M_Z^2}{s})^2}\right]$$
(16)
where $`A\frac{\mathrm{cos}^22\theta _W}{4\mathrm{sin}^42\theta _W}0.15`$, and we have not included the vector coupling of the electron to the $`Z`$, which is proportional to $`\frac{1}{4}\mathrm{sin}^2\theta _W`$. For $`\sqrt{s}200`$ GeV, this gives a cross-section of $`0.6`$ picobarns times the phase space factor. The angular distribution is the usual $`\mathrm{sin}^2\theta `$ distribution for scalar particles.
## III U(2) Model
In this section we consider relaxing the assumption that a symmetry forbids tree level FCNC. We do this by giving the model an explicit flavor symmetry. A very successful and elegant flavor symmetry that has been considered in the literature is the (horizontal) U(2) flavor symmetry . In the U(2) model, which only involves the two Higgs doublets of the MSSM, the matter fields of the first and second generations transform as the components of a doublet, while the third generation transforms as a singlet, i.e. if we denote the matter fields by $`\psi `$, then $`\psi =\psi _a+\psi _3`$, where $`a=1,2`$. The two Higgs doublets transform as singlets. The minimal model contains three flavon fields which are responsible, through their vacuum expectation values (vevs), for the breaking of the flavor symmetry. The flavons consist of a doublet $`\varphi `$, a triplet $`S`$, and a (antisymmetric) singlet $`A`$. The breaking occurs in two steps
$`U(2)\stackrel{ϵ}{}U(1)\stackrel{ϵ^{}}{}nothing,`$ (17)
where $`ϵ`$ is the vev of $`\varphi `$ and $`S`$, while $`ϵ^{}`$ is the vev of $`A`$. We will consider the unified version in which the model is embedded in a SU(5) GUT. In this case, the flavon fields also transform under SU(5), and a new flavon $`\mathrm{\Sigma }`$ is introduced. The flavons and their transformations are
$`\varphi (\mathrm{𝟏},\mathrm{𝟐}),S(\mathrm{𝟕𝟓},\mathrm{𝟑}),`$ (18)
$`A(\mathrm{𝟏},\mathrm{𝟏}),\mathrm{\Sigma }(\mathrm{𝟐𝟒},\mathrm{𝟏}),`$ (19)
where the numbers in parenthesis correspond to the transformations under SU(5) and U(2) respectively. They acquire the following vevs:
$`{\displaystyle \frac{S}{M_f}}`$ $``$ $`\left(\begin{array}{cc}0& 0\\ 0& ϵ\end{array}\right),{\displaystyle \frac{\varphi }{M_f}}\left(\begin{array}{c}0\\ ϵ\end{array}\right),`$ (24)
$`{\displaystyle \frac{A}{M_f}}`$ $``$ $`ϵ^{ab}ϵ^{},{\displaystyle \frac{\mathrm{\Sigma }}{M_f}}ϵ.`$ (26)
where $`M_fϵM_{GUT}`$ is the flavor scale. This model successfully reproduces the observed quark mass ratios and CKM angles, as well as the lepton mass ratios . Let’s now consider the possibility of having three Higgs generations and letting them transform non-trivially under the flavor symmetry . The matter fields are in $`\overline{\mathrm{𝟓}}`$ $`(\overline{F})`$, and $`\mathrm{𝟏𝟎}`$’s $`(T)`$ of SU(5). Then, their transformation properties are $`(\mathrm{𝟏𝟎},\mathrm{𝟐}\mathrm{𝟏})`$ and $`(\overline{\mathrm{𝟓}},\mathrm{𝟐}\mathrm{𝟏})`$ respectively, where again the first term in the parenthesis corresponds to SU(5) and the second to U(2). There are six Higgs doublets with components $`H`$ and $`\overline{H}`$ transforming as $`(\mathrm{𝟓},\mathrm{𝟐}\mathrm{𝟏})`$ and $`(\overline{\mathrm{𝟓}},\mathrm{𝟐}\mathrm{𝟏})`$ respectively.
The Yukawa part of the superpotential is
$`W_Y`$ $`=`$ $`T_3H_3T_3+T_3\overline{H}_3\overline{F}_3\xi +{\displaystyle \frac{1}{M_f}}[T_3H_a\varphi ^aT_3+T_3\varphi ^aH_3T_a+(T_3\varphi ^a\overline{H}_3\overline{F}_a+T_a\varphi ^a\overline{H}_3\overline{F}_3`$ (27)
$`+`$ $`T_3\varphi ^a\overline{H}_a\overline{F}_3+T_a(S^{ab}+A^{ab})\overline{H}_3\overline{F}_b+T_a(S^{ab}+A^{ab})\overline{H}_b\overline{F}_3+T_3(S^{ab}+A^{ab})\overline{H}_a\overline{F}_b)\xi ]`$ (28)
$`+`$ $`{\displaystyle \frac{1}{M_f^2}}[T_a(S^{ab}\mathrm{\Sigma }+A^{ab}\mathrm{\Sigma }+\varphi ^a\varphi ^b)H_3T_b+T_3(S^{ab}\mathrm{\Sigma }+A^{ab}\mathrm{\Sigma }+\varphi ^a\varphi ^b)H_aT_b`$ (29)
$`+`$ $`(T_3\overline{H}_a\varphi ^a\varphi ^b\overline{F}_b+T_a\varphi ^a\varphi ^b\overline{H}_b\overline{F}_3+T_a(S^{ab}+A^{ab})\varphi ^c\overline{H}_b\overline{F}_c)\xi ]`$ (30)
$`+`$ $`{\displaystyle \frac{1}{M_f^3}}\left[T_a\varphi ^a\varphi ^b\varphi ^cH_bT_c+T_a\varphi ^a\varphi ^b\varphi ^c\overline{H}_b\overline{F}_c\xi \right],`$ (31)
where $`\xi m_b/m_t`$.
The Yukawa coupling matrices can now be obtained from Eq. (27) and Eq. (24), their textures are given by
$`Y_U\left(\begin{array}{ccc}0& ϵϵ^{}& 0\\ ϵϵ^{}& ϵ^2& ϵ\\ 0& ϵ& 1\end{array}\right)H_3+\left(\begin{array}{ccc}0& ϵϵ^{}& ϵϵ^{}\\ ϵϵ^{}& ϵ^3& ϵ^2\\ ϵϵ^{}& ϵ^2& ϵ\end{array}\right)H_2+\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& ϵϵ^{}\\ 0& ϵϵ^{}& 0\end{array}\right)H_1,`$ (41)
$`Y_D\left(\begin{array}{ccc}0& ϵ^{}& 0\\ ϵ^{}& ϵ& ϵ\\ 0& ϵ& 1\end{array}\right)\xi \overline{H}_3+\left(\begin{array}{ccc}0& ϵϵ^{}& ϵ^{}\\ ϵϵ^{}& ϵ^2& ϵ\\ ϵ^{}& ϵ& ϵ\end{array}\right)\xi \overline{H}_2+\left(\begin{array}{ccc}0& 0& 0\\ 0& ϵϵ^{}& ϵ^{}\\ 0& ϵ^{}& 0\end{array}\right)\xi \overline{H}_1,`$ (51)
where $`O(1)`$ coefficients have been omitted. Using the values $`ϵ0.02`$ and $`ϵ^{}0.004`$ obtained from fitting the quark masses and mixing angles , one can calculate the contribution to FCNC.
A comprehensive analysis of flavor-changing processes in models with tree-level FCNC was performed in Ref. and Ref. . In Ref. , it was argued that it is natural to take the flavor-changing coupling $`q_iq_j\varphi `$ to be the geometric mean of the $`q_i`$ and $`q_j`$ Yukawa couplings; this ansatz (referred to as “Model III”) has been used in a large number of analyses of FCNC. In Model III, the strongest bound comes from $`K^0\overline{K}^0`$ mixing; the mass of the exchanged scalar has to be greater than approximately a TeV. Much weaker bounds came from processes involving heavier quarks. In many of these analyses, there has been a greater focus on FCNC between the second and third generations; authors have generally assumed that the bound from $`K^0\overline{K}^0`$ mixing can be suppressed by some unspecified mechanism or a small amount of fine tuning. In the model of this Section the Yukawa couplings are different; the $`dsH_2`$ coupling is $`O(ϵϵ^{}\xi )`$ which is $`\sqrt{ϵ^{}}`$ times the geometric mean of the down and strange Yukawa couplings. This significantly weakens the bound; we find a bound on the $`H_2`$ mass of $`17`$ GeV from $`K^0\overline{K}^0`$ mixing. In the case of $`B^0\overline{B}^0`$ mixing we obtain a bound of $`100`$ GeV (with $`f_B200`$ MeV). The bound coming from $`D^0\overline{D}^0`$ mixing is of $`120`$ GeV. Bounds on processes involving leptons and b-quarks are much weaker. Since $`D^0\overline{D}^0`$ mixing is negligible in the standard model, the first signature of this model could come from $`D^0\overline{D}^0`$ mixing.
A similar analysis to the one described in the previous sections could now be made. In this case however, the “extra” Higgs bosons do acquire vevs and couple to fermions, as can be seen from Eqs. (41$``$51). We do not perform a detailed study of the phenomenology of the additional Higgs scalars.
## IV Conclusion
We study the supersymmetric standard model with more than one generation of Higgs doublets. We follow Ref. where a symmetry that forbids tree-level FCNC has been assumed. A result of this assumption is that the additional Higgs bosons do not mix with the standard Higgs bosons, do not couple to fermions, and that there is a mass degeneracy among the lightest neutral bosons. In this paper, it is noted that the charged scalar is heavier than the neutral scalar, for most of parameter space, by between a few hundred MeV and a few GeV, and that standard searches for heavy leptons with nearly degenerate neutrinos may be sensitive to these particles. Lastly, we relax the assumption that tree-level FCNC are eliminated, and discuss a flavor symmetry based on U(2) showing how the level of FCNC processes can be related to quark masses, and how the FCNC coupling between the first and second generations is automatically suppressed relative to the conventional ansatz, significantly weakening bounds from $`K^0\overline{K}^0`$ mixing
Acknowledgments
We thank Christopher D. Carone for useful conversations and for reading the manuscript. We thank the National Science Foundation for support under Grants Nos. PHY-9800741 and PHY-9900657
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# Interaction induced delocalisation for two particles in a periodic potential
\[
## Abstract
We consider two interacting particles evolving in a one-dimensional periodic structure embedded in a magnetic field. We show that the strong localization induced by the magnetic field for particular values of the flux per unit cell is destroyed as soon as the particles interact. We study the spectral and the dynamical aspects of this transition.
\]
As shown by Anderson in 1958$`^{\text{[1]}}`$, a quantum particle in a disordered potential may be trapped in spatially localized eigenstates. A natural question which arises is whether such a dramatic effect operates in a many-body interacting system. This issue has been strongly revived by a series of experiments on $`Si`$ MOSFETs$`^{\text{[2]}}`$ which seem to indicate a metal-insulator transition in two dimensions. Since then, a lot of experimental and theoretical activity has been dedicated to this highly controversial topic$`^{\text{[3]}}`$. Given the complexity of the full many-body systems, various groups have studied the already non trivial two-particle problem$`^{\text{[4, 5, 6]}}`$. In particular, D. L. Shepelyansky has shown convincingly that some two-particle eigenstates exhibit much larger localization length than single particle ones. Note however that in other situations such as in quasiperiodic systems, interactions may generate strongly localized two-particle eigenstates$`^{\text{[7]}}`$. In addition, it seems that for repulsive interactions these interesting effects do not appear in the vicinity of the two-particle ground state, and this leaves open questions for the many electron case.
Recently, it has been shown that an extreme localization mechanism induced by the magnetic field can lead to a complete confinement of the particle motion inside Aharonov-Bohm cages$`^{\text{[8]}}`$. Contrary to the Anderson localization, this phenomenon occurs in pure two-dimensional systems, i.e. without disorder, and is due to a subtle interplay between the structure geometry and the magnetic field. Two series of experiments have confirmed the existence of these Aharonov-Bohm cages. In the first one$`^{\text{[9]}}`$, superconducting wire networks with the adapted structure exhibit a striking reduction of the critical current for the predicted values of the magnetic field. The second experiments$`^{\text{[10]}}`$ measure the magneto-resistance oscillations in two-dimensional mesoscopic structures with a small number of conduction channels and large electronic mean free path. A clearcut dip at half a flux quantum per loop is observed and confirms the presence of this peculiar localization process.
It is therefore very natural to ask whether these cages survive for interacting particles. In this letter, we present an exacly solvable model for two interacting particles under a magnetic field. To simplify, we deal with a quasi-one-dimensional model which exhibits Aharonov-Bohm cages. We show that for half a flux quantum per loop, dispersive two-particle bound states appear even for repulsive local interaction. In this system, the two-particle ground state is non dispersive but the first dispersive band is rather close in energy. Slightly away from these remarkable fluxes, these bound states survive until they merge in a two-particle continuum. Finally, we are led to speculate that a finite repulsive local interaction is able to turn the fully localized non interacting system into a strongly correlated metal provided the electron density is large enough.
We consider a one-dimensional chain of square loops with periodic boundary conditions displayed in Fig. 1 embedded in a uniform perpendicular magnetic field $`𝐁`$, which is a bipartite periodic structure with three sites per unit cell. As we shall see, the various characteristics of this system are similar to those discussed in Ref. for two-dimensional tilings. Hereafter, we fix the total polarization to zero which is equivalent to consider two particles with opposite spin ($``$ and $``$).
Let us consider the standard Hubbard hamiltonian :
$$H=\underset{i,j,\sigma =,}{}t_{ij}c_{i,\sigma }^{}c_{j,\sigma }+U\underset{i}{}n_{i,}n_{i,}\text{,}$$
(1)
where $`c_{i,\sigma }^{}`$ (resp. $`c_{i,\sigma }`$) denotes the creation (resp. annihilation) operator of a fermion with spin $`\sigma `$, $`n_{i,\sigma }=c_{i,\sigma }^{}c_{i,\sigma }`$ the density of spin $`\sigma `$ fermion on site $`i`$, and $`\mathrm{}`$ stands for nearest neighbor pairs. Note that, since the particle considered here are fermions, the interaction term $`U`$ is only efficient in the singlet sector where the orbital part of the wave function is symmetric. When $`𝐁=0`$, the hopping term $`t_{ij}=1`$ if $`i`$ and $`j`$ are nearest neighbors and $`0`$ otherwise. In the presence of a magnetic field$`^{\text{[11]}}`$, $`t_{ij}`$ is multiplied by a phase factor $`e^{i\gamma _{ij}}`$ involving the vector potential $`𝐀`$ :
$$\gamma _{ij}=\frac{2\pi }{\varphi _0}_i^j𝐀.\text{d}𝐥\text{,}$$
(2)
where $`\varphi _0=hc/e`$ is the flux quantum. For convenience, we choose a gauge in which only one hopping term per unit cell is modified (see in Fig. 1). The whole spectrum only depends on the reduced flux $`f=\varphi /\varphi _0`$ where $`\varphi =Ba^2/2`$ is the magnetic flux through an elementary square ($`a`$ is the unit cell vector length).
Let us first analyze the one-particle problem. In this case, the translation invariance of the system along the chain direction allows one to straightforwardly compute the one-particle spectrum that consists of three bands :
$$\epsilon _\alpha (k)=2\alpha \sqrt{1+\mathrm{cos}(\gamma /2)\mathrm{cos}(ka)}k[0,2\pi /a]\text{,}$$
(3)
where $`\alpha =0,\pm 1`$ is the band index and $`\gamma =2\pi f`$. The weight of each band in the normalized density of states equals $`1/3`$. The existence of a non dispersive band at $`\epsilon =0`$ is simply due to the bipartite character of the structure, its degeneracy being equal to the difference between the number of sites of each family. The most striking feature is that for $`f=1/2`$ (half a flux quantum per unit square), the spectrum is made up of three non dispersive bands. As discussed in Ref. , this property leads to a complete lock-in of any wave packet spreading inside the so-called Aharonov-Bohm cages. One thus has a transition induced by the magnetic field.
For $`U=0`$, the two-particle spectrum is the addition of the one-particle spectra so that the eigenenergies are labelled by four quantum numbers :
$$\epsilon _{\alpha _{},\alpha _{}}(k_{},k_{})=\epsilon _\alpha _{}(k_{})+\epsilon _\alpha _{}(k_{})\text{,}$$
(4)
where $`\alpha _\sigma =0,\pm 1`$ (resp. $`k_\sigma `$) is the band index (resp. the wave vector) of the spin $`\sigma `$ particle. Thus, for $`f=1/2`$, the spectrum consists of five non dispersive bands corresponding to $`\epsilon =0,\pm 2,\pm 4`$ and the space evolution of any two-particle wave function is confined in an Aharonov-Bohm cage that is merely the superposition of each one-particle cage.
We now address the interacting case where $`U0`$. The main question is whether or not the latter system remains frozen when the particles are interacting. In other words, can a (local) interaction term destroy the cages and autorize any propagation ? In general, a two-particle problem with on-site interaction in a $`D`$-dimensional structure can be viewed as a single particle one in a $`2D`$-dimensional structure with a local potential in the hyperplane corresponding to a double occupancy of a site. Taking advantage of the translation invariance, this problem can then be mapped onto a (continuous) family of $`D`$-dimensional problems with a finite number of impurity sites which are often easier to solve$`^{\text{[12]}}`$. The same approach could in principle be used here but given the very special nature of the non interacting system for $`f=1/2`$, it is most easily carried out by using the minimally extended one-particle eigenstates displayed in Fig. 2.
These states have non vanishing amplitude only on a finite number of sites and thus reflect the absence of propagation at $`f=1/2`$. We denote them by $`|i,\epsilon _\alpha `$, where $`i`$ is the cell number centered on each 4-fold coordinated site and $`\epsilon _\alpha =0,\pm 2`$.
The two-particle state space is then given by the tensor product of the one-particle state :
$$|i,\alpha ;i^{},\alpha ^{}=|i,\epsilon _\alpha _{}|i^{},\epsilon _\alpha ^{}_{}\text{,}$$
(5)
for all $`i,i^{},\alpha ,\alpha ^{}`$. It is worth stressing that since the $`|i,\epsilon _\alpha `$ are confined eigenstates of the one-particle hamiltonian, most of these two-particle eigenstates are not affected by $`U`$. This local interaction term only acts on states for which the two particles have a non vanishing probability to be on the same site, i.e. such that $`|ii^{}|1`$. As a result, the number of states which are sensitive to $`U`$ scales linearly with the total number of cells $`N`$ whereas in the generic case ($`f1/2`$), it is proportional to $`N^2`$. To proceed further, let us remark that the states sensitive to $`U`$ are space-symmetric :
$$|i,\alpha ;i^{},\alpha ^{}_S=\frac{1}{\sqrt{2}}\left(|i,\alpha ;i^{},\alpha ^{}+|i^{},\alpha ^{};i,\alpha \right)\text{,}$$
(6)
if $`ii^{}`$ or $`\alpha \alpha ^{}`$, and $`|i,\alpha ;i,\alpha _S=|i,\alpha ;i,\alpha `$. Moreover, since the problem is invariant under a translation of the center of mass, it is convenient to build Bloch waves :
$`|\phi _0(\alpha ,\alpha ^{},K)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{N}}}{\displaystyle \underset{n=0}{\overset{N1}{}}}e^{iKna}|n,\alpha ;n,\alpha ^{}_S`$ (7)
$`|\phi _1(\alpha ,\alpha ^{},K)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{N}}}{\displaystyle \underset{n=0}{\overset{N1}{}}}e^{iKna}|n,\alpha ;n+1,\alpha ^{}_S\text{,}`$ (8)
where $`K=2\pi j/Na(j=0,N1)`$ is a wave vector lying in the first Brillouin zone. In the following, we shall always considerthe limit where $`N`$ tends to infinity. One then has to calculate the matrix element of the hamiltonian in each irreducible representation labelled by $`K`$, and restrict this analysis to the $`(15\times 15)`$ subspace generated by $`|\phi _0`$ and $`|\phi _1`$. For any $`(\alpha ,\alpha ^{},\beta ,\beta ^{},K_l,K_m)`$ one has :
$`\phi _0(\alpha ,\alpha ^{},K_l)|H|\phi _1(\beta ,\beta ^{},K_m)`$ $`=`$ $`0`$ (9)
$`\phi _1(\alpha ,\alpha ^{},K_l)|H|\phi _1(\beta ,\beta ^{},K_m)`$ $`=`$ $`\lambda _{\alpha ,\alpha ^{},\beta ,\beta ^{}}\delta _{l,m}\text{,}`$ (11)
where $`\lambda _{\alpha ,\alpha ^{},\beta ,\beta ^{}}`$ is a $`K_{l,m}`$-independent scalar and where $`\delta _{l,m}`$ is the usual Kronecker symbol. This implies that the eigenvalues of the $`(9\times 9)`$ subspace generated by the $`|\phi _1`$ are non dispersive. These eigenvalues are given by the five $`U`$-dependent roots of the characteristic polynomial :
$$P(\epsilon ,U)=\epsilon ^5U\epsilon ^420\epsilon ^3+16U\epsilon ^2+64\epsilon 24U\text{,}$$
(12)
and four $`U`$-independent values $`\pm 2,0`$ (two-fold degenerated) resulting of additional symmetry between $`\alpha `$ and $`\alpha ^{}`$. The non dispersive part of the spectrum$`^{\text{[13]}}`$ is shown in Fig. 3 for $`U0`$.
We emphasize that since the spectrum of the Hubbard hamiltonian (in a bipartite structure) is odd under the transformation $`UU`$$`^{\text{[14]}}`$, we can restrict our analysis to the repulsive case. In the large $`U`$ limit, the five roots of $`P`$ tends toward $`\pm \sqrt{8\pm 2\sqrt{10}},U`$ ; this latter value simply corresponding to a situation where the two particles are localized on the same site (anti-bonding state).
A much more interesting component of the spectrum is provided in the $`(6\times 6)`$ subspace$`^{\text{[15]}}`$ generated by the $`|\phi _0`$. In this subspace, the eigenvalues are given by the roots of the following characteristic polynomial :
$`Q(\epsilon ,U,K)`$ $`=`$ $`\epsilon ^62U\epsilon ^5+(U^220)\epsilon ^4+28U\epsilon ^3+`$ (15)
$`8\left(8U^2\right)\epsilon ^24U(143\mathrm{cos}(Ka))\epsilon +`$
$`4U^2(2+\mathrm{cos}(Ka))\text{.}`$
Contrary to the previous case, these eigenvalues obviously depend on $`K`$, and the associated non degenerated eigenstates are extended (Bloch-like). This dispersive component of the spectrum is displayed in Fig. 4.
In the large $`U`$ limit, the asymptotic eigenvalues are given by $`U`$ (twice degenerated) and $`\pm \sqrt{4\pm 2\sqrt{2\mathrm{cos}(Ka)}}`$, and there still remains a $`K`$-dependent component in the spectrum.
The physical consequences of this dispersion are important. Indeed, let us consider a generic two-particle initial state having a non zero overlap with one of the $`|\phi _0`$$`^{\text{[16]}}`$. The emergence of dispersive states for $`U0`$ indicates that it is now possible for this wave packet to spread over the whole system whereas it was completely trapped inside the Aharonov-Bohm cage in the non interacting case. Moreover, since the dispersive eigenstates are extended, the propagation is ballistic.
For other values of the reduced flux, the full solution of the two-particle problem cannot be cast in a simple analytic form, although it is possible to reduce it into a scattering problem for one particle moving on a chain in the presence of three static impurities. Nevertheless, let us analyze the neighborhood of the half-flux parametrized by $`\delta f=|f1/2|`$. For $`f1/2`$, all the single-particle eigenstates become extended, except those corresponding to the flat band at $`\epsilon =0`$. The non dispersive two-particle states which are insensitive to $`U`$ at $`f=1/2`$ evolve into several two-particle continua, whose band width scales as $`\delta f`$ for small $`\delta f`$. The non dispersive states which are sensitive to $`U`$ become dispersive with a band width scaling as $`(\delta f)^2/U`$ (see inset in Fig. 5). The corresponding wave functions still exhibit a binding of the two particles. Finally, as displayed in Fig. 5, the dispersive states at $`f=1/2`$ remain dispersive. All these dispersive bound states evolve smoothly as a function of $`\delta f`$ until their energies merge in the two-particle continuum, which occurs for $`\delta fU`$ at small $`U`$. As $`\delta f`$ further increases, the total number of bound states gradually decreases.
Since, in this study, the effect of the interaction term between the two particles is clearly to induce a delocalization process, one can expect that for finite density of particles, a subtle correlated conducting state will emerge. Note however that in this system, the energy of two particles is minimal when they do not interact, either because they are far apart or because their orbital wave function is antisymmetric. We therefore conjecture that the many body ground state will remain localized up to a filling factor equal to 1/3. At this point, the lowest flat band is completely filled with a fully polarized electron sea. The next electron will have an opposite spin and is thus likely to delocalize along the chain. This problem obviously deserves futher investigations. Let us also remark that, here, we consider a quasi-one-dimensional model to simplify the calculations, but a similar physics is expected for other tight-binding models which exhibit single-particle confinement inside Aharonov-Bohm cages. Nevertheless, in more general systems, we do not expect to see non dispersive states sensitive to $`U`$. In this context, an interesting open question is whether it is possible to find situations where the ground state of the two-particle spectrum is dispersive, even for repulsive interaction.
How could this interaction induced delocalization be observed experimentally? The two-particle system might be accesible in a Josephson junction or a quantum dot array in the Coulomb blockade regime. Such experimental situations shall be described by a model with on-site disorder, but, very likely, some two-particle states would still exhibit a much larger localization length than single particle ones. Ongoing experiments on ballistic semi-conducting networks with special two-dimensional geometry$`^{\text{[10]}}`$ may also manifest some interaction effects on this Aharonov-Bohm localization. In these structures, the interaction strength can be varied by changing the electronic density or by polarizing the system with a tilted magnetic field. Clearly, a better understanding of the many electron case needs to be achieved.
We would like to thank Cl. Aslangul, B. Delamotte, D. Mouhanna and M. Tissier for fruitful discussions.
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# 1 Introduction
## 1 Introduction
For more than two decades superstring theories have been studied as the most promising candidates for a unified theory of all the interactions including gravity. The theories have the potential to predict the space-time dimensionality, the gauge group, the matter content, and so on, from first principles. The existence of infinitely many perturbative vacua implies, however, that an understanding of nonperturbative effects is crucial to extract information about the real vacuum of the theory. Recent proposals for nonperturbative formulations of superstring theories may therefore be of analogous importance for our understanding of non-perturbative aspects of string theory as lattice gauge theory has been in understanding nonperturbative dynamics of gauge theories. The IIB matrix model , which is conjectured to be a nonperturbative formulation of type IIB superstring theory (for a review, see Ref. ), takes the form of a large $`N`$ reduced model , and in the same way that Monte Carlo studies of lattice gauge theory clarified many important nonperturbative aspects of the strong interaction, Monte Carlo studies of the IIB matrix model might illuminate the nonperturbative dynamics of superstring theories<sup>1</sup><sup>1</sup>1Monte Carlo studies of Matrix Theory would be technically more involved due to a lattice discretization of the time direction .. A number of numerical studies have already been done to pursue that direction. Earlier numerical studies of world-sheet perturbative aspects of superstring-like theories can be found in Refs. .
In this paper, we study the IIB matrix model at large $`N`$ by Monte Carlo simulations. In particular, we investigate the possibility of a spontaneous symmetry breakdown of Lorentz invariance using the low-energy effective theory of the IIB matrix model proposed by Ref. . The authors of also study this issue, using an approach where the bosonic and fermionic matrices are both decomposed into diagonal elements and off-diagonal elements, and the off-diagonal elements are integrated out first perturbatively. Such a perturbative expansion is valid when the bosonic diagonal elements, which can be regarded as coordinates of $`N`$ points in ten-dimensional flat space-time, are well separated from each other. In other words, what one obtains after integrating over the off-diagonal elements perturbatively can be considered as a low-energy effective theory of the IIB matrix model. Note that this perturbative expansion has nothing to do with the string perturbative expansion with respect to the world-sheet topology. Therefore, even at the one-loop level, the low-energy effective theory is expected to include nonperturbative effects of superstring theory, provided, of course, that the IIB matrix model conjecture is true.
In fact, in order to obtain the low-energy effective action only for the bosonic diagonal elements, one still has to perform the integration over the fermionic diagonal elements, which is nontrivial, since their action turns out to be quartic. Although the explicit form of the final low-energy effective theory has not been derived, the theory was shown to be described by some complicated branched-polymer like system, typically involving a “double-tree” structure, in a flat ten-dimensional space time. Thus, even at the one-loop level<sup>2</sup><sup>2</sup>2The validity of the one-loop approximation for studying the low-energy dynamics of supersymmetric large $`N`$ reduced models has been demonstrated in Ref. through the study of the four-dimensional version of the IIB matrix model. We will come back to this point is Section 3., the low-energy effective theory contains highly nontrivial dynamics. It was further argued that the double-tree structure of the one-loop low-energy effective theory might be responsible for a collapse of the distribution of the bosonic diagonal elements. The first Monte Carlo results of a branched-polymer system with a double-tree structure was reported in Ref. .
Here, we write down explicitly the low-energy effective theory of the IIB matrix model in terms of bosonic variables only. Instead of integrating over both bosonic and fermionic off-diagonal elements first, we leave the bosonic off-diagonal elements unintegrated. The action for the fermionic diagonal elements is then still quadratic and can be integrated explicitly, yielding a Pfaffian. Integration over the bosonic off-diagonal elements as well as the bosonic diagonal elements can be done by Monte Carlo simulation. In other words, the bosonic off-diagonal elements play the rôle of auxiliary variables, which enable us to simulate the complicated branched-polymer like system describing the dynamics of the bosonic diagonal elements.
The Pfaffian induced by the integration over the fermionic diagonal elements is generically complex. In general, when the action of a theory has a non-zero imaginary part, the number of configurations needed to extract any information increases exponentially with the system size, except in a few situations, where alternative sampling methods can be invented . This notorious technical problem in Monte Carlo simulations is known as the “sign problem”. In fact, the problem exists already in Monte Carlo simulations of the original IIB matrix model and it is inherited by the low-energy effective theory. In the present work, we simply use the absolute value of the Pfaffian in order to avoid the sign problem and examine only the effect of the modulus of the Pfaffian. Our results suggest that there is no spontaneous symmetry breakdown (SSB) of Lorentz invariance. From this we conclude that if the SSB ever occurs in the IIB matrix model, the phase of the Pfaffian must play a crucial rôle. We also study the six-dimensional version of the IIB matrix model for comparison. The conclusion is the same, but we find an intriguing difference in the finite $`N`$ effects.
The paper is organized as follows. In Section 2, we describe the definition of the IIB matrix model and review some important properties relevant for this work. In Section 3, we derive the low-energy effective theory of the IIB matrix model and explain the model we investigate by Monte Carlo simulations. In Section 4, we present our results for the distribution of the bosonic diagonal elements. In particular, we discuss the possibility of SSB of Lorentz invariance. Section 5 is devoted to a summary and conclusions. In Appendix A, we explain the details of the algorithm we use for the Monte Carlo simulation. In Appendix B, we present some systematic studies for optimizing the parameters involved in the algorithm.
## 2 The IIB matrix model
The IIB matrix model is formally a zero-volume limit of ten-dimensional pure $`𝒩=1`$ supersymmetric Yang-Mills theory. The action, therefore, is given by
$`Z_{\mathrm{IIB}}`$ $`=`$ $`{\displaystyle \text{d}A\text{e}^{S_\mathrm{b}}Z_\mathrm{f}[A]};Z_\mathrm{f}[A]={\displaystyle \text{d}\psi \text{e}^{S_\mathrm{f}}},`$ (2.1)
$`S_\mathrm{b}`$ $`=`$ $`{\displaystyle \frac{1}{4g^2}}\text{tr}([A_\mu ,A_\nu ]^2),`$ (2.2)
$`S_\mathrm{f}`$ $`=`$ $`{\displaystyle \frac{1}{2g^2}}\text{tr}\left(\psi _\alpha (\stackrel{~}{\mathrm{\Gamma }}_\mu )_{\alpha \beta }[A_\mu ,\psi _\beta ]\right).`$ (2.3)
$`A_\mu `$ ($`\mu =1,\mathrm{},10`$) and $`\psi _\alpha `$ ($`\alpha =1,\mathrm{},16`$) are $`N\times N`$ traceless Hermitian matrices, which can be expanded in terms of the generators $`t^a`$ of SU($`N`$) as
$$(A_\mu )_{ij}=\underset{a=1}{\overset{N^21}{}}A_\mu ^a(t^a)_{ij};(\psi _\alpha )_{ij}=\underset{a=1}{\overset{N^21}{}}\psi _\alpha ^a(t^a)_{ij},$$
(2.4)
where $`A_\mu ^a`$ is a real variable and $`\psi _\alpha ^a`$ is a real Grassmann variable. We assume that the generators $`t^a`$ are normalized as $`\text{tr}(t^at^b)=\delta _{ab}`$. The measure $`\text{d}\psi `$ in (2.1) is defined by
$`\text{d}\psi `$ $`=`$ $`{\displaystyle \underset{\alpha =1}{\overset{16}{}}}{\displaystyle \underset{a=1}{\overset{N^21}{}}}\text{d}\psi _\alpha ^a`$ (2.5)
$`=`$ $`{\displaystyle \underset{\alpha =1}{\overset{16}{}}}\left[{\displaystyle \underset{i<j}{}}\{2\text{d}\text{Re}(\psi _\alpha )_{ij}\text{d}\text{Im}(\psi _\alpha )_{ij}\}{\displaystyle \underset{i=1}{\overset{N}{}}}\{\text{d}(\psi _\alpha )_{ii}\}\delta \left({\displaystyle \frac{1}{\sqrt{N}}}{\displaystyle \underset{i=1}{\overset{N}{}}}(\psi _\alpha )_{ii}\right)\right],`$
and similarly for $`\text{d}A`$. The model (2.1) appears after a Wick rotation, so that the metric has Euclidean signature. The $`16\times 16`$ matrices $`\stackrel{~}{\mathrm{\Gamma }}_\mu `$ are defined by
$$\stackrel{~}{\mathrm{\Gamma }}_\mu =𝒞\mathrm{\Gamma }_\mu ,$$
(2.6)
where $`\mathrm{\Gamma }_\mu `$ are ten-dimensional gamma matrices after Weyl projection, and the unitary matrix $`𝒞`$ is a charge conjugation matrix satisfying
$$𝒞\mathrm{\Gamma }_\mu 𝒞^{}=(\mathrm{\Gamma }_\mu )^{};𝒞^{}=𝒞.$$
(2.7)
Due to (2.7), the matrices $`\stackrel{~}{\mathrm{\Gamma }}_\mu `$ are symmetric.
An explicit representation of the gamma matrices is given by
$`\mathrm{\Gamma }_1=i\sigma _2\sigma _2\sigma _2\sigma _2;\mathrm{\Gamma }_2=i\sigma _2\sigma _2\mathrm{𝟏}\sigma _1;\mathrm{\Gamma }_3=i\sigma _2\sigma _2\mathrm{𝟏}\sigma _3;`$ (2.8)
$`\mathrm{\Gamma }_4=i\sigma _2\sigma _1\sigma _2\mathrm{𝟏};\mathrm{\Gamma }_5=i\sigma _2\sigma _3\sigma _2\mathrm{𝟏};\mathrm{\Gamma }_6=i\sigma _2\mathrm{𝟏}\sigma _1\sigma _2;`$
$`\mathrm{\Gamma }_7=i\sigma _2\mathrm{𝟏}\sigma _3\sigma _2;\mathrm{\Gamma }_8=i\sigma _1\mathrm{𝟏}\mathrm{𝟏}\mathrm{𝟏};\mathrm{\Gamma }_9=i\sigma _3\mathrm{𝟏}\mathrm{𝟏}\mathrm{𝟏};`$
$`\mathrm{\Gamma }_{10}=\mathrm{𝟏}\mathrm{𝟏}\mathrm{𝟏}\mathrm{𝟏},`$
for which the charge conjugation matrix $`𝒞`$ becomes a unit matrix and therefore $`\stackrel{~}{\mathrm{\Gamma }}_\mu =\mathrm{\Gamma }_\mu `$.
The model has a manifest ten-dimensional Lorentz invariance, by which we actually mean an SO(10) invariance. $`A_\mu `$ transforms as a vector and $`\psi _\alpha `$ transforms as a Majorana-Weyl spinor. The model is manifestly supersymmetric, and also has a SU($`N`$) symmetry
$$A_\mu VA_\mu V^{};\psi _\alpha V\psi _\alpha V^{},$$
(2.9)
where $`V\text{SU}(N)`$. All these symmetries are inherited from the super Yang-Mills theory before taking the zero-volume limit. In particular, the SU($`N`$) symmetry (2.9) is a remnant of the local gauge symmetry.
The fermion integral $`Z_\mathrm{f}[A]`$ in (2.1) can be obtained explicitly as
$$Z_\mathrm{f}[A]=\text{Pf},$$
(2.10)
where
$`_{a\alpha ,b\beta }`$ $`=`$ $`i{\displaystyle \frac{1}{g^2}}f_{abc}(\stackrel{~}{\mathrm{\Gamma }}_\mu )_{\alpha \beta }A_\mu ^c,`$
$`f_{abc}`$ $`=`$ $`i\text{tr}(t^a[t^b,t^c]).`$ (2.11)
The real totally-antisymmetric tensor $`f_{abc}`$ gives the structure constants of SU($`N`$) and the matrix $`_{a\alpha ,b\beta }`$ is a $`16(N^21)\times 16(N^21)`$ anti-symmetric matrix, regarding each of $`(a\alpha )`$ and $`(b\beta )`$ as a single index.
The convergence of the integration over the bosonic matrices in (2.1) is nontrivial since the integration domain for the Hermitian matrices is non-compact. Even for finite $`N`$ there is a potential danger of divergence when the eigenvalues of $`A_\mu `$ become large. This issue has been addressed in Ref. using one-loop perturbative arguments which pointed to the finiteness of the IIB matrix model given by (2.1) for arbitrary $`N`$. This conclusion is in agreement with an exact result available for $`N=2`$ and a numerical result obtained for $`N=3`$ . Thus it is conceivable that the above conclusion, obtained by one-loop arguments, holds in general.
Since the model is well-defined without any cutoff, the parameter $`g`$, which is the only parameter of the model, can be absorbed by rescaling the variables,
$`A_\mu `$ $`=`$ $`g^{1/2}X_\mu ,`$ (2.12)
$`\psi _\alpha `$ $`=`$ $`g^{3/4}\mathrm{\Psi }_\alpha .`$ (2.13)
Therefore, $`g`$ is a scale parameter rather than a coupling constant, i.e. the $`g`$ dependence of physical quantities is completely determined on dimensional grounds<sup>3</sup><sup>3</sup>3The scale parameter $`g`$ should be tuned appropriately as one sends $`N`$ to infinity so that all the correlation functions of Wilson loops have a finite large $`N`$ limit. Whether such a limit really exists or not is one of the important dynamical issues, which was addressed in Ref. for the four-dimensional version of the IIB matrix model.. In what follows, we take $`g=1`$ without loss of generality.
For comparison, we also study the six-dimensional version of the IIB matrix model. In this regard we recall that pure $`𝒩=1`$ supersymmetric Yang-Mills theory can be also defined in 3D, 4D and 6D, as well as in 10D. Hence, by taking a zero-volume limit of these theories, we arrive at supersymmetric large $`N`$ reduced models which are $`D=3,4,6`$ versions of the IIB matrix model. Using the one-loop argument mentioned above, one concludes that the model is ill-defined for $`D=3`$, but well-defined for $`D=4,6,10`$, irrespectively of $`N`$. For $`D=4`$, Monte Carlo simulations up to $`N=48`$ further confirms this statement . The effective action induced by fermions (logarithm of the fermion integral $`Z_\mathrm{f}[A]`$) is real for $`D=4`$. It is complex in general for $`D=6`$ and $`D=10`$, however, which causes the sign problem in Monte Carlo simulations.
## 3 Low-energy effective theory
In this section, we derive the low-energy effective theory of the IIB matrix model (2.1) along the lines discussed in Ref. . We first decompose the $`N\times N`$ Hermitian matrices $`A_\mu `$ and $`\psi _\alpha `$ as
$`A_\mu `$ $`=`$ $`\widehat{x}_\mu +\widehat{a}_\mu ,`$
$`\psi _\alpha `$ $`=`$ $`\widehat{\xi }_\alpha +\widehat{\phi }_\alpha ,`$ (3.1)
where $`\widehat{x}_\mu `$ and $`\widehat{\xi }_\alpha `$ represent the diagonal parts, while $`\widehat{a}_\mu `$ and $`\widehat{\phi }_\alpha `$ represent the off-diagonal parts. We also introduce $`N`$ ten-dimensional vectors $`x_i`$ ($`i=1,\mathrm{},N`$) through $`x_{i\mu }=(\widehat{x}_\mu )_{ii}`$, and $`N`$ ten-dimensional Majorana-Weyl spinors $`\xi _i`$ ($`i=1,\mathrm{},N`$) through $`\xi _{i\alpha }=(\widehat{\xi }_\alpha )_{ii}`$. For the off-diagonal elements, we use the notations $`a_{\mu ij}=(\widehat{a}_\mu )_{ij}`$ and $`\phi _{\alpha ij}=(\widehat{\phi }_\alpha )_{ij}`$, where $`ij`$. Using the decomposition (3.1), the actions (2.2) and (2.3) can be written as
$`S_\mathrm{b}`$ $`=`$ $`\text{tr}\left({\displaystyle \frac{1}{2}}\widehat{a}_\nu [\widehat{x}_\mu ,[\widehat{x}_\mu ,\widehat{a}_\nu ]]{\displaystyle \frac{1}{2}}[\widehat{x}_\mu ,\widehat{a}_\mu ]^2+[\widehat{x}_\mu ,\widehat{a}_\nu ][\widehat{a}_\mu ,\widehat{a}_\nu ]+{\displaystyle \frac{1}{4}}[\widehat{a}_\mu ,\widehat{a}_\nu ]^2\right),`$ (3.2)
$`S_\mathrm{f}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\stackrel{~}{\mathrm{\Gamma }}_\mu )_{\alpha \beta }\text{tr}\left(\widehat{\phi }_\alpha [\widehat{x}_\mu ,\widehat{\phi }_\beta ]\widehat{\phi }_\alpha [\widehat{\xi }_\beta ,\widehat{a}_\mu ]\widehat{a}_\mu [\widehat{\xi }_\alpha ,\widehat{\phi }_\beta ]+\widehat{\phi }_\alpha [\widehat{a}_\mu ,\widehat{\phi }_\beta ]\right).`$ (3.3)
The one-loop approximation amounts to keeping only the quadratic terms in $`a`$ and $`\phi `$ in the above expressions, neglecting the higher-order terms, i.e., the O($`a^3`$) term and the O($`a^4`$) term in $`S_\mathrm{b}`$ and the O($`a\phi ^2`$) term in $`S_\mathrm{f}`$. The quadratic term in $`a`$ in (3.2) has zero modes due to the fact that the original model (2.1) has the SU($`N`$) invariance (2.9). We thus have to “fix the gauge” properly. Following Ref. , we choose the “gauge-fixing” term (which is the reduced model counterpart of the Feynman gauge in ordinary gauge theory) and the corresponding Faddeev-Popov ghost term as
$`S_{\mathrm{g}.\mathrm{f}.}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\text{tr}([\widehat{x}_\mu ,\widehat{a}_\mu ]^2),`$
$`S_{\mathrm{gh}}`$ $`=`$ $`\text{tr}([\widehat{x}_\mu ,\widehat{b}][\widehat{x}_\mu +\widehat{a}_\mu ,\widehat{c}]).`$ (3.4)
$`\widehat{a}_\mu `$ in the ghost action $`S_{\mathrm{gh}}`$ can be neglected within the one-loop approximation. Integration over the ghost fields $`\widehat{b}`$, $`\widehat{c}`$ then yields $`\{\mathrm{\Delta }(x)\}^2`$, where $`\mathrm{\Delta }(x)`$ is defined as
$$\mathrm{\Delta }(x)=\underset{i<j}{}(x_ix_j)^2.$$
(3.5)
Therefore, the partition function at the one-loop approximation can be written as
$$Z_{\mathrm{IIB}}^{(1\mathrm{loop})}=\text{d}x\text{d}a\text{e}^{S_\mathrm{G}}\{\mathrm{\Delta }(x)\}^2Z_\mathrm{f}^{(1\mathrm{loop})}[x,a],$$
(3.6)
where
$$S_\mathrm{G}=\underset{i<j}{}(x_ix_j)^2|a_{\mu ij}|^2,$$
(3.7)
and the one-loop approximated fermion integral is defined as
$`Z_\mathrm{f}^{(1\mathrm{loop})}[x,a]`$ $`=`$ $`{\displaystyle \text{d}\xi \text{d}\phi \text{e}^{S_\mathrm{f}^{(2)}}},`$ (3.8)
$`S_\mathrm{f}^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\stackrel{~}{\mathrm{\Gamma }}_\mu )_{\alpha \beta }\text{tr}\left(\widehat{\phi }_\alpha [\widehat{x}_\mu ,\widehat{\phi }_\beta ]\widehat{\phi }_\alpha [\widehat{\xi }_\beta ,\widehat{a}_\mu ]\widehat{a}_\mu [\widehat{\xi }_\alpha ,\widehat{\phi }_\beta ]\right).`$ (3.9)
In what follows we calculate $`Z_\mathrm{f}^{(1\mathrm{loop})}[x,a]`$ explicitly. We first note that $`S_\mathrm{f}^{(2)}`$ can be written as
$$S_\mathrm{f}^{(2)}=\frac{1}{2}\underset{ij}{}\stackrel{~}{\phi }_{\alpha ji}(\stackrel{~}{\mathrm{\Gamma }}_\mu )_{\alpha \beta }(x_{i\mu }x_{j\mu })\stackrel{~}{\phi }_{\beta ij}\frac{1}{2}\underset{ij}{}\xi _{i\alpha }M_{i\alpha ,j\beta }^{}\xi _{j\beta },$$
(3.10)
where $`\stackrel{~}{\phi }_{\alpha ij}`$ in (3.10) are defined by
$$\stackrel{~}{\phi }_{\alpha ij}=\phi _{\alpha ij}\frac{(x_{i\rho }x_{j\rho })}{(x_ix_j)^2}(\stackrel{~}{\mathrm{\Gamma }}_\rho ^{}\stackrel{~}{\mathrm{\Gamma }}_\sigma )_{\alpha \beta }a_{\sigma ij}(\xi _{i\beta }\xi _{j\beta }),$$
(3.11)
and $`M_{i\alpha ,j\beta }^{}`$ is a $`16N`$ $`\times `$ $`16N`$ matrix given as
$`M_{i\alpha ,j\beta }^{}`$ $`=`$ $`{\displaystyle \frac{(x_{i\rho }x_{j\rho })}{(x_ix_j)^2}}(\stackrel{~}{\mathrm{\Gamma }}_\mu \stackrel{~}{\mathrm{\Gamma }}_\rho ^{}\stackrel{~}{\mathrm{\Gamma }}_\sigma )_{\alpha \beta }(a_{\mu ji}a_{\sigma ij}a_{\mu ij}a_{\sigma ji})\text{for}ij,`$ (3.12)
$`M_{i\alpha ,i\beta }^{}`$ $`=`$ $`{\displaystyle \underset{ji}{}}M_{i\alpha ,j\beta }^{}.`$ (3.13)
Integration over $`\stackrel{~}{\phi }`$ can now be done yielding $`\{\mathrm{\Delta }(x)\}^8`$. We then integrate out $`\xi _{N\alpha }`$ using the delta functions in (2.5), yielding a factor of $`1/N^8`$ followed by a replacement
$$\xi _{N\alpha }\underset{j=1}{\overset{N1}{}}\xi _{j\alpha }$$
(3.14)
in (3.10). The integration over the $`\xi _{i\alpha }`$ ($`i=1,\mathrm{},(N1)`$) yields $`\text{Pf}M`$. The matrix $`M`$ is a $`16(N1)`$ $`\times `$ $`16(N1)`$ complex matrix defined as
$$M_{i\alpha ,j\beta }=M_{i\alpha ,j\beta }^{}M_{N\alpha ,j\beta }^{}M_{i\alpha ,N\beta }^{}+M_{N\alpha ,N\beta }^{},$$
(3.15)
where indices $`i`$ and $`j`$ run from 1 to $`N1`$. Note that there are identities
$$M_{j\alpha ,i\beta }^{}=M_{i\alpha ,j\beta }^{},M_{i\beta ,j\alpha }^{}=M_{i\alpha ,j\beta }^{},$$
(3.16)
and similarly for $`M_{i\alpha ,j\beta }`$. This means in particular that $`M_{i\alpha ,j\beta }`$ is an anti-symmetric matrix, regarding each of $`(i\alpha )`$ and $`(j\beta )`$ as a single index. Thus the one-loop approximated fermion integral is obtained as
$$Z_\mathrm{f}^{(1\mathrm{loop})}[x,a]=\{\mathrm{\Delta }(x)\}^8\frac{1}{N^8}\text{Pf}M.$$
(3.17)
We have checked numerically that indeed
$$Z_\mathrm{f}[A]Z_\mathrm{f}^{(1\mathrm{loop})}[x,a];(A_\mu )_{ij}=x_{i\mu }\delta _{ij}+a_{\mu ij},$$
(3.18)
holds when the $`x_i`$’s are well-separated and $`a_{\mu ij}`$ are generated with the distribution $`\text{e}^{S_\mathrm{G}}`$, where $`S_\mathrm{G}`$ is given by (3.7). Note that the size of the matrix $`M`$ in (3.17) is of O($`N`$), whereas the size of the matrix $``$ in (2.10) is of O($`N^2`$). The huge reduction is essentially because the integration over the fermionic off-diagonal elements has been done explicitly yielding $`\{\mathrm{\Delta }(x)\}^8`$ in (3.17). This is the substantial gain from using the one-loop approximation.
First we note that the model (3.6) as it stands has a singularity<sup>4</sup><sup>4</sup>4This can be seen more clearly by rescaling $`a_{\mu ij}`$ as in (A.1). Then, one finds that all the $`x_{i\mu }`$-dependence of the partition function is contained in the Pfaffian, which has the singularity. at $`x_{i\mu }=0`$ due to the singularity in (3.12). Therefore if one simulates the model (3.6), the distribution of $`x_i`$ collapses to the origin. We recall that such an ultraviolet singularity does not exist in the original IIB matrix model. On the other hand, the one-loop approximation is valid only when $`x_i`$’s are well separated from each other. Therefore, we have to introduce a UV cutoff to the distribution of $`x_i`$ in order to make the model meaningful. As a UV cutoff, we introduce in the action the term given as
$$S_{\mathrm{cut}}=\underset{i<j}{}f(\sqrt{(x_ix_j)^2}),$$
(3.19)
where the function $`f(x)`$ is taken to be
$$f(r)=\{\begin{array}{cc}\frac{\kappa }{2\mathrm{}^2}(r\mathrm{})^2\hfill & \text{for}r<\mathrm{}\hfill \\ 0\hfill & \text{for}r\mathrm{}.\hfill \end{array}$$
(3.20)
The dimensionless “spring constant” $`\kappa `$ should be taken to be large enough to prevent the $`x_i`$’s from coming closer to each other than the cutoff $`\mathrm{}`$ (See Figs. 1 and 2.). Thus we arrive at the low-energy effective theory of the IIB matrix model (2.1)
$$Z_{\mathrm{LEET}}=\text{d}x\text{d}a\text{e}^{S_\mathrm{G}S_{\mathrm{cut}}}\{\mathrm{\Delta }(x)\}^{10}\text{Pf}M,$$
(3.21)
which is written in terms of bosonic variables only. We have omitted the irrelevant constant factor $`\frac{1}{N^8}`$ in (3.17).
In fact, the UV cutoff parameter $`\mathrm{}`$ can be scaled away from the theory (3.21) by rescaling the variables as $`x_{i\mu }\mathrm{}x_{i\mu }`$ and $`a_{\mu ij}\frac{1}{\mathrm{}}a_{\mu ij}`$. This means that the dependence of the results on the UV cutoff $`\mathrm{}`$ is determined completely on dimensional grounds. In particular, dimensionless quantities are independent of $`\mathrm{}`$. When we are interested in dimensionful quantities, and in particular in their large $`N`$ behavior, we have to know the $`N`$-dependence of the UV cutoff $`\mathrm{}`$. In Ref. , it was argued that $`\mathrm{}`$ should be taken to be $`N`$-independent, based on a reasonable assumption that the ultraviolet behavior of the space-time structure of the IIB matrix model is controlled by the SU(2) matrix model. The issue is also addressed in the four-dimensional version of the IIB matrix model in Ref. . There, an $`N`$-independent UV cutoff was shown to be generated dynamically by treating the full model nonperturbatively instead of making perturbative expansions around diagonal matrices. We therefore take $`\mathrm{}=1`$ for all $`N`$ in the following.
The Pfaffian $`\text{Pf}M`$ in (3.21) is complex in general. This poses the notorious “sign problem”, when one tries to study the model by Monte Carlo simulation. Note, however, that the problem is simply inherited from the original IIB matrix model (2.1), as can be seen from (3.18). In the present work, we take the absolute value $`|\text{Pf}M|`$ and throw away the phase by hand. This corresponds to taking the absolute value $`|\text{Pf}|`$ of the Pfaffian in the IIB matrix model and studying its low-energy effective theory.
To summarize, the model we simulate is given by the partition function
$$Z=\text{d}x\text{d}a\text{e}^{S_\mathrm{G}S_{\mathrm{cut}}}\{\mathrm{\Delta }(x)\}^{10}|\text{Pf}M|.$$
(3.22)
$`S_\mathrm{G}`$ and $`S_{\mathrm{cut}}`$ are given by (3.7) and (3.19) respectively, and $`M`$ is defined through (3.12), (3.13) and (3.15). Note that the model (3.22) still has the 10D Lorentz invariance. We would like to investigate whether the 10D Lorentz invariance breaks down, say, to 4D Lorentz invariance. The model corresponding to the six-dimensional version of the IIB matrix model can be obtained similarly.
The validity of the one-loop approximation for studying the low-energy dynamics of supersymmetric large $`N`$ reduced models has been addressed in Ref. by studying the four-dimensional version of the IIB matrix model without any approximations. It was found that the large $`N`$ behavior of the extent of the distribution of $`x_i`$ agrees with the prediction from the one-loop, low-energy effective theory. We therefore expect that a low-energy phenomenon such as SSB of Lorentz invariance can be studied with the low-energy effective theory.
On the other hand we note that from a technical point of view the one-loop approximation offers certain advantages which are essential in the present work. We recall that even in the four-dimensional version of the IIB matrix model, the largest $`N`$ one can achieve for the full model (without the one-loop approximation) is $`N=48`$ using supercomputers. In order to detect even a small trend of SSB of Lorentz invariance from 10D to 4D, say, one would expect that $`N`$ should be larger than $`4^4=256`$. Due to the fact that the size of the matrix $`M`$ in (3.22) is order of $`N`$ smaller than the size of the matrix $``$ in (2.10), the one-loop approximation enables us to reach $`N=512`$ with reasonable effort. Details of the algorithm used in the Monte Carlo simulations are presented in the Appendix A. The algorithm is a variant of the Hybrid Monte Carlo algorithm which is one of the standard algorithms used in the study of systems with dynamical fermions. The computational effort of the algorithm is estimated to be O($`N^3`$)$``$O($`N^{7/2}`$), which should be compared with an estimate O($`N^5`$) for the algorithm used for the full model in Ref. .
## 4 Results for the distribution of $`x_i`$
First we look at the distribution $`\rho (r)`$ of the distance $`r`$ between $`x_i`$’s, where the distance between two arbitrary points $`x_ix_j`$ is measured by $`\sqrt{(x_ix_j)^2}`$. In Fig. 1, we plot the results for $`D=10`$ with $`N=192`$, 256, 384 and 512. Fig. 2 shows the results for $`D=6`$ with $`N=256`$, 512 and 768. We note that in both cases the distribution at small $`r`$ falls off rapidly below $`r\mathrm{}`$, where $`\mathrm{}=1`$ is the UV cutoff introduced in (3.20). The small penetration into the $`r<\mathrm{}`$ region is due to $`\kappa `$ being finite. However, the results show that the values of $`\kappa `$ we have taken ($`\kappa =300`$ for $`D=10`$ and $`\kappa =100`$ for $`D=6`$) are large enough to make the penetration reasonably small.
In order to see the spontaneous breakdown of Lorentz invariance, we consider the moment of inertia tensor of $`N`$ points $`x_i`$ ($`i=1,\mathrm{},N`$) in a flat $`D`$-dimensional space-time<sup>5</sup><sup>5</sup>5Such a quantity has also been studied in Refs. . . It can be defined as
$$T_{\mu \nu }=\frac{2}{N(N1)}\underset{i<j}{}(x_{i\mu }x_{j\mu })(x_{i\nu }x_{j\nu }),$$
(4.1)
which is a $`D\times D`$ real symmetric matrix. The $`D`$ eigenvalues $`\lambda _1>\lambda _2>\mathrm{}>\lambda _D>0`$ of the matrix $`T`$ represent the principal moments of inertia. We measure $`\lambda _\mu `$ for each configuration and take an average $`\lambda _\mu `$ over all the configurations generated by the Monte Carlo simulation. If Lorentz invariance is not spontaneously broken, $`\lambda _\mu `$ must be all equal in the large $`N`$ limit, representing an isotropic distribution of $`x_i`$. We therefore search for a trend which differs from such a large $`N`$ behavior.
We first note that if the system that describes the dynamics of $`x_i`$ were a simple branched polymer in a flat $`D`$-dimensional space-time, spontaneous breakdown of Lorentz invariance would certainly not occur, and moreover, the large $`N`$ behavior of $`\lambda _\mu `$ could be expected to be $`\lambda _\mu N^{1/2}`$. This is due to the fact that the Hausdorff dimension of a branched polymer is $`d_\mathrm{H}=4`$ irrespectively of the dimension $`D`$ of the space-time in which it is embedded. For this we note that the extent of the distribution of $`x_i`$ is given by
$$R=\sqrt{T_{\mu \mu }}=\sqrt{\underset{\mu =1}{\overset{D}{}}\lambda _\mu },$$
(4.2)
which is related to the number $`N`$ of points through $`N(R/\mathrm{})^{d_\mathrm{H}}`$. The $`\mathrm{}`$ is the UV cutoff taken to be $`\mathrm{}=1`$. Therefore, we find $`RN^{1/d_\mathrm{H}}N^{1/4}`$, which leads to the announced large $`N`$ behavior $`\lambda _\mu N^{1/2}`$.
In Fig. 3 we plot the normalized eigenvalues $`\lambda _\mu /N^{1/2}`$ against $`1/N`$ for $`D=10`$ with $`N=192,256,384,512`$. We find that the smallest (normalized) eigenvalue $`\lambda _{10}/N^{1/2}`$ is almost constant and the larger ones are monotonously decreasing towards the same constant. Thus the observed large $`N`$ behavior suggests that there is no spontaneous breakdown of Lorentz invariance and the large $`N`$ behavior of the extent of the $`x`$-distribution is consistent with a simple branched-polymer prediction. In particular, we see no trend for a gap developing between the fourth and the fifth largest eigenvalues, which could have been observed if the 10D Lorentz invariance were broken down to a four-dimensional one. In Fig. 4 we present the results for $`D=6`$ with $`N=192,256,512,768`$. The qualitative behavior is the same as in $`D=10`$.
The results for $`D=10`$ might suggest that the leading finite $`N`$ effect is given by $`1/N`$. For example, a linear extrapolation in $`1/N`$ using the data for $`N=384,512`$ leads to almost the same results for all the ten (normalized) eigenvalues at $`N=\mathrm{}`$. It is therefore tempting to speculate that finite $`N`$ effects in the IIB matrix model is given by a $`1/N`$ expansion. Such a form of finite $`N`$ effects, if it is true, is in remarkable contrast to the “bosonic” case, in which the fermions in the IIB matrix model (2.1) is omitted by hand. In that case, the large $`N`$ behavior of correlation functions was determined analytically to all orders in the $`1/D`$ expansion and finite $`N`$ effects were found to be given by a $`1/N^2`$ expansion.
For $`D=6`$, the linear extrapolation gives almost the same results up to the third smallest eigenvalue, but not for the larger ones. Thus, as far as reproducing an isotropic distribution of $`x_i`$ in the large $`N`$ limit is concerned, $`D=6`$ seems to have larger finite $`N`$ effects. This rather counter-intuitive result might be understood by considering the branched-polymer representation of the low-energy effective theory . The attractive potential between two $`x_i`$’s connected by a bond is given by $`r^{3(D2)}`$, where $`r=\sqrt{(x_ix_j)^2}`$. Therefore, it is much stronger for $`D=10`$ than for $`D=6`$. Since the branched-polymer system in both cases is not a simple maximal tree but something more complicated, typically involving a double-tree structure, it may well be that the stronger the attractive potential is, the more the distribution of $`x_i`$ tends to become isotropic.
## 5 Summary and Discussion
In this paper, we studied the IIB matrix model (and its six-dimensional version) at large $`N`$ by using the low-energy effective theory developed in Ref. . Unlike in the four-dimensional version studied in Ref. , the fermion integral yields a Pfaffian (or a determinant in the 6D case) which is complex in general. In the present paper we omitted the phase of the Pfaffian by hand in order to avoid the sign problem in the Monte Carlo simulations, and we studied the effect of the modulus only. We have seen that the distribution of $`x_i`$ becomes more and more isotropic as we increase $`N`$. Based on this observation we conclude that if the SSB ever occurs in the original IIB matrix model, the phase of the Pfaffian induced by fermions must play a crucial rôle.
It is interesting to compare with the situation in perturbative superstring theory. Also in that case the fermionic degrees of freedom can be integrated out explicitly , leaving us with an effective bosonic string theory, where the action consists of three parts: the ordinary bosonic string action, an extrinsic curvature term, and (in Euclidean space-time) a purely imaginary Wess-Zumino-like term. This means that if we can neglect the imaginary part, the superstring theory is equivalent to a bosonic string theory with extrinsic curvature. While this equivalence has not been directly disproved, it does not have much support, either, from numerical simulations and analytical calculations .
In Ref. a deformation of the IIB matrix model by introducing an integer parameter $`\nu `$ which couples to the phase of the Pfaffian has been considered. The original IIB matrix model corresponds to $`\nu =1`$. The deformed model continues to be well-defined, and preserves Lorentz invariance, the SU($`N`$) symmetry, and the cluster property. In this language the present work corresponds to the study of the case $`\nu =0`$ using the low-energy effective theory. The opposite extreme, the limit $`\nu =\mathrm{}`$, has been studied analytically in Ref. , and the spontaneous breakdown of Lorentz invariance has been discovered. Of particular interest is the fact that different conclusions have been obtained for $`\nu =\mathrm{}`$ and for $`\nu =0`$, which already implies that a phase transition should occur in between. If the original IIB matrix model ($`\nu =1`$) belongs to the same phase as $`\nu =\mathrm{}`$, we can investigate the dynamics of the IIB matrix model by studying the $`\nu =\mathrm{}`$ model using Monte Carlo simulation as suggested in Ref. . It would be extremely interesting if we can obtain a flat four-dimensional space-time in that way. We hope that the Monte Carlo technique developed in the present work is useful for such future studies and eventually enables us to explore the dynamics of the IIB matrix model.
## Acknowledgments
We would like to thank R.J. Szabo and G. Vernizzi for carefully reading the manuscript. J. A. and K.N. A. acknowledge the support of MaPhySto, a center for mathematical physics and stochastics, funded by the Danish National Research Foundation. K.N. A. acknowledges partial support from a postdoctoral fellowship from the Greek Institute of National Fellowships (IKY). J. N. is supported by the Japan Society for the Promotion of Science as a Postdoctoral Fellow for Research Abroad. The computation has been carried out on VPP500 at High Energy Accelerator Research Organization (KEK), VPP700E at The Institute of Physical and Chemical Research (RIKEN), SR8000 at Computer Centre, University of Tokyo, and SX4 at Research Center for Nuclear Physics (RCNP) of Osaka University. J. N. is grateful to T. Ishikawa, S. Ohta, A.I. Sanda, H. Toki, S. Uehara and Y. Watanabe for many kind helps in using the supercomputers. This work is supported by the Supercomputer Project (No.99-53) of KEK.
## Appendix A: The algorithm for the simulation
In this appendix, we explain the algorithm<sup>6</sup><sup>6</sup>6Ref. gives an overview of effective algorithms for dynamical fermions. we use for the Monte Carlo simulation of our model defined by (3.22).
We first simplify the model by introducing rescaled variables $`b_{\mu ij}`$ as
$$b_{\mu ij}=\sqrt{(x_ix_j)^2}a_{\mu ij}.$$
(A.1)
The model we have to simulate then becomes
$$Z=\text{d}x\text{d}b\text{e}^{S_\mathrm{G}S_{\mathrm{cut}}}|\text{Pf}M|,$$
(A.2)
where
$$S_\mathrm{G}=\underset{i<j}{}|b_{\mu ij}|^2,$$
(A.3)
and (3.12) can be rewritten in terms of $`b_{ij}`$ as
$$M_{i\alpha ,j\beta }^{}=\frac{(x_{i\rho }x_{j\rho })}{\{(x_ix_j)^2+(ϵ\mathrm{})^2\}^2}(\stackrel{~}{\mathrm{\Gamma }}_\mu \stackrel{~}{\mathrm{\Gamma }}_\rho ^{}\stackrel{~}{\mathrm{\Gamma }}_\sigma )_{\alpha \beta }(b_{\mu ji}b_{\sigma ij}b_{\mu ij}b_{\sigma ji})\text{for}ij,$$
(A.4)
where we have also introduced the dimensionless regularization parameter $`ϵ`$ to reduce numerical instabilities. The regularization parameter $`ϵ`$ should be taken to be small enough not to affect the system. For both, 6D and 10D, cases, we took $`ϵ=0.1`$. We note also that with this regularization, the model is well-defined without introducing the cutoff term (3.19). However, by simulating such a model, we find that all the $`x_i`$’s become densely packed within the extent $`ϵ\mathrm{}`$ as we increase $`N`$. Therefore, we still need the cutoff term (3.19) in order to make our model valid as a low-energy effective theory of the original IIB matrix model. For convenience we combine our dynamical variables $`x_{i\mu }`$ and $`b_{\mu ij}`$ ($`ij`$, $`b_{\mu ji}=b_{\mu ij}^{}`$) into a Hermitian matrix $`B_\mu `$ as $`(B_\mu )_{ij}=x_{i\mu }\delta _{ij}+b_{\mu ij}`$.
The algorithm we use to simulate the model (A.2) is a variant of the Hybrid Monte Carlo algorithm . The first step of the Hybrid Monte Carlo algorithm is to apply molecular dynamics . We introduce a conjugate momentum for $`B_{\mu ij}`$ as $`X_{\mu ij}`$, which satisfies $`X_{\mu ji}=(X_{\mu ij})^{}`$; $`X_\mu `$ are Hermitian matrices. The partition function can be rewritten as<sup>7</sup><sup>7</sup>7Unlike in the standard Hybrid Monte Carlo algorithm , we do not introduce the so-called pseudo-fermions. Note, in this regard, that our system (A.2) is different from ordinary field theories with dynamical fermions in the following respects. 1) The number of fermion flavors should be strictly one in order to respect supersymmetry. 2) The size of the matrix $`M`$ is much smaller than the system size. 3) The matrix $`M`$ is not sparse.
$$Z=\text{d}X\text{d}B\text{e}^H,$$
(A.5)
where $`H`$ is the “Hamiltonian” defined by
$$H=\frac{1}{2}\underset{\mu ij}{}|X_{\mu ij}|^2+S_\mathrm{G}[B]+S_{\mathrm{cut}}[B]\frac{1}{2}\mathrm{ln}|detM|.$$
(A.6)
The update of $`X_{\mu ij}`$ can be done by just generating $`X_{\mu ij}`$ with the probability distribution $`\mathrm{exp}(\frac{1}{2}|X_{\mu ij}|^2)`$. In order to update $`B_{\mu ij}`$, we solve the Hamilton equation given by
$`{\displaystyle \frac{\text{d}B_{\mu ij}(\tau )}{\text{d}\tau }}`$ $`=`$ $`{\displaystyle \frac{H}{X_{\mu ij}}}=X_{\mu ji}`$ (A.7)
$`{\displaystyle \frac{\text{d}X_{\mu ij}(\tau )}{\text{d}\tau }}`$ $`=`$ $`{\displaystyle \frac{H}{B_{\mu ij}}}={\displaystyle \frac{1}{4}}\text{Tr}\left({\displaystyle \frac{D}{B_{\mu ij}}}D^1\right){\displaystyle \frac{S_\mathrm{G}}{B_{\mu ij}}}{\displaystyle \frac{S_{\mathrm{cut}}}{B_{\mu ij}}},`$ (A.8)
where we have defined $`D=M^{}M`$. Note that when taking the derivatives in (A.8), $`(B_\mu )_{ij}`$ and $`(B_\mu )_{ji}`$ should be treated as independent complex variables. Along the “classical trajectory” given by the Hamilton equation,
(i) $`H`$ is invariant,
(ii) the motion is reversible,
(iii) the phase-volume is preserved; i.e.
$$\frac{(B(\tau ),X(\tau ))}{(B(0),X(0))}=1,$$
(A.9)
where $`(B(\tau ),X(\tau ))`$ is the point on the trajectory after evolution for fixed $`\tau `$ from $`(B(0),X(0))`$. Therefore, generating a new sets of $`(B,X)`$ by solving the Hamilton equation for a fixed “time” interval $`\tau `$ satisfies the detailed balance. This procedure, together with the proceeding generation of $`X_{\mu ij}`$ with the Gaussian distribution, is called “one trajectory”, which corresponds to “one sweep” in ordinary Monte Carlo simulations.
When solving the Hamilton equation numerically, we have to discretize the equation. A discretization which preserves the properties (ii) and (iii) is known. The property (i) cannot be preserved and yields a small violation of the Hamiltonian conservation. In order to still satisfy the detailed balance exactly, we can perform a Metropolis accept/reject as the end of each trajectory.
We introduce a short-hand notation for discretized $`X_\mu (\tau )`$ and $`B_\mu (\tau )`$ as
$$X_\mu ^{(r)}=X_\mu (r\mathrm{\Delta }\tau );B_\mu ^{(s)}=B_\mu (s\mathrm{\Delta }\tau ).$$
(A.10)
Given the configuration $`(B_\mu ^{(0)},X_\mu ^{(0)})`$, the configuration $`(B_\mu ^{(\nu )},X_\mu ^{(\nu )})`$ after the evolution for a fixed time $`\tau =\nu \mathrm{\Delta }\tau `$ is obtained by solving the discretized Hamilton equation
$`B_{\mu ij}^{(\frac{1}{2})}`$ $`=`$ $`B_{\mu ij}^{(0)}+{\displaystyle \frac{\mathrm{\Delta }\tau }{2}}X_{\mu ji}^{(0)}`$
$`B_{\mu ij}^{(n+\frac{1}{2})}`$ $`=`$ $`B_{\mu ij}^{(n\frac{1}{2})}+\mathrm{\Delta }\tau X_{\mu ji}^{(n)}`$
$`B_{\mu ij}^{(\nu )}`$ $`=`$ $`B_{\mu ij}^{(\nu \frac{1}{2})}+{\displaystyle \frac{\mathrm{\Delta }\tau }{2}}X_{\mu ji}^{(\nu )}`$ (A.11)
$`X_{\mu ij}^{(m+1)}`$ $`=`$ $`X_{\mu ij}^{(m)}\mathrm{\Delta }\tau {\displaystyle \frac{H}{B_{\mu ij}}}(B_\mu ^{(m+\frac{1}{2})}),`$ (A.12)
where $`n=1,2,\mathrm{},\nu 1`$ and $`m=0,1,\mathrm{},\nu 1`$. Note that the first and the final steps in the evolution (Appendix A: The algorithm for the simulation) of $`B_\mu `$ are treated with special care. This particular discretization, which is called as “leap-frog discretization”, preserves the properties (ii) and (iii). At the end of the trajectory, we make a Metropolis accept/reject with the probability $`P=\mathrm{min}(1,\text{e}^{\mathrm{\Delta }H})`$, where $`\mathrm{\Delta }H`$ is the difference of the Hamiltonian $`H`$ defined in (A.6) for the configurations $`(B_\mu ^{(0)},X_\mu ^{(0)})`$ and $`(B_\mu ^{(\nu )},X_\mu ^{(\nu )})`$. Optimization of $`\mathrm{\Delta }\tau `$ and $`\nu `$ is discussed in Appendix B.
In the evolution (A.12) of $`X_\mu `$, one needs to calculate
$$\text{Tr}\left(\frac{D}{B_{\mu ij}}D^1\right)=\text{Tr}\left(\frac{M}{B_{\mu ij}}M^1\right)+\text{Tr}\left(\frac{M^{}}{B_{\mu ij}}M^1\right).$$
(A.13)
If we write the first term as $`T_{\mu ij}`$, the second term can be written as $`T_{\mu ji}^{}`$. In what follows, we calculate $`T_{\mu ij}`$ explicitly. We first note that
$$T_{\mu ij}=\underset{k,l=1}{\overset{N1}{}}\underset{\alpha \beta }{}\frac{M_{k\alpha ,l\beta }}{B_{\mu ij}}C_{l\beta ,k\alpha }=\underset{k,l=1}{\overset{N}{}}\underset{\alpha \beta }{}\frac{M_{k\alpha ,l\beta }^{}}{B_{\mu ij}}C_{l\beta ,k\alpha }^{},$$
(A.14)
where $`C_{i\alpha ,j\beta }`$ is defined as
$$C_{i\alpha ,j\beta }=(M^1)_{i\alpha ,j\beta },$$
(A.15)
and $`C_{i\alpha ,j\beta }^{}`$ is defined as
$`C_{i\alpha ,j\beta }^{}=C_{i\alpha ,j\beta };C_{i\alpha ,N\beta }^{}={\displaystyle \underset{k=1}{\overset{N1}{}}}C_{i\alpha ,k\beta }`$
$`C_{N\alpha ,j\beta }^{}={\displaystyle \underset{k=1}{\overset{N1}{}}}C_{k\alpha ,j\beta };C_{N\alpha ,N\beta }^{}={\displaystyle \underset{k,l=1}{\overset{N1}{}}}C_{k\alpha ,l\beta }.`$ (A.16)
The indices $`i,j`$ in (A.15) and (A.16) run from 1 to $`N1`$. We further rewrite $`T_{\mu ij}`$ as
$$T_{\mu ij}=\underset{k=1}{\overset{N}{}}\underset{lk}{\overset{N}{}}\underset{\alpha \beta }{}\frac{M_{k\alpha ,l\beta }^{}}{B_{\mu ij}}C_{l\beta ,k\alpha }^{\prime \prime },$$
(A.17)
where we have introduced $`C_{i\alpha ,j\beta }^{\prime \prime }`$ through
$$C_{i\alpha ,j\beta }^{\prime \prime }=C_{i\alpha ,j\beta }^{}C_{j\alpha ,j\beta }^{}.$$
(A.18)
Then we calculate $`\frac{M_{k\alpha ,l\beta }^{}}{B_{\mu ij}}`$ explicitly, which yields
$$T_{\mu ij}=\underset{\alpha \beta }{}\frac{(x_{i\rho }x_{j\rho })}{\{(x_ix_j)^2+(ϵ\mathrm{})^2\}^2}\{(\stackrel{~}{\mathrm{\Gamma }}_\sigma \stackrel{~}{\mathrm{\Gamma }}_\rho ^{}\stackrel{~}{\mathrm{\Gamma }}_\mu )_{\alpha \beta }(\stackrel{~}{\mathrm{\Gamma }}_\mu \stackrel{~}{\mathrm{\Gamma }}_\rho ^{}\stackrel{~}{\mathrm{\Gamma }}_\sigma )_{\alpha \beta }\}b_{\sigma ji}C_{i\beta ,j\alpha }^{\prime \prime \prime }$$
(A.19)
for $`ij`$ and
$`T_{\mu ii}`$ $`=`$ $`{\displaystyle \underset{ki}{\overset{N}{}}}{\displaystyle \underset{\alpha \beta }{}}\left[{\displaystyle \frac{\delta _{\mu \rho }}{\{(x_ix_k)^2+(ϵ\mathrm{})^2\}^2}}{\displaystyle \frac{4(x_{i\mu }x_{k\mu })(x_{i\rho }x_{k\rho })}{\{(x_ix_k)^2+(ϵ\mathrm{})^2\}^3}}\right]`$ (A.20)
$`(\stackrel{~}{\mathrm{\Gamma }}_\sigma \stackrel{~}{\mathrm{\Gamma }}_\rho ^{}\stackrel{~}{\mathrm{\Gamma }}_\tau )_{\alpha \beta }(b_{\sigma ki}b_{\tau ik}b_{\sigma ik}b_{\tau ki})C_{i\beta ,k\alpha }^{\prime \prime \prime },`$
where we defined $`C_{i\alpha ,j\beta }^{\prime \prime \prime }`$ through
$`C_{i\alpha ,j\beta }^{\prime \prime \prime }`$ $`=`$ $`C_{i\alpha ,j\beta }^{\prime \prime }+C_{j\alpha ,i\beta }^{\prime \prime }`$ (A.21)
$`=`$ $`C_{i\alpha ,j\beta }^{}C_{j\alpha ,j\beta }^{}C_{i\alpha ,i\beta }^{}+C_{j\alpha ,i\beta }^{}.`$
Let us comment on the required computational effort of our algorithm. The dominant part comes from calculating the inverse in (A.15), which requires $`O(n^3)`$ arithmetic operations, where $`n`$ is the size of the matrix to be inverted. In the present case $`n`$ is of O($`N`$). In order to keep the acceptance rate at the Metropolis accept/reject procedure reasonably high, one has to decrease the step size $`\mathrm{\Delta }\tau `$ as one goes to larger $`N`$. We have seen from simulations that the step size should be taken to be $`\mathrm{\Delta }\tau \frac{1}{\sqrt{N}}`$, which is consistent with a general formula $`\mathrm{\Delta }\tau V^{1/4}`$ in Ref. , where the system size $`V`$ should be replaced by O($`N^2`$) in our case. Accordingly, the number of steps for one trajectory increases as $`\nu \sqrt{N}`$. Therefore, the required computational effort is estimated to be $`O(N^{7/2})`$.
In fact, one can reduce the computational effort by omitting the Metropolis accept/reject procedure, since in that case one can keep the step size $`\mathrm{\Delta }\tau `$ fixed to a small constant for all $`N`$. The required computational effort becomes $`O(N^3)`$. The price one has to pay is that the algorithm then suffers from a systematic error due to the small violation of the Hamiltonian conservation. The systematic error can be estimated to be O($`(\mathrm{\Delta }\tau )^2`$), see Appendix B.
## Appendix B: Optimization of the algorithm
In this appendix, we first discuss the optimization of the parameters in the algorithm with the Metropolis accept/reject procedure. We then move on to the case when the Metropolis accept/reject procedure is omitted. The algorithm and the parameters used for each run is also described.
To start with, one can actually generalize the algorithm described in Appendix A by taking the step size $`\mathrm{\Delta }\tau _x`$ for the diagonal elements to be different<sup>8</sup><sup>8</sup>8This is not the case for the full model , where $`\mathrm{\Delta }\tau `$ should be taken universally for each element of $`A_{\mu ij}`$ in order to respect the SU($`N`$) invariance. In the present case, since the SU($`N`$) “gauge” invariance is fixed by the “gauge-fixing” (3.4), we only have to respect the invariance under permutation of the SU($`N`$) indices. from the step size $`\mathrm{\Delta }\tau `$ for the off-diagonal elements of $`B_{\mu ij}`$ and $`X_{\mu ij}`$. We have fixed the ratio $`\mathrm{\Delta }\tau _x/\mathrm{\Delta }\tau `$ to be the ratio of the mean magnitudes of $`x_{i\mu }`$ and $`b_{\mu ij}`$, which we found from simulations to be approximately $`1`$ for $`D=6`$ and $`1/2`$ for $`D=10`$, respectively.
We still have two parameters in the algorithm, i.e., $`\mathrm{\Delta }\tau `$ and $`\nu `$ (the number of molecular dynamics steps for each trajectory), which can be optimized in a standard way . The key point of the optimization is that the autocorrelation time (in units of accepted trajectory) depends only on $`\nu \mathrm{\Delta }\tau =\tau `$, but not on $`\mathrm{\Delta }\tau `$ and $`\nu `$ separately. This allows us to perform the optimization in two steps. First, one fixes $`\nu \mathrm{\Delta }\tau =\tau `$ and optimize $`\mathrm{\Delta }\tau `$ so that the effective speed of motion in the configuration space, given by acceptance rate times $`\mathrm{\Delta }\tau `$, is maximized. Second, using the optimized $`\mathrm{\Delta }\tau `$ for each $`\tau `$, one minimizes a typical autocorrelation time in units of molecular dynamics step with respect to $`\tau `$.
In Fig. 5, we show the acceptance rate times $`\mathrm{\Delta }\tau `$ as a function of $`\mathrm{\Delta }\tau `$ for $`D=6`$, $`N=16`$ with fixed $`\tau =1.5`$. We find that the optimal $`\mathrm{\Delta }\tau `$ is $`0.0375`$. For $`D=6,N=32`$ with fixed $`\tau =1.5`$, we find that the optimal $`\mathrm{\Delta }\tau `$ is $`0.02142`$, which is smaller than for $`N=16`$ as expected. For both cases, the acceptance rate at the optimal $`\mathrm{\Delta }\tau `$ is found to be $`5060\%`$.
Using the optimal $`\mathrm{\Delta }\tau `$ obtained in the above way for each $`\tau `$, we minimize a typical autocorrelation time (in units of molecular dynamics step) with respect to $`\tau `$. Here, we measure the autocorrelation time of the extent $`R`$ of the $`x_i`$-distribution defined in (4.2) and plot it as a function of $`\tau `$. Fig. 6 shows the result for $`D=6`$ with $`N=16`$. We see that it has a minimum around $`\tau 1.5`$. Similar experiments for $`N=32`$ showed that the optimal $`\tau `$ is almost independent of $`N`$.
When we omit the Metropolis accept/reject procedure, the algorithm suffers from a systematic error, as is explained in Appendix A. Here we fix $`\nu \mathrm{\Delta }\tau `$ to the optimal $`\tau `$ obtained for the case with the Metropolis accept/reject procedure and study the $`\mathrm{\Delta }\tau `$ dependence of the systematic error. As a quantity which shows a large systematic error, we take $`\lambda _1`$, the expectation value of the largest eigenvalue of the moment of inertia tensor defined by (4.1). In Fig. 7, we plot the result against $`(\mathrm{\Delta }\tau )^2`$ for $`D=6,N=16`$ with $`\tau =1.5`$. The systematic error is seen to vanish as O($`(\mathrm{\Delta }\tau )^2`$), which can be also understood theoretically by using the analysis described in Ref. .
Finally, let us comment on the algorithms and the parameters we used in our simulations at large $`N`$. The runs for $`D=6`$ with $`N=192,256`$ were made by the algorithm including the Metropolis accept/reject procedure. The parameters are $`\mathrm{\Delta }\tau =0.0048,\nu =60`$ for $`N=192`$, and $`\mathrm{\Delta }\tau =0.004,\nu =50`$ for $`N=256`$. The numbers of configurations used for the measurements are 560 and 1732 for $`N=192`$ and $`N=256`$, respectively. (The optimization described in this Appendix was not completed when we started these runs. Accordingly, we needed significantly larger numbers of configurations compared to the cases below.) For the other cases, we omitted the Metropolis accept/reject procedure in order to obtain a sufficient statistics. The parameters are $`\mathrm{\Delta }\tau =0.0075,\nu =200`$ for $`D=10`$ with $`N=192,256,384,512`$, and $`\mathrm{\Delta }\tau =0.005,\nu =300`$ for $`D=6`$ with $`N=512,768`$. The numbers of configurations used for the measurements are 88, 40, 84, 44 for $`D=10`$ with $`N=192,256,384,512`$, and 64, 50 for $`D=6`$ with $`N=512,768`$, respectively.
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# Self-dual Codes over the Kleinian Four Group
## 1 Introduction
In this work, we describe a new and natural fourth step in the series of analogies known to exist between binary codes, lattices and vertex operator algebras (see for example \[CS93b, Höh95\]).
Linear codes over the finite field $`𝐅_4`$ are studied in many papers (cf. \[MOSW78, CPS79, Slo79, Slo78, LP90, CS90a, Huf90, Huf91\]), but a developed theory for codes over the Kleinian four-group $`K𝐙_2\times 𝐙_2`$ is missing. It turns out that there is a similar rich theory as one has for binary linear codes. Parts of the results are known from some different viewpoints, but the use of Kleinian codes seems most natural.
We will prove all the results in terms of a theory for Kleinian codes, since this leads to a theory of its own right, although one can deduce most theorems from the corresponding results for self-dual vertex operator algebras or lattices or binary codes. To emphasize this relation, we will give after every theorem a list of references of the analogue<sup>*</sup><sup>*</sup>* An additional asterisk indicates that the theorem can be obtained from the analogues theorems for binary codes, lattices or vertex operator algebras by the relations described in the final section. theorems for binary codes (B), lattices (L) and vertex operator algebras (V).
The second section contains the main definitions and first results. The next section describes the classification of odd and even self-dual codes up to length $`8`$. In the fourth section, we study extremal codes. This are codes with the largest possible minimal weight. The fifth section is about designs for the space $`K^n`$. Section six deals with lexicographic constructions.
In the final section, we explain the relation and discuss some of the analogies with self-dual binary codes, lattices and vertex operator algebras in more detail. Self-dual Kleinian codes of length $`n`$ can be identified with self-dual vertex operator superalgebras of rank $`4n`$ containing a vertex operator algebra of type $`V_{D_4}^n`$. From this viewpoint, Kleinian codes are a special case of codes over a $`3`$-dimensional topological quantum field theory.
Our motivation behind the introduction of Kleinian codes was to have an additional testbed besides binary codes and lattices for the understanding of vertex algebras. Kleinian codes have already found applications as quantum codes and some of the results have been extended to and sharpened for codes of larger length.
## 2 Definitions and basic results
Denote the elements of the Kleinian four group $`K𝐙_2\times 𝐙_2`$ by $`0`$, $`a`$, $`b`$ and $`c`$, where $`0`$ is the neutral element. The automorphism group of $`K`$ is $`S_3`$, the permutation group of the three nonzero elements $`a`$, $`b`$ and $`c`$. A code $`C`$ over $`K`$ of length $`n`$ is a subset of the words of length $`n`$ over the alphabet $`K`$, i.e. consists of vectors $`𝐱=(x_1,\mathrm{},x_n)`$, $`x_iK`$, the codewords of $`C`$. The weight $`\mathrm{wt}(𝐱)`$ of a codeword $`𝐱`$ is the number of nonzero $`x_i`$. The minimal weight of $`C`$ is defined by
$$d=\mathrm{min}\{\mathrm{wt}(𝐱)𝐱C,𝐱0\}.$$
The code $`C`$ is called linear if $`C`$ is a subgroup of the abelian group $`K^n𝐙_2^{2n}`$. A linear code has $`4^k`$ elements with $`k\frac{1}{2}𝐙`$ and we denote $`k`$ the dimension of the code. All codes in this article are assumed to be linear. A code of length $`n`$, dimension $`k`$ and minimal weight $`d`$ is shortly denoted as a $`[n,k,d]`$\- or $`[n,k]`$-code. Let now $`C`$ be a $`[n,k]`$-code.
An important part of the structure which makes the theory of Kleinian codes interesting is the scalar product $`(.,.):K^n\times K^n𝐅_2`$, $`(𝐱,𝐲)=_{i=1}^nx_iy_i`$, where the symmetric bilinear dot product $`:K\times K𝐅_2`$ is defined by $`ab=ba=1`$, $`ac=ca=1`$, $`bc=cb=1`$ and zero otherwise. The dual code $`C^{}`$ is defined by
$$C^{}=\{𝐱K^n(𝐱,𝐲)=0\text{for all }𝐲C\}$$
and has type $`[n,nk]`$.
We call $`C`$ self-orthogonal if $`CC^{}`$ and self-dual if $`C^{}=C`$.
The direct sum $`CD`$ of a $`[n,k]`$-code $`C`$ and a $`[m,l]`$-code $`D`$ is the direct product subgroup of $`K^nK^m`$ and has type $`[n+m,k+l]`$. If $`C`$ can be written in a nontrivial way as a direct sum, $`C`$ is called decomposable, otherwise indecomposable. Obviously $`(CD)^{}=C^{}D^{}`$. Every code $`C`$ is after a renumbering of the positions a direct sum of indecomposable codes. The isomorphism classes of the components are uniquely determined up to permutation.
The (Hamming) weight enumerator of $`C`$ is the degree $`n`$ polynomial
$$W_C(u,v)=\underset{i=0}{\overset{n}{}}A_iu^{ni}v^i\text{with }A_i=\mathrm{\#}\{𝐱C\mathrm{wt}(𝐱)=i\}.$$
The complete weight enumerator is the polynomial
$$\mathrm{cwe}_C(p,q,r,s)=\underset{i,j,k,l}{}A_{i,j,k,l}p^iq^jr^ks^l,$$
where $`A_{i,j,k,l}`$ is the number of code words in $`C`$ containing at $`i`$, $`j`$, $`k`$ resp. $`l`$ of the $`n`$ positions the element $`0`$, $`a`$, $`b`$ resp. $`c`$. There is the obvious relation $`W_C(u,v)=\mathrm{cwe}_C(u,v,v,v)`$. Finally define for a natural number $`g`$ the poly- or $`g`$-weight enumerator $`W_C^g`$ as a polynomial in $`2^g`$ variables $`t_\nu `$ indexed by $`\nu 𝐅_2^g`$:
$$W_C^g=\underset{𝐱^1,\mathrm{},𝐱^gC}{}\underset{i=1}{\overset{n}{}}t_{(\mathrm{wt}(x_i^1),\mathrm{},\mathrm{wt}(x_i^g))}$$
and similar the complete $`g`$-weight enumerator $`\mathrm{cwe}_C^g`$ as a polynomial in $`4^g`$ variables $`s_\nu `$ where $`\nu K^g`$.
The code $`C`$ is called even if the weights of all codewords are divisible by $`2`$. Note, that a code spanned by an orthogonal system of vectors of even weight is itself even.
The automorphisms of the abelian group $`K^n`$ which are also isometries for the metric $`d(𝐱,𝐲)=\mathrm{wt}(𝐱𝐲)`$ on $`K^n`$ form the semidirect product $`G=S_3^n:S_n`$ consisting of the permutation of the positions together with a permutation of the symbols $`a`$, $`b`$ and $`c`$ at each position. The automorphism group of $`C`$ is the subgroup of $`G`$ sending $`C`$ to itself:
$$\mathrm{Aut}(C)=\{gS_3^n:S_ngC=C\}.$$
Two codes $`C`$ and $`D`$ are called to be equivalent if there is a $`gG`$ with $`gC=D`$. The number of distinct codes equivalent to $`C`$ is
$$\frac{6^nn!}{|\mathrm{Aut}(C)|}.$$
Equivalent codes have the same (poly-) weight enumerator, but not necessarily the same complete (poly-) weight enumerator. If $`C`$ is self-orthogonal, self-dual or even, so it is the equivalent code.
Since $`K`$ is isomorphic to the additive group of the field $`𝐅_4`$, we can interpret every code over $`K`$ as a code over $`𝐅_4`$. Every code linear as $`𝐅_4`$-code is linear as Kleinian code, but not conversely. If $`C`$ is a self-dual Type IV $`𝐅_4`$-code for the hermitian scalar product of $`𝐅_4^n`$, then it is also a even self-dual Kleinian code (cf. \[MOSW78\]). Perfect Kleinian codes are the same as perfect $`𝐅_4`$-codes, the only perfect Kleinian codes which exist are $`1`$ error correcting codes \[Tie73\].
Examples of Kleinian codes:
\- The $`[1,\frac{1}{2},1]`$-code $`\gamma _1=\{(0),(a)\}`$: $`|\mathrm{Aut}(\gamma _1)|=2`$, $`W_{\gamma _1}(u,v)=u+v`$.
\- The $`[2,1,2]`$-code $`ϵ_2=\{(00),(aa),(bb),(cc)\}`$: $`|\mathrm{Aut}(ϵ_2)|=12`$, $`W_{ϵ_2}(u,v)=u^2+3v^2`$.
\- The $`[6,3,4]`$-Hexacode $`𝒞_6`$ spanned by
$$\{(a0a0bb),(a0bba0),(bba0a0),(00aaaa),(aa00aa),(b0b0ca)\}$$
as a Kleinian code. One has $`|\mathrm{Aut}(𝒞_6)|=2^221518=2160`$, $`W_{𝒞_6}(u,v)=u^6+45u^2v^4+18v^6`$.
\- The Hamming code $`_m`$, $`m2`$ of type $`[(4^m1)/3,(4^m1)/3m,3]`$ and the extended Hamming code $`\overline{}_m`$, $`m2`$ of type $`[(4^m1)/3+1,(4^m1)/3m,4]`$.
All examples are linear; the first three codes are self-dual; $`ϵ_2`$ and $`𝒞_6\overline{}_2`$ and $`\overline{}_m`$ are even; besides $`\gamma _1`$, they are equivalent to codes over $`𝐅_4`$; the code $`_m`$ is perfect.
Basic results:
The Hamming weight enumerators of $`C`$ and its dual are related by the following equation.
###### Theorem 1 (generalized Mac-Williams identity (cf. \[Del73\]))
$$W_C^{}(u,v)=\frac{1}{|C|}W_C(u+3v,uv).$$
Analogues: B: cf. \[MS77\]; L: cf. \[Ser73\], Ch. VII, Prop. 16; V: \[Zhu90, Höh95\].
Proof: For a function $`f`$ on $`K^n`$ with values in a ring $`R`$ we define its transformation $`g:K^nR`$ by $`g(𝐱)=_{𝐲K^n}f(𝐲)(1)^{(𝐲,𝐱)}`$. One has the following identity:
$$\frac{1}{|C|}\underset{𝐱C}{}g(𝐱)=\underset{𝐲C^{}}{}f(𝐲).$$
(1)
Proof of (1): $`_{𝐱C}g(𝐱)=_{𝐲K^n}_{𝐱C}f(𝐲)(1)^{(𝐱,𝐲)}=|C|_{𝐲C^{}}f(𝐲)+_{𝐲C^{}}f(𝐲)_{𝐱C}(1)^{(𝐱,𝐲)}`$. We have to show that the second sum vanishes. To this end, choose for given $`𝐲K^nC^{}`$ a $`𝐱^{}C`$ with $`(𝐱^{},𝐲)0`$, i.e., $`(1)^{(𝐱^{},𝐲)}=1`$. We get $`s=_{𝐱C}(1)^{(𝐱,𝐲)}=_{𝐱C}(1)^{(𝐱,𝐲)+(𝐱^{},𝐲)}=s`$, which implies $`s=0`$ and proves (1).
Now let $`f(𝐲)=u^{n\mathrm{wt}(𝐲)}v^{\mathrm{wt}(𝐲)}`$. We obtain for its transformation
$`g(𝐱)`$ $`=`$ $`{\displaystyle \underset{𝐲K^n}{}}f(𝐲)(1)^{(𝐱,𝐲)}`$
$`=`$ $`{\displaystyle \underset{y_1,\mathrm{},y_nK}{}}u^{n\mathrm{wt}(y_1)\mathrm{}\mathrm{wt}(y_n)}v^{\mathrm{wt}(y_1)+\mathrm{}+\mathrm{wt}(y_n)}(1)^{x_1y_1+\mathrm{}+x_ny_n}`$
$`=`$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}\left({\displaystyle \underset{zK}{}}u^{1\mathrm{wt}(z)}v^{\mathrm{wt}(z)}(1)^{x_iz}\right)`$
$`=`$ $`(u+3v)^{n\mathrm{wt}(𝐱)}(uv)^{\mathrm{wt}(𝐱)}.`$
Applying (1) we get for the weight enumerator of $`C^{}`$:
$$W_C^{}(u,v)=\underset{𝐲C^{}}{}f(𝐲)=\frac{1}{|C|}\underset{𝐱C}{}(u+3v)^{n\mathrm{wt}(𝐱)}(uv)^{\mathrm{wt}(𝐱)}=\frac{1}{|C|}W_C(u+3v,uv).$$
x
For the other types of weight enumerators we stay only the results, the proofs are similar.
###### Theorem 2 (Mac-Williams identity for complete weight enumerators)
$$\mathrm{cwe}_C^{}(p,q,r,s)=\frac{1}{|C|}\mathrm{cwe}_C(p+q+r+s,p+qrs,pq+rs,pqr+s).$$
From Theorem 1, we get the following descriptions of the weight enumerators of self-dual codes:
###### Theorem 3
Let $`C`$ be a self-dual $`[n,n/2]`$-code. Then, the weight enumerator $`W_C(u,v)`$ is a weighted homogeneous polynomial of weight $`n`$ in $`u+v`$ and $`v(uv)`$, or equivalently in the weight enumerators of $`\gamma _1`$ and $`ϵ_2`$.
Analogues: B: \[Gle71\]; L: cf. \[CS93b\]; V: \[Höh95\], Ch. 2.
Proof: From Theorem 1, we see that $`W_C`$ is invariant under the group $`H𝐙_2`$ generated by the substitution $`\frac{1}{2}\left(\genfrac{}{}{0pt}{}{13}{11}\right)`$. The ring of invariants has Molien series $`1/\left((1\lambda )(1\lambda ^2)\right)`$. (This is the generating function for the multiplicities of the trivial $`H`$-representation in the symmetric powers of the defining two dimensional representation of $`H`$.) The polynomials $`u+v`$ and $`v(uv)`$ or equivalently $`W_{\gamma _1}`$ and $`W_{ϵ_2}`$ are algebraically independent and generate freely the ring of all invariants. x
###### Theorem 4
Let $`C`$ be an even self-dual $`[n,n/2]`$-code. Then, the weight enumerator $`W_C(u,v)`$ is a weighted homogeneous polynomial of weight $`n`$ in $`u^2+3v^2`$ and $`v^2(u^2v^2)^2`$, or equivalently in the weight enumerators of $`ϵ_2`$ and $`𝒞_6`$.
Analogues: B: \[Gle71\]; L: cf. \[CS93b\]; V: see \[God89\] and \[Höh95\], Ch. 2.
Proof: This follows from the corresponding result for even self-dual codes over $`𝐅_4`$ as proven for example in \[MOSW78\], Th. 13: The group generated by $`S=\frac{1}{2}\left(\genfrac{}{}{0pt}{}{13}{11}\right)`$ and $`T=\left(\genfrac{}{}{0pt}{}{10}{01}\right)`$ has order $`12`$ and the Molien series of the corresponding ring of invariants is $`1/\left((1\lambda ^2)(1\lambda ^6)\right)`$. neuezeile x
## 3 Classification of self-dual codes
Let $`\delta _n`$ be the code consisting of all codewords containing only $`0`$’s and an even number of $`a`$’s. This is the even subcode of $`\gamma _1^n`$. One has $`\mathrm{dim}(\delta _n)=(n1)/2`$, coset representatives of $`\delta _n^{}/\delta _n`$ are given by $`(0^n)`$, $`(a,0^{n1})`$, $`(b^n)`$ and $`(c,b^{n1})`$ and its automorphism group consists for $`n2`$ of the permutation of the positions together with possible interchanging $`b`$ and $`c`$ at every position, i.e., $`\mathrm{Aut}(\delta _n)=S_2^n:S_n`$.
The next theorem describes self-orthogonal codes spanned by vectors of small weight.
###### Theorem 5
Minimal weight $`1`$ subcodes of a self-orthogonal code $`C`$ can be split off: $`CD\gamma _1^l`$, with minimal weight of $`D`$ larger then $`1`$. Self-orthogonal codes generated by weight-$`2`$-vectors are equivalent to direct sums of $`\delta _l`$, $`l2`$, and $`ϵ_2`$.
Analogues: First part: B, L: easy to see; V: cf. \[God89\], \[Höh95\], Th. 2.2.8. Second Part: B: \[PS75\], Th. 6.5; L: cf. \[CS93b\], Ch. 4; V: Cartan, Killing, \[Kac89\], \[FZ92\].
Proof: For the first statement, note that a weight-$`1`$-codeword is equivalent to $`(0,\mathrm{},0,a)`$. Then $`C`$=$`C^{}\gamma _1`$, where $`C^{}`$ is the orthogonal complement in $`C`$ of the $`\gamma _1`$ spanned by $`(0,\mathrm{},0,a)`$.
For the proof of the second statement, decompose first the code generated by the weight-$`2`$-codewords into the direct sum of its indecomposable even components and fix one of them. We have two possibilities:
Case i) There are two weight-$`2`$-codewords containing different nonzero entries at the same position.
In this case the component is equivalent to a code containing the two codewords $`(aa0\mathrm{}0)`$ and $`(bb0\mathrm{}0)`$. They generate a $`ϵ_2`$ subcode and, since $`ϵ_2`$ is self-dual, this is the whole component. (The other possible pairs of weight-$`2`$-codewords are not orthogonal.)
Case ii) The component is equivalent to a code whose weight-$`2`$-codewords have at all positions the value $`0`$ or $`a`$.
Inductively, one sees that the component is equivalent to a $`\delta _l`$, $`l2`$. A possible set of generators is given by $`(aa0\mathrm{}0)`$, $`(0aa0\mathrm{})`$, $`\mathrm{}`$, $`(0\mathrm{}0aa)`$. x
Let $`\overline{C}`$ the subcode of $`C`$ generated by the weight $`1`$ and $`2`$ codewords. We can describe $`C`$ by its gluecode $`\mathrm{\Lambda }\overline{C}^{}/\overline{C}`$. The automorphism group of $`C`$ is given by $`\mathrm{Aut}(C)=G_0.G_1.G_2`$, where $`G_0`$ are the “inner automorphisms” of $`\overline{C}`$, i.e. those which are fixing the components of $`\overline{C}`$ and the cosets $`\mathrm{\Lambda }/\overline{C}`$, $`G_1`$ are the automorphisms of $`C`$ fixing the components of $`\overline{C}`$ modulo $`G_0`$ and $`G_2`$ is the induced permutation group on the components of $`\overline{C}`$.
Denote by $`M(n)`$ resp. $`M_e(n)`$ the number of distinct (but maybe equivalent) self-dual resp. even self-dual Kleinian codes of length $`n`$.
###### Theorem 6 (Massformula)
The mass constants are given by
$$M(n)=\underset{i=1}{\overset{n}{}}(2^i+1)=\underset{C}{}\frac{6^nn!}{|\mathrm{Aut}(C)|}$$
where the sum is over equivalence classes of self-dual codes and
$$M_e(n)=\underset{i=0}{\overset{n1}{}}(2^i+1)=\underset{C}{}\frac{6^nn!}{|\mathrm{Aut}(C)|}$$
where the sum is over equivalence classes of even self-dual codes and $`n`$ is even.
Analogues: B: cf. \[PS75\]; L: \[Min84\]; V: unknown.
Proof: First, we prove the formula for $`M(n)`$. Let $`M(n,k)`$ be the number of self-orthogonal codes of dimension $`k`$ and length $`n`$. There are $`|(C^{}C)/C|=4^{n2k}1`$ different extensions of a self-orthogonal $`[n,k]`$-code $`C`$ to a self-orthogonal $`[n,k+\frac{1}{2}]`$-code $`DC`$ by choosing one extra vector $`xC^{}`$. Every self-orthogonal $`[n,k+\frac{1}{2}]`$-code $`D`$ arises from $`|D\{0\}|=4^{k+1/2}1`$ different codes $`C`$. So we get the recursion
$$M(n,k+\frac{1}{2})=M(n,k)\frac{4^{n2k}1}{4^{k+1/2}1}.$$
Together with $`M(n,0)=1`$ we obtain
$$M(n)=M(n,n/2)=\underset{i=0}{\overset{n1}{}}\frac{4^{ni}1}{2^{i+1}1}=\underset{i=1}{\overset{n}{}}(2^i+1).$$
The second expression for $`M(n)`$ describes the decomposition of all self-dual codes into orbits under the action of $`S_3^n:S_n`$.
To get the mass formula for $`M_e(n)`$, define in a similar way as before $`M_e(n,k)`$ as the number of even self-orthogonal codes of dimension $`k`$ and length $`n`$. The dual code $`C^{}`$ of a even self-orthogonal $`[n,k]`$-code $`C`$ contains $`\frac{1}{2}(4^{nk}+(2)^n)`$ vectors of even weight as one can see from Theorem 1. All vectors in a coset $`C^{}/C`$ have the same weight modulo $`2`$. So we get in a similar way as above the recursion
$$M_e(n,k+\frac{1}{2})=M_e(n,k)\frac{\frac{1}{2}(4^{n2k}+2^{n2k})1}{4^{k+1/2}1}.$$
Starting from $`M_e(n,0)=1`$ we obtain
$$M_e(n)=M_e(n,n/2)=\underset{i=0}{\overset{n1}{}}\frac{2^{2n2i1}+2^{ni1}1}{2^{i+1}1}=\underset{i=0}{\overset{n1}{}}(2^i+1)$$
and again we can express the total number as a sum over the different equivalence classes of codes. x
For the weighted sum of the Hamming weight enumerators one has
###### Theorem 7 (Massformula for Hamming weight enumerators)
$$\underset{C}{}\frac{6^nn!}{|\mathrm{Aut}(C)|}W_C(u,v)=M(n)(1+2^n)^1\left[2^nu^n+(u+3v)^n\right]$$
where the sum is over equivalence classes of self-dual codes.
$$\underset{C}{}\frac{6^nn!}{|\mathrm{Aut}(C)|}W_C(u,v)=M_e(n)(1+2^{n1})^1\left[2^{n1}u^n+\frac{1}{2}\left\{(u+3v)^n+(u3v)^n\right\}\right]$$
where the sum is over equivalence classes of even self-dual codes.
Analogues: B: \[PS75\]; L: \[Sie35\]; V: unknown.
Proof: Let $`𝐱`$ be a nonzero vector (of even weight) of length $`n`$. Similar as in the proof of Theorem 6 one gets for the number of (even) self-dual codes containing $`𝐱`$ the expression
$$\underset{i=1}{\overset{n1}{}}(2^i+1)\text{or }\underset{i=0}{\overset{n2}{}}(2^i+1)\text{ for even codes.}$$
From this and Theorem 6 one obtains the result by summing
$$u^{n\mathrm{wt}(𝐱)}v^{\mathrm{wt}(𝐱)}$$
over all pairs $`(𝐱,C)`$, where $`C`$ is a (even) self-dual code with $`𝐱C`$, and expanding the resulting sum in two different ways. x
We remark, that the average Hamming weight enumerator for even self-dual Kleinian codes is the same as for even formal self-dual $`𝐅_4`$-codes (\[MOSW78\], Th. 24) although the mass constants are different.
We call a self-dual code primitive, if no $`\gamma _1`$ subcode can be split off. A primitive code $`C`$ is the first one in the chain $`C`$, $`C\gamma _1`$, $`\mathrm{}`$
###### Theorem 8 (Relation between even and odd self-dual codes)
There is a $`1:1`$-correspondence between isomorphism classes of pairs $`(C,\delta _k)`$, where $`C`$ is an even self-dual code of even length $`n`$ and $`\delta _k`$ a subcode (a defined above) inside $`C`$ (together with the choice of a class $`[x]`$ in $`\delta _k^{}/\delta _k`$ of minimal weight $`1`$, i.e., if $`k=1`$ we must choose $`xK\{0\}`$) and isomorphism classes of self-dual codes $`D`$ of length $`nk`$.
The code $`D`$ is primitive if and only if the subcode $`\delta _k`$ is maximal, i.e. not contained in a $`\delta _{k+1}`$ subcode (with corresponding gluevectors $`[x]`$).
Analogues: B: \[CP80\]; L: \[CS82b\]; V: \[Höh95\], Ch. 3, and \[Höha\].
Proof: We describe the map from self-dual codes $`D`$ of length $`nk`$ to even self-dual codes of even length $`n`$. Denote by $`\delta _k^0=\delta _k`$, $`\delta _k^1`$, $`\delta _k^2`$ and $`\delta _k^3`$ the four cosets of $`\delta _k`$ inside $`\delta _k^{}`$ such that $`(a0^{k1})\delta _k^1`$.
If $`D`$ is even, let $`C=D(\delta _k^0\delta _k^2)`$. Otherwise we have the decomposition $`D_0^{}=D_0D_1D_2D_3`$ of the orthogonal complement of the even subcode $`D_0`$ of $`D=D_0D_1`$ into four $`D_0`$ cosets. Define
$$C=D_0\delta _k^0D_1\delta _k^1D_2\delta _k^2D_3\delta _k^3.$$
Note that for $`k=1`$ the three cosets $`\delta _k^1`$, $`\delta _k^2`$ and $`\delta _k^3`$ are all equivalent under $`\mathrm{Aut}(\delta _1)=S_3`$. It is then easy to check that this map describes the claimed $`1:1`$-correspondence. x
We call $`D`$ a child of the parent code $`C`$. From Theorem 8, we get the following description of the primitive children of an even self-dual code $`C`$ of length $`n`$: Take a position and choose $`x\{a,b,c\}`$ (up to the action of $`\mathrm{Aut}(C)`$), this gives a self-dual code $`D`$ of length $`n1`$.
\- If the position is not in the support of the subcode $`\overline{C}`$ generated by the weight-$`2`$-codewords, the code $`D`$ is primitive.
\- If the position is in an $`\delta _l`$, $`l2`$, component of $`\overline{C}`$ we have two cases: If $`xa`$ then $`D`$ is again maximal, if $`x=a`$ the primitive child is obtained by deleting the remaining $`l1`$ positions of $`\delta _l`$ from $`D`$.
\- If the position is in a $`ϵ_2`$ component, the primitive child is obtained by deleting the second position of $`ϵ_2`$ from $`D`$.
Every non even self-dual code $`D=D_0D_1`$ of even length $`n`$ determines the two even self-dual “neighbours” $`D_0D_2`$ and $`D_0D_3`$, where $`D_0`$, $`D_1`$, $`D_2`$ and $`D_3`$ are the four cosets of $`D_0`$ in $`D_0^{}`$ as above. We define for every even $`n`$ a “neighbourhood graph” by using the isomorphism classes of even self-dual codes as vertices, the isomorphism classes of non even self-dual codes as edges and “neighbourhood” as incidence relation. An edge corresponding to a non primitive code $`D=D^{}\gamma _1^l`$, $`l1`$, is a loop for the vertex corresponding to the even code determined from $`D^{}`$ through Theorem 8. The edges starting on a vertex $`C`$ correspond to the orbits of $`\mathrm{Aut}(C)`$ on the nonzero elements of $`K^n/C`$. It is easy to see that the neighbourhood graph is connected for all $`n`$. For $`n=2`$, $`4`$ and $`6`$ the graph is shown in Figure 1.
Analogues: L: \[Bor84\]; V: not determined.
###### Theorem 9
The even self-dual codes up to length $`8`$ (together with the subcode $`\overline{C}`$, order of $`G_1.G_2`$, weight enumerator and number of children) are given in Table 1.
Analogues: B: \[CP80, CPS92\]; L: \[Kne57, Nie73\]; V: cf. \[God89\], for $`c=24`$ there is a conjectured list in \[Sch93\].
Proof: Use the list of doubly even self-dual binary codes of length $`4n`$ \[CP80, CPS92\] and the construction A described in Section 7 or use Theorem 5 and classify the possibilities for $`\overline{C}`$ and the gluecodes $`\mathrm{\Lambda }\overline{C}^{}/\overline{C}`$ directly. x
We checked the result additionally with the mass formula for the Hamming weight enumerator.
###### Theorem 10
The non even self-dual codes up to length $`6`$ (together with the parent No., the subcode $`\overline{C}`$, order of $`G_1.G_2`$ and the weight enumerator) are given in Table 2.
Analogues: B: \[Ple72, PS75\]; L: \[CS82b, Bor93\]; V: \[Höh95\], Ch. 3, and \[Höha\].
Proof: Look at the list of even self-dual binary codes of length $`4n`$ \[Ple72, PS75\] or apply Theorem 8 to Theorem 9. x
Again we checked the result by the mass formula for the Hamming weight enumerator.
Remark: There is one self-dual code of length $`5`$ without codewords of weight $`2`$: The shorter Hexacode $`𝒞_5`$. There are two self-dual codes of length $`6`$ without codewords of weight $`2`$: The Hexacode $`𝒞_6`$ (even) and the odd Hexacode $`O_6`$ (non even).
The number of inequivalent (even) self-dual codes of small length $`n`$ can be read off from Table 3.
The number of even codes up to length $`8`$ are obtained from Theorem 9, the number of odd codes up to length $`6`$ from Theorem 10 and 9 and for $`n=7`$ it follows from the number of length $`7`$ children of the even length $`8`$ codes. The lower estimates for larger $`n`$ one obtains from the mass formula.
A complete classification up to $`n=10`$ seems possible, but no interesting new structure is expected.
All the self-dual Kleinian codes classified in this section have a nontrivial automorphism group. In analogy to \[OP92, Ban88\], we expect that this holds only for small length $`n`$ and that rather almost all self-dual and even self-dual codes have trivial automorphism group. What are the smallest (even) self-dual codes with trivial automorphism groups (cf. \[Bac94\] for lattices)?
## 4 Extremal codes
In this section, we study self-dual Kleinian codes of type $`[n,n/2,d]`$ where $`d`$ is as large as possible. Let $`m=[n/2]`$. By Theorem 3 the weight enumerator of a code $`C`$ can be written as
$$W_C(u,v)=\underset{i=0}{\overset{m}{}}a_i(u+v)^{n2i}(v(uv))^i$$
(2)
with unique integral numbers $`a_i`$. There is a unique choice of the numbers $`a_0`$, $`\mathrm{}`$, $`a_m`$ such that the right hand side of (2) equals
$$u^n+0u^{n1}v+\mathrm{}+0u^{nm}v^m+A_{m+1}u^{nm1}v^{m+1}+\mathrm{}+A_nv^n.$$
(3)
We call (3) the extremal weight enumerator and a code with this weight enumerator extremal. So an extremal code has minimal weight $`d[n/2]+1`$.
###### Theorem 11
The minimal distance $`d`$ of a self-dual code $`C`$ of length $`n`$ satisfies
$$d\left[\frac{n}{2}\right]+1.$$
Analogues: B: \[MS73\]; L: \[Sie69\]; V: \[Höh95\], Cor. 5.3.3.
Proof: The proof is parallel to \[MS73\], Cor. 3. In fact it can be considered as “case 5”“Case 4” was defined in \[MOSW78\]. of that paper for the parameters $`w=1`$, $`R=2`$, $`S=1`$ and $`\alpha =1`$. It follows also from the next theorem. x
Let $`C_0`$ be the even subcode of $`C`$ as in the proof of Theorem 8. To study extremal codes in more detail, we need the definition of the shadow $`C^{}`$ of $`C`$: We set $`C^{}=C_0^{}C`$ if $`C`$ is not even and $`C^{}=C`$ otherwise.
###### Lemma 1
If the weight enumerator of $`C`$ is written as
$$W_C(u,v)=P_C(W_{\gamma _1},W_{ϵ_2})=Q_C(W_{\gamma _1},W_{ϵ_2}W_{\gamma _1^2})$$
with weighted homogeneous polynomials $`P_C(x,y)`$ and $`Q_C(x,y)`$, then for the shadow one has
$$W_C^{}(u,v)=P_C(W_{\gamma _1^{}},W_{ϵ_2})=Q_C(W_{\gamma _1^{}},W_{ϵ_2}W_{\gamma _{1}^{}{}_{}{}^{2}}).$$
Proof: We show $`W_C^{}(u,v)=\frac{1}{|C|}W_C(u+3v,(1)(uv))`$ from which the lemma follows.
If $`C=C^{}`$ is even this is Theorem 1. Otherwise, we get from there
$`W_C^{}(u,v)`$ $`=`$ $`W_{C_0^{}}(u,v)W_C(u,v)={\displaystyle \frac{2}{|C|}}W_{C_0}(u+3v,uv)W_C(u,v)`$
$`=`$ $`{\displaystyle \frac{1}{|C|}}\left[W_C(u+3v,uv)+W_C(u+3v,(1)(uv))\right]W_C(u,v)`$
$`=`$ $`{\displaystyle \frac{1}{|C|}}W_C(u+3v,(1)(uv)).\text{x}`$
###### Theorem 12
There are exactly five extremal codes: $`\gamma _1`$, $`ϵ_2`$, $`\delta _3^+`$, the shorter Hexacode $`𝒞_5`$ and the Hexacode $`𝒞_6`$.
Analogues: B: \[MS73, War76\]; L: \[COS78\]; V: \[Höh95\], Th. 5.3.2.
For the corresponding extremal weight enumerators see Table 1 and 2.
Proof: The existence and uniqueness of an extremal code for $`n=1`$, $`2`$, $`3`$, $`5`$ and $`6`$ can directly be read off from Table 1 and 2.
The nonexistence for $`n=4`$ follows also from this tables, so we must prove the nonexistence for $`n>6`$. We can assume $`C`$ is non even since for an even code we will show (Theorem 15) that for the minimal weight $`d`$ one has $`d2[n/6]+2`$. But from $`d[n/2]+1`$, we get $`n=2`$ or $`6`$. Now we are using the shadow $`C^{}`$ of $`C`$. From Lemma 1, we get for its weight enumerator for $`n=7`$, $`8`$, $`\mathrm{}`$, $`11`$:
$$\begin{array}{cccccc}n\hfill & 7& 8& 9& 10& 11\\ & & & & & \\ W_C^{}(u,v)\hfill & \frac{7}{4}u^6v+\mathrm{}& \frac{13}{8}u^8+\mathrm{}& \frac{9}{4}u^8v+\mathrm{}& \frac{23}{8}u^{10}+\mathrm{}& \frac{33}{8}u^{10}v+\mathrm{}\end{array}.$$
Since $`W_C^{}`$ must have non negative integral coefficients, there exists no extremal codes for $`7n11`$. For $`n12`$, the coefficient $`A_{m+2}`$ of $`W_C(u,v)`$ is always negative. We will sketch the proof:
Let $`m=[n/2]`$ and replace $`u`$ by $`1`$. Expanding $`(1+v)^n`$ in powers of $`\varphi =\frac{v(1v)}{(1+v)^2}`$ one gets by the Bürmann Lagrange Theorem
$$(1+v)^n=\underset{k=0}{\overset{m}{}}b_k\varphi ^k+\underset{k=m+1}{\overset{\mathrm{}}{}}b_k\varphi ^k$$
(4)
with
$$b_k=\frac{1}{k!}\frac{d^{k1}}{dv^k}\left[\frac{d(1+v)^n}{dv}\left(\frac{v}{\varphi }\right)^k\right]|_{v=0}.$$
Comparing expansion (4) with (2) and (3) yields $`b_k=a_k`$ for $`k=0`$, $`\mathrm{}`$, $`m`$. Furthermore, $`A_{m+1}=b_{m+1}`$, $`A_{m+2}=b_{m+2}+3(m+1)b_{m+1}n`$. Now one estimates with the saddle-point method $`b_{m+1}`$ and $`b_{m+2}`$ and shows that $`A_{m+2}<0`$ for $`m`$ large enough. The smaller $`n`$ are checked by a direct computation. x
Remarks: Similar as in \[CS90c, CS90b, CS91\] one can refine the bound of Theorem 11 to obtain $`d2[n/5]+\mathrm{O}(1)`$ by using the shadow code.
For the difference $`D_C(u,v)=W_{C_2}(u,v)W_{C_3}(u,v)`$ one has the result
$$D_C(u,v)\{\begin{array}{cc}𝐐[W_ϵ,W_{𝒞_6}],\hfill & \text{if }n\text{ is even,}\hfill \\ v(u^2v^2)𝐐[W_ϵ,W_{𝒞_6}],\hfill & \text{if }n\text{ is odd.}\hfill \end{array}$$
This result can be used as in \[CS90c, CS90b, CS91\] to discuss for small $`n`$ the “weakly” extremal codes meeting the stronger bound for $`d`$. As an example, for $`n=5`$ we obtain $`D_C(u,v)=cv(u^2v^2)(u^2+3v^2)`$.
Instead of looking for codes with large minimal weight, one can ask the same question for the shadow itself. For self-dual codes with shadows of large minimal weight one gets similar results as recently described by N. Elkies and the author:
###### Theorem 13
The minimal weight $`h`$ of the shadow $`C^{}`$ of a self-dual code $`C`$ of length $`n`$ satisfies $`hn`$, with equality if and only if $`C\gamma _1^n`$.
Analogues: B: \[Elk95b\]; L: \[Elk95a\]; V: \[Höh97\], Th. 1.
Proof: Clearly $`hn`$. By Lemma 1, the weight enumerator of $`C^{}`$ is a polynomial $`P_C(W_{\gamma _1^{}},W_{ϵ_2})`$ in the weight enumerators of $`\gamma _1^{}`$ and $`ϵ_2`$, i.e. $`W_C^{}(u,v)`$ is a homogeneous polynomial of weight $`n`$ in $`2v`$ and $`u^2+3v^2`$. So $`h=n`$ implies $`W_C^{}(u,v)=(2v)^n`$; but then $`W_C(u,v)=(u+v)^n`$ and $`C\gamma _1^n`$. x
###### Theorem 14
Let $`C`$ be a self-dual code of length $`n`$ without words of weight $`1`$. Then one has
$`C`$ hat at least $`(n/2)(5n)`$ codewords of weight $`2`$.
The equality holds if and only if $`h(C^{})=n2`$.
In this case the number of codewords of weight $`n2`$ in the shadow is $`2^{n3}n`$.
Analogues: B, L: \[Elk95b\]; V: \[Höh97\], Th. 2.
Proof: Assume first $`h(C^{})n2`$. In the same way as in the proof of Theorem 13 we see that $`P_C(x,y)`$ is a linear combination of $`x^n`$ and $`x^{n2}y`$ and we obtain
$`W_C(u,v)`$ $`=`$ $`(u+v)^n{\displaystyle \frac{n}{2}}(u+v)^{n2}\left((u+v)^2(u^2+3v^2)\right)`$ (5)
$`=`$ $`u^n+0u^{n1}v+{\displaystyle \frac{n}{2}}(5n)u^{n2}v^2+\mathrm{}.`$ (6)
This proves one direction of ii).
Conversely, we can assume $`n<6`$, so the weight enumerator of $`C`$ can be written as
$`W_C(u,v)`$ $`=`$ $`(u+v)^n{\displaystyle \frac{n}{2}}(u+v)^{n2}\left(2uv2v^2\right)+{\displaystyle \frac{A_2(n/2)(5n)}{4}}(u+v)^{n4}\left(2uv2v^2\right)^2.`$
From Lemma 1, we get $`A_2(n/2)(5n)0`$ since $`W_C^{}(u,v)`$ has nonnegative coefficients, and we have i) and the converse of of ii).
Finally, Part iii) follows also from (5) and Lemma 1:
$`W_C^{}(u,v)`$ $`=`$ $`(2v)^n{\displaystyle \frac{n}{2}}(2v)^{n2}\left((2v)^2(u^2+3v^2)\right)`$
$`=`$ $`2^{n3}nu^2v^{n2}+(2^nn\mathrm{\hspace{0.17em}2}^{n3})v^n.\text{x}`$
There are exactly four such codes meeting the bound $`h(C^{})=n2`$, namely $`ϵ_2`$, $`\delta _3^+`$, $`(\delta _2^2)^+`$ and $`𝒞_5`$.
For even codes there are similar definitions and results. The following result was proven for $`𝐅_4`$-codes, but since its proof uses only Theorem 4 it is also true for Kleinian codes.
###### Theorem 15 (see \[MOSW78\])
The minimal distance $`d`$ of an even self-dual code $`C`$ of length $`n`$ satisfies
$$d2\left[\frac{n}{6}\right]+2.$$
Analogues: B: \[MOS75\]; L: \[MOS75\]; V: \[Höh95\], Section 5.2.
Remark: The analogous bound for doubly-even binary codes has recently been improved in \[KL97, KL00\] for large lengths.
An even self-dual code matching this bound is called extremal. The corresponding weight enumerator is called the extremal weight enumerator of length $`n`$. A table of extremal weight enumerators was given in \[MOSW78\], Table 1.
Again from the $`𝐅_4`$ case, the next result follows.
###### Theorem 16 (see \[MOSW78\])
There are no extremal even codes of length $`n136`$.
Analogues: B: \[MOS75\]; L: \[MOS75\]; V: no known bound, cf. \[Höh95\], Section 5.2.
Examples of extremal $`𝐅_4`$-codes are known for $`n=2`$ ($`ϵ_2`$), $`4`$ ($`ϵ_2^2`$), $`6`$ ($`𝒞_6`$), $`8`$ ($`3`$ codes), $`10`$, $`14`$, $`\mathrm{}`$, $`22`$, $`28`$ and $`30`$ (see \[CPS79\]). They are also examples of extremal even Kleinian codes.
There is no extremal $`𝐅_4`$-code of length 12. But there is an extremal even Kleinian code of this length with generator matrix
$$\left(\begin{array}{cc}\mathrm{𝚊𝚊𝚊𝚊𝚊𝚊}\hfill & \mathrm{𝟶𝟶𝟶𝟶𝟶𝟶}\hfill \\ \mathrm{𝚋𝚋𝚋𝚋𝚋𝚋}\hfill & \mathrm{𝟶𝟶𝟶𝟶𝟶𝟶}\hfill \\ \mathrm{𝟶𝟶𝟶𝟶𝟶𝟶}\hfill & \mathrm{𝚊𝚊𝚊𝚊𝚊𝚊}\hfill \\ \mathrm{𝟶𝟶𝟶𝟶𝟶𝟶}\hfill & \mathrm{𝚋𝚋𝚋𝚋𝚋𝚋}\hfill \\ \mathrm{𝚊𝟶𝚋𝚊𝚋𝟶}\hfill & \mathrm{𝚊𝚊𝚊𝚊𝟶𝟶}\hfill \\ \mathrm{𝚊𝚋𝚌𝚌𝚋𝚊}\hfill & \mathrm{𝚋𝚋𝚋𝚋𝟶𝟶}\hfill \\ \mathrm{𝚌𝚊𝚌𝚊𝟶𝟶}\hfill & \mathrm{𝚊𝟶𝚊𝚊𝚊𝟶}\hfill \\ \mathrm{𝚌𝚌𝚊𝟶𝚊𝟶}\hfill & \mathrm{𝚋𝟶𝚋𝚋𝚋𝟶}\hfill \\ \mathrm{𝚌𝚌𝚋𝚊𝚊𝚋}\hfill & \mathrm{𝚊𝟶𝟶𝚊𝚊𝚊}\hfill \\ \mathrm{𝚋𝚌𝚌𝚋𝚊𝚊}\hfill & \mathrm{𝚋𝟶𝟶𝚋𝚋𝚋}\hfill \\ \mathrm{𝚌𝚊𝚊𝚋𝚌𝚋}\hfill & \mathrm{𝚊𝚊𝟶𝟶𝚊𝚊}\hfill \\ \mathrm{𝚋𝟶𝚋𝚊𝚊𝟶}\hfill & \mathrm{𝚋𝚋𝟶𝟶𝚋𝚋}\hfill \end{array}\right)$$
and weight enumerator $`W_C(u,v)=u^{12}+396u^6v^6+1485u^4v^8+1980u^2v^{10}+234v^{12}`$.
Besides the question of the existence of a projective plane of order ten and of a doubly even code of type $`[72,36,16]`$, a $`[24,12,10]`$ self-dual $`𝐅_4`$-code was most wanted. After the first question, also the third question has found a negative answer \[LP90\]. Since Kleinian codes are combinatorial more natural than $`𝐅_4`$-codes, we ask if there is an even self-dual Kleinian code of type $`[24,12,10]`$. This is the smallest open case for extremal even Kleinian codes.
Good even and doubly even self-dual binary codes meeting the Gilbert-Varshamov bound exist, as was shown by using the mass formula for the Hamming weight enumerator \[MST72\]. A similar result holds for lattices (see \[MH73\], Ch. II). We expect the same for self-dual and even self-dual Kleinian codes.
## 5 Constant weight codes and generalized $`t`$-designs
Let $`X_k`$ be the fiber over $`k`$ of the weight map $`\mathrm{wt}:K^n\{0,1,\mathrm{},n\}`$. We can write it as the (not two point) homogenous space $`X_k=G/H=S_3^n:S_n/(S_2^k:S_k\times S_3^{nk}:S_{nk})`$. The $`H`$-module structure of the function space $`L_2(X_k)`$ for general alphabets instead of $`K`$ has been studied in \[Dun76\]. The space $`X_k`$ carries the structure of a symmetric association scheme, called the nonbinary Johnson scheme (cf. \[TAG85\]) as follows: A pair $`(𝐱,𝐲)X_k\times X_k`$ belongs to the relation $`R_{r,s}`$, with $`r`$, $`s\{0,1,\mathrm{},k\}`$, $`rs`$, if $`r=\mathrm{\#}\{ix_i=y_i0\}`$ and $`s=\mathrm{\#}\{ix_i0,y_i0\}`$. This structures allow one to use the usual association scheme methods to study subsets $`YX_k`$ (cf. \[DL98\]). I like to thank C. Bachoc for mentioning the references \[Dun76, TAG85, DL98, Del73, Bac99\] to me.
Here, we use the definition of a generalized $`t`$-designs as in \[Del73\]: An element $`𝐱K^n`$ is said to be covered by an element $`𝐲K^n`$ if each nonzero component $`x_i`$ of $`𝐱`$ is equal to the corresponding component $`y_i`$ of $`𝐲`$. A generalized $`t(n,k,\mu )`$ design (of type $`3`$) is a nonempty subset $`YX_k`$ such that any element of $`X_t`$ is covered by exactly $`\mu `$ elements from $`Y`$. For $`t=2`$, this definition is identical with the notion of a group divisible incomplete block design with $`n`$ groups of $`3`$ elements, blocksize $`k`$ and $`\lambda _1=0`$, $`\lambda _2=\mu `$ introduced in \[BN39\].
As an example, the three codewords of weight $`2`$ in $`ϵ_2`$ form a generalized $`1`$-$`(2,2,1)`$ design. The next result describes a method to obtain generalized $`2`$-designs.
###### Theorem 17
Let $`C`$ be an extremal even code of length $`n=6k`$. Then, the codewords of $`C`$ of fixed non-zero weight form a generalized $`2`$-design.
Analogues: B: \[AM69\]; L: \[Ven84\]; V: unknown.
Proof: This follows from Th. 5.3. in \[Del73\], a generalization of the Assmus and Mattson theorem: By Theorem 15, there are at most $`\frac{1}{2}\left(n(2(n/6)+2)\right)+1=2(n/6)`$ nonzero weights in such a code. Note, that our scalar product on $`K`$ defines a required identification map $`\chi _{(.)}:K\mathrm{Hom}(K,𝐂^{})`$. x
The result applies in particular to the unique extremal even code of length $`6`$, the Hexacode $`𝒞_6`$ and the extremal even code of length $`12`$ given in the last section. The generalized $`2`$-$`(6,4,2)`$ and $`2`$-$`(6,6,2)`$ designs formed by the vectors of the Hexacode of weight $`4`$ and $`6`$ are unique.
In this case, the design property can also be obtained from the following result about $`\mathrm{Aut}(𝒞_6)`$:
###### Theorem 18
The automorphism group of the Hexacode acts transitively on the weight $`2`$ vectors in $`K^6`$.
Analogues: B: \[Mat61, Car31\]; L: \[GS87, HS79\]; V: unknown.
Proof: By computing the double cosets $`(S_2^k:S_k\times S_3^{6k}:S_{6k})S_3^6/\mathrm{Aut}(𝒞_6)`$ for $`k=0`$, $`1`$, $`\mathrm{}`$, $`6`$, we get the orbit decomposition of $`K^6`$ under $`\mathrm{Aut}(𝒞_6)`$ as shown in Table 4.
There is only one orbit for $`k=2`$. x
This gives also the information about the structure of the deep holes and the cocode $`K^6/𝒞_6`$.
###### Theorem 19
The covering radius of the Hexacode $`𝒞_6`$ is $`2`$. There is one type of deep holes in $`K^6`$. Representatives are the $`\mathrm{Aut}(𝒞_6)`$-orbits $`A_2`$, $`B_3`$, $`C_4`$, $`D_4`$, $`C_5`$, $`C_6`$, $`D_6`$ and $`E_6`$. For every deep hole there are exactly three codewords with distance $`2`$. The three orbits $`A_0`$, $`A_1`$ and $`C_6`$ form a complete system of representatives for the cocode $`K^6/𝒞_6`$, representing the cosets of minimal weight $`0`$, $`1`$ and $`2`$, respectively.
Analogues: B: \[CS90d\]; L: partially \[CPS82, BCQ93, Bor98\]; V: unknown.
The $`135`$ deep holes of weight $`2`$ are partioned into $`45`$ sets of “trios”, the members of each trio are representing the same coset in $`K^6/𝒞_6`$. The subcode of $`𝒞_6`$ generated by pairs of members in a trio forms a frame which corresponds to a twisted construction of $`𝒞_6`$ from a $`D_8^{}/D_8`$-code (cf. the end of section 7).
From the next theorem, one deduces immediately that the $`18`$ vectors of weight $`6`$ in the Hexacode are the smallest possible number of elements necessary to form a generalized $`2`$-design with $`n=k=6`$.
###### Theorem 20 (Th. 5 and 6 in \[BC52\])
For the number of elements of a generalized $`2`$-$`(n,k,\lambda )`$ design $`Y`$ of type $`3`$ one has
$$|Y|\{\begin{array}{cc}3n,\hfill & \text{for }k<n\text{,}\hfill \\ 2n+1,\hfill & \text{for }k=n\text{.}\hfill \end{array}$$
Analogues: B: \[RCW75\], L: \[DGS77\], V: unknown.
The set of the $`45`$ weight $`4`$ vectors in the Hexacode has the smallest cardinality for a generalized $`2(6,4,\lambda )`$ design.
By taking the $`253`$ of the $`759`$ vectors of weight $`8`$ in the binary Golay code having first coordinate $`1`$, one gets the essentially only tight $`4`$-design \[Bre79\]. The $`196560`$ vectors of squared length $`4`$ in the Leech lattice form the only tight spherical $`11`$-design \[BD79, BD80\] in dimension greater then $`2`$. This leads to the question: Is there a good notion of tight generalized $`t`$-designs, using a bound generalizing Theorem 20 for its definition, characterizing one of the two designs belonging to the Hexacode?
## 6 Lexicographic codes
The lexicographic code of length $`n`$ and minimal distance $`d`$ is defined by the greedy algorithm: After writing down the elements of $`K^n`$ in lexicographic order one chooses in every step the lexicographic first word which has distance at least $`d`$ to the already chosen codewords.
###### Theorem 21 (Conway-Sloane \[CS86\])
The lexicographic code of length $`2`$ and minimal distance $`2`$ is $`ϵ_2`$. The lexicographic code of length $`6`$ and minimal distance $`4`$ is the Hexacode $`𝒞_6`$.
Analogues: B: \[CS86\]; L: \[CS82a\].
Define self-orthogonal lexicographic codes by restricting the choice of the next codeword to the dual code of the code spanned by the codewords already chosen. This is some analogy to the definition of integral laminated lattices.
###### Theorem 22
The self-orthogonal lexicographic codes with minimal distance $`1`$, $`2`$, $`3`$ and $`4`$ are “periodic” under direct sum. The periods are $`1`$, $`2`$, $`5`$ and $`6`$ with periodicity elements $`\gamma _1`$, $`ϵ_2`$, $`𝒞_5`$ and $`𝒞_6`$ respectively.
Analogues: B: $`c_2`$, $`H_8`$, $`g_{22}`$ and $`g_{24}`$ \[Mon96\]; L: $`𝐙`$, $`E_8`$, $`\mathrm{\Lambda }_{23}`$ and $`\mathrm{\Lambda }_{24}`$ \[PP85, CS83\]; V: $`V_{\mathrm{Fermi}}`$, $`V_{E_8}`$, $`VB^{\mathrm{}}`$ and $`V^{\mathrm{}}`$.
## 7 Relations to binary codes, lattices and vertex operator algebras
In this section, we assume that the reader is familiar with the notation of a vertex operator algebra (VOA) and a vertex operator super algebra (SVOA) (see \[FLM88, FHL93, Kac97\] for an introduction). All (S)VOA’s are assumed to be simple, unitary and “nice” (cf. \[Höh95\], Ch. 1).
All the definitions and results of this work have analogies for binary codes, lattices and VOA’s, although for VOA’s the theory is not completely developed. Analogously to the relation between binary codes and lattices and between lattices and VOA’s one has two constructions (an “untwisted” and a “twisted” one) for binary codes from Kleinian codes.
Construction A: Define a map $`\rho _A`$ from Kleinian codes of length $`n`$ to binary codes of length $`4n`$ by
$$\rho _A(C):=\widehat{C}+d_4^n,$$
where $`\widehat{}:K^n𝐅_2^{4n}`$ is the map induced from $`\widehat{}:K(D_4^{}/D_4)(D_2^{}/D_2)^2𝐅_2^4`$, $`0(0000)`$, $`a(1100)`$, $`b(1010)`$, $`c(0110)`$ and $`d_4^n=\{(0000),(1111)\}^n`$. So every codeword in $`C`$ is replaced with $`2^n`$ binary codewords in $`𝐅_2^{4n}`$.
Construction B: Assume $`n`$ is even. Then
$$\rho _B(C):=\widehat{C}+(d_4^n)_0\widehat{C}+(d_4^n)_0+\{\begin{array}{cc}(1000\mathrm{}\mathrm{1000\; 1000}),\hfill & \text{if }n0(mod4)\text{,}\hfill \\ (1000\mathrm{}\mathrm{1000\; 0111}),\hfill & \text{if }n2(mod4)\text{,}\hfill \end{array}$$
where $`\widehat{}:K^n𝐅_2^{4n}`$ is the map as defined before and $`(d_4^n)_0`$ is the subcode of $`d_4^n`$ consisting of vectors of weight divisible by $`8`$.
###### Lemma 2
If $`C`$ is a linear, self-dual resp. even Kleinian code then $`\rho _A(C)`$ is a linear, even self-dual resp. doubly even binary code. The same is true for $`\rho _B(C)`$ if the length is even. x
###### Lemma 3
For the weight enumerators one has:
$`W_{\rho _A(C)}(x,y)`$ $`=`$ $`W_C(x^4+y^4,2x^2y^2),`$
$`W_{\rho _B(C)}(x,y)`$ $`=`$ $`{\displaystyle \frac{1}{2}}W_C(x^4+y^4,2x^2y^2)+{\displaystyle \frac{1}{2}}(x^4y^4)^n+`$
$`{\displaystyle \frac{2^n}{2}}((x^3y+xy^3)^n+(1)^{n/2}(x^3yxy^3)^n).\text{x}`$
Analogues: B–L: see \[CS93b\], Ch. 7; L–V: cf. \[Höh95\], Ch. 1 and 5.
Remarks:
$`\rho _B(𝒞_6)`$ gives the Golay code. (This is the MOG-construction.)
If we denote the untwisted (twisted) construction from binary codes to lattices and from lattices to VOA’s also with $`\rho _A`$ resp. $`\rho _B`$ (cf. \[DGH98\]) then one has
$$\rho _X(\rho _Y(\rho _Z))=\rho _{\pi (X)}(\rho _{\pi (Y)}(\rho _{\pi (Z)})),\text{with }X\text{}Y\text{}Z\{A,B\}\text{ and }\pi S_3\text{.}$$
Markings and frames:
A marking for a code $`C`$ is the choice of a vector $`(K\{0\})^n`$. Table 4 shows that there exist $`5`$ inequivalent markings for the Hexacode.
For $`i=1`$, $`\mathrm{}`$, $`n`$ we define
$$I_i=\{\begin{array}{cc}\{(4i3,4i2),(4i1,4i)\},\hfill & \text{if }_i=a\text{,}\hfill \\ \{(4i3,4i1),(4i2,4i)\},\hfill & \text{if }_i=b\text{,}\hfill \\ \{(4i3,4i),(4i1,4i2)\},\hfill & \text{if }_i=c\text{.}\hfill \end{array}$$
Then $`I=_{i=1}^nI_i`$ is a marking for the binary code $`\rho _X(C)`$ as defined in \[DGH98\]. As described in \[DGH98\] one gets from $`I`$ a $`D_1`$-frame in $`\rho _X(\rho _Y(C))`$ (or equivalent a $`𝐙_4`$-code, cf. \[CS93a\]) and a Virasoro frame in $`\rho _X(\rho _Y(\rho _Z(C)))`$. Since $`\mathrm{Aut}(K^n)=S_3^n:S_n`$ acts transitively on $`(K\{0\})^n`$ we can assume $`=(aa\mathrm{}a)`$ by replacing $`C`$ with an equivalent code. For this standard marking we define the symmetrized (marked) weight enumerator $`\mathrm{swe}_C`$ as
$$\mathrm{swe}_C(U,V,W)=\mathrm{cwe}_C(U,V,W,W).$$
The symmetrized marked weight enumerator of the above marked binary code $`\rho _X(C)`$ as defined in \[DGH98\] can be obtained from $`\mathrm{swe}_C(U,V,W)`$:
###### Lemma 4
$`\mathrm{smwe}_{\rho _A(C)}(x,y,z)`$ $`=`$ $`\mathrm{swe}_C(x^2+y^2,2xy,2z^2),`$
$`\mathrm{smwe}_{\rho _B(C)}(x,y,z)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{swe}_C(x^2+y^2,2xy,2z^2)+{\displaystyle \frac{1}{2}}(x^2y^2)^n+`$
$`{\displaystyle \frac{1}{2}}2^n((x+y)^n+(1)^{n/2}(xy)^n)z^n.\text{x}`$
Analogues: B–L: \[DGH98\]; L–V: \[DGH98\].
We remark that the symmetrized marked weight enumerator of an even self-dual code belongs to a ring of polynomials with Molien series $`\left(1+\lambda ^4\right)/\left((1\lambda ^2)^2(1\lambda ^6)\right)`$ generated by $`p_2=x^2+2y^2+z^2`$, $`q_2=x^2+4yzz^2`$, $`p_4=x^4+8y^4+6x^2z^2+z^4`$, $`p_6=x^6+6x^2y^4+4y^6+24x^2y^3z+12x^2y^2z^2+6y^4z^2+8y^3z^3+3x^2z^4`$ subject to one relation for $`p_4^2`$.
Now, we describe how codes and lattices can be understood in terms of VOA’s. Let $`V`$ be a rational VOA whose intertwiner algebra is abelian, i.e. the set of irreducible $`V`$-modules form an abelian group $`G`$ under the fusion product (cf. \[DL93\]). The map $`\alpha :G𝐂^{}`$, $`Me^{2\pi ih(M)}`$, where $`h(M)`$ is the conformal weight of the $`V`$-module $`M`$ defines a quadratic form on $`G`$ and can be interpreted as an element of $`H^4(K(G,2),𝐂^{})`$; where $`K(G,2)`$ is the Eilenberg-MacLane space with $`\pi _2(K(G,2))G`$ (see \[Höhb\]). Another description is the following: The monodromy structure of the intertwiner operators of $`V`$ give rise to a three dimensional topological quantum field theory which is example I.1.7.2 of \[Tur94\].
The fusion algebra of $`V^n`$ is $`(V^n)𝐙[G^n]`$. A subgroup $`CG^n`$ is called an even self-orthogonal linear code if $`C`$ is an isotropic subspace of the quadratic space $`(G^n,\alpha ^n)`$. It is proven in \[Höhb\] that (simple) VOA-extensions $`W`$ of $`V^n`$ are in one to one correspondence with such codes $`C`$; in particular, $`W=_{\alpha C}M_\alpha `$ has a unique VOA-structure up to isomorphism extending the VOA-structure of $`V=M_0`$. The uniqueness follows from $`H^3(K(C,2),𝐂^{})=0`$. Similar remarks hold for odd self-orthogonal codes and SVOA’s.
As an example, let $`V`$ be the lattice-VOA $`V_L`$ belonging to an even integral positive definite lattice $`L`$ of rank $`n`$. In this case $`G=L^{}/L`$ with $`\alpha `$ induced from $`e^{2\pi i\frac{(.,.)}{2}}:𝐑^n𝐂^{}`$, where $`(.,.)`$ is the standard scalar product of $`𝐑^n`$. In fact, the triple $`(G,\alpha ,n)`$ is a complete invariant of the genus of $`L`$ (see \[Nik80\]).
Since the VOA belonging to the root lattice $`D_4`$ of $`\mathrm{Spin}(8)`$ has four irreducible modules with the conformal weights $`0`$ and three times $`\frac{1}{2}`$ and one has $`(V_{D_4})𝐙[𝐙_4]`$, we get from the above example the following description of Kleinian codes:
Even (odd) self-dual $`K`$-codes of length $`n`$ are the same as self-dual VOA’s (SVOA’s) of rank $`4n`$ with sub-VOA $`V_{D_4}^n`$, the $`n`$-th tensor product of the VOA associated to the Level-$`1`$-representation of the affine Lie algebra $`\widehat{\mathrm{Spin}}(8)`$. The automorphism group of $`K^n`$ corresponds to the outer automorphism group of $`V_{D_4^n}`$ in the VOA-sense (Triality of $`\mathrm{Spin}(8)`$!); the group algebra $`𝐙[K^n]`$ is the fusion algebra of $`V_{D_4^n}`$.
One has a similar description for binary codes in terms of the lattice-VOA $`V_{A_1}^n`$.
For $`V`$ be the (non rational) Heisenberg-VOA $`V_h`$ of rank $`1`$ on has $`G^n=𝐑^n`$, $`\alpha =e^{2\pi i\frac{(.,.)}{2}}`$. Isotropic subspaces are even integral lattices, i.e., we have a $`1:1`$-correspondence between rank $`n`$ VOA’s containing the Heisenberg-VOA $`V_h^nV_{h^n}`$ and even integral lattices.
The description of (marked/framed) Kleinian codes, binary codes and lattices in terms of VOA’s is summarized in the next table.
$$\begin{array}{cccccc}\text{Object}\hfill & \text{Rank}& \text{Sub-VOA (framed)}& \text{Group}& \text{Sub-VOA}& \text{Group}\\ & & & & & \\ K\text{-codes}\hfill & 4n& V_{D_4,}^n& 2^n:S_n& V_{D_4}^n& S_3^n:S_n\\ \text{binary codes}\hfill & 2n& V_{D_2}^n& 2^n:S_n& V_{A_1}^{2n}& S_{2n}\\ \text{lattices}\hfill & n& V_{D_1}^n& 2^n:S_n& V_h^n& \mathrm{SO}(n)\\ \text{VOA’s}\hfill & n/2& L_{1/2}(0)^n& \text{cf. }\text{[GH]}& \mathrm{Vir}_n& \mathrm{`}\mathrm{`}\mathrm{Aut}((\mathrm{Vir}_n))\mathrm{"}\end{array}$$
Construction A (including marking/frames) can now be completely understood in terms of VOA’s as indicated in following table of inclusions:
$$\begin{array}{ccccccccc}& \text{K–B}& & & \text{B–L}& & & \text{L–V}& \\ & & & & & & & & \\ V_{D_4,}^n& & V_{D_2}^{2n}& V_{D_2}^n& & V_{D_1}^{2n}& V_{D_1}^n& & L_{1/2}(0)^{2n}\\ & & & & & & & & \\ V_{D_4}^n& & V_{A_1}^{2n}& V_{A_1}^{2n}& & V_h^{2n}& V_h^n& & \mathrm{Vir}_n\end{array}$$
For all four theories one has analogous basic objects. We display their relations in Table 5.
Final Remarks: The way from Kleinian codes over binary codes and lattices to VOA’s is not canonically given. There is no way to see what is the next step. But in the other direction there is in some sense always a canonical choice: Consider the self-dual objects of rank $`24`$. There are always two objects without “roots”: An even and an odd one.<sup>§</sup><sup>§</sup>§ In the case of vertex operator algebras the uniqueness of the moonshine module $`V^{\mathrm{}}`$ and the odd moonshine module $`VO^{\mathrm{}}`$ is only a conjecture. Look at the even subobject of the odd one. Exactly one of its $`4`$ modules contains “roots”. Take the direct sum of the even subobject and the “root”-module and consider inside the subobject generated by the “roots”. It is a direct product of indecomposable objects. The next step is now represented by “Codes” over the modules of one such indecomposable object.
There is one more such step before Kleinian codes, namely codes over the $`3`$-dimensional topological quantum field theory belonging the vertex operator algebra $`V_{D_8}`$.
Some historical comments and further developments:
I found the structure of Kleinian codes as developed in this paper by searching for an analogue of the shorter Moonshine module in autumn 1995. This was motivated by the work on Virasoro frames inside the Moonshine module. The weight enumerator of the shorter Hexacode (which is not a $`𝐅_4`$-code) dropped out. Compare the last paragraphs above.
A first outline of this paper was distributed during the first two month of 1996 including all the results but most proofs not yet written up in Kleinian code language. Some other preliminary versions, but now without Section 5, were distributed in summer 1996. The only exception to this is the extremal code of length 12. I tried to find such a code by hand (cf. letter to Hirzebruch \[Höh96\]), but without success. Back in Germany in October 1996, it popped up on the screen of my old AT-286 PC after a few minutes (or hours) by running a simple back-tracking algorithm. This code was also found in \[CRSS98\], where the authors applied the theory of Kleinian codes to quantum codes. This paper became the stimulus of a lot of research on quantum codes. It seems that only a late 1996 preprint found the widest distribution. I am sorry about the delay in publishing the paper. I like to thank C. Bachoc, J.-L. Kim and V. Pless for comments on the final version.
Since that time, Kleinian codes have been investigated further. In the following, I will give an overview.
Section 2: The invariant ring for the complete weight enumerator of even self-dual Kleinian codes has been given in \[RS98a\].
Section 3: Examples of cyclic self-dual codes for all odd length have been given by M. Ran and J. Snyders in \[RS00\].
It was pointed out to me by J.-L. Kim that the papers \[GHKPa, BG\] are answering partially my question for the smallest codes with trivial automorphism group: There is at least one such code of length $`12`$ (called $`QC\mathrm{\_}12g`$ in \[GHKPa\]; non even) and there are at least $`273`$ such extremal even codes of length $`14`$ (see \[BG\]). Since all the even codes of length $`8`$ and $`10`$ without weight $`2`$ vectors are extremal, it follows from Section 3 and \[BG\] that the answer for even self-dual codes must be $`12`$ or $`14`$.
Section 4: The upper bound of Theorem 11 has been sharpened by E. Rains in \[Rai98\] to $`d2[n/6]+2+e`$ with $`e=1`$ for $`n5(mod6)`$ and $`e=0`$ else. For $`6|n`$, a code meeting this bound is even. An analogue sharpened bound for binary codes can also be found in \[Rai98\] and for odd lattices in \[RS98b\].
In \[GHKPb, GHKPa\], Gaborit, Huffman, Kim and Pless classified self-dual Kleinian codes with minimal weight reaching the above bound for length $`8`$, $`9`$ and $`11`$ (there are $`5`$, $`8`$ resp. $`1`$ such codes). They also proved the uniqueness of the extremal even code of length $`12`$. There is no such code for length $`13`$ (see \[RS98a\]). For even codes, the length $`10`$ has been settled in \[BG\] ($`19`$ codes), where also partial results for length $`14`$ and $`18`$ are obtained.
Section 5: C. Bachoc (see \[Bac\]) has proven Theorem 17 and some extensions of it for all $`n`$ by using discrete harmonic analysis on $`X_k`$. Interestingly, this approach works only for alphabets with $`2`$, $`3`$ and $`4`$ elements and a unique choice of group structure and bilinear form. For four elements, one gets our scalar product on $`K`$. The binary analogue was studied before in \[Bac99\]. This approach forms the direct analogue to the approach of B. Venkov for lattices \[Ven84\].
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# The rotating detector and vacuum fluctuations
## I Introduction
The key point of special relativity is that the Poincaré group is the symmetry group of all physical systems. The definitions of the Lorentz and Poincaré groups are based as groups of mappings that leave invariant the flat metric. The natural consequence is that an inertial observer (in an inertial reference frame) can assign a time and space location to any event occurring in space-time, using light clocks, etc. To obtain a “physical” interpretation of the Poincaré mapping we have to derive general relations between the space-time measurements made by different observers who are in different inertial frames.
So far, we have been considering only classes of measuring devices in inertial reference frames. However, suppose we are to make measurements with a device in non-inertial frames, as for example in a rotating disc. To discuss such measurements and to show how they can be incorporated into a space-time description, entails that distance and time measurements made with some arbitrary set of measuring devices can always be made to correspond to the coordinates of space-time by means of a suitable space-time mapping. In other words, in order to compare measurements made by inertial and non-inertial (e.g. rotating) observers, we must present the mapping that relates the measurements made with the two different sets of devices.
For example, in a Galilean scenario it is possible to relate the space and time measurements made in a rotating frame to those in an inertial one by the mapping:
$`t`$ $`=`$ $`t^{},`$ (1)
$`r`$ $`=`$ $`r^{},`$ (2)
$`\theta `$ $`=`$ $`\theta ^{}\mathrm{\Omega }t^{},`$ (3)
$`z`$ $`=`$ $`z^{},`$ (4)
where $`\mathrm{\Omega }`$ is the constant angular velocity around the $`z`$ axis of the inertial frame. In the above, the cylindrical coordinate system $`x^\mu =\{t^{},r^{},\theta ^{},z^{}\}`$ is adapted to an inertial observer and the rotating coordinate system $`x^\mu =\{t,r,\theta ,z\}`$ is the one adapted to the rotating observer. Although some authors tried to construct the counterpart of this mapping incorporating special relativity , the final answer for this question is still open.
In two recent papers it is assumed that the mapping which relates the inertial frame with the rotating one is given by:
$`t`$ $`=`$ $`t^{}\mathrm{cosh}\mathrm{\Omega }r^{}r^{}\theta ^{}\mathrm{sinh}\mathrm{\Omega }r^{},`$ (5)
$`r`$ $`=`$ $`r^{},`$ (6)
$`\theta `$ $`=`$ $`\theta ^{}\mathrm{cosh}\mathrm{\Omega }r^{}{\displaystyle \frac{t^{}}{r^{}}}\mathrm{sinh}\mathrm{\Omega }r^{},`$ (7)
$`z`$ $`=`$ $`z^{}.`$ (8)
Such group of transformations was presented a long time ago by Trocheries and also Takeno . In Takeno’s derivation, three assumptions were made: (i) the transformation laws constitute a group; (ii) for small velocities we must recover the usual linear velocity law ($`v=\mathrm{\Omega }r`$); and (iii) the velocity composition law is also in agreement with special relativity. In fact, the above transformation predicts that the velocity of a point at distance $`r`$ from the axis is given by $`v(r)=\mathrm{tanh}(\mathrm{\Omega }r)`$.
It is our purpose, in this work, to investigate how does an uniformly rotating observer see the Minkowski vacuum. In this direction we performed the quantization of a scalar field as an observer who rotates uniformly around some fixed point would do it. We assume the Trocheries-Takeno transformations (5-8) to compare measurements made in the inertial and in the rotating frames.
The canonical quantization of a scalar field in the rotating frame, related to the inertial one through coordinate transformations (1-4), was made by Denardo and Percacci and also by Letaw and Pfautsch . To compare the quantizations performed in the inertial and rotating frames, the authors calculated the Bogolubov transformation between the inertial modes $`\psi _i(t^{},r^{},\theta ^{},z^{})`$ and the modes adapted to the rotating frame, $`\overline{\psi }_j(t,r,\theta ,z)`$. Since they found that the Bogolubov coefficients $`\beta _{ij}`$ are null, they conclude that the rotating vacuum or no-particle state, as defined by the rotating observer, is just the Minkowski vacuum $`|0,M`$.
Another way to compare different quantizations is to study the vacuum activity of a quantum field, and this is performed introducing a measuring device which couples with the quantum field via an interaction Lagrangian. (In the following discussion we will take into account the detector model due to Unruh and De Witt .) Using first-order perturbation theory, it is possible to calculate the probability of excitation per unit proper time (excitation rate) of such a detector, that is, the probability per unit time that the detector, travelling in a given world-line and initially prepared in its ground state, to wind up in an excited state when it interacts with the field in a given state . As an example, consider that the detector is in an inertial frame and the field is prepared in the Minkowski vacuum state. In this situation the detector will remain in the ground state (null excitation rate). This is easily understood, because there are no inertial particles in the Minkowski vacuum state. On the other hand, if the detector is put in a world-line of an observer with constant proper acceleration and the field is prepared again in the Minkowski vacuum, the detector has a non-null probability to suffer a transition to an excited state. This is the so-called Unruh-Davies effect and the quantitative result is in agreement with the fact that the Minkowski vacuum state is seen as a thermal state by the accelerated observer, with temperature proportional to its proper acceleration . The construction of a quantum field theory with the implementation of the Fock space in Rindler’s manifold leads to define the Rindler vacuum state $`|0,R`$. For completeness, in the situation when the field is prepared in such a state and the detector is uniformly accelerated, the detector remains inert. Again one is tempted to conclude that this is so because there are no Rindler particles in $`|0,R`$ to be detected in the uniformly accelerated frame.
The agreement between the response of a detector and canonical quantum field theory seems not to occur for more general situations. Indeed, as was shown by Letaw and Pfautsch, if the detector is put in a uniformly rotating world-line and the field is prepared in the Minkowski vacuum state $`|0,M`$, it is found a non-null excitation rate, in spite of the fact that $`|0,M`$ is considered as the vacuum state for a rotating observer, as discussed above. The rotating detector is excited even though there are no particles as an orbiting observer would define them (see also ).
Recently Davies et al (see also ) solved this paradox, still assuming the Galilean coordinate transformations (1-4) between the inertial and rotating frames. First of all note that the world-line of an observer in the rotating frame is an integral curve of the Killing vector $`\xi =(1\mathrm{\Omega }^2r^2)^{\frac{1}{2}}/t`$, which is timelike only for $`\mathrm{\Omega }r<1`$. Therefore, for a given angular velocity $`\mathrm{\Omega }`$ there will be a maximum value of the radial coordinate $`r_{max}=1/\mathrm{\Omega }`$ (the light cylinder) for which an observer a distance $`r>r_{max}`$ will be moving faster than light. The Bogolubov coefficient $`\beta `$ is a scalar product over a spacelike hypersurface, where the radial coordinate ranges over $`0r<\mathrm{}`$ and the notion of Bogolubov transformations outside the light cylinder becomes obscure. In order to circumvent this problem, Davies et al introduced a perfectly conducting cylinder with radius $`a<r_{max}`$ and they prove that the response of the detector vanishes when it is put a distance $`r<a`$ from the rotation axis. They conclude that “a rotating particle detector corotating with a rotating vacuum state registers the absence of quanta”, although they continue to regard the Minkowski vacuum as the rotating vacuum state.
The aim of this paper is two-fold. The first one is to present an exact solution of the Klein-Gordon equation in the Trocheries-Takeno coordinate system and to show that the Bogolubov coefficients $`\beta `$ between inertial and rotating modes are not zero. This fact proves that there is a Trocheries-Takeno vacuum state adapted to rotating observers. The second one is to analyse the behavior of an apparatus device, a detector which is coupled with the scalar field travelling in inertial or rotating world-lines, interacting with the field in the Minkowski or the Trocheries-Takeno vacuum states.
We organize this paper as follows. In Section II we second quantize a massless scalar field in Takeno’s rotating coordinate system, and also compare this quantization with the usual one in the inertial frame via the calculation of the Bogolubov coefficients. In Section III we introduce the measuring apparatus – the Unruh-De Witt detector . We calculate its response function when it is rotating and the field is prepared in two different states: the Trocheries-Takeno vacuum state and the Minkowski vacuum state. We also consider the case of an inertial detector interacting with the field in the rotating vacuum. Conclusions are given in Section IV. In this paper we use $`\mathrm{}=c=k_B=1`$.
## II Canonical quantization in the inertial and rotating frames
In this section we will make a comparison between the quantizations of the massless Klein-Gordon field performed in the inertial frame and in the rotating one (Trocheries-Takeno), when the two coordinate systems are related by the mapping (5-8). Such a comparison will be made by calculating the Bogolubov transformation between the inertial and rotating modes, solutions of the respective Klein-Gordon equations. For the quantization in the inertial frame one chooses cylindrical coordinates on $`t^{}=`$ constant hypersurfaces and writes the Klein-Gordon equation in terms of them. We just quote the results of refs. . Positive-frequency modes (with respect to inertial time $`t^{}`$), solutions of the Klein-Gordon equation, are found to be:
$$v_{q^{}m^{}k^{}}(t^{},r^{},\theta ^{},z^{})=N_1e^{ik^{}z^{}+im^{}\theta ^{}}e^{i\omega ^{}t^{}}J_m^{}(q^{}r^{}),$$
(9)
where $`\omega ^{}{}_{}{}^{2}=q^{}{}_{}{}^{2}+k^{}^2`$, $`J_m^{}(q^{}r^{})`$ are Bessel functions well-behaved at the origin and $`N_1=q^{\frac{1}{2}}[2\pi (2\omega ^{})^{1/2}]^1`$ is a normalization factor. In the above, $`m^{}=0,\pm 1,\pm 2,\pm 3,\mathrm{}`$, $`0q^{}<\mathrm{}`$ and $`\mathrm{}<k^{}<\mathrm{}`$. In this way the field is expanded as:
$$\varphi (t^{},r^{},\theta ^{},z^{})=\underset{m^{}}{}𝑑q^{}𝑑k^{}\left[b_{q^{}m^{}k^{}}v_{q^{}m^{}k^{}}(t^{},r^{},\theta ^{},z^{})+b_{q^{}m^{}k^{}}^{}v_{q^{}m^{}k^{}}^{}(t^{},r^{},\theta ^{},z^{})\right],$$
(10)
where the coefficients $`b_{q^{}m^{}k^{}}`$ and $`b_{q^{}m^{}k^{}}^{}`$ are, respectively, the annihilation and creation operators of the inertial quanta of the field and satisfy the usual commutation rule $`[b_i,b_j^{}]=\delta _{ij}`$. In the above, the modes $`v_i`$ and $`v_i^{}`$ are called, respectively, positive and negative frequency modes with respect to the Killing vector $`/t^{}`$. It is important to stress that in stationary geometries, such as the Minkowski space-time, the definition of positive and negative frequency modes has no ambiguities. The Minkowski vacuum state is then defined by
$$b_{q^{}m^{}k^{}}|0,M=0,q^{},m^{},k^{}.$$
(11)
Now we shall consider the quantization in the rotating frame. Assuming the mapping (5-8) to connect measurements made in the rotating frame and those made in the inertial one, the line element in the rotating coordinates assumes the non-stationary form :
$$ds^2=dt^2(1+P)dr^2r^2d\theta ^2dz^2+2Qdrd\theta +2Sdtdr,$$
(12)
where $`P,Q`$ and $`S`$ are given by:
$`P`$ $`=`$ $`({\displaystyle \frac{Y}{r^2}}+4\mathrm{\Omega }\theta t)\mathrm{sinh}^2\mathrm{\Omega }r{\displaystyle \frac{\mathrm{\Omega }}{r}}(t^2+r^2\theta ^2)\mathrm{sinh}2\mathrm{\Omega }r+\mathrm{\Omega }^2Y,`$ (13)
$`Q`$ $`=`$ $`r\theta \mathrm{sinh}^2\mathrm{\Omega }r{\displaystyle \frac{1}{2}}t\mathrm{sinh}2\mathrm{\Omega }r+\mathrm{\Omega }rt,`$ (14)
$`S`$ $`=`$ $`{\displaystyle \frac{t}{r}}\mathrm{sinh}^2\mathrm{\Omega }r{\displaystyle \frac{1}{2}}\theta \mathrm{sinh}2\mathrm{\Omega }r\mathrm{\Omega }r\theta ,`$ (15)
with $`Y=(t^2r^2\theta ^2)`$. Note that this metric presents no event horizons. In order to implement the canonical quantization first we have to solve the Klein-Gordon equation in the Trocheries-Takeno coordinate system:
$$\mathrm{}\varphi (t,r,\theta ,z)=0.$$
(16)
It is possible to show that a complete set, basis in the space of solutions of the Klein-Gordon equation is given by $`\{u_{qmk},u_{qmk}^{}\}`$, in which
$$u_{qmk}(t,r,\theta ,z)=N_2e^{ikz}\mathrm{exp}\left[i\left(m\mathrm{cosh}\mathrm{\Omega }r+\omega r\mathrm{sinh}\mathrm{\Omega }r\right)\theta i\left(\frac{m}{r}\mathrm{sinh}\mathrm{\Omega }r+\omega \mathrm{cosh}\mathrm{\Omega }r\right)t\right]J_m(qr),$$
(17)
where $`\omega ^2=q^2+k^2`$ and $`N_2`$ is a normalization factor. Again, $`m=0,\pm 1,\pm 2,\pm 3,\mathrm{}`$, $`0q<\mathrm{}`$ and $`\mathrm{}<k<\mathrm{}`$. One sees that these modes are well-behaved throughout the whole manifold. Making use of the transformations (5-8) one can show that these modes are of positive frequency by using the criterium of di Sessa , which states that a given mode is of positive frequency if it vanishes in the limit $`(t^{})i\mathrm{}`$, where $`t^{}`$ is the inertial time coordinate, while $`u_j^{}`$ are modes of negative frequency. In this way, the field operator is expanded in terms of these modes as:
$$\varphi (t,r,\theta ,z)=\underset{m}{}𝑑q𝑑k\left(a_{qmk}u_{qmk}(t,r,\theta ,z)+a_{qmk}^{}u_{qmk}^{}(t,r,\theta ,z)\right),$$
(18)
where the coefficients $`a_{qmk}`$ and $`a_{qmk}^{}`$ are, respectively, the annihilation and creation operators of the Trocheries-Takeno quanta of the field. The vacuum state defined by the rotating observer is thus the Trocheries-Takeno vacuum state $`|0,T`$ and it is given by
$$a_{qmk}|0,T=0,q,m,k.$$
(19)
The many-particle states of the theory can be obtained through successive applications of the creation operators on the vacuum state.
We are now ready to compare both quantizations by using the Bogolubov transformations . Since both sets of modes are complete, one can expand modes (9) in terms of modes (17) and vice-versa, the coefficients of this expansion being called Bogolubov coefficients. For instance:
$$u_i(x)=\underset{j}{}\alpha _{ij}v_j(x)+\beta _{ij}v_j^{}(x)$$
(20)
and conversely:
$$v_j(x)=\underset{i}{}\alpha _{ij}^{}u_i(x)\beta _{ij}u_i^{}(x),$$
(21)
where $`\alpha _{jj^{}}=(u_j,v_j^{})`$ and $`\beta _{jj^{}}=(u_j,v_j^{}^{})`$, where the scalar product is defined by:
$$(\varphi ,\chi )=i_\mathrm{\Sigma }𝑑\mathrm{\Sigma }^\mu \sqrt{g}\left[\varphi _\mu \chi ^{}\chi ^{}_\mu \varphi \right],$$
(22)
with $`d\mathrm{\Sigma }^\mu =n^\mu d\mathrm{\Sigma }`$, where $`n^\mu `$ is a future-oriented unit vector orthogonal to the spacelike hypersurface $`\mathrm{\Sigma }`$. As $`\mathrm{\Sigma }`$ we will choose the hypersurface $`t^{}=0`$, $`t^{}`$ being the inertial time. The relevant coefficient for our present analysis is the $`\beta `$ coefficient since it is this coefficient which gives the content of rotating particles in the Minkowski vacuum , and for such a calculation we need to express the rotating modes $`u_j(x)`$ in terms of the inertial coordinates, using transformations (5-8). In this way:
$`u_{qmk}(t^{},r^{},\theta ^{},z^{})`$ $`=`$ $`N_2e^{ikz^{}}\mathrm{exp}\left(i\theta ^{}\left[m\mathrm{cosh}2\mathrm{\Omega }r^{}+\omega r^{}\mathrm{sinh}2\mathrm{\Omega }r^{}\right]\right)`$ (23)
$`\times `$ $`\mathrm{exp}\left(it^{}\left[\omega \mathrm{cosh}2\mathrm{\Omega }r^{}+{\displaystyle \frac{m}{r^{}}}\mathrm{sinh}2\mathrm{\Omega }r^{}\right]\right)J_m(qr^{})`$ (24)
and the Bogolubov coefficient $`\beta _{jj^{}}`$ is written as:
$`\beta _{jj^{}}`$ $`=`$ $`+i{\displaystyle _0^{\mathrm{}}}r^{}𝑑r^{}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑z^{}{\displaystyle _0^{2\pi }}𝑑\theta ^{}\left[u_j(x^{}){\displaystyle \frac{v_j^{}(x^{})}{t^{}}}v_j^{}(x^{}){\displaystyle \frac{u_j(x^{})}{t^{}}}\right]`$ (25)
$`=`$ $`N_1N_2{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑z^{}e^{i(k+k^{})z^{}}{\displaystyle _0^{\mathrm{}}}r^{}𝑑r^{}\left[\omega ^{}\omega \mathrm{cosh}2\mathrm{\Omega }r^{}{\displaystyle \frac{m}{r^{}}}\mathrm{sinh}2\mathrm{\Omega }r^{}\right]J_m^{}(q^{}r^{})J_m(qr^{})`$ (26)
$`\times `$ $`{\displaystyle _0^{2\pi }}𝑑\theta ^{}\mathrm{exp}\left(i\theta ^{}\left[m^{}+m\mathrm{cosh}2\mathrm{\Omega }r^{}+\omega r^{}\mathrm{sinh}2\mathrm{\Omega }r^{}\right]\right).`$ (27)
The first integral is easily evaluated to a delta function $`2\pi \delta (k+k^{})`$, while the third one gives us:
$$_0^{2\pi }𝑑\theta ^{}\mathrm{exp}\left(i\theta ^{}A_{m,m^{}}(r^{},\omega )\right)=\left(iA_{m,m^{}}(r^{},\omega )\right)^1\left[\mathrm{exp}\left(2\pi iA_{m,m^{}}(r^{},\omega )\right)1\right],$$
where
$$A_{m,m^{}}(r^{},\omega )=m^{}+m\mathrm{cosh}2\mathrm{\Omega }r^{}+\omega r^{}\mathrm{sinh}2\mathrm{\Omega }r^{}.$$
(28)
Thus, we obtain
$`\beta _{jj^{}}`$ $`=`$ $`2\pi N_1N_2\delta (k+k^{}){\displaystyle _0^{\mathrm{}}}r^{}𝑑r^{}\left(\omega ^{}\omega \mathrm{cosh}2\mathrm{\Omega }r^{}{\displaystyle \frac{m}{r^{}}}\mathrm{sinh}2\mathrm{\Omega }r^{}\right)J_m^{}(q^{}r^{})J_m(qr^{})`$ (29)
$`\times `$ $`\left(iA_{m,m^{}}(r^{},\omega )\right)^1\left[\mathrm{exp}\left(2\pi iA_{m,m^{}}(r^{},\omega )\right)1\right].`$ (30)
The resulting expression is difficult to evaluate, but nonetheless it is non-zero. (In the appendix we give an indirect proof that it is non-zero.) This means that the two vacua considered are non-equivalent, i.e., $`|0,M|0,T`$, which means that the Minkowski vacuum $`|0,M`$ contains rotating quanta, i.e., Trocheries-Takeno particles .
## III Detector excitation rate
We now pass to consider the probability of excitation of a detector which is moving in a circular path at constant angular velocity $`\mathrm{\Omega }`$ and at a distance $`R_0`$ from the rotation axis, interacting with the scalar field. The initial state of the detector is its ground state and for the initial state of the field we will consider the two vacuum states: the usual Minkowski vacuum state and also the Trocheries-Takeno vacuum state. The interaction with the field may cause transitions between the energy levels of the detector and if it is found, after the interaction, in an excited state, one can say that it has detected a vacuum fluctuation of the field .
As a detector we shall be considering mainly the detector model of Unruh-De Witt , which is a system with two internal energy eigenstates with monopole matrix element between these two states different from zero. According to standard theory , the probability of excitation per unit proper time of such a system (normalized by the selectivity of the detector), or simply, its excitation rate, is given by:
$$R(E)=_{\mathrm{}}^{\mathrm{}}𝑑\mathrm{\Delta }te^{iE\mathrm{\Delta }t}G^+(x,x^{}),$$
(31)
where $`\mathrm{\Delta }t=tt^{}`$, $`E>0`$ is the difference between the excited and ground state energies of the detector and $`G^+(x,x^{})`$ is the positive-frequency Wightman function calculated along the detector’s trajectory. Let us note that the positive-frequency Wightman function is given by
$$G^+(x,x^{})=0|\varphi (x)\varphi (x^{})|0,$$
(32)
where $`|0`$ is the vacuum state of the field, which can either be $`|0,M`$ or $`|0,T`$. Let us consider first the second possibility.
If one splits the field operator in its positive and negative frequency parts with respect to the Trocheries-Takeno time coordinate $`t`$, as $`\varphi (x)=\varphi ^+(x)+\varphi ^{}(x)`$, where $`\varphi ^+(x)`$ contains only annihilation operators and $`\varphi ^{}(x)`$ contains only creation operators (see Eq.(18)), and also considers $`|0`$ as the Trocheries-Takeno vacuum state, i.e., $`|0=|0,T`$ then, using Eq.(19), one finds that:
$$G_T^+(x,y)=\underset{i}{}u_i(x)u_i^{}(y),$$
(33)
where the subscript $`T`$ stands for the Wightman function calculated in the Trocheries-Takeno vacuum state. Considering now the modes given by Eq.(17) and that we are interested in the situation where the detector is at rest in the Trocheries-Takeno frame, i.e., $`\theta =`$ constant, $`z=`$ constant and $`r=R_0=`$ constant, one finds:
$$G_T^+(x,y)=\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}_0^{\mathrm{}}𝑑q_{\mathrm{}}^{\mathrm{}}𝑑kN_2^2e^{i[\frac{m}{R_0}\mathrm{sinh}\mathrm{\Omega }R_0+\omega \mathrm{cosh}\mathrm{\Omega }R_0]\mathrm{\Delta }t}J_m^2(qR_0).$$
(34)
Putting the above expression in Eq.(31), we find:
$$R_T^{(r)}(E,R_0)=\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}_0^{\mathrm{}}𝑑q_{\mathrm{}}^{\mathrm{}}𝑑kN_2^2J_m^2(qR_0)_{\mathrm{}}^{\mathrm{}}𝑑\mathrm{\Delta }te^{i[E+\frac{m}{R_0}\mathrm{sinh}\mathrm{\Omega }R_0+\omega \mathrm{cosh}\mathrm{\Omega }R_0]\mathrm{\Delta }t}.$$
(35)
(In the above, the subscript $`T`$ stands for the Takeno vacuum and the superscript $`(r)`$ stands for the rotating world-line followed by the detector.) The last integral gives us $`2\pi \delta \left(E+\frac{m}{R_0}\mathrm{sinh}\mathrm{\Omega }R_0+\omega \mathrm{cosh}\mathrm{\Omega }R_0\right)`$, for which the argument is non-null only if $`m<0`$; we can take the summation index to run for $`m=1,2,3,\mathrm{}`$, leaving us with
$$R_T^{(r)}(E,R_0)=2\pi \underset{m=1}{\overset{\mathrm{}}{}}_0^{\mathrm{}}𝑑q_{\mathrm{}}^{\mathrm{}}𝑑kN_2^2J_m^2(qR_0)\delta \left(E\frac{m}{R_0}\mathrm{sinh}\mathrm{\Omega }R_0+\omega \mathrm{cosh}\mathrm{\Omega }R_0\right).$$
(36)
The above expression predicts excitation for the detector, and depends in a non-trivial way on the position $`R_0`$ where it is put. So we once again arrive at the same confrontation between canonical quantum field theory and the detector formalism, which was settled by Letaw and Pfautsch and Padmanabhan: how is it possible for the orbiting detector to be excited in the rotating vacuum? However a crucial distinction exists between our present analysis and the above-mentioned works: we state, as proved in the last section, that the rotating vacuum is not the Minkowski vacuum. We now analyse the two independent origins of the non-null excitation rate, Eq.(36).
Note that the present situation of an Unruh-De Witt detector being excited when put in an orbiting world-line interacting with the field in the rotating (Trocheries-Takeno) vacuum is to be contrasted with the two following situations. In fact, this same detector is not excited whether it is in an inertial world-line and interacting with the field in the inertial (Minkowski) vacuum or when it is uniformly accelerated and interacting with the field in the accelerated (Rindler) vacuum . However note that both Minkowski and Rindler space-times are static, differently of the Trocheries-Takeno metric, and differently also of the rotating metric obtained using the galilean transformation, which are examples of non-static metrics. Recall that using instead the galilean transformation to a rotating frame it was also found by Letaw and Pfautsch a non-null excitation rate for the orbiting detector in the rotating vacuum, considered by them as $`|0,M`$. Therefore the excitation in the present case may be attributed to the non-staticity of the Trocheries-Takeno metric.
The other origin of the excitation found above for the detector can be atributed to the detector model we adopted . Indeed, note that splitting the field operator in its positive and negative frequency parts with respect to rotating time $`t`$ in Eq.(32), one can express the Wightman function as:
$`G^+(x,x^{})`$ $`=`$ $`0|\varphi (x)\varphi (x^{})|0`$ (37)
$`=`$ $`0|\varphi ^+(x)\varphi ^+(x^{})|0+0|\varphi ^+(x)\varphi ^{}(x^{})|0`$ (38)
$`+`$ $`0|\varphi ^{}(x)\varphi ^+(x^{})|0+0|\varphi ^{}(x)\varphi ^{}(x^{})|0.`$ (39)
In the case where $`|0=|0,T`$, because of Eq.(19) only the second term above is non-vanishing, corresponding to the emission (creation) of a Trocheries-Takeno quantum with simultaneous excitation of the detector and this is the term responsible for the non-vanishing excitation rate, Eq.(36). In the context of quantum optics photodetection is regarded as photoabsorption processes only , a detector being able to be excited only when it absorbs (annihilates) a quantum of the field. In this way terms like the second one above are discarded and only terms like the third one are taken into account, such a procedure being called the rotating-wave approximation . Therefore purely absorptive detectors (Glauber model) always give vanishing excitation rate in the vacuum state of the field. From this discussion, it is clear that the Glauber detector model will not be excited when put in the orbiting world-line and the field is in the Trocheries-Takeno vacuum state. Another context in which the inclusion of the antiresonant term (second one above) plays a crucial role is that of accelerated observers, where the thermal character of the Minkowski vacuum as seen by a Rindler observer is not revealed if one uses the Glauber correlation function, but only if one uses the Wightman one. In effect, the Wightman correlation function includes the vacuum fluctuations that are omitted in the Glauber function , and to these very vacuum fluctuations can be attributed the non-vanishing excitation rate Eq.(36). Because of this feature the model of Unruh-De Witt is also called a fluctuometer .
We now discuss the other case of putting the detector in an orbiting trajectory and preparing the scalar field in the usual inertial vacuum $`|0,M`$. Writing $`|0,M`$ for $`|0`$ in Eq.(32), it is easy to show that the positive frequency Wightman function is given by:
$$G_M^+(x^{},y^{})=\underset{j}{}v_j(x^{})v_j^{}(y^{}),$$
(40)
where $`M`$ stands for the Minkowski vacuum state. As the rate of excitation Eq.(31) is given in terms of the detector’s proper time, we shall express Eq.(40) in terms of the rotating coordinates, using the inverse of Takeno’s transformations. Let us begin with $`G_M^+(x^{},y^{})`$, in inertial coordinates, with identifications $`r_1^{}=r_2^{}=R_0`$ and $`z_1^{}=z_2^{}`$, as demanded for this case:
$$G_M^+(x^{},y^{})=\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}_0^{\mathrm{}}𝑑q_{\mathrm{}}^{\mathrm{}}𝑑kN_1^2e^{i\omega (t_1^{}t_2^{})+im(\theta _1^{}\theta _2^{})}J_m^2(qR_0).$$
(41)
The inverse of Takeno’s transformations read
$`t^{}`$ $`=`$ $`t\mathrm{cosh}\mathrm{\Omega }r+r\theta \mathrm{sinh}\mathrm{\Omega }r,`$ (42)
$`r^{}`$ $`=`$ $`r,`$ (43)
$`\theta ^{}`$ $`=`$ $`\theta \mathrm{cosh}\mathrm{\Omega }r+{\displaystyle \frac{t}{r}}\mathrm{sinh}\mathrm{\Omega }r,`$ (44)
$`z^{}`$ $`=`$ $`z.`$ (45)
Using the above in Eq.(41) and taking note of the fact that the detector is at rest in the rotating frame, i.e., $`\theta _1=\theta _2`$, we see that in this manner the Minkowski Wightman function is a function of the difference in proper time $`\mathrm{\Delta }t=t_1t_2`$, which allows us to calculate the rate of excitation of the orbiting detector when the field is in the Minkowski vacuum:
$$R_M^{(r)}(E,R_0)=2\pi \underset{m=1}{\overset{\mathrm{}}{}}_0^{\mathrm{}}𝑑q_{\mathrm{}}^{\mathrm{}}𝑑kN_1^2J_m^2(qR_0)\delta \left(E\frac{m}{R_0}\mathrm{sinh}\mathrm{\Omega }R_0+\omega \mathrm{cosh}\mathrm{\Omega }R_0\right).$$
(46)
The result above is very much like Eq.(36), with the exception that in the above it appears the normalization of the inertial modes $`N_1`$ instead of $`N_2`$.
Finally, let us suppose that it is possible to prepare the field in the rotating vacuum and the detector is in an inertial world-line and let us calculate the excitation rate in this situation:
$$R_T^{(i)}(E,R_0)=_{\mathrm{}}^{\mathrm{}}𝑑\mathrm{\Delta }t^{}e^{iE\mathrm{\Delta }t^{}}G_T^+(x^{},y^{}),$$
(47)
where the superscript $`(i)`$ stands for the inertial world-line followed by the detector, $`\mathrm{\Delta }t^{}`$ is the difference in proper time in the inertial frame, and $`G_T^+(x^{},y^{})`$ is given by Eq.(33), but written now in terms of the inertial coordinates. It is not difficult to write $`G_T^+(x^{},y^{})`$ in terms of the inertial coordinates, recalling that now the detector is not at rest in the rotating frame. We have therefore the result that:
$`R_T^{(i)}(E,R_0)`$ $`=`$ $`2\pi {\displaystyle \underset{m=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle _0^{\mathrm{}}}𝑑q{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑kN_2^2J_m^2(qR_0)`$ (48)
$`\times `$ $`\delta \left(E\left(\omega \mathrm{\Omega }R_0{\displaystyle \frac{m}{R_0}}\right)\mathrm{sinh}(2\mathrm{\Omega }R_0)(m\mathrm{\Omega }\omega )\mathrm{cosh}(2\mathrm{\Omega }R_0)\right).`$ (49)
In order to study the activity of the Trocheries-Takeno vacuum, we calculated the rate of excitation of an Unruh-De Witt detector in two different situations: when it is put in the orbiting and in the inertial world-lines. Since in the first case we found a non-null rate, contrary to the idea that the orbiting detector co-rotating with the rotating vacuum should not perceive anything, it can be considered as a noise of the rotating vacuum, being it perceived regardless of the state of motion of the detector. This amounts to say that the inertial detector will also measure this noise, and we normalize the rate in this situation by subtracting from it the value Eq.(36), resulting in a normalized excitation rate for the inertial detector in interaction with the field in the rotating vacuum:
$`_T^{(i)}(E,R_0)=2\pi {\displaystyle \underset{m=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle _0^{\mathrm{}}}dq{\displaystyle _{\mathrm{}}^{\mathrm{}}}dkN_2^2J_m^2(qR_0)\times `$ (50)
$`\left[\delta \left(E\left(\omega \mathrm{\Omega }R_0{\displaystyle \frac{m}{R_0}}\right)\mathrm{sinh}(2\mathrm{\Omega }R_0)(m\mathrm{\Omega }\omega )\mathrm{cosh}(2\mathrm{\Omega }R_0)\right)\delta \left(E{\displaystyle \frac{m}{R_0}}\mathrm{sinh}\mathrm{\Omega }R_0+\omega \mathrm{cosh}\mathrm{\Omega }R_0\right)\right].`$ (51)
It remains to be clarified the meaning of $`R_0`$ in the excitation above. When studying the quantization in the rotating frame, one has to choose the world-line followed by the rotating observer, and this is parametrized by two quantities: the angular velocity $`\mathrm{\Omega }`$ and the distance $`R_0`$ from the rotation axis. The vacuum state which appears in such a quantization is thus also indexed by these parameters and this is the origin of $`R_0`$ in the above. A similar dependence of the excitation rate of a detector on a geometrical parameter appears, for instance, in the well-known Unruh-Davies effect: a uniformly accelerated detector interacting with the field in the Minkowski vacuum state will absorb particles in the same way as if it were inertial and interacting with the field in a thermal bath, with a temperature that depends on the proper acceleration of the detector.
## IV Summary and discussions
In this work we quantize a massless scalar field in a uniformly rotating frame and compare this quantization with the usual one in an inertial frame. As a difference with regard to previous works, we assume a coordinate transformation between both frames that takes into account the finite velocity of light and is valid in the whole manifold. In doing so, a material point orbiting around the axis of rotation never exceeds the speed of light, no matter how far it is from the axis. The metric, when written in rotating coordinates, presents no event horizons, although it is non-static and non-stationary. We recourse to a criterium of di Sessa to define positive and negative frequency modes in the rotating frame and the field is quantized along these lines. Such a quantization entails a vacuum state and by using the Bogolubov transformations we were able to show that this vacuum state is inequivalent to the Minkowski one. This is the main result of the paper. This results in that the Minkowski vacuum is seen as a many-rotating-particle state by a uniformly rotating observer, although it can not be seen as a thermal state (of Takeno particles), as in the case of a uniformly accelerated observer (for Rindler particles) and thus we cannot assign a temperature to it.
We obtain the response function of an Unruh-De Witt detector in three different situations: the orbiting detector is interacting with the field prepared in the rotating and in the Minkowski vacuum states, and finally the detector is travelling in an inertial world-line and the field is prepared in the rotating vacuum. In the first case it is found that the detector gets excited and we attributed this excitation to two different causes: Firstly, we are using the Unruh-De Witt detector model instead of a purely absorptive detector as the Glauber’s model, and secondly the Trocheries-Takeno metric is non-static. Because the rotating vacuum excites even a rotating detector, we consider this as a noise which will be measured by any other state of motion of the detector. In this way, when calculating the reponse of the inertial detector in the presence of the rotating vacuum we subtract from it this noise.
###### Acknowledgements.
This work was partially suported by Conselho Nacional de Desenvolvimento Científico e Tecnológico (CNPq) of Brazil.
## V Appendix
What interests us here is the number of Trocheries-Takeno particles in a given state $`j=(q,m,k)`$ that is present in the Minkowski vacuum, given in terms of the Bogolubov coefficients :
$$0,M\left|a_j^{}a_j\right|0,M=\underset{j^{}}{}\left|\beta _{jj^{}}\right|^2.$$
(52)
We see that this is given by a sum of the squared modulus of the various coefficients. If we can show that at least one of these coefficients is non-null, so we prove that there is a non-vanishing content of Trocheries-Takeno particles in the Minkowski vacuum.
We now give an indirect proof that the Bogolubov coefficient $`\beta `$ is in fact non-zero. The calculations in the text show that:
$`\beta _{jj^{}}`$ $`=`$ $`2\pi N_1N_2\delta (k+k^{}){\displaystyle _0^{\mathrm{}}}r^{}𝑑r^{}\left(\omega ^{}\omega \mathrm{cosh}2\mathrm{\Omega }r^{}{\displaystyle \frac{m}{r^{}}}\mathrm{sinh}2\mathrm{\Omega }r^{}\right)J_m^{}(q^{}r^{})J_m(qr^{})`$ (53)
$`\times `$ $`\left(iA_{m,m^{}}(r^{},\omega )\right)^1\left[\mathrm{exp}\left(2\pi iA_{m,m^{}}(r^{},\omega )\right)1\right],`$ (54)
where
$$A_{m,m^{}}(r,\omega )=m^{}+m\mathrm{cosh}2\mathrm{\Omega }r^{}+\omega r^{}\mathrm{sinh}2\mathrm{\Omega }r^{}.$$
(55)
Let us study the high-velocity limit, $`\mathrm{\Omega }\mathrm{}`$, of $`\beta _{00}`$, which is for $`m=0`$ and $`m^{}=0`$:
$$\beta _{00}=2\pi iN_1(m^{}=0)N_2(m=0)\delta (k+k^{})_0^{\mathrm{}}r𝑑r\frac{\omega ^{}\omega \mathrm{cosh}(2\mathrm{\Omega }r)}{\omega r\mathrm{sinh}(2\mathrm{\Omega }r)}J_0(qr)J_0(q^{}r)\left[e^{2\pi i\omega r\mathrm{sinh}(2\mathrm{\Omega }r)}1\right].$$
(56)
As $`\frac{1}{\mathrm{sinh}x}0`$ and $`\frac{\mathrm{cosh}x}{\mathrm{sinh}x}1`$ in the limit $`x\mathrm{}`$, we find:
$$\underset{\mathrm{\Omega }\mathrm{}}{lim}\beta _{00}=2\pi iN_1(m^{}=0)N_2(m=0)\delta (k+k^{})_0^{\mathrm{}}𝑑rJ_0(qr)J_0(q^{}r)\left[1e^{2\pi i\omega r\mathrm{sinh}(2\mathrm{\Omega }r)}\right].$$
(57)
Calling $`K(q,q^{},\omega ,\mathrm{\Omega })`$ the integral above, note that:
$$K(q,q^{},\omega ,\mathrm{\Omega })=\mathrm{}\left\{K(q,q^{},\omega ,\mathrm{\Omega })\right\}+i\mathrm{}\left\{K(q,q^{},\omega ,\mathrm{\Omega })\right\},$$
(58)
where the real and imaginary parts read:
$$\mathrm{}\left\{K(q,q^{},\omega ,\mathrm{\Omega })\right\}=_0^{\mathrm{}}𝑑rJ_0(qr)J_0(q^{}r)\left(1\mathrm{cos}\left[2\pi \omega r\mathrm{sinh}(2\mathrm{\Omega }r)\right]\right)$$
(59)
and
$$\mathrm{}\left\{K(q,q^{},\omega ,\mathrm{\Omega })\right\}=_0^{\mathrm{}}𝑑rJ_0(qr)J_0(q^{}r)\mathrm{sin}\left[2\pi \omega r\mathrm{sinh}(2\mathrm{\Omega }r)\right].$$
(60)
Note that in Eq.(59) above, the second integral is not capable to make $`\mathrm{}\left\{K(q,q^{},\omega ,\mathrm{\Omega })\right\}`$ vanish if the first one is non-zero. And note that this is indeed the case, since one can find the first of them in Gradshteyn , in terms of hypergeometric functions:
$$_0^{\mathrm{}}𝑑rJ_0(qr)J_0(q^{}r)=\frac{1}{q+q^{}}F[\frac{1}{2};\frac{1}{2};1;\frac{4qq^{}}{(q+q^{})^2}],$$
(61)
and it is non-zero.
We now present a different route to prove that $`\beta `$ is different from zero. We start from Eq.(54); as for $`m=m^{}=0`$ the integrands go like $`\frac{1}{\mathrm{sinh}(2\mathrm{\Omega }r)}`$, and $`\frac{1}{\mathrm{sinh}(2\mathrm{\Omega }r)}`$ diverges for $`r0`$, let us calculate $`\beta _{0,1}`$, which is for $`m=0`$ and $`m^{}=1`$:
$`\beta _{01}`$ $`=`$ $`2\pi iN_1(m^{}=1)N_2(m=0)\delta (k+k^{}){\displaystyle _0^{\mathrm{}}}r𝑑r{\displaystyle \frac{\left(\omega ^{}\omega \mathrm{cosh}2\mathrm{\Omega }r\right)}{\left(1+r\omega \mathrm{sinh}2\mathrm{\Omega }r\right)}}J_0(qr)J_1(q^{}r)`$ (62)
$`\times `$ $`\left[\mathrm{exp}\left(2\pi i(1+r\omega \mathrm{sinh}2\mathrm{\Omega }r)\right)1\right].`$ (63)
So we call $`I`$ the integral above:
$$I=_0^{\mathrm{}}r𝑑r\frac{\left(\omega ^{}\omega \mathrm{cosh}2\mathrm{\Omega }r\right)}{\left(1+r\omega \mathrm{sinh}2\mathrm{\Omega }r\right)}J_0(qr)J_1(q^{}r)\left[\mathrm{exp}\left(2\pi i(1+r\omega \mathrm{sinh}2\mathrm{\Omega }r)\right)1\right].$$
(64)
The first integral is the only one which is $`\omega ^{}`$dependent, i.e., it is a function of $`\omega ^{}`$. So, if we prove that it is different from zero, we prove that $`\beta _{01}0`$, since the second integral is not sufficient to make it zero. Let us call it $`I(\omega ^{})`$:
$$I(\omega ^{})=\omega ^{}_0^{\mathrm{}}r𝑑r\frac{J_0(qr)J_1(q^{}r)}{\left(1+r\omega \mathrm{sinh}2\mathrm{\Omega }r\right)}\left[\mathrm{exp}\left(2\pi ir\omega \mathrm{sinh}2\mathrm{\Omega }r\right)1\right].$$
(65)
It is possible to separate $`I(\omega ^{})`$ in its real and imaginary parts and, using the same reasoning as for Eq.(59), we can concentrate in the second integral above, which we call $`I_1(\omega ^{})`$:
$$I_1(\omega ^{})=\omega ^{}_0^{\mathrm{}}r𝑑r\frac{J_0(qr)J_1(q^{}r)}{\left(1+r\omega \mathrm{sinh}2\mathrm{\Omega }r\right)}.$$
(66)
There is not an analytical expression for $`I_1(\omega ^{})`$, but using the Maple one can calculate particular values, such as:
$$_0^{\mathrm{}}r𝑑r\frac{J_0(r)J_1(r)}{\left(1+r\mathrm{sinh}r\right)}=0.183096;$$
(67)
so we proved, very indirectly, that $`\beta 0`$.
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# 1 (a) The upper hystogram shows the yearly number of papers on MMs published from 1973 till 1983 [83P1]. (b) The middle hystogram shows the number of papers of the present bibliography as a function of the year of publication. (c) The lower hystogram shows the yearly number of papers found in the SLAC database with the search command find title monopole# or title dyon# or keyword magnetic monopole as a function of the year of entry in the database.
DFUB 2000-9
Bologna, May 2000
Magnetic Monopole Bibliography
G. Giacomelli<sup>1</sup>, M. Giorgini<sup>1</sup>, T. Lari<sup>1</sup>, M. Ouchrif$`^{\mathrm{𝟏}\mathbf{,}\mathrm{𝟐}}`$,
L. Patrizii<sup>1</sup>, V. Popa<sup>1,3</sup>, P. Spada<sup>1</sup> and V. Togo<sup>1</sup>
<sup>1</sup>Dipartimento di Fisica dell’Università di Bologna and INFN, Sezione di Bologna, I-40127 Bologna, Italy
<sup>2</sup>Faculty of Sciences, University Mohamed I, B.P. 524 Oujda, Morocco
<sup>3</sup>Institute of Space Sciences, Bucharest R-76900, Romania
## Abstract
We present a bibliography compilation on magnetic monopoles updated to include references till the end of year 1999. It is intended to contain nearly all the experimental papers on the subject and only the theoretical papers which have some specific experimental implications.
1. Introduction
Even though Maxwell’s equations formally allow the existence of magnetic monopoles (MMs), interest in this kind of objects arose only in 1931 after the paper of P. A. M. Dirac \[31D1\], in which it was shown that magnetic charges can be introduced in the framework of quantum mechanics and that the product of the basic electric charge and of the basic magnetic charge is quantized according to the Dirac relation $`eg=\frac{n\mathrm{}c}{2}`$, where $`n`$ is an integer. Such a particle is called magnetic monopole if it carries only a magnetic charge, and dyon if it carries both magnetic and electric charges (a monopole bound with a nucleus behaves effectively as a dyon). Dirac could not constrain the monopole mass; rough estimates indicated that the MM mass should be larger than several GeV.
Many types of searches for magnetic monopoles with masses not much larger than the proton mass were performed at each new accelerator and in bulk matter. Very many theoretical studies on MMs have been published.
The other date of fundamental importance in the history of monopoles is 1974. In that year ’t Hooft \[74H1\] and Polyakov \[74P1\] demonstrated that Grand Unified Theories (GUT) of the electroweak and strong interactions implied the existence of MMs with masses of the order of $`10^{17}\text{ GeV/c}^2`$ ($`m_Mm_X/\alpha `$ where $`m_X`$ is the mass of the carrier of the unified force and $`\alpha `$ is the unified coupling constant). These masses are too large for monopoles to be produced at present or future high energy accelerators or somewhere in the present universe. They could have been produced immediately after the big bang, either as topological defects or in very high energy collisions such as $`e^+e^{}M\overline{M}`$, immediately after the phase transition at the end of the GUT epoch; thus GUT monopoles could be present in the cosmic radiation, since the lightest monopole should be stable, due to conservation of magnetic charge. From 1974 to the present time a very large number of theoretical studies were made on magnetic monopoles; also many experimental searches were performed.
The present paper gives a bibliography of publications on monopoles. The bibliography is intended to contain nearly all the experimental papers on the subject and only the theoretical papers which have specific experimental implications. With some exceptions only papers published in international refereed journals have been included.
The publications on MMs comprise many different subjects.
1) Theoretical works on Dirac MMs (for ex. \[31D1\], \[66S1\], \[76W1\], \[77K2\], \[97I2\], …). The possible existence of bound states of magnetic monopoles with electrons and/or nuclei has also been investigated (for ex. \[51M1\], \[77K3\], \[83B1\], \[83B2\], \[84B1\], …).
2) Theoretical works on GUT MMs (for ex. \[74H1\], \[74P1\], \[84P1\], …).
3) Papers on the catalysis of baryon decay by GUT monopoles, such as \[80R1\], \[82C2\].
4) Papers on the cosmological production of MMs (for ex. \[76K1\], \[79P1\], \[80L2\], \[80E1\], \[81G2\], …) and papers which derived limits on MM fluxes from astrophysical considerations (for ex. \[70P1\], \[82K2\], \[85A2\], \[85B2\], …).
5) Theoretical works on supermassive MMs based on other theories, like monopoles from superstrings (\[87L1\]), intermediate symmetry breaking monopoles (\[83L1\], \[84L2\]), or lighter monopoles of electroweak nature (\[97C1\], \[97C4\]).
6) Studies of the energy losses of monopoles in matter and on the possible techniques to detect them (for ex. \[78A1\], \[82A4\], \[83A2\], \[83D3\], \[84B2\], \[85B5\], \[87F1\], \[89P1\], \[91O1\], \[97A1\], …).
7) Searches at high energy accelerators; the searches are either direct (detection of monopoles immediately after their production in high-energy collisions) or indirect (for example when a piece of matter is exposed to a beam for a long time, and then later analyzed). Examples of direct searches are \[75G1\], \[82K1\], \[83A7\], \[83M1\], \[87G2\], \[90B3\], \[00B2\], … Examples of indirect searches are \[61F1\], \[63A1\], \[66A1\], \[74C1\], \[75E1\], \[78C1\], \[83B7\], ….
8) Searches for possible effects of virtual monopoles, such as \[95D2\], \[95A2\], \[97K2\], …
9) Direct searches for MMs in the cosmic radiation. Examples of such searches are those by MACRO \[91B1\], \[91P1\], \[94A1\], \[95A1\], \[95M2\], \[97A2\], \[97A3\] and by other experiments \[82B2\], \[83C1\], \[86P1\], \[90B1\], \[90B2\], \[90B4\], \[90G1\], \[91O1\], \[92T1\].
10) Indirect searches for monopoles in the cosmic radiation; the experiments look for monopoles trapped in matter or for the effects due to the passage of MMs in the past. Examples are \[63G2\], \[73R1\], \[83E1\], \[84P3\], \[86P1\], \[87E1\], \[89A1\], \[90G1\], \[95J1\], …
11) Review papers on various aspects of monopoles and on their experimental search; examples are \[81G3\], \[82G1\], \[83P1\], \[86G2\], \[94G1\], \[96B3\], …
Several MM bibliographies have been made in the past, see \[73S1\], \[77C1\], \[80R1\], \[82C4\], \[84G1\], \[94G1\]. The bibliography in \[84G1\] is more complete than the present one for the period before 1984. The bibliography of \[80R1\] covers essentially all the papers dealing with classical Dirac monopoles.
A bibliography on the experimental limits and on astrophysical bounds is presented regularly every two years by the Particle Data Group (see \[96P1\], \[94P2\], \[92P1\]).
The present bibliography covers the period before May 2000 and is an update of DFUB 98-9 \[98G6\].
Fig. 1a presents an hystogram with the number of papers on MMs and dyons published each year from 1973 to 1983 \[83P1\]; Fig. 1b shows the number of papers in the present bibliography (until December 1999) as a function of the year of publication; Fig. 1c shows an hystogram of the number of yearly papers in the SLAC database which meet one of the following conditions:
\- they have monopole or monopoles in the title;
\- they have dyon or dyons or dyonic in the title;
\- they have the keyword magnetic monopole assigned.
For the first and second hystogram the year is that of publication, for the third it is the year of receipt by the SLAC library; this third hystogram has many more entries, mostly theoretical papers, conference proceedings and unpublished reports.
The peak rates in the hystograms correspond to the periods immediately after the paper of ’t Hooft and Polyakov (1974) and after the 1982 Cabrera candidate event.
Fig. 2 shows a compilation of the 90% C.L. limits on MMs in the cosmic radiation; the limits apply to an isotropic flux of bare $`g=g_D`$ massive magnetic monopoles for a catalysis cross section smaller than few mb.
In the literature one finds references to many types of MMs. For completeness we recall the simplest definitions of many of these MMs.
$``$ GUT monopoles are the MMs connected with the Grand Unification of the electroweak and strong interactions and have masses of $`10^{16}÷10^{17}`$ GeV. They appear in the early universe at the end of the GUT epoch \[74H1\], \[74P1\], \[83F2\], \[84P3\], \[85K2\]; ’t Hooft-Polyakov monopoles are GUT MMs \[90K1\], \[92B1\], \[93D1\], \[98K1\]; SO(3)-Z2 monopoles are particular GUT MMs \[98G4\].
$``$ QCD monopoles are MMs with a colour charge \[98S1\], \[99G1\].
$``$ BPS monopoles are MMs appearing in the Bogomolny and Prasad-Sommerfield limit \[93B2\], \[97B4\], \[97C3\], \[97S1\], \[98B2\], \[99I2\], \[99L1\].
$``$ Kaluza-Klein monopoles are MMs connected with the unification of the GUT interaction with the gravitational interaction; they have typical masses $`>10^{19}`$ GeV \[83S3\], \[86S3\], \[97B3\], \[98B3\].
$``$ Non abelian monopoles are MMs which appear in non abelian theories (including GUT, Kaluza-Klein,…, monopoles) \[94M1\], \[94M2\], \[99C1\], \[00L1\].
$``$ Classical (Dirac) monopoles are abelian MMs hypothesized by Dirac in 1931; they could have relatively low masses \[31D1\], \[33T1\], \[51M1\], \[59B1\], \[72B2\], \[78Z1\], \[90A4\], \[97H1\].
$``$ Intermediate mass monopoles are MMs connected with an intermediate symmetry breaking scale; they could have masses of the order of $`10^{10}`$ GeV \[84L2\], \[95D2\], \[00B1\], \[00B2\].
$``$ Wu-Yang monopoles are particular solutions in Yang-Mills theories \[76W1\], \[98D2\].
$``$ Monopoles in String Theories \[93S1\], \[94B2\], \[99H1\].
$``$ Constituent monopoles are MMs formed by the superposition of $`n`$ BPS MMs (for a SU(n) gauge group) \[98K2\], \[99K1\].
$``$ Complex monopoles are solutions of Yang-Mills theories when adding a Chern-Simons term \[98T1\], \[99H2\], \[99T1\].
For further possibilities, see other references, such as \[98L2\], \[99F1\], \[99T2\].
We gratefully acknowledge many members of the MACRO Collaboration, in particular all the members of the Bologna group. We apologize for possible omissions.
2. References on Magnetic Monopoles
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\[00B3\] V. A. Balkanov et al., (Search for fast magnetic monopoles in a deep-sea experiment at Lake Baikal) Bull. Russ. Acad. Sci. Phys. 63 (1999) 485, Izv. Ross. Akad. Nauk, Fiz. 63 (1999) 598.
\[00D1\] N. M. Davies, V. V. Khoze (On Affleck-Dine-Seiberg superpotential and magnetic monopoles in supersymmetric QCD) JHEP 0001:015 (2000), hep-th/9911112.
\[00D2\] Zh. Dzhilkibaev (Search for fast monopoles in the Baikal experiment) Zeuthen 1998, Simulation and analysis methods for large neutrino telescopes (1998) 364.
\[00G1\] L. Gamberg, K. A. Milton (Eikonal scattering of monopoles and dyons in dual QED) hep-ph/0005016.
\[00H1\] Y. Hosotani, J. Bjoraker (Monopoles and dyons in the pure Einstein-Yang-Mills theory) gr-qc/0001105.
\[00J1\] O. Jahn (Instantons and monopoles in general abelian gauges) J. Phys. A33 (2000) 2997, hep-th/9909004.
\[00J2\] R. Jeannerot, S. Khalil, G. Lazarides, Q. Shafi (Inflation and monopoles in supersymmetric SU(4)C X SU(2)(L) X SU(2)(R)) hep-ph/0002151.
\[00K1\] G. R. Kalbfleisch et al., (Improved experimental limits on the production of magnetic monopoles) hep-ex/0005005.
\[00K2\] B. Kol, M. Kroyter, (On the spatial structure of monopoles) hep-th/0002118.
\[00L1\] N. F. Lepora, (On the spectrum and representation theory of nonabelian monopoles) hep-th/0002163.
\[00L2\] S. L. Liebling (Static gravitational global monopoles) Phys. Rev. D61:024030 (2000), gr-qc/9906014.
\[00M1\] D. Maison (Gravitational global monopoles with horizons) gr-qc/9912100.
\[00W1\] S. D. Wick, T. W. Kephart, T. J. Weiler, P. L. Biermann (Signatures for a cosmic flux of magnetic monopoles) submitted to Astropart. Phys., astro-ph/0001233.
Other MM bibliographies can be found in: \[73S1\], \[77C1\], \[80R1\], \[82C4\], \[84G1\], \[94G1\], \[98G6\].
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# Contents
## 1 Prologue
Continuum strong QCD. The phrase comes to us via Ref. , although that is likely not its debut. In our usage it embraces all continuum nonperturbative methods and models that seek to provide an intuitive understanding of strong interaction phenomena, and particularly those where a direct connection with QCD can be established, in one true limit or another. The community of continuum strong QCD practitioners is a large one and that is helped by the appealing simplicity of models, which facilitates insightful contributions that are not labour intensive.
Our goal is to present an image of the gamut of strong QCD phenomena currently at the forefront of the interaction between experiment and theory, and of contemporary methods and their application. Our portrayal will not be complete but ought to serve as a conduit into this field. We judge that the best means to pursue our aim is to focus on a particular approach, making connections and comparisons with others when possible and helpful. All methods have strengths and weaknesses, and, when modelling is involved, they are not always obviously systematic. That is the cost of leaving perturbation theory behind. They can also yield internally consistent results that nevertheless are without fidelity to QCD. Vigilance is therefore necessary but herein lies another benefit of a leagued community.
Everyone has a favourite tool and the Dyson-Schwinger equations \[DSEs\] are ours: they provide the primary medium for this discourse. The framework is appropriate here because the last decade has seen something of a renaissance in its phenomenological application and we chronicle that herein. Additionally, the DSEs have been applied simultaneously to phenomena as apparently unconnected as low-energy $`\pi \pi `$ scattering, $`BD^{}`$ decays and the equation of state for a quark gluon plasma, and hence they provide a single framework that serves to conduct us through a wide range of hadron physics. Focusing on one approach is no impediment to a broad perspective: the continuum, nonperturbative studies complement each other, with agreement between them being a strong signal of a real effect. In this and their flexibility they provide a foil to the phlegmatic progress of numerical simulations of lattice-QCD, and a fleet guide to unanticipated phenomena.
To close this short introduction we present a list of abbreviations. They are defined in the text when first introduced, however, that point can be difficult to locate.
| BSA | | Bethe-Salpeter amplitude | BSE | | Bethe-Salpeter equation |
| --- | --- | --- | --- | --- | --- |
| DCSB | | dynamical chiral symmetry breaking | DSE | | Dyson-Schwinger equation |
| EOS | | equation of state | LHC | | large hadron collider |
| QCD | | quantum chromodynamics | QC<sub>2</sub>D | | two-colour QCD |
| QED | | quantum electrodynamics | QED<sub>3</sub> | | three-dimensional QED |
| QGP | | quark gluon plasma | RHIC | | relativistic heavy ion collider |
## 2 Dyson-Schwinger Equations
In introducing the DSEs we find useful the brief text book discussions in Chap. 10 of Ref. and Chap. 2 of Ref. , and particularly their use in proving the renormalisability of quantum electrodynamics \[QED\] in Ref. , Chap. 19. However, these sources do not cover the use of DSEs in contemporary nuclear and high-energy physics. Reference ameliorates that with a wide ranging review of the theoretical and phenomenological applications extant when written, and it provides a good foundation for us.
The DSEs are a nonperturbative means of analysing a quantum field theory. Derived from a theory’s Euclidean space generating functional, they are an enumerable infinity of coupled integral equations whose solutions are the $`n`$-point Schwinger functions \[Euclidean Green functions\], which are the same matrix elements estimated in numerical simulations of lattice-QCD. In theories with elementary fermions, the simplest of the DSEs is the gap equation, which is basic to studying dynamical symmetry breaking in systems as disparate as ferromagnets, superconductors and QCD. The gap equation is a good example because it is familiar and has all the properties that characterise each DSE: its solution is a $`2`$-point function \[the fermion propagator\] while its kernel involves higher $`n`$-point functions; e.g., in a gauge theory, the kernel is constructed from the gauge-boson $`2`$-point function and fermion–gauge-boson vertex, a $`3`$-point function; a weak-coupling expansion yields all the diagrams of perturbation theory; and solved self-consistently, the solution of the gap equation exhibits nonperturbative effects unobtainable at any finite order in perturbation theory; e.g, dynamical symmetry breaking.
The coupling between equations; i.e., the fact that the equation for a given $`m`$-point function always involves at least one $`n>m`$-point function, necessitates a truncation of the tower of DSEs in order to define a tractable problem. One systematic and familiar truncation is a weak coupling expansion to reproduce perturbation theory. However, that precludes the study of nonperturbative phenomena and hence something else is needed for the investigation of strongly interacting systems, bound state phenomena and phase transitions.
In analysing the ferromagnetic transition, the Hartree-Fock approximation yields qualitatively reliable information and in QED and QCD its analogue: the rainbow truncation, has proven efficacious. However, a priori it can be difficult to judge whether a given truncation will yield reliable results and a systematic improvement is not always obvious. It is here that some model-dependence enters but that is not new, being typical in the study of strongly-interacting few- and many-body systems. To proceed with the DSEs one just employs a truncation and explores its consequences, applying it to different systems and constraining it, where possible, by comparisons with experimental data and other theoretical approaches on their common domain. In this way a reliable truncation can be identified and then attention paid to understanding the keystone of its success and improving its foundation. This pragmatic approach has proven rewarding in strong QCD, as we shall describe.
### 2.1 DSE Primer
Lattice-QCD is defined in Euclidean space because the zero chemical potential \[$`\mu =0`$\] Euclidean QCD action defines a probability measure, for which many numerical simulation algorithms are available. The Gaussian distribution:
$$𝒦_t(q,q^{})=(2\pi t)^{3/2}\mathrm{exp}[\left(qq^{}\right)^2/(2t)],$$
(2.1.1)
defines the simplest probability measure: $`dq^{}𝒦_t(q,q^{})`$. $`𝒦_t`$ is positive and normalisable, which allows its interpretation as a probability density. An heuristic exposition of probability measures in quantum field theory can be found in Ref. , Chap. 6, while Ref. , Chaps. 3 and 6, provides a more rigorous discussion in the context of quantum mechanics and quantum field theory.
Working in Euclidean space, however, is more than simply pragmatic: Euclidean lattice field theory is currently a primary candidate for the rigorous definition of an interacting quantum field theory and that relies on it being possible to define the generating functional via a proper limiting procedure. The moments of the measure; i.e., vacuum expectation values of the fields, are the $`n`$-point Schwinger functions and the quantum field theory is completely determined once all its Schwinger functions are known. The time-ordered Green functions of the associated Minkowski space theory can be obtained in a well-defined fashion from the Schwinger functions. This is one reason why we employ a Euclidean formulation. Another is a desire to maintain contact with perturbation theory where the renormalisation group equations for QCD and their solutions are best understood .
To make clear our conventions: for $`4`$-vectors $`a`$, $`b`$:
$$ab:=a_\mu b_\nu \delta _{\mu \nu }:=\underset{i=1}{\overset{4}{}}a_ib_i,$$
(2.1.2)
so that a spacelike vector, $`Q_\mu `$, has $`Q^2>0`$; our Dirac matrices are Hermitian and defined by the algebra
$$\{\gamma _\mu ,\gamma _\nu \}=2\delta _{\mu \nu };$$
(2.1.3)
and we use
$$\gamma _5:=\gamma _1\gamma _2\gamma _3\gamma _4$$
(2.1.4)
so that
$$\mathrm{tr}\left[\gamma _5\gamma _\mu \gamma _\nu \gamma _\rho \gamma _\sigma \right]=4\epsilon _{\mu \nu \rho \sigma },\epsilon _{1234}=1.$$
(2.1.5)
The Dirac-like representation of these matrices is:
$$\stackrel{}{\gamma }=\left(\begin{array}{cc}0& i\stackrel{}{\tau }\\ i\stackrel{}{\tau }& 0\end{array}\right),\gamma _4=\left(\begin{array}{cc}\tau ^0& 0\\ 0& \tau ^0\end{array}\right),$$
(2.1.6)
where the $`2\times 2`$ Pauli matrices are:
$$\tau ^0=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),\tau ^1=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),\tau ^2=\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right),\tau ^3=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).$$
(2.1.7)
Using these conventions the \[unrenormalised\] Euclidean QCD action is
$$S[\overline{q},q,A]=d^4x\left[\frac{1}{4}F_{\mu \nu }^aF_{\mu \nu }^a+\frac{1}{2\xi }A^aA^a+\underset{f=1}{\overset{N_f}{}}\overline{q}_f\left(\gamma +m_f+ig\frac{1}{2}\lambda ^a\gamma A^a\right)q_f\right],$$
(2.1.8)
where: $`F_{\mu \nu }^a=_\mu A_\nu ^a_\nu A_\mu ^agf^{abc}A_\mu ^bA_\nu ^s`$; $`N_f`$ is the number of quark flavours; $`m_f`$ are the current-quark masses; $`\{\lambda ^a:a=1,\mathrm{},8\}`$ with $`[\lambda ^a,\lambda ^b]=2if^{abc}\lambda ^c`$ are the Gell-Mann matrices for $`SU(3)`$ colour; and $`\xi `$ is the covariant gauge fixing parameter. The generating functional follows:
$$𝒵[\overline{\eta },\eta ,J]=𝑑\mu (\overline{q},q,A,\overline{\omega },\omega )\mathrm{exp}d^4x\left[\overline{q}\eta +\overline{\eta }q+J_\mu ^aA_\mu ^a\right],$$
(2.1.9)
with sources: $`\overline{\eta }`$, $`\eta `$, $`J`$, and a functional integral measure
$`d\mu (\overline{q},q,A,\overline{\omega },\omega ):=`$
$`{\displaystyle \underset{x}{}}{\displaystyle \underset{\varphi }{}}𝒟\overline{q}_\varphi (x)𝒟q_\varphi (x){\displaystyle \underset{a}{}}𝒟\overline{\omega }^a(x)𝒟\omega ^a(x){\displaystyle \underset{\mu }{}}𝒟A_\mu ^a(x)\mathrm{exp}(S[\overline{q},q,A]S^g[\overline{\omega },\omega ,A]),`$
where $`\varphi `$ represents both the flavour and colour index of the quark field, and $`\overline{\omega }`$ and $`\omega `$ are scalar, Grassmann \[ghost\] fields that are a necessary addition in covariant gauges, and most other gauges too. \[Without gauge fixing the action is constant along trajectories of gauge-equivalent gluon field configurations, which leads to a gauge-orbit volume-divergence in the continuum generating functional.\] The normalisation
$$𝒵[\overline{\eta }=0,\eta =0,J=0]=1$$
(2.1.11)
is implicit in the measure.
The ghosts only couple directly to the gauge field:
$$S_g[\overline{\omega },\omega ,A]=d^4x\left[_\mu \overline{\omega }^a_\mu \omega ^agf^{abc}_\mu \overline{\omega }^a\omega ^bA_\mu ^c\right],$$
(2.1.12)
and restore unitarity in the subspace of transverse \[physical\] gauge fields. Practically, the normalisation means that ghost fields are unnecessary in the calculation of gauge invariant observables using lattice-regularised QCD because the gauge-orbit volume-divergence in the generating functional, associated with the uncountable infinity of gauge-equivalent gluon field configurations in the continuum, is rendered finite by the simple expedient of only summing over a finite number of configurations.
It is not necessary to employ the covariant gauge fixing condition in constructing the measure although it does yield DSEs with a simple form. Indeed a rigorous definition of the measure may require a gauge fixing functional that either completely eliminates Gribov copies or restricts the functional integration domain to a subspace without them. Concerned, as they are, with unobservable degrees of freedom, such modifications of the measure cannot directly affect the colour-singlet $`n`$-point functions that describe physical observables. However, they have the capacity to modify the infrared behaviour of coloured $`n`$-point functions and thereby influence the manner in which we intuitively understand observable effects.
The DSEs are derived from the generating functional using the elementary observation that, with sensible Euclidean space boundary conditions, the integral of a total derivative is zero:
$$0=\frac{\delta }{\delta A}𝑑\mu (\overline{q},q,A,\overline{\omega },\omega )\mathrm{exp}d^4x\left[\overline{q}\eta +\overline{\eta }q+J_\mu ^aA_\mu ^a\right].$$
(2.1.13)
Examples of such derivations: the QED gap equation, the equation for the photon vacuum polarisation and that for the fermion-fermion scattering matrix \[a $`4`$-point function\], are presented in Ref. , Chap. 10. The scattering matrix is important because it lies at the heart of two-body bound state studies in quantum field theory, being a keystone of the Bethe-Salpeter equation. Section 2.1 of Ref. repeats the first two derivations and also makes explicit the effects of renormalisation. Herein we simply begin with the relevant DSE leaving its derivation as an exercise. That too can be side-stepped if a Minkowski space version of the desired DSE is at hand. Then one need merely employ the transcription rules:
Configuration Space
1. $`{\displaystyle ^M}d^4x^Mi{\displaystyle ^E}d^4x^E`$
2. $`/i\gamma ^E^E`$
3. $`/Ai\gamma ^EA^E`$
4. $`A_\mu B^\mu A^EB^E`$
5. $`x^\mu _\mu x^E^E`$
Momentum Space
1. $`{\displaystyle ^M}d^4k^Mi{\displaystyle ^E}d^4k^E`$
2. $`/ki\gamma ^Ek^E`$
3. $`/Ai\gamma ^EA^E`$
4. $`k_\mu q^\mu k^Eq^E`$
5. $`k_\mu x^\mu k^Ex^E`$
These rules are valid in perturbation theory; i.e., the correct Minkowski space integral for a given diagram will be obtained by applying these rules to the Euclidean integral: they take account of the change of variables and rotation of the contour. However, for the diagrams that represent DSEs, which involve dressed $`n`$-point functions whose analytic structure is not known a priori, the Minkowski space equation obtained using this prescription will have the right appearance but it’s solutions may bear no relation to the analytic continuation of the solution of the Euclidean equation.
### 2.2 Core Issues
The use of DSEs as a unifying, phenomenological tool in QCD has grown much since the publication of Ref. , and Ref. reviews applications to electromagnetic and hadronic interactions of light mesons and the connection to the successful Nambu-Jona-Lasinio and Global-Colour models. In most of their widespread applications these two models yield results kindred to those obtained with the rainbow-ladder DSE truncation.
Quark DSE \[Gap Equation\] The idea that the fermion DSE, or gap equation, can be used to study dynamical chiral symmetry breaking \[DCSB\] perhaps began with a study of the electron propagator . Since that analysis there have been many studies of dynamical mass generation in QED, and those reviewed in Ref. establish the context: there is some simplicity in dealing with an Abelian theory. Continuing research has led to an incipient understanding of the role played by multiplicative renormalisability and gauge covariance in constraining the dressed-fermion-photon vertex , and progress in this direction provides intuitive guidance for moving beyond the rainbow truncation. However, the studies have also made plain the difficulty in defining the chiral limit of a theory without asymptotic freedom. Renormalised, strong coupling, quenched QED yields a scalar self energy for the electron that is not positive-definite: damped oscillations appear after the renormalisation point . \[Quenched $`=`$ bare photon propagator in the gap equation. It is the simplest and most widely explored truncation of the gauge sector.\] This pathology can be interpreted as a signal that four fermion operators $`[\overline{\psi }(x)\psi (x)]^2`$ have acquired a large anomalous dimension and have thus become relevant operators for the range of gauge couplings that support DCSB. That hypothesis has implications for the triviality of the theory, and is reviewed and explored in Refs. . Nonperturbative QED remains an instructive challenge.
The study of DCSB in QCD is much simpler primarily because the chiral limit is well-defined. We discuss that now using the renormalised quark-DSE:
$`S(p)^1`$ $`=`$ $`Z_2(i\gamma p+m_{\mathrm{bare}})+Z_1{\displaystyle _q^\mathrm{\Lambda }}g^2D_{\mu \nu }(pq){\displaystyle \frac{\lambda ^a}{2}}\gamma _\mu S(q)\mathrm{\Gamma }_\nu ^a(q,p),`$ (2.2.1)
where $`D_{\mu \nu }(k)`$ is the renormalised dressed-gluon propagator, $`\mathrm{\Gamma }_\nu ^a(q;p)`$ is the renormalised dressed-quark-gluon vertex, $`m_{\mathrm{bare}}`$ is the $`\mathrm{\Lambda }`$-dependent current-quark bare mass that appears in the Lagrangian and $`_q^\mathrm{\Lambda }:=^\mathrm{\Lambda }d^4q/(2\pi )^4`$ represents mnemonically a translationally-invariant regularisation of the integral, with $`\mathrm{\Lambda }`$ the regularisation mass-scale. The final stage of any calculation is to remove the regularisation by taking the limit $`\mathrm{\Lambda }\mathrm{}`$. Using a translationally invariant regularisation makes possible the preservation of Ward-Takahashi identities, which is crucial in studying DCSB, for example. One implementation well-suited to a nonperturbative solution of the DSE is Pauli-Villars regularisation, which has the quark interacting with an additional massive gluon-like vector boson: mass$`\mathrm{\Lambda }`$, that decouples as $`\mathrm{\Lambda }\mathrm{}`$ . An alternative is a numerical implementation of dimensional regularisation, which, although more cumbersome, can provide the necessary check of scheme-independence .
In Eq. (2.2.1), $`Z_1(\zeta ^2,\mathrm{\Lambda }^2)`$ and $`Z_2(\zeta ^2,\mathrm{\Lambda }^2)`$ are the quark-gluon-vertex and quark wave function renormalisation constants, which depend on the renormalisation point, $`\zeta `$, and the regularisation mass-scale, as does the mass renormalisation constant
$$Z_m(\zeta ^2,\mathrm{\Lambda }^2)=Z_4(\zeta ^2,\mathrm{\Lambda }^2)/Z_2(\zeta ^2,\mathrm{\Lambda }^2),$$
(2.2.2)
with the renormalised mass given by
$$m(\zeta ):=m_{\mathrm{bare}}(\mathrm{\Lambda })/Z_m(\zeta ^2,\mathrm{\Lambda }^2).$$
(2.2.3)
Although we have suppressed the flavour label, $`S`$, $`\mathrm{\Gamma }_\mu ^a`$ and $`m_{\mathrm{bare}}`$ depend on it. However, one can always use a flavour-independent renormalisation scheme, which we assume herein, and hence all the renormalisation constants are flavour-independent .
The solution of Eq. (2.2.1) has the form
$$S(p)^1=i\gamma pA(p^2,\zeta ^2)+B(p^2,\zeta ^2)=\frac{1}{Z(p^2,\zeta ^2)}\left[i\gamma p+M(p^2,\zeta ^2)\right].$$
(2.2.4)
The functions $`A(p^2,\zeta ^2)`$, $`B(p^2,\zeta ^2)`$ embody the effects of vector and scalar quark-dressing induced by the quark’s interaction with its own gluon field. The ratio: $`M(p^2,\zeta ^2)`$, is the quark mass function and a pole mass would be the solution of
$$m_{\mathrm{pole}}^2M^2(p^2=m_{\mathrm{pole}}^2,\zeta ^2)=0.$$
(2.2.5)
A widely posed conjecture is that confinement rules out a solution of this equation . We discuss this further below.
Equation (2.2.1) must be solved subject to a renormalisation \[boundary\] condition, and because the theory is asymptotically free it is practical and useful to impose the requirement that at a large spacelike $`\zeta ^2`$
$$S(p)^1|_{p^2=\zeta ^2}=i\gamma p+m(\zeta ),$$
(2.2.6)
where $`m(\zeta )`$ is the renormalised current-quark mass at the scale $`\zeta `$. By “large” here we mean $`\zeta ^2\mathrm{\Lambda }_{\mathrm{QCD}}^2`$ so that in quantitative, model studies extensive use can be made of matching with the results of perturbation theory. It is the ultraviolet stability of QCD; i.e., the fact that perturbation theory is valid at large spacelike momenta, that makes possible a straightforward definition of the chiral limit. It also provides the starkest contrast to strong coupling QED.
Multiplicative renormalisability in gauge theories entails that
$$\frac{A(p^2,\zeta ^2)}{A(p^2,\overline{\zeta }^2)\text{}}=\frac{Z_2(\zeta ^2,\mathrm{\Lambda }^2)}{Z_2(\overline{\zeta }^2,\mathrm{\Lambda }^2)\text{}}=A(\overline{\zeta }^2,\zeta ^2)=\frac{1}{A(\zeta ^2,\overline{\zeta }^2)\text{}}$$
(2.2.7)
and beginning with Ref. this relation has been used efficaciously to build realistic Ansätze for the fermion–photon vertex in quenched QED. A systematic approach to such nonperturbative improvements is developing and these improvements continue to provide intuitive guidance in QED, where they complement the perturbative calculation of the vertex . They are also useful in exploring model dependence in QCD studies.
At one loop in QCD perturbation theory
$$Z_2(\zeta ^2,\mathrm{\Lambda }^2)=\left[\frac{\alpha (\mathrm{\Lambda }^2)}{\alpha (\zeta ^2)}\right]^{\gamma _F/\beta _1},\gamma _F=\frac{2}{3}\xi ,\beta _1=\frac{1}{3}N_f\frac{11}{2},$$
(2.2.8)
and at this order the running strong coupling is
$$\alpha (\zeta ^2)=\frac{\pi }{\text{}\frac{1}{2}\beta _1\mathrm{ln}\left[\zeta ^2/\mathrm{\Lambda }_{\mathrm{QCD}}^2\right]}.$$
(2.2.9)
In Landau gauge: $`\xi =0`$, so $`Z_21`$ at one loop order. This, plus the fact that Landau gauge is a fixed point of the renormalisation group \[Eq. (2.2.20)\], makes it the most useful covariant gauge for model studies. It also underlies the quantitative accuracy of Landau gauge, rainbow truncation estimates of the critical coupling in strong QED . In a self consistent solution of Eq. (2.2.1), $`Z_21`$ even in Landau gauge but, at large $`\zeta ^2`$, the $`\zeta `$-dependence is very weak. However, as will become evident, in studies of realistic QCD models this dependence becomes significant for $`\zeta ^2<1`$$`2`$GeV<sup>2</sup>, and is driven by the same effect that causes DCSB.
The dressed-quark mass function: $`M(p^2,\zeta ^2)=B(p^2,\zeta ^2)/A(p^2,\zeta ^2)`$, is independent of the renormalisation point; i.e., with $`\zeta \overline{\zeta }`$
$$M(p^2,\zeta ^2)=M(p^2,\overline{\zeta }^2):=M(p^2),p^2:$$
(2.2.10)
it is a function only of $`p^2/\mathrm{\Lambda }_{\mathrm{QCD}}^2`$, which is another constraint on models. At one loop order the running \[or renormalised\] mass
$$m(\zeta )=M(\zeta ^2)=\frac{\widehat{m}}{\left(\text{}\frac{1}{2}\mathrm{ln}\left[\zeta ^2/\mathrm{\Lambda }_{\mathrm{QCD}}^2\right]\right)^{\gamma _m}},\gamma _m=12/(332N_f),$$
(2.2.11)
where $`\widehat{m}`$ is the renormalisation point independent current-quark mass, and the mass renormalisation constant is, Eq. (2.2.2),
$$Z_m(\zeta ^2,\mathrm{\Lambda }^2)=\left[\frac{\alpha (\mathrm{\Lambda }^2)}{\alpha (\zeta ^2)}\right]^{\gamma _m}.$$
(2.2.12)
The mass anomalous dimension, $`\gamma _m`$, is independent of the gauge parameter to all orders perturbation theory and for two different quark flavours the ratio: $`m_{f_1}(\zeta )/m_{f_2}(\zeta )=\widehat{m}_{f_1}/\widehat{m}_{f_2}`$, which is independent of the renormalisation point and of the renormalisation scheme. The chiral limit is unambiguously defined by
$$\mathrm{𝐜𝐡𝐢𝐫𝐚𝐥}\mathrm{𝐥𝐢𝐦𝐢𝐭}:\widehat{m}=0.$$
(2.2.13)
In this case there is no perturbative contribution to the scalar piece of the quark self energy; i.e., $`B(p^2,\zeta ^2)0`$ at every order in perturbation theory and in fact there is no scalar mass-like divergence in the calculation of the self energy. This is manifest in the quark DSE, with Eq. (2.2.1) yielding, in addition to the perturbative result: $`B(p^2,\zeta ^2)0`$, a solution $`M(p^2)=B(p^2,\zeta ^2)/A(p^2,\zeta ^2)0`$ that is power-law suppressed in the ultraviolet: $`M(p^2)1/p^2`$, guaranteeing convergence of the associated integral without subtraction. Dynamical chiral symmetry breaking is
$$\mathrm{𝐃𝐂𝐒𝐁}:M(p^2)0\mathrm{when}\widehat{m}=0.$$
(2.2.14)
As we shall see, in QCD this is possible if and only if the quark condensate is nonzero. The criteria are equivalent.
The solution of the quark DSE depends on the anatomy of the dressed-gluon propagator, which in concert with the dressed-quark-gluon vertex encodes in Eq. (2.2.1) all effects of the quark-quark interaction. In a covariant gauge the renormalised propagator is
$$D_{\mu \nu }(k)=\left(\delta _{\mu \nu }\frac{k_\mu k_\nu }{k^2}\right)\frac{d(k^2,\zeta ^2)}{k^2}+\xi \frac{k_\mu k_\nu }{k^4},$$
(2.2.15)
where $`d(k^2,\zeta ^2)=1/[1+\mathrm{\Pi }(k^2,\zeta ^2)]`$, with $`\mathrm{\Pi }(k^2,\zeta ^2)`$ the renormalised gluon vacuum polarisation for which the conventional renormalisation condition is
$$\mathrm{\Pi }(\zeta ^2,\zeta ^2)=0;\mathrm{i}.\mathrm{e}.,d(\zeta ^2,\zeta ^2)=1.$$
(2.2.16)
For the dressed gluon propagator, multiplicative renormalisability entails
$$\frac{d(k^2,\zeta ^2)}{d(k^2,\overline{\zeta }^2)\text{}}=\frac{Z_3(\overline{\zeta }^2,\mathrm{\Lambda }^2)}{Z_3(\zeta ^2,\mathrm{\Lambda }^2)\text{}}=d(\zeta ^2,\overline{\zeta }^2)=\frac{1}{d(\overline{\zeta }^2,\zeta ^2)\text{}},$$
(2.2.17)
and at one loop in perturbation theory
$$Z_3(\zeta ^2,\mathrm{\Lambda }^2)=\left[\frac{\alpha (\mathrm{\Lambda }^2)}{\alpha (\zeta ^2)}\right]^{\gamma _1/\beta _1},\gamma _1=\frac{1}{3}N_f\frac{1}{4}(133\xi ).$$
(2.2.18)
The gauge parameter is also renormalisation point dependent; i.e., the renormalised theory has a running gauge parameter. However, because of Becchi-Rouet-Stora \[BRST or gauge\] invariance, there is no new dynamical information in that: its evolution is completely determined by the gluon wave function renormalisation constant
$$\xi (\zeta ^2)=Z_3^1(\zeta ^2,\mathrm{\Lambda }^2)\xi _{\mathrm{bare}}(\mathrm{\Lambda }).$$
(2.2.19)
One can express $`\xi (\zeta ^2)`$ in terms of a renormalisation point invariant gauge parameter: $`\widehat{\xi }`$, which is an overall multiplicative factor in the formula and hence
$$\mathrm{𝐋𝐚𝐧𝐝𝐚𝐮}\mathrm{𝐆𝐚𝐮𝐠𝐞}:\widehat{\xi }=0\xi (\zeta ^2)0$$
(2.2.20)
at all orders in perturbation theory; i.e., Landau gauge is a fixed point of the renormalisation group.
The renormalised dressed-quark-gluon vertex has the form
$$\mathrm{\Gamma }_\nu ^a(k,p)=\frac{\lambda ^a}{2}\mathrm{\Gamma }_\nu (k,p);$$
(2.2.21)
i.e., the colour matrix structure factorises. It is a fully amputated vertex, which means all the analytic structure associated with elementary excitations has been eliminated. To discuss this further we introduce the notion of a particle-like singularity. It is one of the form: $`P=kp`$,
$$\frac{1}{(P^2+b^2)^x},x(0,1].$$
(2.2.22)
If the vertex possesses such a singularity then it can be expressed in terms of non-negative spectral densities, which is impossible if $`x>1`$. $`x=1`$ is the ideal case of an isolated $`\delta `$-function distribution in the spectral densities and hence an isolated free-particle pole. $`x(0,1)`$ corresponds to an accumulation at the particle pole of branch points associated with multiparticle production, as occurs with the electron propagator in QED because of photon dressing.
The vertex is a fully amputated $`3`$-point function. Hence, the presence of such a singularity entails the existence of a flavour singlet composite (quark-antiquark bound state) with colour octet quantum numbers and mass $`m=b`$. The bound state amplitude follows immediately from the associated homogeneous Bethe-Salpeter equation, which the singularity allows one to derive. Such an excitation must not exist as an asymptotic state, that would violate the observational evidence of confinement, and any modelling of $`\mathrm{\Gamma }_\mu ^a(k,p)`$ ought to be consistent with this.
Expressing the Dirac structure of $`\mathrm{\Gamma }_\nu (k,p)`$ requires $`12`$ independent scalar functions:
$$\mathrm{\Gamma }_\nu (k,p)=\gamma _\nu F_1(k,p,\zeta )+\mathrm{},$$
(2.2.23)
which for our purposes it is not necessary to reproduce fully. A pedagogical discussion of the perturbative calculation of $`\mathrm{\Gamma }_\nu (k,p)`$ can be found in Ref. while Refs. explore its nonperturbative structure and properties. We only make $`F_1(k,p,\zeta )`$ explicit because the renormalisability of QCD entails that it alone is ultraviolet divergent. Defining
$$f_1(k^2,\zeta ^2):=F_1(k,k,\zeta ),$$
(2.2.24)
the conventional renormalisation boundary condition is
$$f_1(\zeta ^2,\zeta ^2)=1,$$
(2.2.25)
which is practical because QCD is asymptotically free. Multiplicative renormalisability entails
$$\frac{f_1(k^2,\zeta ^2)}{f_1(k^2,\overline{\zeta }^2)\text{}}=\frac{Z_1(\zeta ^2,\mathrm{\Lambda }^2)}{Z_1(\overline{\zeta }^2,\mathrm{\Lambda }^2)\text{}}=f_1(\overline{\zeta }^2,\zeta ^2)=\frac{1}{f_1(\zeta ^2,\overline{\zeta }^2)\text{}},$$
(2.2.26)
and at one loop order
$$Z_1(\mu ^2,\mathrm{\Lambda }^2)=\left[\frac{\alpha (\mathrm{\Lambda }^2)}{\alpha (\mu ^2)}\right]^{\gamma _\mathrm{\Gamma }/\beta _1},\gamma _\mathrm{\Gamma }=\frac{1}{2}[\frac{3}{4}(3+\xi )+\frac{4}{3}\xi ].$$
(2.2.27)
Dynamical Chiral Symmetry Breaking At this point each element in the quark DSE, Eq. (2.2.1), is defined, with some of their perturbative properties elucidated, and the question is how does that provide an understanding of DCSB? It is best answered using an example, in which the model-independent aspects are made clear.
The quark DSE is an integral equation and hence its elements must be known at all values of their momentum arguments, not just in the perturbative domain but also in the infrared. While the gluon propagator and quark-gluon vertex each satisfy their own DSE, that couples the quark DSE to other members of the tower of equations and hinders rather than helps in solving the gap equation. Therefore here, as with all applications of the gap equation, one employs Ansätze for the interaction elements \[$`D_{\mu \nu }(k)`$ and $`\mathrm{\Gamma }_\nu (k,p)`$\], constrained as much and on as large a domain as possible. This approach has a long history in exploring QCD and we illustrate it using the model of Ref. .
The renormalised dressed-ladder truncation of the quark-antiquark scattering kernel \[$`4`$-point function\] is
$$\overline{K}(p,q;P)_{tu}^{rs}=g^2(\zeta ^2)D_{\mu \nu }(pq)\left[\text{}\mathrm{\Gamma }_\mu ^a(p_+,q_+)S(q_+)\right]_{tr}\left[\text{}S(q_{})\mathrm{\Gamma }_\nu ^a(q_{},p_{})\right]_{su},$$
(2.2.28)
where $`p_\pm =p\pm P/2`$, $`q_\pm =q\pm P/2`$, with $`P`$ the total momentum of the quark-antiquark pair, and although we use it now we have suppressed the $`\zeta `$-dependence of the Schwinger functions. From Eqs. (2.2.16-2.2.18) it follows that for $`Q^2:=(pq)^2`$ large and spacelike
$$d(Q^2,\zeta ^2)=\frac{Z_3(\zeta ^2,\mathrm{\Lambda }^2)}{Z_3(Q^2,\mathrm{\Lambda }^2)}d(\zeta ^2,\zeta ^2)=\left[\frac{\alpha (Q^2)}{\alpha (\zeta ^2)}\right]^{\gamma _1/\beta _1}D_{\mu \nu }(pq)=\left[\frac{\alpha (Q^2)}{\alpha (\zeta ^2)}\right]^{\gamma _1/\beta _1}D_{\mu \nu }^{\mathrm{free}}(pq).$$
(2.2.29)
Using this and analogous results for the other Schwinger functions then on the kinematic domain for which $`Q^2p^2q^2`$ is large and spacelike \[$`g^2(\zeta ^2):=4\pi \alpha (\zeta ^2)`$\]
$$\overline{K}(p,q;P)_{tu}^{rs}4\pi \alpha (Q^2)D_{\mu \nu }^{\mathrm{free}}(pq)\left[\text{}\frac{1}{2}\lambda ^a\gamma _\mu S^{\mathrm{free}}(q_+)\right]_{tr}\left[\text{}S^{\mathrm{free}}(q_{})\frac{1}{2}\lambda ^a\gamma _\nu \right]_{su},$$
(2.2.30)
because Eqs. (2.2.8), (2.2.18) and (2.2.27) yield
$$\frac{2\gamma _F}{\beta _1}+\frac{\gamma _1}{\beta _1}\frac{2\gamma _\mathrm{\Gamma }}{\beta _1}=1.$$
(2.2.31)
This is one way of understanding the origin of an often used Ansatz in studies of the gap equation; i.e., making the replacement
$$g^2D_{\mu \nu }(k)4\pi \alpha (k^2)D_{\mu \nu }^{\mathrm{free}}(k)$$
(2.2.32)
in Eq. (2.2.1), and using the “rainbow truncation:”
$$\mathrm{\Gamma }_\nu (q,p)=\gamma _\nu .$$
(2.2.33)
Equation (2.2.32) is often described as the “Abelian approximation” because the left- and right-hand-sides \[r.h.s.\] are equal in QED. In QCD, equality between the two sides cannot be obtained easily by a selective resummation of diagrams. As reviewed in Ref. , Eqs. (5.1-5.8), it can only be achieved by enforcing equality between the renormalisation constants for the ghost-gluon vertex and ghost wave function: $`\stackrel{~}{Z}_1=\stackrel{~}{Z}_3`$. A mutually consistent constraint, which follows formally from $`\stackrel{~}{Z}_1=\stackrel{~}{Z}_3`$, is to enforce the Abelian Ward identity: $`Z_1=Z_2`$. At one-loop this corresponds to neglecting the contribution of the 3-gluon vertex to $`\mathrm{\Gamma }_\nu `$, in which case $`\gamma _\mathrm{\Gamma }\frac{2}{3}\xi =\gamma _F`$. This additional constraint provides the basis for extensions of Eq. (2.2.33); i.e., using Ansätze for $`\mathrm{\Gamma }_\nu `$ that are consistent with the QED vector Ward-Takahashi identity; e.g., Ref. .
Arguments such as these inspire the following Ansatz for the kernel in Eq. (2.2.1:
$$Z_1_q^\mathrm{\Lambda }g^2D_{\mu \nu }(pq)\frac{\lambda ^a}{2}\gamma _\mu S(q)\mathrm{\Gamma }_\nu ^a(q,p)_q^\mathrm{\Lambda }𝒢((pq)^2)D_{\mu \nu }^{\mathrm{free}}(pq)\frac{\lambda ^a}{2}\gamma _\mu S(q)\frac{\lambda ^a}{2}\gamma _\nu ,$$
(2.2.34)
with the ultraviolet behaviour of the “effective coupling:” $`𝒢(k^2)`$, fixed by that of the strong running coupling. Since it is not possible to calculate $`Z_1`$ nonperturbatively without analysing the DSE for the dressed-quark-gluon vertex, this Ansatz absorbs it in the model effective coupling.
Equation (2.2.34) is a model for the product of the dressed-propagator and dressed-vertex and its definition is complete once the behaviour of $`𝒢(k^2)`$ in the infrared is specified; i.e., for $`k^2<1`$-$`2`$GeV<sup>2</sup>. Reference used
$$\frac{𝒢(k^2)}{k^2}=8\pi ^4D\delta ^4(k)+\frac{4\pi ^2}{\omega ^6}Dk^2\mathrm{e}^{k^2/\omega ^2}+4\pi \frac{\gamma _m\pi }{\frac{1}{2}\mathrm{ln}\left[\tau +\left(1+k^2/\mathrm{\Lambda }_{\mathrm{QCD}}^2\right)^2\right]}(k^2),$$
(2.2.35)
with $`(k^2)=[1\mathrm{exp}(k^2/[4m_t^2])]/k^2`$ and $`\tau =\mathrm{e}^21`$. For $`N_f=4`$, $`\mathrm{\Lambda }_{\mathrm{QCD}}^{N_f=4}=0.234\mathrm{GeV}`$. The qualitative features of Eq. (2.2.35) are plain. The first term is an integrable infrared singularity and the second is a finite-width approximation to $`\delta ^4(k)`$, normalised such that it has the same $`d^4k`$ as the first term. In this way the infrared strength is split into the sum of a zero-width and a finite-width piece. The last term in Eq. (2.2.35) is proportional to $`\alpha (k^2)/k^2`$ at large spacelike-$`k^2`$ and has no singularity on the real-$`k^2`$ axis.
There are ostensibly three parameters in Eq. (2.2.35): $`D`$, $`\omega `$ and $`m_t`$. However, in Ref. the authors fixed $`\omega =0.3`$GeV$`(=1/[.66\mathrm{fm}])`$ and $`m_t=0.5`$GeV$`(=1/[.39\mathrm{fm}])`$, and only varied $`D`$ and the renormalised $`u=d`$\- and $`s`$-current-quark masses in an attempt to obtain a good description of low-energy $`\pi `$\- and $`K`$-meson properties, using a renormalisation point $`\zeta =19`$GeV that is large enough to be in the perturbative domain. \[The numerical values of $`\omega `$ and $`m_t`$ are chosen so as to ensure that $`𝒢(k^2)4\pi \alpha (k^2)`$ for $`k^2>2`$GeV<sup>2</sup>. Minor variations in $`\omega `$ and $`m_t`$ can be compensated by small changes in $`D`$.\] Such a procedure could self-consistently yield $`D=0`$, which would indicate that agreement with observable phenomena precludes an infrared enhancement in the effective interaction. However, that was not the case and a good fit required
$$D=(0.884\mathrm{GeV})^2,$$
(2.2.36)
with renormalised current-quark masses
$$\begin{array}{cc}m_{u,d}(\zeta )=3.74\mathrm{MeV},& m_s(\zeta )=82.5\mathrm{MeV},\end{array}$$
(2.2.37)
which are in the ratio $`1`$:$`\mathrm{\hspace{0.17em}22}`$, and yielded, in MeV,
$$\begin{array}{ccccc}& m_\pi & m_K& f_\pi & f_K\\ & & & & \\ \mathrm{Calc}.\text{[20]}\hfill & 139& 497& 131& 154\\ \mathrm{Expt}.\text{[34]}\hfill & 139& 496& 131& 160\end{array}$$
(2.2.38)
and other quantities to be described below. An explanation of how this fit was accomplished requires a discussion of the homogeneous Bethe-Salpeter equation, which we postpone. Here we instead focus on describing the properties of the DSE solution obtained with these parameter values.
Using Eqs. (2.2.1-2.2.3) and (2.2.34) the gap equation can be written
$`S(p,\zeta )^1`$ $`=`$ $`Z_2i\gamma p+Z_4m(\zeta )+\mathrm{\Sigma }^{}(p,\mathrm{\Lambda }),`$ (2.2.39)
with the regularised quark self energy
$`\mathrm{\Sigma }^{}(p,\mathrm{\Lambda })`$ $`:=`$ $`{\displaystyle _q^\mathrm{\Lambda }}𝒢((pq)^2)D_{\mu \nu }^{\mathrm{free}}(pq){\displaystyle \frac{\lambda ^a}{2}}\gamma _\mu S(q){\displaystyle \frac{\lambda ^a}{2}}\gamma _\nu .`$ (2.2.40)
When $`\widehat{m}0`$ the renormalisation condition, Eq. (2.2.6), is straightforward to implement. Writing
$$\mathrm{\Sigma }^{}(p,\mathrm{\Lambda }):=i\gamma p\left(A^{}(p^2,\mathrm{\Lambda }^2)1\right)+B^{}(p^2,\mathrm{\Lambda }^2),$$
(2.2.41)
which emphasises that these functions depend on the regularisation mass-scale, $`\mathrm{\Lambda }`$, Eq. (2.2.6) entails
$$Z_2(\zeta ^2,\mathrm{\Lambda }^2)=2A^{}(\zeta ^2,\mathrm{\Lambda }^2)\mathrm{and}m(\zeta )=Z_2(\zeta ^2,\mathrm{\Lambda }^2)m_{\mathrm{bm}}(\mathrm{\Lambda }^2)+B^{}(\zeta ^2,\mathrm{\Lambda }^2)$$
(2.2.42)
so that
$$A(p^2,\zeta ^2)=1+A^{}(p^2,\mathrm{\Lambda }^2)A^{}(\zeta ^2,\mathrm{\Lambda }^2),B(p^2,\zeta ^2)=m(\zeta )+B^{}(p^2,\mathrm{\Lambda }^2)B^{}(\zeta ^2,\mathrm{\Lambda }^2).$$
(2.2.43)
Multiplicative renormalisability requires that having fixed the solutions at a single renormalisation point, $`\zeta `$, their form at another point, $`\overline{\zeta }`$, is given by
$$S^1(p,\overline{\zeta })=i\gamma pA(p^2,\overline{\zeta }^2)+B(p^2,\overline{\zeta }^2)=\frac{Z_2(\overline{\zeta }^2,\mathrm{\Lambda }^2)}{Z_2(\zeta ^2,\mathrm{\Lambda }^2)}S^1(p,\zeta ).$$
(2.2.44)
This feature is evident in the solutions obtained in Ref. . It means that, in evolving the renormalisation point to $`\overline{\zeta }`$, the “1” in Eqs. (2.2.43) is replaced by $`Z_2(\overline{\zeta }^2,\mathrm{\Lambda }^2)/Z_2(\zeta ^2,\mathrm{\Lambda }^2)`$, and the “$`m(\zeta )`$” by $`m(\overline{\zeta })`$; i.e., the “seeds” in the integral equation evolve according to the QCD renormalisation group. This is why Eq. (2.2.34) is called a “renormalisation-group-improved rainbow truncation.”
Turning to the chiral limit, it follows from Eqs. (2.2.2), (2.2.3), (2.2.11) and (2.2.13) that for $`\widehat{m}=0`$
$$Z_2(\zeta ^2,\mathrm{\Lambda }^2)m_{\mathrm{bare}}(\mathrm{\Lambda }^2)=0,\mathrm{\Lambda }.$$
(2.2.45)
Hence, as remarked on page 2.2.13, there is no subtraction in the equation for $`B(p^2,\zeta ^2)`$; i.e., Eq. (2.2.43) becomes
$$B(p^2,\zeta ^2)=B^{}(p^2,\mathrm{\Lambda }^2),\underset{\mathrm{\Lambda }\mathrm{}}{lim}B^{}(p^2,\mathrm{\Lambda }^2)<\mathrm{},$$
(2.2.46)
which is only possible if the mass function is at least $`1/p^2`$-suppressed. This is not the case in quenched strong coupling QED, where the mass function behaves as $`\mathrm{cos}(\mathrm{const}.\mathrm{ln}[p^2/\zeta ^2])/(p^2/\zeta ^2)^{1/2}`$ , and that is the origin of the complications indicated on page 2.2 .
In Fig. 2.1 we present the renormalised dressed-quark mass function, $`M(p^2)`$, obtained by solving Eq. (2.2.39) using the model and parameter values of Ref. , Eqs. (2.2.34-2.2.37), and also in the chiral limit and with typical heavy-quark current-mass values. In the presence of explicit chiral symmetry breaking Eq. (2.2.11) describes the form of $`M(p^2)`$ for $`p^2>\mathrm{O}(1\mathrm{GeV}^2)`$. In the chiral limit, however, the ultraviolet behaviour is given by
$$M(p^2)\stackrel{\mathrm{large}p^2}{=}\frac{2\pi ^2\gamma _m}{3}\frac{\left(\overline{q}q^0\right)}{p^2\left(\frac{1}{2}\mathrm{ln}\left[p^2/\mathrm{\Lambda }_{\mathrm{QCD}}^2\right]\right)^{1\gamma _m}},$$
(2.2.47)
where $`\overline{q}q^0`$ is the renormalisation-point-independent vacuum quark condensate. This behaviour too is characteristic of the QCD renormalisation group and exhibits the power-law suppression anticipated on page 2.2.13. These results for the large $`p^2`$ behaviour of the mass function are model independent; i.e., they arise only because the DSE truncation is consistent with the QCD renormalisation group at one loop. (It has long been known that the truncation defined by Eq. (2.2.34) yields results in agreement with the QCD renormalisation group at one loop; e.g., Refs. .)
The gauge invariant expression for the renormalisation-point-dependent vacuum quark condensate was derived in Ref. :
$$\overline{q}q_\zeta ^0:=Z_4(\zeta ^2,\mathrm{\Lambda }^2)N_c\mathrm{tr}_\mathrm{D}_q^\mathrm{\Lambda }S^0(q,\zeta ),$$
(2.2.48)
where $`\mathrm{tr}_D`$ identifies a trace over Dirac indices only and the superscript “$`0`$” indicates the quantity was calculated in the chiral limit. Substituting Eq. (2.2.47) into Eq. (2.2.48), recalling that $`Z_4=Z_m`$ in Landau gauge and using Eq. (2.2.12), yields the one-loop expression
$$\overline{q}q_\zeta ^0=\left(\frac{1}{2}\mathrm{ln}\zeta ^2/\mathrm{\Lambda }_{\mathrm{QCD}}^2\right)^{\gamma _m}\overline{q}q^0.$$
(2.2.49)
Employing Eq. (2.2.11), this exemplifies the general result that
$$m(\zeta )\overline{q}q_\zeta ^0=\widehat{m}\overline{q}q^0;$$
(2.2.50)
i.e., that this product is renormalisation point invariant and, importantly, it shows that the behaviour expressed in Eq. (2.2.47) is exactly that required for consistency with the gauge invariant expression for the quark condensate. A model, such as Ref. , in which the scalar projection of the chiral limit dressed-quark propagator falls faster than $`1/p^4`$, up to $`\mathrm{ln}`$-corrections, is only consistent with this quark condensate vanishing, and it is this condensate that appears in the current algebra expression for the pion mass .
Equation (2.2.47) provides a reliable means of calculating the quark condensate because corrections are suppressed by powers of $`\mathrm{\Lambda }_{\mathrm{QCD}}^2/\zeta ^2`$. Analysing the asymptotic form of the numerical solution one finds
$$\overline{q}q^0=(0.227\mathrm{GeV})^3.$$
(2.2.51)
Using Eq. (2.2.49) one can define a one-loop evolved condensate
$$\overline{q}q_\zeta ^0|_{\zeta =1\mathrm{GeV}}:=\left(\mathrm{ln}\left[1/\mathrm{\Lambda }_{\mathrm{QCD}}\right]\right)^{\gamma _m}\overline{q}q^0=(0.241\mathrm{GeV})^3.$$
(2.2.52)
This can be directly compared with the value of the quark condensate employed in contemporary phenomenological studies : $`(0.236\pm 0.008\mathrm{GeV})^3`$. The authors of Ref. noted that increasing $`\omega 1.5\omega `$ in $`𝒢(k^2)`$ increases the calculated value in Eq. (2.2.52) by $`10`$%; i.e., the magnitude of the condensate is correlated with the degree of infrared enhancement/strength in the effective interaction. That is unsurprising because it has long been known that there is a critical coupling for DCSB; i.e., the kernel in the gap equation must have an integrated strength that exceeds some critical value . This is true in all fermion-based studies of DCSB.
The renormalisation-point-invariant current-quark masses corresponding to the $`m_f(\zeta )`$ in Fig. 2.1 are obtained in the following way: using Eq. (2.2.48), direct calculation from the chiral limit numerical solution gives
$$\overline{q}q_{\zeta =19\mathrm{GeV}}^0=(0.275\mathrm{GeV})^3,$$
(2.2.53)
and hence from the values of $`m_f^\zeta m_f(\zeta )`$ listed in Fig. 2.1 and Eqs. (2.2.50), (2.2.51), in MeV,
$$\begin{array}{cccc}\widehat{m}_{u,d}=6.60,& \widehat{m}_s=147,& \widehat{m}_c=\mathrm{1\hspace{0.17em}030},& \widehat{m}_b=\mathrm{6\hspace{0.17em}760},\end{array}$$
(2.2.54)
from which also follow one-loop evolved values in analogy with Eq. (2.2.52):
$$\begin{array}{cccc}m_{u,d}^{1\mathrm{GeV}}=5.5,& m_s^{1\mathrm{GeV}}=130,& m_c^{1\mathrm{GeV}}=860,& m_b^{1\mathrm{GeV}}=\mathrm{5\hspace{0.17em}700}.\end{array}$$
(2.2.55)
Figure 2.1 highlights a number of qualitative aspects of the quark mass function. One is the difference in the ultraviolet between the behaviour of $`M(p^2)`$ in the chiral limit and in the presence of explicit chiral symmetry breaking. In the infrared, however, the $`u,d`$-quark mass function and the chiral limit solution are almost indistinguishable. The chiral limit solution is nonzero only because of the nonperturbative DCSB mechanism whereas the $`u,d`$-quark mass function is purely perturbative at $`p^2>20`$GeV<sup>2</sup>. Hence the evolution to coincidence of the chiral-limit and $`u,d`$-quark mass functions makes clear the transition from the perturbative to the nonperturbative domain. It is on this nonperturbative domain that $`A(p^2,\zeta ^2)`$ differs significantly from one. \[This behaviour and that of the light-quark mass function depicted in Fig. 2.1 have recently been confirmed in numerical simulations of lattice-QCD .\] A concomitant observation is that the DCSB mechanism has a significant effect on the propagation characteristics of $`u,d,s`$-quarks. However, as evident in the figure, that is not the case for the $`b`$-quark. Its large current-quark mass almost entirely suppresses momentum-dependent dressing, so that $`M_b(p^2)`$ is nearly constant on a substantial domain. This is true to a lesser extent for the $`c`$-quark.
To quantify the effect of the DCSB mechanism on massive quark propagation characteristics the authors of Ref. introduced a single measure: $`_f:=M_f^E/m_f^\zeta `$, where $`M_f^E`$ is the Euclidean constituent-quark mass, defined as the solution of
$$(M_f^E)^2M^2(p^2=(M_f^E)^2,\zeta ^2)=0.$$
(2.2.56)
In this exemplifying model the Euclidean constituent-quark mass takes the values listed in the figure, which have magnitudes and ratios consistent with contemporary phenomenology; e.g., Refs. , and
$$\begin{array}{cccccc}f& u,d& s& c& b& \\ & & & & & \\ _f& 150& 10& 2.2& 1.2& \end{array}.$$
(2.2.57)
These values are representative and definitive: for light-quarks $`_{q=u,d,s}10`$-$`100`$, while for heavy-quarks $`_{Q=c,b}1`$, and highlight the existence of a mass-scale characteristic of DCSB: $`M_\chi `$. The propagation characteristics of a flavour with $`m_f^\zeta M_\chi `$ are significantly altered by the DCSB mechanism, while for flavours with $`m_f^\zeta M_\chi `$ momentum-dependent dressing is almost irrelevant. It is apparent and unsurprising that $`M_\chi 0.2`$GeV$`\mathrm{\Lambda }_{\mathrm{QCD}}`$. This feature of the dressed-quark mass function provides the foundation for a constituent-quark-like approximation in the treatment of heavy-meson decays and transition form factors .
To recapitulate. The quark DSE describes the phenomena of DCSB and the concomitant dynamical generation of a momentum dependent quark mass function, with the renormalisation-group-improved rainbow-truncation yielding model-independent results for the momentum dependence of the mass function in the ultraviolet. For light quarks, defined by $`_f1`$, the magnitude of the mass function in the infrared; i.e., for $`p^2<1`$-$`2`$GeV<sup>2</sup>, is determined by the behaviour of the effective quark-quark interaction on the same domain, and so is the value of the vacuum quark condensate. An infrared enhancement in this effective interaction is required to describe observable light-meson phenomena. We have only provided a single illustration but it is supported by, e.g., Refs. and the observation that with insufficient infrared integrated strength in the kernel of Eq. (2.2.1) the quark condensate vanishes , which is a poor starting point for light-hadron phenomenology. The question now arises, where does this strength come from?
Gluon DSE Guidance here comes from studies of the DSE satisfied by the dressed-gluon propagator, which is depicted in Fig. 2.2. However, as we now describe, these studies are inconclusive.
Early analyses used the ghost-free axial gauge: $`nA^a=0`$, $`n^2>1`$, in which case the second equation in the figure is absent and two independent scalar functions: $`F_1`$, $`F_2`$, are required to fully specify the dressed-gluon propagator, cf. the covariant gauge expression in Eq. (2.2.15), which requires only one function. In the absence of interactions: $`F_1(k^2)=1/k^2`$, $`F_2(k^2)0`$. These studies employed an Ansatz for the three-gluon vertex that doesn’t possess a particle-like singularity and neglected the coupling to the quark DSE. They also assumed $`F_20`$, even nonperturbatively, and ignored it in solving the DSE. The analysis then yielded
$$F_1(k^2)\stackrel{k^20}{}\frac{1}{k^4};$$
(2.2.58)
i.e., a marked infrared enhancement that can yield an area law for the Wilson loop and hence confinement, and DCSB as described above without fine-tuning. This effect is driven by the gluon vacuum polarisation, diagram three in the first line of Fig. (2.2). A similar result was obtained in Ref. . However, a possible flaw in these analyses was identified in Ref. , which argued from properties of the spectral density in ghost-free gauges that $`F_2`$ cannot be zero but acts to cancel the enhancement in $`F_1`$. \[Preserving $`F_2`$ yields a coupled system of equations for the gluon propagator that is at least as complicated as that obtained in covariant gauges, which perhaps outweighs the apparent benefit of eliminating ghost fields in the first place.\]
There have also been analyses of the gluon DSE using Landau gauge and those of Refs. are unanimous in arriving at the covariant gauge analogue of Eq. (2.2.58), again driven by the gluon vacuum polarisation diagram. In these studies Ansätze were used for the dressed-three-gluon vertex, all of which were free of particle-like singularities. However, these studies too have weaknesses: based on an anticipated dominance of the gluon-vacuum polarisation, truncations were implemented so that only the third and fifth diagrams on the r.h.s. of the first equation in Fig. (2.2) were retained. In covariant gauges there is a priori no reason to neglect the ghost loop contribution, diagram six, although perturbatively its contribution is estimated to be small .
Another class of Landau gauge studies are described in Refs. , which propose solving the DSEs via rational polynomial Ansätze for the one-particle irreducible components of the Schwinger functions appearing in Fig. 2.2; i.e., the self energies and vertices. This method attempts to preserve aspects of the organising principle of perturbation theory in truncating the DSEs. In concrete calculations, for simplicity, only the first, third and sixth diagrams on the r.h.s. of the first equation in Fig. 2.2 survive, the last \[fermion\] equation is neglected and the leading order solution of the ghost equation has the appearance of the massless free propagator: $`1/k^2`$. The analysis suggests that a consistent solution for the dressed-gluon propagator is one that vanishes at $`k^2=0`$; i.e., in Eq. (2.2.15)
$$d(k^2)\frac{k^4}{k^4+\gamma ^4}.$$
(2.2.59)
However, the associated polynomial Ansatz for the dressed-three-gluon vertex exhibits particle-like singularities and this is characteristic of the method. The question of how this can be made consistent with the absence of coloured bound states in the strong interaction spectrum is currently unanswered.
Proponents of the result in Eq. (2.2.59) claim support from studies of “complete” gauge fixing; i.e., in the outcome of attempts to construct a Faddeev-Popov-like determinant that eliminates Gribov copies or ensures that the functional integration domain for the gauge field is restricted to a subspace without them. Fixing a so-called “minimal Landau gauge,” which enforces a constraint of integrating only over gauge field configurations inside the Gribov horizon; i.e., on the simplest domain for which the Faddeév-Popov operator is invertible, the dressed-gluon $`2`$-point function is shown to vanish at $`k^2=0`$. However, the approach advocated in Refs. makes no use of the additional ghost-like fields necessary to restrict the integration domain.
Thus far in this discussion of the gluon DSE we have reported nothing qualitatively new and a more detailed review of the studies described can be found in Ref. , Sec. 5.1. What about contemporary studies?
The direct approach to solving the Landau gauge gluon DSE, pioneered in Refs. , has been revived by two groups: $`𝒜`$, Refs. ; and $``$, Refs. , with the significant new feature that nonperturbative effects in the ghost sector are admitted; i.e., a nonperturbative solution of the DSE for the ghost propagator is sought in the form
$$G^{ab}(k)=\delta ^{ab}\frac{\varpi (k^2)}{k^2}[\mathrm{without}\mathrm{interactions},\varpi (k^2)1].$$
(2.2.60)
These studies analyse a truncated gluon-ghost DSE system, retaining only the third and sixth loop diagrams in the first equation of Fig. (2.2) and also the second equation. Superficially this is the same complex of equations as studied in Refs. . However, the procedure for solving it is different, arguably less systematic but also less restrictive. The difference between the groups is that $`𝒜`$ employ Ansätze for the dressed-ghost-gluon and dressed-three-gluon vertices constructed so as to satisfy the relevant Slavnov-Taylor identities while $``$ simply use the bare, perturbative vertices. Nevertheless they agree in the conclusion that in this case the infrared behaviour of the gluon DSE’s solution is determined by the ghost loop alone: it overwhelms the gluon vacuum polarisation contribution. That is emphasised in Ref. , which eliminates every loop diagram in truncating the first equation of Fig. 2.2 except the ghost loop and still recovers the behaviour of Ref. . That behaviour is
$$\varpi (k^2)\frac{1}{(k^2)^\kappa },d(k^2)(k^2)^{2\kappa }\mathrm{for}k^2<\mathrm{\Lambda }_{\mathrm{QCD}}^2,\mathrm{with}0.8<\kappa 1.$$
(2.2.61)
Exact evaluation of the angular integrals that arise when solving the integral equations gives the integer valued upper bound, $`\kappa =1`$ . This corresponds to a dressed-gluon $`2`$-point function that vanishes at $`k^2=0`$, although the suppression is very sudden with the propagator not peaking until $`k^2\mathrm{\Lambda }_{\mathrm{QCD}}^2`$, where
$$(d(k^2)/k^2)|_{k^2=\mathrm{\Lambda }_{\mathrm{QCD}}^2}100/\mathrm{\Lambda }_{\mathrm{QCD}}^2;$$
(2.2.62)
i.e., it is very much enhanced over the free propagator. \[See, e.g., Ref. , Fig. 12.\] $`\kappa =1`$ also yields a dressed-ghost propagator that exhibits a dipole enhancement analogous to that of Eq. (2.2.58). A \[renormalisation group invariant\] strong running coupling consistent with this truncations is:
$$\alpha (k^2):=\frac{1}{4\pi }g^2\varpi ^2(k^2)d(k^2)$$
(2.2.63)
and its value at $`k^2=0`$ is fixed by the numerical solutions:
$$\begin{array}{ccc}& 𝒜& \\ & & \\ \alpha (k^2=0)\hfill & 9.5& 4\text{or}\mathrm{\hspace{0.33em}12}\end{array}.$$
(2.2.64)
\[NB. Group $`𝒜`$ approximates the angular integrals and uses vertex Ansätze. Group $``$ uses bare vertices and in Ref. approximates the angular integrals to obtain $`\alpha (0)12`$, while in Ref. the integrals are evaluated exactly, which yields $`\alpha (0)=\frac{4}{3}\pi 4.2`$.\]
The qualitative common feature is that the Grassmannian ghost loops act to suppress the dressed-gluon propagator in the infrared. That may also be said of Refs. . \[Indications that the quark loop, diagram seven in Fig. 2.2, acts to oppose an enhancement of the type in Eq. (2.2.58) may here, with hindsight, be viewed as suggestive.\] One aspect of ghost fields is that they enter because of gauge fixing via the Faddeev-Popov determinant. Hence, while none of the groups introduce the additional Faddeev-Popov contributions advocated in Refs. , they nevertheless do admit ghost contributions, and in their solution the number of ghost fields does not have a qualitative impact. Reference also obtains a dressed-propagator for the Faddeev-Popov fields with a $`k^2=0`$ dipole singularity. It contributes to the action via the term employed to restrict the gauge field integration domain, in which capacity the dipole singularity can plausibly drive an area law for Wilson loops.
Schwinger functions are the primary object of study in numerical simulations of lattice-QCD and Refs. report contemporary estimates of the lattice Landau gauge dressed-gluon $`2`$-point function. They are consistent with a finite although not necessarily vanishing value of $`d(k^2=0)`$. However, simulations of the dressed-ghost $`2`$-point function find no evidence of a dipole singularity, with the ghost propagator behaving as if $`\varpi (k^2)=1`$ in the smallest momentum bins . \[NB. Since the quantitative results from groups $`𝒜`$ and $``$ differ and exhibit marked sensitivity to details of the numerical analysis, any agreement between the DSE results for $`\varpi (k^2)`$ or $`d(k^2)`$ and the lattice data on some subdomain can be regarded as fortuitous.\]
The behaviour in Eqs. (2.2.61) also entails the presence of particle-like singularities in extant Ansätze for the dressed-ghost-gluon, dressed-three-gluon and dressed-quark-gluon vertices that are consistent with the relevant Slavnov-Taylor identities. \[$`\kappa =1`$ corresponds to an ideal simple pole singularity.\] Hence while this behaviour may be consistent with the confinement of elementary excitations, as currently elucidated it also predicts the existence of coloured bound states in the strong interaction spectrum. Furthermore, while it does yield a strong running coupling with $`\alpha (k^2=0)>1`$, that makes DCSB dependent on fine tuning , and quantitative calculations based on the present numerical solutions give a quark condensate only $`5`$% of the value in Eq. (2.2.52. Notwithstanding these remarks, the studies of Refs. and subsequently Refs. are laudable. They have focused attention on a previously unsuspected qualitative sensitivity to truncations in the gauge sector.
To recapitulate. It is clear from Refs. that DCSB requires the effective interaction in the quark DSE to be strongly enhanced at $`k^2\mathrm{\Lambda }_{\mathrm{QCD}}^2`$. \[Remember too that modern lattice simulations confirm the pattern of behaviour exhibited by quark DSE solutions obtained with such an enhanced interaction.\] Studies of QCD’s gauge sector indicate that gluon-gluon and/or gluon-ghost dynamics can generate such an enhancement. However, the qualitative nature of the mechanism and its strength remains unclear: is it the gluon vacuum polarisation or that of the ghost that is the driving force? It is a contemporary challenge to explore and understand this.
Finally, in discussing aspects of the gauge sector one might consider whether instanton configurations play a role? Instantons are solutions of the classical equation of motion for the Euclidean gauge field. As such they form a set of measure zero in the gauge field integration space. Nonetheless, they can form the basis for a semi-classical approximation to the gauge field action and models based on this notion have been phenomenologically successful . In this context we note that the DSEs in Fig. (2.2) are derived nonperturbatively and their self-consistent solution includes the effects of all field configurations. Hence instanton-like configurations may contribute to the form of the solution. However, a successful description of observable phenomena does not require that their contribution be quantified in a particular truncation. Nevertheless, one might estimate that a dilute liquid of instantons, each with radius $`\overline{\rho }1/(0.6`$GeV), could significantly effect the propagation characteristics of gluons only on the domain of intermediate momenta: $`k^2(0.6\mathrm{GeV})^2`$. Therefore they cannot qualitatively affect the infrared aspects discussed in this subsection. They also make no contribution in the perturbative domain: $`k^2>1`$$`2`$GeV<sup>2</sup>, where the perturbative matching inherent in the DSEs is a strength that makes possible a unification of infrared and ultraviolet phenomena, such as in the behaviour of bound state elastic and transition form factors; e.g., Refs. .
Confinement Confinement is the failure to directly observe coloured excitations in a detector: neither quarks nor gluons nor coloured composites. The contemporary hypothesis is stronger; i.e., coloured excitations cannot propagate to a detector. To ensure this it is sufficient that coloured $`n`$-point functions violate the axiom of reflection positivity , which is guaranteed if the Fourier transform of the momentum-space $`n`$-point Schwinger function is not a positive-definite function of its arguments. Reflection positivity is one of a set of five axioms that must be satisfied if the given $`n`$-point function is to have a continuation to Minkowski space and hence an association with a physical, observable state. If an Hamiltonian exists for the theory but a given $`n`$-point function violates reflection positivity then the space of observable states, which is spanned by the eigenstates of the Hamiltonian, does not contain anything corresponding to the excitation(s) described by that Schwinger function. \[The violation of reflection positivity is not a necessary condition for confinement . A text-book counterexample is massless two-dimensional QED but in this case confinement of electric charge and DCSB both arise as a peculiar consequence of the number of dimensions.\]
The free boson propagator does not violate reflection positivity:
$$\mathrm{\Delta }(x):=\frac{d^4k}{(2\pi )^4}\mathrm{e}^{ikx}\frac{1}{k^2+m^2}=\frac{1}{4\pi ^2x}_0^{\mathrm{}}𝑑\mathrm{}J_1(\mathrm{}x)\frac{\mathrm{}^2}{\mathrm{}^2+m^2}=\frac{m}{4\pi ^2x}K_1(mx).$$
(2.2.65)
Here $`x:=(xx)^{1/2}>0`$, $`J_1`$ is an oscillatory Bessel function of the first kind and $`K_1`$ is the monotonically decreasing, strictly convex-up, non-negative modified Bessel function of the second kind. The same is true of the free fermion propagator:
$$S(x)=\frac{d^4k}{(2\pi )^4}\mathrm{e}^{ikx}\frac{mi\gamma k}{k^2+m^2}=(m\gamma )\mathrm{\Delta }(x)=\frac{m^2}{4\pi ^2x}\left[K_1(mx)+\frac{\gamma x}{x}K_2(mx)\right],$$
(2.2.66)
which is also positive definite. The spatially averaged Schwinger function is a particularly insightful tool . Consider the fermion and let $`T=x_4`$ represent Euclidean “time,” then
$$\sigma _S(T):=d^3x\mathrm{tr}_DS(\stackrel{}{x},T)=\frac{1}{\pi }_0^{\mathrm{}}𝑑\mathrm{}\frac{m}{\mathrm{}^2+m^2}\mathrm{cos}(\mathrm{}T)=\frac{1}{2}\mathrm{e}^{mT}.$$
(2.2.67)
Hence the free fermion’s mass can be easily obtained from the large $`T`$ behaviour of the spatial average:
$$mT=\underset{T\mathrm{}}{lim}\mathrm{ln}\sigma _S(T).$$
(2.2.68)
\[The boson analogy is obvious.\] This is just the approach used to determine bound state masses in simulations of lattice-QCD.
For contrast, consider the dressed-gluon $`2`$-point function in Eq. (2.2.59):
$$D(x):=\frac{d^4k}{(2\pi )^4}\mathrm{e}^{ikx}\frac{k^2}{k^4+\gamma ^4}=\frac{1}{4\pi ^2x}_0^{\mathrm{}}𝑑\mathrm{}J_1(\mathrm{}x)\frac{\mathrm{}^4}{\mathrm{}^4+\gamma ^4}=\frac{\gamma }{4\pi ^2x}\left(\frac{d}{dz}\mathrm{ker}(z)\right)|_{z=\gamma x},$$
(2.2.69)
where ker$`(z)`$ is the oscillatory Thomson function. $`D(x)`$ is not positive definite and hence a dressed-gluon $`2`$-point function that vanishes at $`k^2=0`$ violates the axiom of reflection positivity and is therefore not observable; i.e., the excitation it describes is confined. At asymptotically large Euclidean distances
$`D(x)`$ $`\stackrel{x\mathrm{}}{}`$ $`{\displaystyle \frac{\gamma ^{1/2}}{x^{3/2}}}\mathrm{e}^{\gamma x/\sqrt{2}}\left[\mathrm{cos}(\frac{1}{\sqrt{2}}\gamma x+\frac{\pi }{8})+\mathrm{sin}(\frac{1}{\sqrt{2}}\gamma x+\frac{\pi }{8})\right].`$ (2.2.70)
Comparing this with Eq. (2.2.65) one identifies a mass as the coefficient in the exponential: $`m_D=\gamma /\sqrt{2}`$. \[NB. At large $`x`$, $`K_1(x)\mathrm{exp}(x)/\sqrt{x}`$.\] By an obvious analogy, the coefficient in the oscillatory term is the lifetime : $`\tau =1/m_D`$. Both the mass and lifetime are tied to the dynamically generated mass-scale $`\gamma `$, which, using
$$\frac{z}{z^2+\gamma ^4}=\frac{1}{2}\frac{1}{z+i\gamma ^2}+\frac{1}{2}\frac{1}{zi\gamma ^2},$$
(2.2.71)
is just the displacement of the complex conjugate poles from the real-$`k^2`$ axis. It is a general result that the Fourier transform of a real function with complex conjugate poles is not positive definite. Hence the existence of such poles in a $`n`$-point Schwinger function is a sufficient condition for the violation of reflection positivity and thus for confinement. The spatially averaged Schwinger function is also useful here.
$$D(T):=d^3xD(\stackrel{}{x},T)=\frac{1}{\pi }_0^{\mathrm{}}𝑑\mathrm{}\frac{\mathrm{}^2}{\mathrm{}^4+\gamma ^4}\mathrm{cos}(\mathrm{}T)=\frac{1}{2\gamma }\mathrm{e}^{\frac{1}{\sqrt{2}}\gamma T}\mathrm{cos}(\frac{1}{\sqrt{2}}\gamma T+\frac{\pi }{4}),$$
(2.2.72)
and, generalising Eq. (2.2.68), one can define a $`T`$-dependent mass:
$$m(T)T:=\mathrm{ln}D(T)=\mathrm{ln}(2\gamma )+\frac{1}{\sqrt{2}}\gamma T\mathrm{ln}\left[\mathrm{cos}(\frac{1}{\sqrt{2}}\gamma T+\frac{\pi }{4})\right].$$
(2.2.73)
It exhibits periodic singularities whose frequency is proportional to the dynamical mass-scale that is responsible for the violation of reflection positivity. If a dressed-fermion $`2`$-point function has complex conjugate poles it too will be characterised by a $`T`$-dependent mass that exhibits such behaviour.
This reflection positivity criterion has been employed to very good effect in three dimensional QED . First, some background. QED<sub>3</sub> is confining in the quenched truncation . That is evident in the classical potential
$$V(r):=_{\mathrm{}}^{\mathrm{}}𝑑x_3\frac{d^3k}{(2\pi )^3}\mathrm{e}^{i\stackrel{}{k}\stackrel{}{x}+ik_3x_3}\frac{e^2}{k^2}=\frac{e^2}{2\pi }\mathrm{ln}(e^2r),r^2=x_1^2+x_2^2,$$
(2.2.74)
which describes the interaction between two static sources. \[NB. $`e^2`$ has the dimensions of mass in QED<sub>3</sub>.\] It is a logarithmically growing potential, showing that the energy required to separate two charges is infinite. Furthermore, $`V(r)`$ is just a one-dimensional average of the spatial gauge-boson $`2`$-point Schwinger function and it is not positive definite, which indicates that the photon is also confined.
If now, however, the photon vacuum polarisation tensor is evaluated at order $`e^2`$ using $`N_f`$ massless fermions then, using the notation of Eq. (2.2.15), the photon propagator is characterised by
$$\frac{d(k^2)}{k^2}=\frac{1}{k^2+\stackrel{~}{\alpha }k},\text{from }\mathrm{\Pi }(k^2)=\frac{\stackrel{~}{\alpha }}{k},\stackrel{~}{\alpha }=N_fe^2/8,$$
(2.2.75)
and one finds
$$V(r)=\frac{e^2}{4}\left[\text{H}_0(\stackrel{~}{\alpha }r)N_0(\stackrel{~}{\alpha }r)\right],$$
(2.2.76)
where H$`{}_{0}{}^{}(x)`$ is a Struve function and $`N_0(x)`$ a Neumann function, both of which are related to Bessel functions. In this case $`V(r)`$ is positive definite, with the limiting cases
$$V(r)\stackrel{r0}{}\mathrm{ln}(\stackrel{~}{\alpha }r),V(r)\stackrel{r\mathrm{}}{=}\frac{e^2}{2\pi }\frac{1}{\stackrel{~}{\alpha }r},$$
(2.2.77)
and confinement is lost in QED<sub>3</sub>. That is easy to understand: pairs of massless fermions cost no energy to produce and can propagate to infinity so they are very effective at screening the interaction.
With $`d(k^2)=1/[1+\mathrm{\Pi }(k^2)]`$ and sensible, physical constraints on the form of $`\mathrm{\Pi }(k^2)`$, such as boundedness and vanishing in the ultraviolet, one can show that
$$V(r)\stackrel{r\mathrm{}}{=}\frac{e^2}{2\pi }\frac{1}{1+\mathrm{\Pi }(0)}\mathrm{ln}(e^2r)+\mathrm{const}.+h(r),$$
(2.2.78)
where $`h(r)`$ falls-off at least as quickly as $`1/r`$. Hence, the existence of a confining potential in QED<sub>3</sub> just depends on the value of the vacuum polarisation at the origin. In the quenched truncation, $`\mathrm{\Pi }(0)=0`$ and the theory is logarithmically confining. With massless fermions, $`1/[1+\mathrm{\Pi }(0)]=0`$ and confinement is absent. Finally, when the vacuum polarisation is evaluated from a loop of massive fermions, whether that mass is obtained dynamically via the gap equation or simply introduced as an external parameter, one obtains $`\mathrm{\Pi }(0)<\mathrm{}`$ and hence a confining theory.
In Ref. the QED<sub>3</sub> gap equation is solved for all four cases and the fermion propagator analysed. The results are summarised by Fig. 2.3. In the quenched theory, Eq. (2.2.74), the dressed-fermion $`2`$-point function exhibits exactly those periodic singularities that, via Eq. (2.2.73), are indicative of complex conjugate poles. Hence this feature of the $`2`$-point function, tied to the violation of reflection positivity, is a clear signal of confinement in the theory. That is emphasised further by a comparison with the theory that is unquenched via massless fermions in the vacuum polarisation, Eq. (2.2.75). As we have described, that theory is not confining and in this case $`\sigma _S(T)`$ has the noninteracting, unconfined free particle form in Eq. (2.2.67). The difference could not be more stark. The remaining two cases exhibit the periodic singularities that signal confinement, just as they should based on Eq. (2.2.78).
At this point we note that any concern that the presence of complex conjugate singularities in coloured $`n`$-point functions leads to a violation of causality is misguided. Microscopic causality only constrains the commutativity of operators, and products thereof, that represent elements in the space of observable particle states; i.e., the space spanned by eigenstates of the Hamiltonian. Since Schwinger functions that violate reflection positivity do not have a continuation into that space there can be no question of violating causality. It is only required that $`𝒮`$-matrix elements that describe colour-singlet to colour-singlet transitions should satisfy the axioms, including reflection positivity.
The violation of reflection positivity by coloured $`n`$-point functions is a sufficient condition for confinement. However, it is not necessary, as the example of planar, two-dimensional QCD shows . There the fermion two-point function exhibits particle-like singularities but the colour singlet meson bound state amplitudes, obtained from a Bethe-Salpeter equation, vanish at momenta coincident with the constituent-fermion mass shell. This excludes the pinch singularities that would otherwise lead to bound state break-up and liberation of the constituents. It is a realisation of confinement via a failure of the cluster decomposition property \[CDP\] . The CDP is a requirement that the difference between the vacuum expectation value of a product of fields and all products of vacuum expectation values of subsets of these fields must vanish faster than any power. \[This is modified slightly in theories, like QED, with a massless, asymptotic state: the photon.\] It can be understood as a statement about charge screening and its failure means that, irrespective of the separation between sources, the interaction between them is never negligible. That is an appealing, intuitive representation of confinement. Failure of the CDP is an implicit basis for confinement in the bulk of QCD potential models with; e.g., Refs. providing contemporary illustrations.
Confinement is a more contentious issue than DCSB and its origin and realisation less-well understood. However, in this subsection we have described a perspective that is common to many authors, and an interested reader will find variations and more expansive discussions of various points in; e.g., Refs. , and complementary perspectives in; e.g., Refs. . It is, of course, because confinement is poorly understood that its study and modelling are important. To presume otherwise is a misapprehension.
### 2.3 Phenomenological Applications
When Ref. was written the only application of DSEs to observable phenomena consisted in the oft-repeated calculation of well-known quantities, such as the pion mass and decay constant. References represent a departure from that, and are progenitors of the wide-ranging application of DSEs to observables newly accessible at the current generation of experimental facilities. Many of these applications are reviewed in Refs. and herein we only describe three recent, significant developments.
Light Mesons The model used to illustrate renormalisation and DCSB in Sec. 2.2 has been applied to the calculation of vector meson masses and decay constants , and to elucidate the role of vector mesons in connection with the electromagnetic pion form factor . These studies are important because, in concert with Ref. , they complete a DSE description of those light-mesons in the strong interaction spectrum that are most often produced in reactions involving hadrons. In so doing they illustrate the efficacy of the renormalisation-group-improved rainbow-ladder truncation for flavour non-singlet pseudoscalar and vector mesons composed of light-quarks ($`u`$, $`d`$, $`s`$), and thereby that of the systematic, Ward-Takahashi identity preserving truncation scheme introduced in Ref. .
The renormalised homogeneous Bethe-Salpeter equation for a bound state of a dressed-quark and dressed-antiquark with total momentum $`P`$ is
$$[\mathrm{\Gamma }_H(k;P)]_{tu}=_q^\mathrm{\Lambda }[\chi (q;P)]_{rs}K_{tu}^{rs}(q,k;P),\chi (q;P):=𝒮(q_+)\mathrm{\Gamma }_H(q;P)𝒮(q_+),$$
(2.3.1)
with: $`\mathrm{\Gamma }_H(k;P)`$ the Bethe-Salpeter amplitude \[BSA\], where $`H`$ specifies the flavour structure of the meson; $`𝒮(p):=\mathrm{diag}[S_u(p),S_d(p),S_s(p)]`$; $`q_+=q+\eta _PP`$, $`q_{}=q(1\eta _P)P`$; and $`r,\mathrm{},u`$ represent colour-, Dirac- and flavour-matrix indices. \[$`\eta _P[0,1]`$ is the momentum partitioning parameter. It appears in Poincaré covariant treatments because, in general, the definition of the relative momentum is arbitrary. Physical observables, such as the mass, must be independent of $`\eta _P`$ but that is only possible if the Bethe-Salpeter amplitude depends on it. $`\eta _P=1/2`$ for charge-conjugation eigenstates.\]
In Eq. (2.3.1), $`K_{tu}^{rs}(q,k;P)`$ is the renormalised, fully-amputated quark-antiquark scattering kernel, which, as we have seen, also appears implicitly in Eq. (2.2.1) because it is the kernel of the inhomogeneous DSE satisfied by $`\mathrm{\Gamma }_\nu (q;p)`$. $`K_{tu}^{rs}(q,k;P)`$ is a $`4`$-point Schwinger function, obtained as the sum of a countable infinity of skeleton diagrams. It is two-particle-irreducible, with respect to the quark-antiquark pair of lines and does not contain quark-antiquark to single gauge-boson annihilation diagrams, such as would describe the leptonic decay of a pseudoscalar meson. \[A connection between the fully-amputated quark-antiquark scattering amplitude: $`M=K+K(𝒮𝒮)K+\mathrm{}`$, and the Wilson loop is discussed in Ref. .\] The complexity of $`K_{tu}^{rs}(q,k;P)`$ is one reason why quantitative studies of the quark DSE currently employ Ansätze for $`D_{\mu \nu }(k)`$ and $`\mathrm{\Gamma }_\nu (k,p)`$. However, as illustrated by Ref. , the complexity of $`K_{tu}^{rs}(q,k;P)`$ does not prevent one from analysing aspects of QCD in a model independent manner and proving general results that provide useful constraints on model studies.
Equation (2.3.1) is an eigenvalue problem and solutions exist only for particular, separated values of $`P^2`$. The eigenvector associated with each eigenvalue: $`\mathrm{\Gamma }_H(k;P)`$, the BSA, is a one-particle-irreducible, fully-amputated quark-meson vertex. In the flavour non-singlet pseudoscalar channels the solutions having the lowest eigenvalues correspond to the $`\pi `$\- and $`K`$-mesons, while in the vector channels they correspond to the $`\omega `$-, $`\rho `$\- and $`\varphi `$ mesons.
Following Ref. , the renormalised inhomogeneous BSE for the axial-vector vertex, consistent with the renormalisation–group-improved quark DSE, Eqs. (2.2.39) and (2.2.40), is
$$\mathrm{\Gamma }_{5\mu }^l(k;P)=Z_2\frac{1}{2}\lambda ^l\gamma _5\gamma _\mu _q^\mathrm{\Lambda }𝒢((kq)^2)D_{\alpha \beta }^{\mathrm{free}}(kq)\frac{\lambda ^a}{2}\gamma _\alpha 𝒮(q_+)\mathrm{\Gamma }_{5\mu }^l(q;P)𝒮(q_{})\frac{\lambda ^a}{2}\gamma _\beta ,$$
(2.3.2)
where $`\{\frac{1}{2}\lambda _F^l:l=1,\mathrm{},8\}`$ are the generators of $`SU(3)_{\mathrm{flavour}}`$. It is straightforward to verify that the axial-vector Ward-Takahashi identity is satisfied; i.e.,
$$P_\mu \mathrm{\Gamma }_{5\mu }^l(k;P)=𝒮^1(k_+)\frac{1}{2}\lambda ^li\gamma _5+\frac{1}{2}\lambda ^li\gamma _5𝒮^1(k_{})M_\zeta i\mathrm{\Gamma }_5^l(k;P)i\mathrm{\Gamma }_5^l(k;P)M_\zeta ,$$
(2.3.3)
where $`M_\zeta =\mathrm{diag}[m_u(\zeta ),m_d(\zeta ),m_s(\zeta )]`$ and the renormalised pseudoscalar vertex satisfies its own inhomogeneous BSE:
$$\mathrm{\Gamma }_5^l(k;P)=Z_4\frac{1}{2}\lambda ^l\gamma _5_q^\mathrm{\Lambda }𝒢((kq)^2)D_{\mu \nu }^{\mathrm{free}}(kq)\frac{\lambda ^a}{2}\gamma _\mu 𝒮(q_+)\mathrm{\Gamma }_5^l(q;P)𝒮(q_{})\frac{\lambda ^a}{2}\gamma _\nu .$$
(2.3.4)
\[NB. The product $`M_\zeta \mathrm{\Gamma }_5^l(k;P)`$ is renormalisation point independent.\]
The pseudoscalar mesons appear as poles in both the axial-vector and pseudoscalar vertices and equating pole residues yields the homogeneous BSE
$$\mathrm{\Gamma }_H(k;P)+_q^\mathrm{\Lambda }𝒢((kq)^2)D_{\mu \nu }^{\mathrm{free}}(kq)\frac{\lambda ^a}{2}\gamma _\mu 𝒮(q_+)\mathrm{\Gamma }_H(q;P)𝒮(q_{})\frac{\lambda ^a}{2}\gamma _\nu =0.$$
(2.3.5)
As is characteristic of homogeneous equations, the normalisation of the solution is not fixed by this equation. The canonical normalisation enforces a requirement that the bound state contribution to the fully-amputated quark-antiquark scattering amplitude: $`M`$, have unit residue. In this rainbow-ladder truncation that condition is expressed via
$$2P_\mu =_q^\mathrm{\Lambda }\mathrm{tr}\left[\overline{\mathrm{\Gamma }}_H(q;P)\frac{𝒮(q_+)}{P_\mu }\mathrm{\Gamma }_H(q;P)𝒮(q_{})+\overline{\mathrm{\Gamma }}_H(q;P)𝒮(q_+)\mathrm{\Gamma }_H(q;P)\frac{𝒮(q_{})}{P_\mu }\right],$$
(2.3.6)
where
$$\overline{\mathrm{\Gamma }}_H(k,P)^\mathrm{t}:=C^1\mathrm{\Gamma }_H(k,P)C,$$
(2.3.7)
with $`C=\gamma _2\gamma _4`$ the charge conjugation matrix:
$$C\gamma _\mu ^\mathrm{t}C^{}=\gamma _\mu ;[C,\gamma _5]=0,$$
(2.3.8)
and $`X^\mathrm{t}`$ denotes the matrix transpose of $`X`$. The general form of a pseudoscalar BSA is
$$\mathrm{\Gamma }_H(k;P)=𝒯^H\gamma _5\left[iE_H(k;P)+\gamma PF_H(k;P)\text{}+\gamma kkPG_H(k;P)+\sigma _{\mu \nu }k_\mu P_\nu H_H(k;P)\right],$$
(2.3.9)
where $`𝒯^H`$ is a matrix that describes the flavour content of the meson; e.g., $`𝒯^{\pi ^+}=\frac{1}{2}(\lambda _F^1+i\lambda _F^2)`$ and, for bound states of constituents with equal current-quark masses, the scalar functions $`E`$, $`F`$, $`G`$ and $`H`$ are even under $`kPkP`$. \[NB. Since the homogeneous BSE is an eigenvalue problem, $`E_H(k;P)=E_H(k^2,kP|P^2)`$; i.e., $`P^2`$ is not a variable, instead it labels the solution. The same is true of each function.\]
Equation (2.3.5) also describes vector mesons, as can be shown by considering the inhomogeneous equation for the renormalised vector vertex. In general, twelve independent scalar functions are required to express the Dirac structure of a vector vertex. However, a vector meson bound state is transverse:
$$P_\mu \mathrm{\Gamma }_\mu ^H(k^2,kP|P^2=m_H^2)=0,$$
(2.3.10)
where $`m_H`$ is the bound state’s mass, and this constraint reduces to eight the number of independent scalar functions. One therefore has
$$\mathrm{\Gamma }_\mu ^H(k;P)=𝒯^H\underset{l=1}{\overset{8}{}}O_\mu ^lF_l(k;P),$$
(2.3.11)
with eight orthonormalised matrix covariants
$$\begin{array}{cccccc}\hfill O_\mu ^1& =& \gamma _\mu ^\mathrm{T},\hfill & O_\mu ^2\hfill & =& \frac{6}{\sqrt{5}}\left(\widehat{k}_\mu ^\mathrm{T}\gamma ^\mathrm{T}\widehat{k}\frac{1}{3}\gamma _\mu ^\mathrm{T}\widehat{k}^\mathrm{T}\widehat{k}^\mathrm{T}\right),\hfill \\ \hfill O_\mu ^3& =& 2\widehat{k}^\mathrm{T}\gamma \widehat{P},\hfill & O_\mu ^4\hfill & =& i\sqrt{2}\left(\gamma _\mu ^\mathrm{T}\gamma \widehat{k}^\mathrm{T}\gamma \widehat{P}+\widehat{k}_\mu ^\mathrm{T}\gamma \widehat{P}\right),\hfill \\ \hfill O_\mu ^5& =& 2\widehat{k}_\mu ^\mathrm{T},\hfill & O_\mu ^6\hfill & =& \frac{i}{\sqrt{2}}\left(\gamma _\mu ^\mathrm{T}\gamma ^\mathrm{T}\widehat{k}\gamma ^\mathrm{T}\widehat{k}\gamma _\mu ^\mathrm{T}\right),\hfill \\ \hfill O_\mu ^7+\frac{1}{\sqrt{2}}O_\mu ^8& =& i\sqrt{\frac{3}{5}}[1+(\widehat{k}\widehat{P})^2]\left(\gamma _\mu ^\mathrm{T}\gamma \widehat{P}\gamma \widehat{P}\gamma _\mu ^\mathrm{T}\right),\hfill & O_\mu ^8\hfill & =& 2i\sqrt{\frac{6}{5}}\widehat{k}_\mu ^\mathrm{T}\gamma ^\mathrm{T}\widehat{k}\gamma \widehat{P},\hfill \end{array}$$
(2.3.12)
where: $`\gamma _\mu ^\mathrm{T}:=\gamma _\mu +\gamma \widehat{P}\widehat{P}_\mu `$, $`\widehat{P}\widehat{P}=1`$; and $`\widehat{k}\widehat{k}=1`$, $`\widehat{k}_\mu ^\mathrm{T}:=\widehat{k}_\mu +\widehat{k}\widehat{P}\widehat{P}_\mu `$. With this decomposition the magnitudes of the invariant functions, $`F_l`$, are a direct measure of the relative importance of a given Dirac covariant in the BSA; e.g., one expects $`F_1`$ to be the function with the greatest magnitude for $`J^{PC}=1^{}`$ bound states.
To calculate the meson masses, one first solves Eqs. (2.2.39) and (2.2.40) for the renormalised dressed-quark propagator. This numerical solution for $`S(p)`$ is used in the pseudoscalar BSE, Eq. (2.3.5) under the substitution of Eq. (2.3.9), which is a coupled set of four homogeneous equations, one set for each meson; and the vector BSE, Eq. (2.3.5) with Eq. (2.3.11), a coupled set of eight equations for each meson. Solving the equations is a challenging numerical exercise, requiring careful attention to detail, and two complementary methods were both used in Refs. . While the numerical methods were identical, the authors of Ref. used a simplified version of the effective interaction:
$$\frac{𝒢(k^2)}{k^2}=\frac{4\pi ^2}{\omega ^6}Dk^2\mathrm{e}^{k^2/\omega ^2}+4\pi \frac{\gamma _m\pi }{\frac{1}{2}\mathrm{ln}\left[\tau +\left(1+k^2/\mathrm{\Lambda }_{\mathrm{QCD}}^2\right)^2\right]}(k^2),$$
(2.3.13)
and varied the single parameter $`D`$ along with the current-quark masses: $`\widehat{m}_u=\widehat{m}_d`$ and $`\widehat{m}_s`$, in order to reproduce the observed values of $`m_\pi `$, $`m_K`$ and $`f_\pi `$. All other calculated results are predictions in the sense that they are unconstrained. (NB. The Poincaré invariant four-dimensional BSE is solved directly, eschewing the commonly used artefice of a three-dimensional reduction, which introduces spurious effects when imposing compatibility with Goldstone’s theorem and also leads to a misinterpretation of a model’s parameters .)
Before reporting the results it is necessary to introduce the formulae for the meson decay constants: $`f_H`$, which completely describe the strong interaction contribution to a meson’s weak or electromagnetic decay. Following Ref. , the pseudoscalar meson decay constant is given by
$$\frac{1}{\sqrt{2}}f_HP_\mu :=0|\overline{𝒬}(𝒯^H)^\mathrm{t}\gamma _\mu \gamma _5𝒬|H(P)=Z_2\mathrm{tr}_k^\mathrm{\Lambda }\left(𝒯^H\right)^\mathrm{t}\gamma _5\gamma _\mu 𝒮(k_+)\mathrm{\Gamma }_H(k;P)𝒮(k_{}),$$
(2.3.14)
where here $`𝒬=\mathrm{column}(u,d,s)`$. The factor of $`Z_2`$ on the r.h.s. ensures that $`f_H`$ is gauge invariant, and independent of the renormalisation point and regularisation mass-scale; i.e., that $`f_H`$ is truly an observable. Equation (2.3.14) is the pseudovector projection of the unamputated Bethe-Salpeter wave function, $`\chi (k;P)`$, calculated at the origin in configuration space. As such, it is one field theoretical generalisation of the “wave function at the origin,” which describes the decay of bound states in quantum mechanics. The analogous expression for vector mesons is
$`\frac{1}{\sqrt{2}}f_Hm_Hϵ_\mu ^\lambda (P)`$ $`:=`$ $`0|\overline{𝒬}(𝒯^H)^\mathrm{t}\gamma _\mu 𝒬|H(P)`$ (2.3.15)
$``$ $`\frac{1}{\sqrt{2}}f_Hm_H=\frac{1}{3}Z_2\mathrm{tr}{\displaystyle _k^\mathrm{\Lambda }}\left(𝒯^H\right)^\mathrm{t}\gamma _\mu 𝒮(k_+)\mathrm{\Gamma }_\mu ^H(k;P)𝒮(k_{}),`$
where $`ϵ_\mu ^\lambda (P)`$ is the vector meson’s polarisation vector: $`Pϵ^\lambda (P)=0`$.
A best-fit is obtained with
$$D=(1.12\mathrm{GeV})^2,$$
(2.3.16)
which is a 60% increase over Eq. (2.2.36), as expected because the single Gaussian term in Eq. (2.3.13) must here replace the sum of the first two terms in Eq. (2.2.35); and renormalised current-quark masses
$$\begin{array}{cc}m_{u,d}^{1\mathrm{GeV}}=5.5\mathrm{MeV},& m_s^{1\mathrm{GeV}}=124\mathrm{MeV},\end{array}$$
(2.3.17)
which are little changed from the values used in Ref. , Eq. (2.2.55). \[NB. $`\omega `$ and $`m_t`$ are unchanged. See the discussion preceding Eq. (2.2.36).\] The results are presented in Table 2.1, and are characterised by a root-mean-square error over predicted quantities of just $`3.6`$%. We emphasise that this is obtained with a one-parameter model of the effective interaction.
Experimental values of the pseudoscalar meson decay constants are obtained directly via observation of their prominent $`\beta `$-decay mode. However, for the vector mesons this mode is not easily accessible and to proceed we note that the rainbow-ladder truncation predicts ideal flavour mixing. Using this and isospin symmetry, one can relate the $`f_\rho `$ matrix element to that describing $`\rho e^+e^{}`$ decay:
$$\frac{m_\rho ^2}{g_\rho }ϵ_\mu ^\lambda (P):=\frac{1}{\sqrt{2}}0|\overline{𝒬}Q_e(𝒯^{\rho ^0})^\mathrm{t}\gamma _\mu 𝒬|\rho _\lambda ^0(P)=0|\overline{𝒬}(𝒯^\rho ^{})^\mathrm{t}\gamma _\mu 𝒬|\rho _\lambda ^{}(p)=\frac{1}{\sqrt{2}}f_\rho m_\rho ϵ_\mu ^\lambda (P),$$
(2.3.18)
with $`Q_e:=\mathrm{diag}[2/3,1/3,1/3]`$; i.e., the quark’s electromagnetic charge matrix. $`\mathrm{\Gamma }_{\rho ^0e^+e^{}}=6.77\pm 0.32`$keV $``$ $`g_\rho =5.03\pm 0.12`$ and hence $`f_\rho =216\pm 5`$MeV. For the $`\varphi `$-meson
$$\frac{m_\varphi ^2}{g_\varphi }ϵ_\mu ^\lambda (P):=\frac{1}{3}0|\overline{s}\gamma _\mu s|\varphi _\lambda (P):=\frac{1}{3}f_\varphi m_\varphi ϵ_\mu ^\lambda (P)$$
(2.3.19)
and hence $`\mathrm{\Gamma }_{\varphi e^+e^{}}=1.37\pm 0.05`$keV $``$ $`g_\varphi =12.9\pm 0.2`$ or $`f_\varphi =238\pm 4`$MeV. $`f_K^{}`$ follows from
$$\mathrm{\Gamma }_{\tau K^{}\nu _\tau }/\mathrm{\Gamma }_{\tau \rho \nu _\tau }=0.051f_K^{}=1.042f_\rho .$$
(2.3.20)
A number of other important observations are recorded in Refs. . First: The calculated values of observable quantities are independent of the momentum partitioning parameter: $`\eta _P`$, when all of the Dirac covariants, and their complete momentum dependence, are retained; i.e., Poincaré invariance is manifest. Second: For pseudoscalar mesons the leading $`\gamma _5`$-covariant is dominant but the pseudovector components also play an important role; e.g., $`f_K`$ is $`<30`$% smaller without them. For the vector mesons, while $`F_1`$ is dominant, $`F_{2\mathrm{}5}`$ are also important; e.g., $`m_\rho `$ is $`>20`$% larger without them. Third: Flavour nonsinglet pseudoscalar mesons obey a mass formula , exact in QCD,
$$f_Hm_H^2=_H^\zeta r_H^\zeta ,_H^\zeta =\mathrm{tr}_F[M_\zeta \{𝒯^H,(𝒯^H)^\mathrm{t}\}],$$
(2.3.21)
where
$$r_H^\zeta =i\sqrt{2}Z_4\mathrm{tr}_k^\mathrm{\Lambda }\left(𝒯^H\right)^\mathrm{t}\gamma _5𝒮(k_+)\mathrm{\Gamma }_H(k;P)𝒮(k_{}):=2i\overline{q}q_\zeta ^H\frac{1}{f_H}$$
(2.3.22)
is the gauge-invariant and cutoff-independent residue of the pion pole in the pseudoscalar vertex. As a residue, it is an analogue of $`f_H`$ and describes the pseudoscalar projection of the unamputated Bethe-Salpeter wave function calculated at the origin in configuration space. For small current-quark masses, Eq. (2.3.21) yields the “Gell-Mann–Oakes–Renner” relation as a corollary. However, it is also valid for heavy-quarks and predicts $`m_H\widehat{m}_Q`$ in the heavy-meson domain, which is verified in the strong interaction spectrum. Fourth: The behaviour at large $`k^2`$ \[ultraviolet relative momenta\] is model independent, determined as it is by the one-loop improved strong running coupling. Fifth: For the pseudoscalar mesons, the asymptotic behaviour of the subdominant pseudovector amplitudes \[$`F_H`$, $`G_H`$\] is crucial for convergence of the integral describing $`f_H`$. For the vector mesons, the same is true of $`F_{2\mathrm{}5}`$. This is also a model-independent result because of the fourth point.
This subsection illustrates the reliability of the rainbow-ladder truncation for light vector and flavour nonsinglet pseudoscalar mesons. That is not an accident but rather, as elucidated in Ref. , it is the result of cancellations between vertex corrections and crossed-box contributions at each higher order in the quark-antiquark scattering kernel. There are two other classes of light-meson: scalar and axial-vector. A separable model that expresses characteristics of the rainbow-ladder truncation has been employed successfully in calculating the masses and decay constants of $`u`$,$`d`$-quark axial-vector mesons . Hence, while a more sophisticated study is still lacking, and is indeed required, this suggests that the truncation can provide a good approximation in this sector; i.e., that the conspiratorial cancellations are also effective here.
In the scalar channel, however, the rainbow-ladder truncation is not certain to provide a reliable approximation because the cancellations described above do not occur . This is entangled with the phenomenological difficulties encountered in understanding the composition of scalar resonances below $`1.4`$GeV . For the isoscalar-scalar meson the problem is exacerbated by the presence of timelike gluon exchange contributions to the kernel, which are the analogue of those diagrams expected to generate the $`\eta `$-$`\eta ^{}`$ mass splitting in BSE studies . If the rainbow-ladder truncation is employed one obtains ideal flavour mixing and degenerate isospin partners, and; e.g.,
$$\begin{array}{ccccc}& \hfill \text{Ref.}& (u/d)_{I=0,1}& u\overline{s}& s\overline{s}\\ & & & & \\ \mathrm{calculated}\mathrm{mass}_{\mathrm{in}\mathrm{GeV}}\hfill & \hfill \text{[47]}& 0.59& 0.90& 1.20\\ & \hfill \text{[52]}& 0.67& & \\ & \hfill \text{[114]}& 0.71& 1.18& 1.54\\ & \hfill \text{[121]}& 0.59& & \\ & & & & \\ \mathrm{averaged}\mathrm{mass}\hfill & & 0.64\pm 0.06& 1.04\pm 0.20& 1.37\pm 0.24\end{array}$$
(2.3.23)
($`\mathrm{\hspace{0.17em}0}^{++}`$ mesons containing at least one $`s`$-quark were not considered in Refs. .) Each model represented in this compilation was constrained to accurately describe $`\pi `$\- and $`K`$-meson observables, and the standard deviation about their averaged vector meson masses is a factor of $`2`$$`5`$ smaller than that exhibited in the last row here.
In Eq. (2.3.23) we do not compare directly with observed masses because of the uncertainty in identifying the members of the scalar nonet. We only note that: 1) the $`u\overline{s}`$ channel is least affected by those corrections to the rainbow-ladder truncation that can alter the feature of ideal flavour mixing and hence it may be appropriate to identify the $`u\overline{s}`$ scalar with the $`K_0^{}(1430)`$, in which case the mass is underestimated by $`<30`$%; and 2) an analysis of $`\pi \pi `$ data identifies an isoscalar-scalar with
$$m_{0_{I=0}^{++}}0.46\mathrm{GeV},\mathrm{\Gamma }_{0_{I=0}^{++}\pi \pi }0.220.47\mathrm{GeV}.$$
(2.3.24)
This supports ideal flavour mixing but, because $`\mathrm{\Gamma }_\sigma /m_\sigma `$ is large, suggests an accurate calculation of this state’s mass will require a Bethe-Salpeter kernel that explicitly includes couplings to the important $`\pi \pi `$ mode. In contrast, that coupling can be handled perturbatively in the $`\omega `$-$`\rho `$ sector . If this identification is correct then the mass estimate in Eq. (2.3.23) is $`<40`$% too large.
The discussion here makes plain that the constituent-quark-like rainbow-ladder scalar bound states are significantly altered by corrections to the BSE’s kernel. That is good because it is consistent with observation: understanding the scalar meson nonet is a complex problem. (NB. The difficulties to be anticipated are illustrated; e.g., in Ref. .) Using the DSEs this complexity is expressed in unanswered questions. For example, why do timelike gluon exchange contributions to the kernel in the isoscalar-scalar channel not force a deviation from ideal flavour mixing and, returning to the pseudoscalar sector, what effect do they play in the $`\eta `$-$`\eta ^{}`$ mass splitting?
There are other contemporary questions. For example, which improvements to the rainbow-ladder kernel are necessary in order to study bound states containing at least one heavy-quark? The extent to which the cancellations elucidated in Ref. persist as the current-quark mass evolves to values larger than $`M_\chi `$ is not known. Even though the rainbow-ladder truncation can provide an acceptable estimate of heavy-meson masses; e.g., Ref. , improvements are necessary and potentials derived from dual-QCD models; e.g., Refs. , have been applied more exhaustively . (Heavy-heavy-mesons are also amenable to study via heavy-quark expansions ). Other composites admitted by QCD can also be considered. The rainbow-ladder truncation has recently been applied to “exotic” light-mesons, predicting a $`J^{PC}=1^+`$ meson with a mass of $`1.4`$$`1.5`$GeV . \[“exotic” because such a $`J^{PC}`$ value is unobtainable in the $`q\overline{q}`$ constituent quark model.\] However, there is a dearth of DSE applications to the glueball spectrum, which has been explored using other methods; e.g., Refs. . A gluon-sector analogue of the rainbow-ladder truncation is an obvious starting point and such studies must precede any exploration of hybrid quark-gluon states .
That applications in all these areas are actively being pursued and contemplated is an indication of a healthy discipline.
Electromagnetic Pion Form Factor and Vector Dominance We have seen that the renormalisation group improved rainbow-ladder truncation of the quark-antiquark scattering kernel: $`K`$, provides a good understanding of colour singlet, mesonic spectral functions. That is a key success since describing these correlation functions is a core problem in QCD . However, it is only a single step and one must proceed from this foundation to the study of scattering observables, which delve deeper into hadron structure.
The best such observables to study are elastic electromagnetic form factors, because the probe is well understood, and the simplest “target” for a theorist is the pion, as long as the theoretical framework accurately describes DCSB. The electromagnetic pion form factor is a much studied observable but here, to be concrete, we focus on the application of DSEs to this problem, and in that connection Ref. is a pilot. As our exemplar we choose Ref. because it is a direct application of the effective interaction described in the previous subsection, Eq. (2.3.13), and it is the most complete study to date.
In the isospin-symmetric limit the renormalised impulse approximation to the pion’s electromagnetic form factor is
$`(p_1+p_2)_\mu F_\pi (q^2):=\mathrm{\Lambda }_\mu (p_1,p_2)`$
$`=2\mathrm{tr}{\displaystyle _k^\mathrm{\Lambda }}\overline{\mathrm{\Gamma }}_\pi (k;p_2)𝒮(k_{++})iQ_e\mathrm{\Gamma }_\mu ^\gamma (k_{++},k_+)𝒮(k_+)\mathrm{\Gamma }_\pi (kq/2;p_1)𝒮(k_{}),`$
$`k_{\alpha \beta }:=k+\alpha p_1/2+\beta q/2`$ and $`p_2:=p_1+q`$. Here, $`\mathrm{\Gamma }_\pi (k;P)`$ is the pion BSA, which has the form in Eq. (2.3.9), and $`𝒮(k)=\mathrm{diag}[S_{u=d}(k),S_{u=d}(k)]`$. No renormalisation constants appear explicitly in Eq. (2.3) because the renormalised dressed-quark-photon vertex: $`\mathrm{\Gamma }_\mu ^\gamma `$, satisfies the vector Ward-Takahashi identity:
$$(p_1p_2)_\mu i\mathrm{\Gamma }_\mu ^\gamma (p_1,p_2)=S^1(p_1)S^1(p_2).$$
(2.3.26)
Importantly, this also ensures current conservation:
$$(p_1p_2)_\mu \mathrm{\Lambda }_\mu (p_1,p_2)=0,$$
(2.3.27)
and, using Eq. (2.3.6), the correct normalisation of the form factor:
$$F(q^2=0)=1;$$
(2.3.28)
i.e., combining the rainbow-ladder truncation of the quark-antiquark scattering kernel with the impulse approximation yields a minimal, consistent approximation . \[NB. As the $`\pi ^0`$ is a charge conjugation eigenstate, $`F_{\pi ^0}(q^2)0`$, $`q^2`$. The impulse approximation yields this result.\]
The only quantity in Eq. (2.3) not already known is $`\mathrm{\Gamma }_\mu ^\gamma `$. It is the solution of an inhomogeneous BSE, which in rainbow-ladder truncation is
$$\mathrm{\Gamma }_\mu ^\gamma (k;P)=Z_2\gamma _\mu _q^\mathrm{\Lambda }𝒢((kq)^2)D_{\alpha \beta }^{\mathrm{free}}(kq)\frac{\lambda ^a}{2}\gamma _\alpha 𝒮(q_+)\mathrm{\Gamma }_\mu ^\gamma (k;P)𝒮(q_{})\frac{\lambda ^a}{2}\gamma _\beta .$$
(2.3.29)
It is straightforward to verify that Eq. (2.3.26) is satisfied, and owing to this the general solution of Eq. (2.3.29) involves only eight independent scalar functions and can be expressed as
$$\mathrm{\Gamma }_\mu ^\gamma (k;P)=\mathrm{\Gamma }_\mu ^{\mathrm{BC}}(k;P)+\underset{l=1}{\overset{8}{}}O_\mu ^lF_l^\gamma (k;P),$$
(2.3.30)
where the matrices $`O_\mu ^l`$ are given in Eq. (2.3.12) and
$$\mathrm{\Gamma }_\mu ^{\mathrm{BC}}(k;P)=i\mathrm{\Sigma }_A(k_+^2,k_{}^2)\gamma _\mu +(k_++k_{})_\mu \left[\frac{1}{2}i\gamma (k_++k_{})\mathrm{\Delta }_A(k_+^2,k_{}^2)+\mathrm{\Delta }_B(k_+^2,k_{}^2)\right];$$
(2.3.31)
$$\mathrm{\Sigma }_F(k_+^2,k_{}^2)=\frac{1}{2}[F(k_+^2)+F(k_{}^2)],\mathrm{\Delta }_F(k_+^2,k_{}^2)=\frac{F(k_+^2)F(k_{}^2)}{k_+^2k_{}^2},$$
(2.3.32)
with $`F=A,B`$; i.e., the scalar functions in the dressed-quark propagator. $`\mathrm{\Gamma }_\mu ^{\mathrm{BC}}(k;P)`$ saturates the vector Ward-Takahashi identity, and the remaining terms are transverse and hence do not contribute to the r.h.s. of Eq. (2.3.26).
The importance of determining the dressed-quark-photon vertex from Eq. (2.3.29) was recognised in Ref. , where a solution was obtained using the simple model of $`𝒢(k^2)`$ introduced in Ref. . As is readily anticipated, the dressed-vertex exhibits isolated simple poles at timelike values of $`P^2`$. Each pole corresponds to a $`1^{}`$ bound state, $`P^2=`$ mass-squared, and its matrix-valued residue is proportional to the bound state amplitude. For $`Q^2`$ in the neighbourhood of any one of these poles the behaviour of the pion form factor is primarily determined by the manifestation of that bound state in the $`1^{}`$ spectral density. Vector meson dominance, in any of its forms, is an Ansatz to be used for extrapolating outside of these neighbourhoods. A direct solution of Eq. (2.3.29) obviates the need for such an expedient and also any need to fabricate an interpretation of an off-shell bound state.
Using the interaction of Eq. (2.3.13), with its single parameter fixed as discussed in connection with Eq. (2.3.16), and solving numerically for the renormalised dressed-quark propagator, pion BSA and dressed-quark-photon vertex, Ref. obtains the pion form factor depicted in Fig. 2.4 with
$$r_\pi ^2:=6\frac{dF_\pi (y)}{dy}|_{y=0}r_\pi =0.68\mathrm{fm}\mathrm{cf}.r_\pi ^{\mathrm{Obs}}=0.663\pm 0.006.$$
(2.3.33)
The complete calculation describes the data for both spacelike and timelike momenta, plainly exhibiting the evolution to a simple pole corresponding to the $`\rho `$-meson. Here the $`\rho `$ is described by a simple pole on the real-$`P^2`$ axis because the rainbow-ladder truncation excludes the $`\pi \pi `$ contribution in the kernel. However, that can be included perturbatively and estimates show it yields no-more-than a $`15`$% increase in $`r_\pi `$. \[NB. As illustrated by Refs. , the extension to $`K`$-meson form factors is straightforward but with the interesting new feature that, unlike the neutral pion’s elastic form factor, $`F_{K^0}(Q^2)0`$, $`Q^2>0`$, because the neutral kaons are not charge-conjugation eigenstates.\]
Proponents of vector meson dominance Ansätze have historically claimed support in the accuracy of the estimate $`r_\pi ^26/m_\rho ^2`$ and in this connection one can ask what alternative does the direct calculation describe? To address this Ref. observed that, for $`Q^2[m_\rho ^2,0.2\mathrm{GeV}^2]`$, an interpolation of the solution of Eq. (2.3.29) is provided by
$$\mathrm{\Gamma }_\mu ^\gamma (k;Q)\mathrm{\Gamma }_\mu ^{\mathrm{BC}}(k;Q)\frac{1}{g_\rho }\frac{Q^2}{Q^2+m_\rho ^2}\underset{l=1}{\overset{5}{}}O_\mu ^lF_l^\rho (k^2|Q^2=m_\rho ^2),$$
(2.3.34)
where
$$F_l^\rho (k^2|Q^2=m_\rho ^2):=\frac{2}{\pi }_1^1𝑑x\sqrt{1x^2}F_l^\rho (k^2,ikm_\rho x|Q^2=m_\rho ^2)$$
(2.3.35)
are the leading Chebyshev moments of the five dominant scalar functions in the $`\rho `$-meson BSA and, following Ref. , $`1/g_\rho `$ is the residue of the $`\rho `$-meson pole in the photon vacuum polarisation. Substituting Eq. (2.3.34) into Eq. (2.3) yields
$$F_\pi (Q^2)F_\pi ^{\mathrm{BC}}(Q^2)\frac{g_{\rho \pi \pi }}{g_\rho }F_{\rho \pi \pi }(Q^2)\frac{Q^2}{Q^2+m_\rho ^2},$$
(2.3.36)
where $`g_{\rho \pi \pi }F_{\rho \pi \pi }(Q^2)`$ is the impulse approximation to the $`\rho \pi \pi `$ vertex, $`F_{\rho \pi \pi }(Q^2=m_\rho ^2)=1`$. Implicit in Eq. (2.3.34) is a clear but not unique definition of an off-shell $`\rho `$-$`\overline{q}q`$ correlation, which yields Eq. (2.3.36) as a calculable approximation to the form factor wherein
$$(r_\pi ^\rho )^2:=\frac{6}{m_\rho ^2}\frac{g_{\rho \pi \pi }}{g_\rho }F_{\rho \pi \pi }(0).$$
(2.3.37)
This is simply the vector meson dominance result corrected for nonuniversality of the strong and electromagnetic $`\rho `$-meson couplings; i.e., $`g_\rho g_{\rho \pi \pi }`$, and, via $`F_{\rho \pi \pi }(0)`$, for the nonpointlike nature of the off-shell $`\rho `$-$`\overline{q}q`$ correlation. $`F_{\rho \pi \pi }(Q^2)`$ was not calculated in Ref. and therefore here we make an estimate. Typically $`F_{\rho \pi \pi }(0)0.5`$ and using experimental values for the observables \[including $`g_{\rho \pi \pi }=6.05`$\] we find
$$(r_\pi ^\rho )^20.5r_\pi ^2,$$
(2.3.38)
which is a significant suppression with respect to the “naive” vector meson dominance estimate.
This illustration demonstrates that the pion charge radius is indeed influenced by the $`\rho `$-pole’s contribution to the $`1^{}`$ spectral density. However, that is unsurprising: the bound state poles are a significant feature of the dressed-vertex. More important is a realisation that the separation in Eq. (2.3.34) is completely arbitrary and hence so is the fraction of the charge radius attributed to the “off-shell $`\rho `$-meson.” One can shift any amount of strength between the two terms and yet maintain an accurate interpolation of the dressed-vertex. For example,
$$\frac{Q^2}{Q^2+m_\rho ^2}=\frac{Q^2}{Q^2+m_\rho ^2}\mathrm{e}^{\omega (1+Q^2/m_\rho ^2)}+\frac{Q^2}{Q^2+m_\rho ^2}\left[1\mathrm{e}^{\omega (1+Q^2/m_\rho ^2)}\right],$$
(2.3.39)
is an apparently nugatory rearrangement. However, only the first term has a pole at $`Q^2+m_\rho ^2=0`$, and substituting Eq. (2.3.39) into Eq. (2.3.34) and repeating the analysis one arrives at
$$(r_\pi ^\rho )^2=\mathrm{e}^\omega \frac{6}{m_\rho ^2}\frac{g_{\rho \pi \pi }}{g_\rho }F_{\rho \pi \pi }(0)\stackrel{\omega =1}{}\mathrm{\hspace{0.17em}0.2}r_\pi ^2.$$
(2.3.40)
It is obvious now that $`r_\pi ^\rho `$ can be made arbitrarily small while preserving $`r_\pi r_\pi ^{\mathrm{Obs}}`$. \[NB. Eq. (2.3.40) and the procedure leading to it are no more contrived than Eq. (2.3.37) because Eq. (2.3.39) can be interpreted as “unfreezing” the vector meson BSA; i.e., of allowing the bound state correlation to be suppressed off-shell: $`F_l^\rho (k^2|Q^2=m_\rho ^2)\mathrm{exp}(\omega [1+Q^2/m_\rho ^2])F_l^\rho (k^2|Q^2=m_\rho ^2)`$.\]
For large spacelike $`Q^2`$ the calculation of $`F_\pi (Q^2)`$ directly from $`𝒢(k^2)`$ is a computational challenge and for practical reasons the study in Ref. was restricted to $`Q^21`$GeV. The effective interaction doesn’t appear explicitly in the impulse approximation, Eq. (2.3), only the dressed-quark propagator, pion Bethe-Salpeter amplitude and dressed-quark-photon vertex. Hence the numerical analysis is markedly simplified if algebraic approximations to these Schwinger functions are employed.
This was the approach adopted in Ref. , which has since been used in a wide range of applications; e.g., Refs. . It is efficacious because the vector and axial-vector Ward-Takahashi identities can be used to motivate Ansätze for $`\mathrm{\Gamma }_\pi `$ and $`\mathrm{\Gamma }_\mu ^\gamma `$ that are expressed solely in terms of $`S`$. In the present context, algebraic models for $`S`$ and $`\mathrm{\Gamma }_\pi `$ have been developed that encode the important qualitative aspects of the DSE and BSE solutions in Ref. . These forms, along with $`\mathrm{\Gamma }_\mu ^\gamma `$ determined via Eq. (2.3.31), make possible a calculation of $`F_\pi (Q^2)`$, $`Q^2>0`$, and an analytic analysis of the asymptotic behaviour. One finds that the pion’s pseudoscalar covariant can alone provide a quantitative description of $`F_\pi (Q^2)`$ for $`Q^2<5`$GeV<sup>2</sup>. However, beyond this point the pseudovector covariants become important. It is these terms in Eq. (2.3.9) that ensure
$$Q^2F_\pi (Q^2)=\mathrm{const}.,$$
(2.3.41)
up to calculable $`\mathrm{ln}(Q/\mathrm{\Lambda }_{\mathrm{QCD}})^p`$-corrections, and this behaviour is unmistakable for $`Q^2>10`$GeV<sup>2</sup>. The anomalous dimension: $`p`$, is determined by that of $`F_\pi `$ and $`G_\pi `$, which is a model independent result. \[NB. If the pseudovector covariants are neglected, $`Q^4F_\pi (Q^2)=\mathrm{const}.`$. Further, the dressed-quark propagator obtained in Ref. necessarily yields results inconsistent with Eq. (2.3.41).\] This application is an emphatic demonstration of the DSEs’ ability to provide a single and simultaneous description of the infrared and ultraviolet aspects of observables.
That is also much in evidence in the study of anomalous processes. As first observed in connection with processes like $`K^+K^{}\pi ^+\pi ^0\pi ^{}`$, the systematically truncated DSEs yield the anomalies of current algebra without any model dependence; i.e., the anomalies remain a feature of the global aspects of DCSB . For the $`\pi ^0\gamma \gamma `$ process this feature was verified in Refs. and for the $`\gamma \pi ^{}\pi \pi `$ transition form factor, in Ref. . Consequently the DSEs provide a single framework wherein the value of such transition form factors is fixed and model-independent at the soft-pion, zero momentum transfer point, and their evolution is calculable on the entire range of momentum transfer, reproducing the ultraviolet behaviour anticipated from perturbative QCD. These features have been elucidated and exploited in Refs. .
Describing Baryons Hitherto we have described the application of DSEs to the study of mesonic observables, which requires and illustrates a contemporary understanding of the two-body problem in quantum field theory. Baryons, as a three-body problem, pose a greater challenge and historically they have been described using constituent-quark-like models, which remain an important contemporary tool; e.g., Refs. . We identify a beginning of progress with a direct assault on the baryon problem in a realisation that field theoretical models of the strong interaction admit the construction of a meson-diquark auxiliary-field effective action and thereby a description of baryons as loosely-bound quark-diquark composites.
A head-on DSE approach begins with a relativistic Faddeev equation, which can be derived by exploiting the fact that single gluon exchange between two quarks is attractive in the colour-antitriplet channel, whether or not the gluon is dressed. Indeed, using a rainbow-ladder truncation of the quark-quark scattering kernel, one obtains nonpointlike, colour-antitriplet diquark \[quark-quark\] bound states in the strong interaction spectrum . However, as demonstrated in Ref. , this is a flaw of the rainbow-ladder truncation: higher order terms in the kernel ensure that the quark-quark scattering matrix does not exhibit the simple poles that correspond to asymptotic states. Nevertheless, studies with improved kernels do support a physical interpretation of the spurious rainbow-ladder diquark masses. Denoting the mass by $`m_{qq}`$, then $`\mathrm{}_{qq}:=1/m_{qq}`$ represents the range over which a true diquark correlation in this channel can persist inside a baryon. In this sense they are “pseudo-particle” masses and can be used to estimate which two-body correlations should be retained in solving the Faddeev equation. (NB. Gluon mediated interactions are repulsive in the colour-sextet quark-quark channel, just as they are in the colour-octet meson channel .)
Reference tabulates calculated values of these pseudo-particle masses, from which we extract:
$$\begin{array}{ccccccccc}(qq)_{J^P}\hfill & (ud)_{0^+}& (us)_{0^+}& (uu)_{1^+}& (us)_{1^+}& (ss)_{1^+}& (uu)_1^{}& (us)_1^{}& (ss)_1^{}\\ & & & & & & & & \\ m_{qq}(\mathrm{GeV})\hfill & 0.74& 0.88& 0.95& 1.05& 1.13& 1.47& 1.53& 1.64\end{array}$$
(2.3.42)
The mass ordering is characteristic and model-independent, and indicates that an accurate study of the nucleon should retain the scalar and pseudovector correlations: $`(ud)_{0^+}`$, $`(uu)_{1^+}`$, because for these $`m_{qq}<m_N`$. This expectation is verified in calculations, where one finds that including the $`(uu)_{1^+}`$ correlation yields a nucleon whose mass is a welcome $`33`$% less-than that of a quark$`+`$scalar-diquark-only nucleon. We note that $`m_{(ud)_{0^+}}/m_{(uu)_{1^+}}=0.78`$ cf. $`0.76=m_N/m_\mathrm{\Delta }`$ and hence one might anticipate that the presence of diquark correlations in baryons is likely to provide a straightforward explanation of the $`N`$-$`\mathrm{\Delta }`$ mass-splitting. Lattice estimates, where available , agree with these results.
Reference is an extensive study of the octet and decuplet baryon spectrum based on a quark-diquark Faddeev equation. It represents the nucleon as a composite of a quark and pointlike diquark, which are bound together by a repeated exchange of roles between the dormant and diquark-participant quarks, as depicted in Fig. 2.5, and employs the confined-particle representation of the dressed-quark and -diquark propagators advocated in Ref. . Seven parameters appear in the model and, unfortunately, variations are permitted in some of them that allow the calculated results to override intuitive expectations. For example, the $`N`$-$`\mathrm{\Delta }`$ mass-splitting is fitted by adjusting the relative strength of the quark-diquark couplings in the scalar and pseudovector channels whilst simultaneously enforcing $`m_{(ud)_{0^+}}=m_{(uu)_{1^+}}`$. Bethe-Salpeter equation studies show that this is erroneous; i.e., these masses cannot be equal, and they and the couplings are not independent. Hence this means of generating the $`N`$-$`\mathrm{\Delta }`$ mass-splitting is unlikely to be completely correct. Nevertheless, the importance of the study is a demonstration that an accurate description of the spectrum is possible, and indeed no calculated mass is more-than $`1`$% away from its experimental value. Improvements are now possible and they will build on the lessons Ref. provides; e.g, replacing undetermined parameters with values obtained in precursor calculations. This is an important frontier.
Even before such thorough relativistic bound state calculations, the notion that diquark correlations in baryons could be significant found support in the analysis of scattering observables . As remarked earlier, these observables delve deep into hadron structure and so are a perennial focus of theory and experiment. With a growing understanding of the nature of diquark correlations, detailed nucleon models are now being applied in the nonperturbative evaluation of nucleon structure functions; e.g., Ref. , and nucleon form factors . As an exemplar, we review a calculation of the nucleon’s scalar form factor, which also yields the nucleon $`\sigma `$-term . This is an important application because it illustrates the only method that allows an unambiguous off-shell extrapolation in the estimation of meson-nucleon form factors and thereby provides an analogue for the discussion of vector meson dominance, Eqs. (2.3.34-2.3.40).
Underlying Ref. is the observation that nucleon propagation is described by a $`6`$-point Schwinger function
$$G_{\alpha \alpha ^{}}^{\tau \tau ^{}}(RR^{}):=\mathrm{\Psi }_\alpha ^\tau (R)\overline{\mathrm{\Psi }}_\alpha ^{}^\tau ^{}(R^{}),$$
(2.3.43)
where
$`\mathrm{\Psi }_\alpha ^\tau (R)`$ $`:=`$ $`{\displaystyle \underset{i=1}{\overset{3}{}}d^4x_i\psi (x_iR;\alpha _i,\tau _i;\alpha ,\tau )\epsilon _{abc}q_a(x_1;\alpha _1,\tau _1)q_b(x_2;\alpha _2,\tau _2)q_c(x_3;\alpha _3,\tau _3)},`$ (2.3.44)
$`\overline{\mathrm{\Psi }}_\alpha ^\tau (R)`$ $`:=`$ $`{\displaystyle \underset{i=1}{\overset{3}{}}d^4x_i\psi (x_iR;\alpha _i,\tau _i;\alpha ,\tau )^{}\epsilon _{abc}\overline{q}_a(x_1;\alpha _1,\tau _1)\overline{q}_b(x_2;\alpha _2,\tau _2)\overline{q}_c(x_3;\alpha _3,\tau _3)},`$ (2.3.45)
with: $`q_a(x_1;\alpha _1,\tau _1)`$, etc., quark Grassmann variables; $`\alpha _i`$, $`\alpha `$ quark and nucleon Dirac subscripts; $`\tau _i`$, $`\tau `$ the isospin analogues; and $`\epsilon _{abc}`$ ensuring colour neutrality. In these expressions $`\psi (x_iR;\alpha _i,\tau _i;\alpha ,\tau )`$ describes the distribution of quarks in the nucleon and; e.g., it can represent a nucleon Faddeev amplitude. The electromagnetic interaction of this nucleon is described by the current
$$J_\mu (R^{}R_0,R_0R)=\overline{\mathrm{\Psi }}(R^{})\overline{q}(R_0)iQ_e\gamma _\mu q(R_0)\mathrm{\Psi }(R).$$
(2.3.46)
To proceed, Ref. writes
$$\psi (x_iR;\alpha _i,\tau _i;\alpha ,\tau )=\underset{i=1}{\overset{3}{}}\frac{d^4p_i}{(2\pi )^4}\psi (p_i;\alpha _i,\tau _i;\alpha ,\tau )\mathrm{exp}\left[i\underset{i=1}{\overset{3}{}}p_i(x_iR)\right]$$
(2.3.47)
and employs a product Ansatz for the nucleon amplitude
$$\psi (p_i;\alpha _i,\tau _i;\alpha ,\tau )=\delta ^{\tau \tau _3}\delta _{\alpha \alpha _3}\psi (p_1+p_2,p_3)\mathrm{\Delta }(p_1+p_2)\mathrm{\Gamma }_{\alpha _1\alpha _2}^{\tau _1\tau _2}(p_1,p_2),$$
(2.3.48)
where $`\psi (\mathrm{}_1,\mathrm{}_2)`$ is a Bethe-Salpeter-like amplitude characterising the relative-momentum dependence of the correlation between diquark and quark \[$`\mathrm{\Phi }`$ in Fig. 2.5\], $`\mathrm{\Delta }(K)`$ describes the pseudo-particle propagation characteristics of the diquark, and
$`\mathrm{\Gamma }_{\alpha _1\alpha _2}^{\tau _1\tau _2}(p_1,p_2)`$ $`=`$ $`(Ci\gamma _5)_{\alpha _1\alpha _2}(i\tau ^2)^{\tau _1\tau _2}\mathrm{\Gamma }(p_1,p_2)`$ (2.3.49)
represents the momentum-dependence, and spin and isospin character of the diquark correlation; i.e., it corresponds to a Bethe-Salpeter-like amplitude for what here is a nonpointlike diquark. Complete antisymmetrisation is not explicit in this product Ansatz but that is effected in $`\mathrm{\Psi }`$ via the contraction with the Grassmann elements, Eq. (2.3.44).
The impulse approximation to the electromagnetic current in this model can now be obtained directly from Eq. (2.3.46). This $`8`$-point Schwinger function expresses a product of eight Grassmann variables and, via an analogue of Wick’s theorem, that can be reduced to a sum of products of four $`2`$-point dressed-quark Schwinger functions and the dressed-quark-photon vertex, with momenta and indices correlated via Eq. (2.3.48). Using the $`12`$ particle exchange symmetry exhibited explicitly by Eq. (2.3.48) one arrives at the result depicted in Fig. 2.6.
The product Ansatz is completely specified once explicit forms for the functions are given and Ref. employs
$`\psi (\mathrm{}_1,\mathrm{}_2)`$ $`=`$ $`{\displaystyle \frac{1}{𝒩_\mathrm{\Psi }}}(\mathrm{}^2/\omega _\psi ^2),\mathrm{}:=\frac{1}{3}\mathrm{}_1\frac{2}{3}\mathrm{}_2,`$ (2.3.50)
$`\mathrm{\Gamma }(q_1,q_2)`$ $`=`$ $`{\displaystyle \frac{1}{𝒩_\mathrm{\Gamma }}}(q^2/\omega _\mathrm{\Gamma }^2),q:=\frac{1}{2}q_1\frac{1}{2}q_2,`$ (2.3.51)
$`\mathrm{\Delta }(K)`$ $`=`$ $`{\displaystyle \frac{1}{m_\mathrm{\Delta }^2}}(K^2/\omega _\mathrm{\Gamma }^2),`$ (2.3.52)
$`(y)`$ $`=`$ $`{\displaystyle \frac{1\mathrm{e}^y}{y}},`$ (2.3.53)
whose parameters were fixed by fitting the proton’s charge form factor on $`Q^2[0,3]`$GeV<sup>2</sup> :
$$\begin{array}{cccc}& \omega _\psi & \omega _\mathrm{\Gamma }& m_\mathrm{\Delta }\\ & & & \\ \text{in GeV}\hfill & \mathrm{\hspace{0.33em}0.20}& 1.4& 0.63\end{array},\begin{array}{cccc}& \mathrm{\hspace{0.33em}1}/\omega _\psi & 1/\omega _\mathrm{\Gamma }& 1/m_\mathrm{\Delta }\\ & & & \\ \text{in fm}\hfill & 0.99& 0.14& 0.31\end{array}.$$
(2.3.54)
The two tables demonstrate the internal consistency of the model. $`d_\mathrm{\Gamma }:=1/\omega _\mathrm{\Gamma }`$ is a measure of the mean separation between the quarks constituting the scalar diquark and $`d_\psi :=1/\omega _\psi `$ is the analogue for the quark-diquark separation. $`d_\mathrm{\Gamma }<d_\psi `$ is necessary if the quark-quark clustering interpretation is to be valid. $`\mathrm{}_{(ud)_{0^+}}=1/m_\mathrm{\Delta }`$ is a measure of the range over which the diquark persists and that must be significantly less than the nucleon’s diameter. \[$`𝒩_\mathrm{\Psi }`$ and $`𝒩_\mathrm{\Gamma }`$ are the nucleon and $`(ud)`$ diquark normalisation constants, which are defined and calculated analogously to Eq. (2.3.6), and ensure composite electric charges of $`1`$ for the proton and $`1/3`$ for the diquark.\]
As is plain in Fig. 2.6, the calculation also involves the dressed-quark propagator and, based on the success of applications such as Refs. , the following algebraic parametrisation is employed in Refs. :
$`S(p)`$ $`=`$ $`i\gamma p\sigma _V(p^2)+\sigma _S(p^2),`$ (2.3.55)
$`\overline{\sigma }_S(x)`$ $`=`$ $`2\overline{m}(2(x+\overline{m}^2))+(b_1x)(b_3x)\left[b_0+b_2(ϵx)\right],`$ (2.3.56)
$`\overline{\sigma }_V(x)`$ $`=`$ $`{\displaystyle \frac{1}{x+\overline{m}^2}}\left[1(2(x+\overline{m}^2))\right],`$ (2.3.57)
$`x=p^2/\lambda ^2`$, $`\overline{m}`$ = $`m/\lambda `$, $`\overline{\sigma }_S(x)=\lambda \sigma _S(p^2)`$ and $`\overline{\sigma }_V(x)=\lambda ^2\sigma _V(p^2)`$. The mass-scale, $`\lambda =0.566`$GeV, and parameter values
$$\begin{array}{ccccc}\overline{m}& b_0& b_1& b_2& b_3\\ & & & & \\ 0.00897& 0.131& 2.90& 0.603& 0.185\end{array},$$
(2.3.58)
were fixed in a least-squares fit to light-meson observables. \[$`ϵ=10^4`$ in (2.3.56) acts only to decouple the large- and intermediate-$`p^2`$ domains.\] This simple form captures the essential features of direct solutions of Eq. (2.2.1). It represents the dressed-quark propagator as an entire function, which is inspired by the algebraic solutions in Refs. and ensures confinement via the means described in Sec. 2.2; and exhibits DCSB with
$$\overline{q}q_{1\mathrm{GeV}^2}=\lambda ^3\frac{3}{4\pi ^2}\frac{b_0}{b_1b_3}\mathrm{ln}\frac{1}{\mathrm{\Lambda }_{\mathrm{QCD}}^2}=(0.221\mathrm{GeV})^3,$$
(2.3.59)
which is calculated directly from Eqs. (2.2.12), (2.2.48) after noting that Eqs. (2.3.56), (2.3.57) yield the chiral limit quark mass function of Eq. (2.2.47) with $`\gamma _m=1`$. This is a general feature; i.e., the parametrisation exhibits asymptotic freedom at large-$`p^2`$ omitting only the additional $`\mathrm{ln}(p^2/\mathrm{\Lambda }_{\mathrm{QCD}}^2)`$-suppression, which is a useful but not necessary simplification. As we see here, this omission introduces model artefacts that are easily identified and accounted for.
Reference obtains a good description of the proton’s charge and magnetic form factors, and also the neutron’s magnetic form factor. The nonpointlike nature of the scalar diquark correlation makes possible a magnetic moment ratio of $`|\mu _n/\mu _p|=0.55`$, cf. $`0.68`$ experimentally. This ratio is always less-than $`0.5`$ when the correlation is pointlike. The neutron’s charge form factor is poorly described. The charge-radius-squared is negative, consistent with the data, but is $`60`$% too large. That defect results primarily from neglecting the contribution of the axial-vector diquark.
The nucleon’s scalar form factor is
$$\sigma (q^2)\overline{u}(P^{})u(P):=P^{}|m(\overline{u}u+\overline{d}d)|P,$$
(2.3.60)
where the nucleon spinors satisfy:
$$\gamma Pu(P)=iMu(P),\overline{u}(P)\gamma P=iM\overline{u}(P),$$
(2.3.61)
with the nucleon mass $`M=0.94`$GeV, $`q=(P^{}P)`$ and $`R=(P^{}+P)`$. The $`\pi N`$ $`\sigma `$-term is just $`\sigma (q^2=0)`$, which is the in-nucleon expectation value of the explicit chiral symmetry breaking term in the QCD Lagrangian. The general form of a fermion-scalar vertex is
$$\mathrm{\Lambda }_\mathrm{𝟏}(q,P)=f_1+i\gamma qf_2+i\gamma Rf_3+i\sigma _{\mu \nu }R_\mu q_\nu f_4,f_i=f_i(q^2,R^2)$$
(2.3.62)
since $`qR=0`$ for elastic processes. However, using Eqs. (2.3.61) the scalar current simplifies:
$`J_\mathrm{𝟏}(P^{},P)`$ $`:=`$ $`\overline{u}(P^{})\mathrm{\Lambda }_\mathrm{𝟏}(q,P)u(P)=s(q^2)\overline{u}(P^{})u(P),`$ (2.3.63)
$`s(q^2)`$ $`=`$ $`f_12Mf_3+q^2f_4.`$ (2.3.64)
The impulse approximation to $`J_\mathrm{𝟏}(P^{},P)`$ is also given by the five diagrams in Fig. 2.6 but with the dressed-quark-photon vertex replaced by the dressed-quark-scalar vertex, which is the solution of an inhomogeneous BSE analogous to Eq. (2.3.29). In Ref. , to hasten an exemplifying result, the scalar vertex equation was solved using the Goldstone-theorem-preserving separable model of Ref. . In that model the BSE assumes the form
$$\mathrm{\Gamma }_\mathrm{𝟏}(k;Q)=\mathrm{𝟏}\frac{4}{3}\frac{d^4q}{(2\pi )^4}\mathrm{\Delta }(kq)\gamma _\mu S(q_+)\mathrm{\Gamma }_\mathrm{𝟏}(q;Q)S(q_{})\gamma _\mu ,$$
(2.3.65)
with the interaction
$$\mathrm{\Delta }(kq)=G(k^2)G(q^2)+kqF(k^2)F(q^2),$$
(2.3.66)
where $`F(k^2)`$, $`G(k^2)`$ are regularised forms of the $`A`$, $`B`$ obtained from Eqs. (2.3.56), (2.3.57). \[NB. The regularisation ensures convergence of necessary integrals in the separable model.\] The solution of Eq. (2.3.65) is
$$\mathrm{\Gamma }_\mathrm{𝟏}(k;Q)=\mathrm{𝟏}+t_1(Q^2)G(k^2)+it_2(Q^2)F(k^2)\frac{kQ\gamma Q}{Q^2}+it_3(Q^2)F(k^2)\gamma k.$$
(2.3.67)
The $`\sigma `$-term is only sensitive to the vertex at $`Q^2=0`$, where the solution reduces to
$$\mathrm{\Gamma }_\mathrm{𝟏}(k;Q)|_{Q^2=0}=\mathrm{𝟏}+t_1(0)G(k^2)+t_3(0)F(k^2)i\gamma k,$$
(2.3.68)
with $`t_1(0)=0.242`$GeV, $`t_3(0)=0.0140`$GeV. Calculation shows that at $`k^2=0`$ the $`t_1`$-term is 6-times larger than the bare term; i.e., it is dominant in the infrared. That is to be expected because it represents the effect of the nonperturbative DCSB mechanism in the solution. This and the other $`t_i`$-terms vanish as $`k^2\mathrm{}`$, which is just a manifestation of asymptotic freedom.
What makes the inhomogeneous scalar BSE important here is that it has a solution for all $`Q^2`$ and that solution exhibits a pole at the $`\sigma `$-meson mass; i.e., in the neighbourhood of $`(Q^2)=m_\sigma ^2=(0.715\mathrm{GeV})^2`$
$$\mathrm{\Gamma }_\mathrm{𝟏}(k;Q)=\mathrm{𝑟𝑒𝑔𝑢𝑙𝑎𝑟}+\frac{n_\sigma m_\sigma ^2}{Q^2+m_\sigma ^2}\mathrm{\Gamma }_\sigma (k;Q),$$
(2.3.69)
where regular indicates terms that are regular in this neighbourhood and $`\mathrm{\Gamma }_\sigma (k;Q)`$ is the canonically normalised $`\sigma `$-meson Bethe-Salpeter amplitude, whose form is exactly that of $`(\mathrm{\Gamma }_\mathrm{𝟏}(k;Q)\mathrm{𝟏})`$ in Eq. (2.3.67). This is qualitatively identical to the behaviour of the solution of Eq. (2.3.29) discussed above in connection with $`F_\pi (Q^2)`$ and elucidated in Refs. . The simple pole appears in each of the functions $`t_i(Q^2)`$ and a pole fit yields, with $`m`$ the current-quark mass,
$$mn_\sigma =3.3\mathrm{MeV}.$$
(2.3.70)
$`n_\sigma m_\sigma ^2`$ is the analogue of the residue of the $`\pi `$-pole in the pseudoscalar vertex, Eq. (2.3.22): $`\sqrt{2}\overline{q}q_\pi /f_\pi `$, and its flow under the renormalisation group is identical. $`mn_\sigma `$ is renormalisation point independent and its value can be compared with
$$\frac{m\overline{q}q_\pi }{\frac{1}{\sqrt{2}}f_\pi }\frac{1}{m_\sigma ^2}=3.6\mathrm{MeV};$$
(2.3.71)
i.e., the magnitude of $`n_\sigma `$ is typical of effects driven by dynamical chiral symmetry breaking. The $`\sigma `$-meson-$`\overline{q}q`$ coupling is defined on-shell
$$g_{\sigma \overline{q}q}:=\mathrm{\Gamma }_\sigma (0;Q)|_{Q^2=m_\sigma ^2}=12.6,$$
(2.3.72)
and its magnitude can be placed in context via a comparison with $`g_{\pi \overline{q}q}=11.8`$, which is defined analogously.
The expectation value in Eq. (2.3.60) is obtained with
$$\mathrm{\Gamma }_m(k;Q)=m\mathrm{\Gamma }_\mathrm{𝟏}(k;Q)$$
(2.3.73)
as the probe vertex in Fig. 2.6, and using the solution for $`\mathrm{\Gamma }_\mathrm{𝟏}`$ described here the calculated $`\sigma `$-term is
$$\sigma /M_N=0.015.$$
(2.3.74)
No parameters were varied to obtain this result, which may be compared with a recent lattice computation : $`\sigma /M_N=0.019\pm 0.05`$, calculated using an extrapolation in the current-quark mass. Alternative extrapolation methods can lead to larger values; e.g., $`\sigma /M_N=0.047`$$`0.059`$ , which are also suggested by some phenomenological analyses . \[The $`\sigma `$-term is not directly accessible via experiment but a value is theoretically inferred by extrapolating $`\pi N`$ scattering data using dispersion relations : $`\sigma /M_N=0.047`$$`0.076`$.\] Simple estimates indicate that the result in Eq. (2.3.74) increases with decreasing $`m_\sigma `$, and a reduction in $`m_\sigma `$ is a likely consequence of using an improved kernel in the inhomogeneous scalar BSE. Thus Eq. (2.3.74) is an excellent first estimate.
The nucleon’s scalar form factor is depicted in Fig. 2.7, where the evolution to the $`\sigma `$-meson pole is evident. Fitting ($`t=q^2`$)
$$\sigma (t)=g_{\sigma NN}\frac{mn_\sigma }{1t/m_\sigma ^2},t[0.1,0.5]\mathrm{GeV}^2,$$
(2.3.75)
which isolates the residue associated with $`\mathrm{\Gamma }_m(k;Q)`$, one obtains the on-shell coupling: $`g_{\sigma NN}=27.3`$. A direct calculation using the solution of the homogeneous Bethe-Salpeter equation yields $`g_{\sigma NN}=27.7`$, in agreement within Monte-Carlo errors. Equation (2.3.75) alone overestimates the magnitude of the calculated $`\sigma (t)`$ everywhere except in the neighbourhood of the pole.
As the lowest-mass pole-solution of a rainbow-ladder BSE, Eq. (2.3.65), the $`\sigma `$-meson described here is not obviously related to the scalar meson introduced in phenomenological nucleon-nucleon potentials and meson exchange models to mock-up two-pion exchange. However, a coupling relevant to such models can be estimated by introducing $`g_\sigma (t)`$:
$$\sigma (t)=:g_\sigma (t)\frac{mn_\sigma }{1t/m_\sigma ^2},$$
(2.3.76)
where a fit to the calculated $`\sigma (t)`$ in Fig. 2.7 yields
$$g_\sigma (t)=1.61+2.61\frac{1}{(1t/\mathrm{\Lambda }_\sigma ^2)^{10}},\mathrm{\Lambda }_\sigma =1.56\mathrm{GeV},$$
(2.3.77)
with the large exponent merely reflecting the rapid evolution from bound state to continuum dominance of the vertex in the spacelike region. $`g_\sigma (t)`$ describes the $`t`$-dependent nucleon to scalar-quark-antiquark-correlation coupling strength and at the mock-$`\sigma `$-mass: $`m_\sigma ^{2\pi }=0.5`$GeV,
$$g_\sigma :=g_\sigma ((m_\sigma ^{2\pi })^2)=9.3,$$
(2.3.78)
which may be compared with a phenomenologically inferred value: $`g_\sigma =10`$ . This gives a microscopic interpretation of the mock-$`\sigma `$. \[It is important to note that $`g_\sigma (4m_\pi ^2)=5.2`$ and hence this comparison is meaningful on a relevant phenomenological domain.\] Further, $`g_\sigma (q^2\mathrm{})=1.61`$ so that $`\sigma (q^2)`$ is well approximated by a single monopole for $`q^2>1`$GeV<sup>2</sup>. However, the residue is very different from the on-shell value. The scalar radius of the nucleon is obtained from
$$r_{\sigma NN}^2:=\frac{6}{\sigma }\frac{d\sigma (q^2)}{dq^2}|_{q^2=0}=(0.89\text{fm})^2,$$
(2.3.79)
which may be compared with an inferred value : $`(1.2\mathrm{fm})^2`$.
Reference calculates a range of non-electromagnetic nucleon form factors and provides a uniformly good description, yielding meson-nucleon form factors that are “soft” and couplings that generally agree well with those employed in meson-exchange models. Where there are discrepancies with these models or with experiment, a plausible cause and means for its amelioration was readily identified. The study demonstrates that it is realistic to hope for useful constraints on meson-exchange models from well-moulded models of hadron structure. The analysis of the nucleon $`\sigma `$-term is particularly important because it illustrates the only method that allows an unambiguous off-shell extrapolation in the estimation of meson-nucleon form factors: one must employ solutions of the inhomogeneous BSEs to describe correlations in a channel when the total momentum is not in the neighbourhood of a bound state pole. Improvements to the simple quark$`+`$scalar-diquark nucleon model are being developed and their application to a wide range of hadronic phenomena is being explored. These studies hold the promise of providing a much-needed standard model of the nucleon.
## 3 Nonzero Density and Temperature
Identifying the onset and properties of a quark gluon plasma \[QGP\] is the primary objective of studies at nonzero temperature and density. Such a plasma is expected to be characterised by light quarks and gluons propagating freely over distances $`10`$-times larger than the proton. The exploration of QCD in this domain requires a knowledge of equilibrium statistical field theory, for which Sec. 6 of Ref. is a primer and Refs. provide a graduate-level introduction. Here the application of DSEs provides our primary medium and we will describe the necessary and related concepts.
Section 2 reviews the application of DSEs at zero density and temperature. It shows that the framework is efficacious on this domain and well constrained. Both are important because little is known about the domain of high density and/or temperature, and much of what we do know comes from the extrapolation of models. Since the DSEs provide a nonperturbative framework that admits the simultaneous study of DCSB and confinement they are well suited to explore the phase transition that yields the QGP. Furthermore, as they also accurately describe bound states, they can be used to explore the response of hadron properties to extremes of density and temperature; i.e., to elucidate signals of QGP formation. These features also mean that the DSEs can address all the phenomena accessible in lattice-QCD simulations and therefore they can be used to confirm the results of such studies. Importantly too the lattice simulations can be used to constrain the model-dependent elements of DSE studies, making possible a reliable DSE extrapolation into the domain of nonzero chemical potential QCD, for example, which is inaccessible in contemporary lattice-QCD.
### 3.1 In-medium Essentials
Equation of State The $`T0`$ analogue of the partition function in Eq. (2.1.9) is
$$𝒵_{(\mu ,T)}[\overline{\eta },\eta ,J]=𝑑\mu _\beta (\overline{q},q,A,\overline{\omega },\omega )\mathrm{exp}_0^\beta 𝑑\tau d^3x\left[\overline{q}\eta +\overline{\eta }q+J_\mu ^aA_\mu ^a\right],$$
(3.1.1)
with $`x_4\tau `$, $`\beta :=1/T`$ and the measure
$`d\mu _\beta (\overline{q},q,A,\overline{\omega },\omega ):=`$
$`{\displaystyle \underset{\stackrel{}{x}\tau [0,\beta ]}{}}{\displaystyle \underset{\varphi }{}}𝒟\overline{q}_\varphi (\stackrel{}{x},\tau )𝒟q_\varphi (\stackrel{}{x},\tau ){\displaystyle \underset{a}{}}𝒟\overline{\omega }^a(\stackrel{}{x},\tau )𝒟\omega ^a(\stackrel{}{x},\tau ){\displaystyle \underset{\mu }{}}𝒟A_\mu ^a(\stackrel{}{x},\tau )\mathrm{exp}(S_\beta [\overline{q},q,A]S_\beta ^g[\overline{\omega },\omega ,A]),`$
where $`S_\beta [\overline{q},q,A]`$, $`S_\beta ^g[\overline{\omega },\omega ,A]`$ are the actions of Eqs. (2.1.8), (2.1.12) but for the obvious bounding of the $`\tau `$-integral. The frame specified by separating the spacetime integral into a product of a compact integral over Euclidean time and the volume integral is not prescribed by the theory. This separation breaks the theory’s original $`O(4)`$ symmetry to $`O(3)`$ and is effected by introducing an arbitrary, normalised, spacelike vector, $`u_\mu `$, with $`up=0`$ defining the $`T=0`$-hyperplane. Here we have employed a conventional choice for the heat bath vector: $`u=(\stackrel{}{0},1)`$. \[NB. The zero temperature theory is recovered by taking the limit $`T0`$; i.e., $`\beta \mathrm{}`$, in Eqs. (3.1.1), (3.1).\] The chemical potential is introduced as a Lagrange multiplier
$$S_\beta [\overline{q},q,A]S_{\mu ,\beta }[\overline{q},q,A]:=S_\beta [\overline{q},q,A]_0^\beta 𝑑\tau d^3x\mu \overline{q}\gamma _4q.$$
(3.1.3)
With $`\mu 0`$, $`S_{\mu ,\beta }[\overline{q},q,A]`$ is not Hermitian and hence $`d\mu _{\mu ,\beta }`$ is not a probability measure. That makes numerical simulations of lattice-regularised $`\mu 0`$ QCD very difficult; in fact, there is currently no satisfactory algorithm . The functional integral in Eq. (3.1.1) is evaluated on the space of Grassmann fields for which $`q(\stackrel{}{x},0)=q(\stackrel{}{x},\beta )`$; i.e., over fields satisfying antiperiodic boundary conditions, and over gauge fields satisfying periodic boundary conditions. These boundary conditions are an essential aspect of the difference between fermion and boson statistics.
The partition function is particularly important in equilibrium statistical field theory. As usual, it yields the DSEs in a straightforward manner, however, more than that, all the thermodynamic functions are obtained from the partition function. The thermodynamic potential and pressure densities are
$$\omega (\mu ,T)=p(\mu ,T)=\frac{1}{\beta V}\mathrm{ln}𝒵_{(\mu ,T)},$$
(3.1.4)
where $`\beta V`$ is the four-volume normalising factor. The stable phase of the system is that in which the pressure is maximal or equivalently the thermodynamic potential energy is minimised. The expression for the pressure is the Equation of State \[EOS\] and from that one immediately obtains the baryon number and entropy densities:
$$\rho (\mu ,T)=\frac{}{\mu }p(\mu ,T)|_{T\mathrm{fixed}},s(\mu ,T)=\frac{}{T}p(\mu ,T)|_{\mu \mathrm{fixed}}$$
(3.1.5)
and also the energy density:
$$ϵ(\mu ,T)=p(\mu ,T)+\mu \rho (\mu ,T)+Ts(\mu ,T).$$
(3.1.6)
The calculation of these quantities in QCD is a contemporary focus; e.g., Refs. .
A simple estimate of the pressure due to dressed-quarks is obtained via the “steepest descent” approximation, which yields
$$p_\mathrm{\Sigma }(\mu ,T)=\frac{1}{\beta V}\left\{\mathrm{TrLn}\left[\beta S^1\right]\frac{1}{2}\mathrm{Tr}\left[\mathrm{\Sigma }S\right]\right\},$$
(3.1.7)
where $`S`$ is the solution of the gap equation with $`\mathrm{\Sigma }`$ the associated self energy, and “Tr” and “Ln” are extensions of “tr” and “$`\mathrm{ln}`$” to matrix-valued functions. Equation (3.1.7) is just the auxiliary field effective action , which yields the free fermion pressure in the absence of interactions; i.e., when $`\mathrm{\Sigma }0`$. At this level of truncation the total pressure receives an additive contribution from dressed-gluons:
$$p_\mathrm{\Delta }(\mu ,T)=\frac{1}{\beta V}\frac{1}{2}\mathrm{TrLn}\left[\beta ^2D_{\mu \nu }^1\right],$$
(3.1.8)
where $`D_{\mu \nu }`$ is the $`T0`$ dressed-gluon $`2`$-point function, and this yields the free-gluon pressure in the absence of interactions.
Propagators As just described, the introduction of temperature to Euclidean QCD breaks the original $`O(4)`$ symmetry to $`O(3)`$. In this case the most general form of the dressed-quark propagator is
$$S(p,u)^1=i\gamma pa(p^2,[up]^2)+b(p^2,[up]^2)+i\gamma uupc(p^2,[up]^2)+i\sigma _{\mu \nu }p_\mu u_\nu d(p^2,[up]^2),$$
(3.1.9)
where $`u_\mu =(\stackrel{}{0},1)`$, and the surviving $`O(3)`$ invariance of Eq. (3.1.9) is apparent upon inspection. Even at $`T0`$, the chiral-limit dressed-quark propagator should satisfy
$$V(\alpha )S(p,u)^1V(\alpha )=S(p,u)^1,V(\alpha )=\mathrm{exp}\left(i\gamma _5\frac{1}{2}\lambda _F^l\alpha ^l\right),$$
(3.1.10)
if chiral symmetry is realised in the Wigner-Weyl mode. Hence $`b0d`$ in this case. Furthermore, since the term involving $`d`$ does not appear in the free fermion action of thermal field theory, there are no perturbative divergences in this function’s evaluation and $`d`$ is therefore power-law suppressed in the ultraviolet. \[cf. The discussion of the chiral limit fermion mass function on page 2.2.13.\] These features, along with the vector exchange nature of the interaction, underly the fact that $`d(p^2,up^2)`$ only plays a minor role in the solution of the QCD gap equation. Hereafter we neglect it and assume the following as the most general form of the dressed-quark $`2`$-point function:
$`S(\stackrel{}{p},\stackrel{~}{\omega }_k)`$ $`=`$ $`{\displaystyle \frac{1}{i\stackrel{}{\gamma }\stackrel{}{p}A(\stackrel{}{p}^2,\stackrel{~}{\omega }_k^2)+B(\stackrel{}{p}^2,\stackrel{~}{\omega }_k^2)+i\gamma _4\stackrel{~}{\omega }_kC(\stackrel{}{p}^2,\stackrel{~}{\omega }_k^2)}},`$ (3.1.11)
$`=`$ $`i\stackrel{}{\gamma }\stackrel{}{p}\sigma _A(\stackrel{}{p}^2,\stackrel{~}{\omega }_k^2)+\sigma _B(\stackrel{}{p}^2,\stackrel{~}{\omega }_k^2)i\gamma _4\stackrel{~}{\omega }_k\sigma _C(\stackrel{}{p}^2,\stackrel{~}{\omega }_k^2),`$ (3.1.12)
where $`\stackrel{~}{\omega }_k=\omega _k+i\mu `$, with $`\mu `$ the chemical potential, and $`\omega _k=(2k+1)\pi T`$, $`kZZ`$, are the fermion Matsubara frequencies, which ensure $`q(\stackrel{}{x},0)=q(\stackrel{}{x},\beta )`$. The scalar functions: $`=A`$, $`B`$, $`C`$, are complex and satisfy
$$(\stackrel{}{p}^2,\stackrel{~}{\omega }_k^2)^{}=(\stackrel{}{p}^2,\stackrel{~}{\omega }_{k1}^2).$$
(3.1.13)
\[NB. The functions $`A`$, $`C`$ cannot be neglected. For $`T=0`$, $`A=C`$ and $`A10`$; e.g., Refs. , and the momentum dependence of $`A`$ can conspire with that of $`B`$ to ensure confinement of dressed-quarks . Furthermore, as will become clear, the $`(\stackrel{}{p}^2,\stackrel{~}{\omega }_k^2)`$-dependence of $`A`$, $`C`$ has a marked effect on the behaviour of bulk thermodynamic quantities.\]
The dressed-quark $`2`$-point function satisfies a gap equation
$$S(p_{\omega _k})^1=Z_2^Ai\stackrel{}{\gamma }\stackrel{}{p}+Z_2(i\gamma _4\stackrel{~}{\omega }_k+m_{\mathrm{bare}})+\mathrm{\Sigma }^{}(p_{\omega _k}),$$
(3.1.14)
which can be derived in the usual way from the partition function, Eq. (3.1.1). Here we have introduced a shorthand notation: $`p_{\omega _k}=(\stackrel{}{p},\stackrel{~}{\omega }_k)`$, and the regularised self energy is
$`\mathrm{\Sigma }^{}(p_{\omega _k})`$ $`=`$ $`i\stackrel{}{\gamma }\stackrel{}{p}\mathrm{\Sigma }_A^{}(p_{\omega _k})+i\gamma _4\omega _k\mathrm{\Sigma }_C^{}(p_{\omega _k})+\mathrm{\Sigma }_B^{}(p_{\omega _k}),`$ (3.1.15)
$`\mathrm{\Sigma }_{}^{}(p_{\omega _k})`$ $`=`$ $`\frac{1}{4}\mathrm{tr}𝒫_{}{\displaystyle _{l,q}^{\overline{\mathrm{\Lambda }}}}\frac{4}{3}g^2D_{\mu \nu }(\stackrel{}{p}\stackrel{}{q},\omega _k\omega _l)\gamma _\mu S(q_{\omega _l})\mathrm{\Gamma }_\nu (q_{\omega _l};p_{\omega _k}),`$ (3.1.16)
where the renormalisation factors are $`𝒫_A:=(Z_1^A/|\stackrel{}{p}|^2)i\stackrel{}{\gamma }\stackrel{}{p}`$, $`𝒫_B:=Z_1`$, $`𝒫_C:=(Z_1/\omega _k)i\gamma _4`$, and $`_{l,q}^{\overline{\mathrm{\Lambda }}}:=T_{l=\mathrm{}}^{\mathrm{}}^{\overline{\mathrm{\Lambda }}}d^3q/(2\pi )^3`$, with $`^{\overline{\mathrm{\Lambda }}}`$ representing a translationally invariant regularisation of the three-dimensional integral and $`\overline{\mathrm{\Lambda }}`$ is the regularisation mass-scale. Since introducing $`\mu `$, $`T`$ does not generate qualitatively new divergences, the regularisation and renormalisation proceeds just as in the absence of the medium so that the renormalised self energies are
$$\begin{array}{ccc}\hfill (p_{\omega _k};\zeta )& =& \xi _{}+\mathrm{\Sigma }_{}^{}(p_{\omega _k};\overline{\mathrm{\Lambda }})\mathrm{\Sigma }_{}^{}(\zeta _{\omega _0}^{};\overline{\mathrm{\Lambda }}),\hfill \end{array}$$
(3.1.17)
with $`\zeta `$ the renormalisation point, $`(\zeta _{\omega _0}^{})^2:=\zeta ^2\omega _0^2`$, $`\xi _A=1=\xi _C`$, and $`\xi _B=m_R(\zeta )`$. \[cf. Eqs. (2.2.39-2.2.43).\]
The regularised self energy is expressed in terms of the renormalised dressed-gluon $`2`$-point function and dressed-quark-gluon $`3`$-point function. At $`T0`$ in Landau gauge the complete expression of the former requires two $`O(3)`$-scalar functions \[cf. Eq. (2.2.15)\]
$$g^2D_{\mu \nu }(p_{\mathrm{\Omega }_k})=P_{\mu \nu }^L(p_{\mathrm{\Omega }_k})\mathrm{\Delta }_F(p_{\mathrm{\Omega }_k})+P_{\mu \nu }^T(p_{\mathrm{\Omega }_k})\mathrm{\Delta }_G(p_{\mathrm{\Omega }_k}),$$
(3.1.18)
where $`p_{\mathrm{\Omega }_k}:=(\stackrel{}{p},\mathrm{\Omega }_k)`$, $`\mathrm{\Omega }_k=2k\pi T`$ is the Matsubara frequency for bosons, and
$`P_{\mu \nu }^T(p_{\mathrm{\Omega }_k})`$ $`:=`$ $`\{\begin{array}{cc}0,\hfill & \mu \mathrm{and}/\mathrm{or}\nu =4,\hfill \\ \delta _{ij}{\displaystyle \frac{p_ip_j}{p^2}},\hfill & \mu ,\nu =i,j=1,2,3,\hfill \end{array}`$ (3.1.21)
$`P_{\mu \nu }^L(p_{\mathrm{\Omega }_k})`$ $`=`$ $`\delta _{\mu \nu }{\displaystyle \frac{p_\mu p_\nu }{p^2}}P_{\mu \nu }^T(p_{\mathrm{\Omega }_k}).`$ (3.1.22)
In the absence of interactions: $`\mathrm{\Delta }_F(p_{\mathrm{\Omega }_k}^2)=\mathrm{\Delta }_G(p_{\mathrm{\Omega }_k})=1/p_{\mathrm{\Omega }_k}^2`$. The QCD analogue of a Debye \[or electric screening\] mass appears as a $`T`$-dependent contribution to $`\mathrm{\Delta }_F`$, the gluon’s longitudinal polarisation function. The complete expression of the renormalised dressed-quark-gluon $`3`$-point function requires at least fifty-four scalar functions \[cf. twelve for $`T=0`$\] but, since the theory is renormalisable, ultraviolet divergences are only encountered in evaluating two of them. It follows from this that sensible, character-preserving truncations are possible.
Order Parameters In order to demarcate the QGP phase it is necessary to identify order parameters; i.e., find operators $`X_i`$ whose expectation values are nonzero in the normal phase: $`X_i0`$, but vanish in the QGP: $`X_i0`$. This is often a difficult task. Having identified such operators then, if $`X_i`$ is discontinuous at the transition point we will describe the transition as first order; otherwise we classify it as second order.
In Sec. 2.2 we described the phenomenon of DCSB and its connection with the purely nonperturbative generation of a scalar term in the dressed-quark self energy when the current-quark mass is zero. One characteristic of this effect is the large magnitude, Eq. (2.2.51), of the vacuum quark condensate that appears in Eqs. (2.2.47), (2.2.48), which has many observable consequences; e.g., it is the primary cause of what is naively an unexpectedly large $`\pi `$-$`\rho `$ mass splitting. DCSB is a defining feature of the normal phase of QCD. If, with increasing $`\mu `$ and/or $`T`$, the condensate vanishes at some particular value of these parameters then a new phase of QCD has been reached. It is a phase in which chiral symmetry is realised explicitly and, as weakly interacting systems cannot support condensate formation, this is plainly a good defining property of a QGP. The discussion of the chiral limit, Eqs. (2.2.45-2.2.52), also makes plain that it is not possible to have a nonzero quark condensate unless the dressed-quark mass function, $`M(p_{\omega _k})`$, is nonzero too. Therefore the magnitude of this function at any single point is also an order parameter appropriate for characterising the chiral symmetry restoring aspect of the QGP transition. In fact, in DSE studies it is by far the most direct, since it appears in all the integrands that describe other quantities.
We also discussed an arguably more significant aspect of cold, sparse QCD, namely confinement. It is easy to build a model that exhibits DCSB without confinement; e.g., Refs. . However, building a covariant model with confinement and an uncomplicated realisation of all aspects of DCSB is a challenging task. As emphasised by the study of QED<sub>3</sub> in Ref. , that challenge is met with a modicum of success using the DSE truncation scheme explored in Ref. and the realisation of confinement via a violation of reflection positivity by coloured $`n`$-point functions. This also exposes a simple deconfinement order parameter appropriate to both light- and heavy-quarks, whose application is illustrated in Fig. 2.3. To make it plain we return to $`D(T)`$ defined in Eq. (2.2.72). When calculated from a Schwinger function with complex conjugate poles, $`D(T)`$ has at least one zero. Denote its position by $`T_0^{z_1}`$ and define
$$\kappa _0:=1/T_0^{z_1}.$$
(3.1.23)
If, with increasing $`\mu `$ and/or $`T`$, one observes $`\kappa _0`$ approaching zero then the \[first\] zero in $`D(T)`$ is moving to $`T=\mathrm{}`$. This evolution of $`\kappa _0`$ corresponds to the limiting case when there is no zero at all, which is just the free, unconfined particle described by $`\sigma _S(T)`$ in Eqs. (2.2.67). Thus $`\kappa _0`$ is a good order parameter for deconfinement : it is nonzero in the confined phase, and remains nonzero until the Schwinger function evolves into one that respects reflection positivity and corresponds to an asymptotic state. It is clear that deconfinement is also a good defining property of a QGP.
Critical Behaviour In second order transitions the length scale associated with correlations in the system diverges as the order parameter vanishes and a range of critical exponents can be defined that characterise the behaviour of macroscopic properties at the transition point. In the context of QCD this can be elucidated via the free energy, which here we write as $`f(t,h)`$, where $`t:=T/T_c1`$ is the reduced temperature, with $`T_c`$ a putative critical temperature, and $`h:=\beta m_R`$ is the source of explicit chiral symmetry breaking measured in units of the temperature. \[We omit the chemical potential for now so that $`f=ϵ=(T^2/V)(\mathrm{ln}𝒵_{(T)}/T)`$.\] $`h`$ plays a role analogous to that of an external magnetic field applied to a ferromagnet. Since correlation lengths diverge it follows that for $`t,h0`$ the free energy is a generalised homogeneous function; i.e.,
$$f(t,h)=\frac{1}{b}f(tb^{y_t},hb^{y_h}).$$
(3.1.24)
As a consequence the “magnetisation” behaves as follows:
$$M(t,h):=\frac{f(t,h)}{h}|_{t\mathrm{fixed}},M(t,h)=b^{y_h1}M(tb^{y_t},hb^{y_h}).$$
(3.1.25)
\[NB. Using the partition function one finds easily that $`M(t,h)`$ is just the vacuum quark condensate. In these formulae it can be replaced by any equivalent order parameter.\] The scaling parameter: $`b`$, is arbitrary and along the trajectory $`|t|b^{y_t}=1`$ one has
$$M(t,h)=|t|^{(1y_h)/y_t}M(\mathrm{sgn}(t),h|t|^{y_h/y_t});\mathrm{i}.\mathrm{e}.,M(t,0)|t|^\beta ,\beta :=\frac{1y_h}{y_t}.$$
(3.1.26)
Alternatively, along the trajectory $`hb^{y_h}=1`$
$$M(t,h)=h^{(1y_h)/y_h}M(th^{y_t/y_h},1);\mathrm{i}.\mathrm{e}.,M(0,h)h^{1/\delta },\delta :=\frac{y_h}{1y_h}.$$
(3.1.27)
Equations (3.1.26), (3.1.27) quantify the behaviour to be expected of an order parameter at a second order transition. They also introduce two critical exponents and, using the renormalisation group, scaling laws can be derived that relate all other such exponents to $`\beta `$ and $`\delta `$ . It is widely conjectured that the values of these exponents are fully determined by the dimension of space and the nature of the order parameter. This is the notion of universality ; i.e., that the critical exponents are independent of a theory’s microscopic details and hence all theories can be grouped into a much smaller number of universality classes according to the values of their critical exponents. If this is the case, the behaviour of a complicated theory near criticality is completely determined by the behaviour of a simpler theory in the same universality class. In mean-field theories
$$\beta ^{\mathrm{MF}}=\frac{1}{2},\delta ^{\mathrm{MF}}=3.0.$$
(3.1.28)
The success of the nonlinear $`\sigma `$-model in describing long-wavelength pion dynamics underlies a conjecture that chiral symmetry restoration at $`T0`$ in 2-flavour QCD is a second order transition with the theory lying in the universality class characterised by the 3-dimensional, $`N=4`$ Heisenberg magnet \[$`O(4)`$ model\] whose critical exponents are :
$$\beta ^H=0.38\pm 0.01,\delta ^H=4.82\pm 0.05.$$
(3.1.29)
As an alternative to Eqs. (3.1.26), (3.1.27), the critical exponents can be determined by studying the pseudocritical behaviour of the chiral and thermal susceptibilities:
$$\chi _h(t,h):=\frac{M(t,h)}{h}|_{t\mathrm{fixed}},\chi _t(t,h):=\frac{M(t,h)}{t}|_{h\mathrm{fixed}},$$
(3.1.30)
which are smooth functions of $`t`$ with maxima at the pseudocritical points: $`t_{\mathrm{pc}}^h`$ and $`t_{\mathrm{pc}}^t`$, where $`t_{\mathrm{pc}}^hh^{1/(\beta \delta )}t_{\mathrm{pc}}^t`$. Since $`\beta \delta >0`$, the pseudocritical points approach the critical point: $`t=0`$, as $`h0^+`$ and
$$\chi _h^{\mathrm{pc}}:=\chi _h(t_{\mathrm{pc}}^h,h)h^{z_h},z_h:=1\frac{1}{\delta },\chi _t^{\mathrm{pc}}:=\chi _t(t_{\mathrm{pc}}^t,h)h^{z_t},z_t:=\frac{1}{\beta \delta }(1\beta ).$$
(3.1.31)
\[The appendix of Ref. provides a demonstration.\] Mean field behaviour corresponds to
$$z_h^{\mathrm{MF}}=\frac{2}{3},z_t^{\mathrm{MF}}=\frac{1}{3}.$$
(3.1.32)
Analysing the susceptibilities can provide more accurate results when numerical noise is anticipated, which is why the method has been used in lattice-QCD analyses .
### 3.2 Quark Gluon Plasma Phase
Temperature We are now in a position to describe the application of DSEs to demarcating the QGP and exploring its properties. For the moment we persist with $`\mu =0`$ and explore the temperature induced transition because that is the domain on which the most complete studies exist \[in all approaches\].
The properties of the class of rainbow-ladder models are elucidated in Ref. . This class is defined by Eq. (2.2.33) with the mutually consistent $`T0`$ constraints: $`Z_1=Z_2`$, $`Z_1^A=Z_2^A`$. \[NB. The rainbow truncation is the leading term in a $`1/N_c`$-expansion of $`\mathrm{\Gamma }_\nu (q_{\omega _l};p_{\omega _k})`$.\] The form of $`D_{\mu \nu }(p_{\mathrm{\Omega }_k})`$ has no bearing on whether or not a model is in the class. To be concrete, Ref. considered models defined by
$`\mathrm{\Delta }_F(p_{\mathrm{\Omega }_k})`$ $`=`$ $`𝒟(p_{\mathrm{\Omega }_k};m_g),\mathrm{\Delta }_G(p_{\mathrm{\Omega }_k})=𝒟(p_{\mathrm{\Omega }_k};0),`$ (3.2.1)
$`𝒟(p_{\mathrm{\Omega }_k};m_g)`$ $`:=`$ $`2\pi ^2D\frac{2\pi }{T}\delta _{0k}\delta ^3(\stackrel{}{p})+𝒟_\mathrm{M}(p_{\mathrm{\Omega }_k};m_g),`$ (3.2.2)
with three choices of $`𝒟_\mathrm{M}(p_{\mathrm{\Omega }_k};m_g)`$. The first term in $`𝒟(p_{\mathrm{\Omega }_k};m_g)`$ is a simple $`T0`$ generalisation of the distribution in Eq. (2.2.35). It represents the long-range piece of the effective interaction and the $`𝒟_\mathrm{M}(p_{\mathrm{\Omega }_k};m_g)`$ are chosen so as not to pollute that. (NB. $`\mathrm{\Delta }_G`$ does not contain a simple $`T`$-dependent mass. Such “magnetic” screening is more complicated than electric screening because of infrared divergences in QCD .) The three choices are:
A) $`𝒟_{\mathrm{M}:=\mathrm{A}}(p_{\mathrm{\Omega }_k};m_g)0`$, with $`D:=\eta ^2/2`$ and $`\eta =1.06`$ fixed by fitting $`\pi `$\- and $`\rho `$-meson masses at $`T=0`$. $`m=12`$MeV yields $`m_\pi =140`$MeV. $`𝒟_\mathrm{A}`$ yields an ultraviolet finite model and hence the renormalisation point and cutoff can be removed simultaneously. Nevertheless it exhibits many features in common with more sophisticated Ansätze .
B) The model of Ref. , obtained with $`D:=(8/9)m_t^2`$ and
$$𝒟_{\mathrm{M}:=\mathrm{B}}(p_{\mathrm{\Omega }_k};m_g)=\frac{16}{9}\pi ^2\frac{1\mathrm{e}^{s_{\mathrm{\Omega }_k}/(4m_t^2)}}{s_{\mathrm{\Omega }_k}},$$
(3.2.3)
$`s_{\mathrm{\Omega }_k}:=p_{\mathrm{\Omega }_k}^2+m_g^2`$, where, following Ref. , $`m_g^2=4\pi ^2\gamma _m(N_c/3+N_f/6)=(8/3)\pi ^2T^2`$ \[$`N_f=3`$\]. The mass-scale $`m_t=0.69\mathrm{GeV}=1/0.29\mathrm{fm}`$ marks the boundary between the perturbative and nonperturbative domains, and was also fixed by fitting $`\pi `$\- and $`\rho `$-meson properties at $`T=0`$. At $`\zeta =9.5`$GeV, $`m_R=1.1`$MeV yields $`m_\pi =140`$MeV. This model adds a Coulomb-like short-range interaction to that of Ref. , thereby marginally improving its ultraviolet behaviour.
C) A minimal $`T0`$ extension of the effective interaction in Eq. (2.2.35) with the only modification being $`k^2s_{\mathrm{\Omega }_k}`$, $`m_g^2=(16/5)\pi ^2T^2`$, in the last two terms. \[Recall that this model’s parameters were fitted with $`N_f=4`$\]. As described above, $`𝒟_\mathrm{C}`$ ensures that the model preserves the one-loop renormalisation group behaviour of QCD. To further explore model-dependence, two parameter sets were employed in Ref. , which differ only in the value of $`\omega `$: $`{}_{1}{}^{}\omega =0.6m_t`$, $`{}_{2}{}^{}\omega =1.2m_t`$. They are equivalent, yielding the same values for pion observables: with $`{}_{1}{}^{}\omega `$, a renormalisation point invariant current-quark mass of $`\widehat{m}=6.6`$MeV yields $`m_\pi =140`$MeV; while with $`{}_{2}{}^{}\omega `$, $`\widehat{m}=5.7`$MeV effects that.
Two chiral order parameters were employed in analysing chiral symmetry restoration
$$𝒳:=B(\stackrel{}{p}=0,\omega _0),𝒳_C:=\frac{B(\stackrel{}{p}=0,\omega _0)}{C(\stackrel{}{p}=0,\omega _0)}.$$
(3.2.4)
They should be equivalent and the onset of that equivalence determines the $`h`$-domain on which Eqs. (3.1.31) are valid. These being bona fide order parameters is useful and particularly important because it means that the lowest Matsubara frequency completely determines the character of the chiral phase transition, as conjectured in Ref. . Furthermore, it was numerically verified in Ref. that in the chiral limit and for $`t0`$: $`f_\pi \overline{q}q𝒳(t,0)`$; i.e., these quantities are all bona fide order parameters.
The results for the critical temperature and exponents are presented in Table 3.1, where the quoted critical temperatures are easily determined using either Eqs. (3.1.26), (3.1.27) or $`t_{\mathrm{pc}}^hh^{1/(\beta \delta )}t_{\mathrm{pc}}^t`$. The models are all constrained by observables in two-flavour QCD and the critical temperature agrees with the current estimate from simulations of lattice-QCD with two dynamical light-quark flavours . (Quenched lattice-QCD exhibits a first order \[deconfining\] transition at a much higher critical temperature: $`T_c^{\mathrm{quench}}=270(5)`$MeV .) Studies using simple versions of the Nambu–Jona-Lasinio model typically yield higher values: $`T_c>200`$MeV . However, that is likely a model artefact. Modifications that make possible a relaxation of the overly restrictive connection between the model’s mass-scale and the dressed-quark mass; such as $`1/N_c`$-\[or meson-loop-\]corrections or replacing the contact interaction by a separable form, can easily effect a reduction of the critical temperature by $`<25`$%; e.g., Refs. . This sort of reduction brings it in-line with the estimate in Table 3.1.
The behaviour of the pseudocritical maxima in the chiral and thermal susceptibilities is illustrated in Fig. 3.1. It is plain that the results exhibit curvature and hence the scaling relations of Eq. (3.1.31) are not satisfied on the entire domain. The domain on which they are valid may be established by defining a “local” critical exponent:
$$z_i:=\frac{\mathrm{ln}\chi _i^{\mathrm{pc}}\mathrm{ln}\chi _{i+1}^{\mathrm{pc}}}{\mathrm{ln}h_i\mathrm{ln}h_{i+1}},$$
(3.2.5)
where $`(h_i,\chi _i^{\mathrm{pc}})`$ and $`(h_{i+1},\chi _{i+1}^{\mathrm{pc}})`$ are adjacent data pairs. $`h`$ lies in the scaling region when $`z_i`$ is independent of the order parameter. Applying this to the $`𝒟_\mathrm{B}`$ model yields the results in Fig. 3.2, which indicate clearly that the scaling relations are not valid until
$$\mathrm{log}_{10}(h/h_u)<7,$$
(3.2.6)
where $`h_u:=m_R/T_c`$ is defined with the current-quark mass that gives $`m_\pi =140`$MeV in this model. Applying the same method to the renormalisation-group improved models, $`𝒟_\mathrm{C}`$, shows that in these cases the scaling relations are only valid for
$`{}_{1}{}^{}\omega :\mathrm{log}_{10}(h/h_u)<5,`$ $`{}_{2}{}^{}\omega :\mathrm{log}_{10}(h/h_u)<6;`$ (3.2.7)
i.e., for current-quark masses six orders of magnitude smaller than those of real $`u`$,$`d`$-quarks.
While these are model studies the result is also likely to be true in QCD; i.e., while the critical temperature is relatively easy to determine, very small current-quark masses may be necessary to accurately calculate the critical exponents from the chiral and thermal susceptibilities. If that is the case, the calculation of these exponents via numerical simulations of lattice-QCD will not be feasible. The lattice-to-lattice variation of the critical exponents described in Refs. , with even a first order transition being possible, could be a signal of this.
It is clear from Table 3.1 that each of the models considered in Ref. is mean field in nature. As the long-range part of the effective interaction is identical in each case and the correlation length diverges as $`t0`$, that may not be surprising. The models are distinguished by the feature that they describe the long-wavelength dynamics of QCD very well, with mesons represented as nonpointlike quark-antiquark composites. Coulomb gauge models also describe mesons as composites and exhibit mean field critical exponents. The long-range part of the interaction there corresponds directly to the regularised Fourier amplitude of a linearly rising potential and therefore it is in no simple way equivalent to the distribution in Eq. (2.2.35). Separable models with : $`𝒟(p_{\omega _k}q_{\omega _l};m_g)g(|\stackrel{}{p}|)g(|\stackrel{}{q}|)`$, where $`g(|\stackrel{}{p}|)`$ is a non-increasing function of its argument, can also provide a good description of long-wavelength pion dynamics. They too describe mesons as composites and exhibit an explicit fermion substructure, and they also have mean field critical exponents, as do models with $`g=g(\omega _k^2+\stackrel{}{p}^2)`$ . \[NB. The Nambu–Jona-Lasinio model can be viewed as a separable model obtained with $`g(|\stackrel{}{p}|)=\theta (1|\stackrel{}{p}|/_3\mathrm{\Lambda })`$, where $`{}_{3}{}^{}\mathrm{\Lambda }`$ is a cutoff parameter. Also, in the rainbow-truncation DSE model of Ref. , which employs an infrared-finite effective interaction, the critical exponent obtained from $`𝒳`$ is $`1/2`$. However, a different exponent is reported to be obtained from $`\overline{q}q`$. We expect that is erroneous, as these must be equivalent order parameters.\]
All these examples are members of the class of rainbow truncation models. Therefore the results indicate that chiral symmetry restoration in this class is a mean field transition. This is an essential consequence of the fermion substructure in the rainbow gap equation: non-trivial critical exponents require the appearance of an infrared divergence in this equation. However, in rainbow truncation the self energy is infrared-regulated by the zeroth fermion Matsubara frequency: $`\omega _0=\pi T0`$ in the vicinity of the transition. For the same reason chiral symmetry restoration at finite-$`T`$ in QCD will be a mean field transition unless $`1/N_c`$-corrections to the vertex are large for $`TT_c`$. There are, however, examples in which that is certainly true . Infrared divergences might be anticipated if bosonic modes dominate the gap equation near $`T_c`$ because the zeroth boson Matsubara frequency vanishes. It is therefore important to note that pseudoscalar meson-like fluctuations are nonperturbative $`1/N_c`$-corrections to the dressed-vertex. They do not contribute to the dressed-gluon $`2`$-point function. Therefore a true determination of the critical exponents must await the accurate incorporation of these $`1/N_c`$-corrections.
We have seen that chiral symmetry restoration in $`2`$-flavour QCD occurs at $`T150`$MeV via a second order transition whose critical exponents are not yet reliably determined. What of deconfinement? In QCD with two dynamical flavours the Wilson-line/Polyakov-loop is not strictly a good order parameter for deconfinement because of flux tube breaking mediated by quark-antiquark pair formation. However, in simulations the Polyakov loop’s susceptibility develops a peak at the same temperature as the chiral susceptibility and this has been interpreted as evidence of coincident deconfinement and chiral symmetry restoration .
$`\kappa _0`$ in Eq. (3.1.23) provides an order parameter that can be used to explore deconfinement via the gap equation. It was exploited in Ref. where the deconfinement transition in model “B)” was studied. In this case the Schwinger function analysed is
$$\sigma _S(r):=\frac{2}{\pi }_0^{\mathrm{}}𝑑pp\mathrm{sin}(pr)\sigma _B(p^2,\omega _0^2),$$
(3.2.8)
where $`\sigma _B`$ is calculated in the chiral limit. Figure 3.3 depicts the behaviour of $`\kappa _0(T)`$ obtained from $`\sigma _S(r)`$ and clearly the deconfinement order parameter is nonzero in the normal phase. $`\kappa _0(T)`$ vanishes at the same temperature as does $`𝒳`$. Hence, in this model, the deconfinement and chiral symmetry restoring transitions are coincident \[to better-than $`1`$%\] and analysis that parallels Ref. shows that deconfinement is also a mean field transition. (This improves over the critical exponents estimated in Ref. .) Whether the transitions are always coincident is unknown, however, that is not unlikely given the role the scalar piece of the dressed-quark self energy plays in determining both transitions.
The rainbow-ladder class of models fails to exhibit any flavour dependence in the transition to a QGP, a flaw it shares; e.g., with random matrix models . It is anticipated that with three light quark flavours the chiral symmetry restoring transition is first order but becomes second order if the mass of the “$`s`$-quark” is raised above a presently-undetermined critical value. This can be expected because for an infinitely heavy $`s`$-quark one recovers the two-flavour theory with its second order transition. (A summary of lattice results is presented in Ref. .) A coupling between flavours is necessary to achieve such a pattern of symmetry breaking and, as seen above, that appears only when $`1/N_c`$-corrections to the vertex are incorporated.
Chemical Potential The introduction of a chemical potential moves us to a domain in which the only available nonperturbative information comes from models. For $`\mu 0`$ the dressed-quark self energies in general acquire an imaginary part driven by the magnitude of $`\mu `$; e.g., Refs. . This effect is not observed in explorations of the Nambu–Jona-Lasinio model or indeed in any model in which the interaction is energy-independent; i.e., instantaneous. That follows from Eq. (3.1.13): an instantaneous interaction cannot yield an energy-dependent solution and hence $`M=M^{}M\text{ }\mathrm{R}`$. It is an unrealistic limitation since the interaction in the normal \[confined\] phase of QCD is not instantaneous.
An extensive exploration of the phase structure and thermodynamics of two-flavour QCD with $`\mu 0T`$ is presented Ref. , which employs the $`𝒟_\mathrm{A}`$ model. $`𝒟_\mathrm{A}`$ is an infrared-dominant model that does not represent well the behaviour of $`D_{\mu \nu }(\stackrel{}{p},\mathrm{\Omega }_k)`$ away from $`|\stackrel{}{p}|^2+\mathrm{\Omega }_k^20`$ and that introduces model-dependent artefacts. However, there is significant merit in its algebraic simplicity and, since the artefacts are easily identified, the model remains useful as a means of elucidating many of the qualitative features of more sophisticated Ansätze. Therefore we use it here as an exemplar.
In this model the gap equation is
$$S(p_{\omega _k})^1=S_0(p_{\omega _k})^1+\frac{1}{4}\eta ^2\gamma _\nu S(p_{\omega _k})\gamma _\nu $$
(3.2.9)
and the simplicity is now apparent: the model allows the reduction of an integral equation to an algebraic equation, which is always a useful step when the goal is to develop an intuitive understanding of complicated phenomena. Substituting Eq. (3.1.11) yields
$`\eta ^2m^2`$ $`=`$ $`B^4+mB^3+\left(4p_{\omega _k}^2\eta ^2m^2\right)B^2m\left(2\eta ^2+m^2+4p_{\omega _k}^2\right)B,`$ (3.2.10)
$`A(p_{\omega _k})`$ $`=`$ $`C(p_{\omega _k})={\displaystyle \frac{2B(p_{\omega _k})}{m+B(p_{\omega _k})}},`$ (3.2.11)
and it is now apparent that neglecting $`A`$, $`C`$, as in Refs. , is a poor approximation in the presence of DCSB; i.e., on the domain where $`B(p_{\omega _k})m`$, which must at least introduce quantitative errors.
In the chiral limit: $`m=0`$, Eq. (3.2.10) reduces to a quadratic equation for $`B(p_{\omega _k})`$, which has two qualitatively distinct solutions. The “Nambu-Goldstone” solution, for which
$`B(p_{\omega _k})`$ $`=`$ $`\{\begin{array}{ccc}\sqrt{\eta ^24p_{\omega _k}^2},\hfill & & \text{ }\mathrm{R}(p_{\omega _k}^2)<\frac{\eta ^2}{4}\hfill \\ 0,\hfill & & \mathrm{otherwise}\hfill \end{array}`$ (3.2.14)
$`C(p_{\omega _k})`$ $`=`$ $`\{\begin{array}{ccc}2,\hfill & & \text{ }\mathrm{R}(p_{\omega _k}^2)<\frac{\eta ^2}{4}\hfill \\ \frac{1}{2}\left(1+\sqrt{1+2\eta ^2/p_{\omega _k}^2}\right),\hfill & & \mathrm{otherwise},\hfill \end{array}`$ (3.2.17)
describes a phase in which: 1) chiral symmetry is dynamically broken, because one has a nonzero quark mass-function, $`B(p_{\omega _k})`$, in the absence of a current-quark mass; and 2) the dressed-quarks are confined, because the dressed-quark $`2`$-point Schwinger function violates reflection positivity. The alternative “Wigner” solution, for which
$$\widehat{B}(p_{\omega _k})0,\widehat{C}(p_{\omega _k})=\frac{1}{2}\left(1+\sqrt{1+2\eta ^2/p_{\omega _k}^2}\right),$$
(3.2.18)
characterises a phase in which chiral symmetry is not broken and the dressed-quarks are not confined.
To explore the phase transition one can consider the relative stability of the confined and deconfined phases, which is measured by their difference in pressure. Using Eq. (3.1.7) alone, because the gluon contributions cancel, that difference is
$`(\mu ,T)`$ $`:=`$ $`p_{\mathrm{\Sigma }_{\mathrm{NG}}}(\mu ,T)p_{\mathrm{\Sigma }_\mathrm{W}}(\mu ,T),`$ (3.2.19)
$`=`$ $`4N_c{\displaystyle _{l,q}^{\overline{\mathrm{\Lambda }}}}\left\{\mathrm{ln}\left[{\displaystyle \frac{|\stackrel{}{p}|^2A^2+\stackrel{~}{\omega }_k^2C^2+B^2}{|\stackrel{}{p}|^2\widehat{A}^2+\stackrel{~}{\omega }_k^2\widehat{C}^2}}\right]+|\stackrel{}{p}|^2\left(\sigma _A\widehat{\sigma }_A\right)+\stackrel{~}{\omega }_k^2\left(\sigma _C\widehat{\sigma }_C\right)\right\},`$ (3.2.20)
which in the particular case considered here yields
$$(\mu ,T)=\eta ^4\mathrm{\hspace{0.17em}2}N_cN_f\frac{\overline{T}}{\pi ^2}\underset{l=0}{\overset{l_{\mathrm{max}}}{}}_0^{\overline{\mathrm{\Lambda }}_l}𝑑yy^2\left\{\text{ }\mathrm{R}\left(2\overline{p}_l^2\right)\text{ }\mathrm{R}\left(\frac{1}{C(\overline{p}_l)}\right)\mathrm{ln}\left|\overline{p}_l^2C(\overline{p}_l)^2\right|\right\},$$
(3.2.21)
with: $`\overline{T}=T/\eta `$, $`\overline{\mu }=\mu /\eta `$; $`\overline{\omega }_{l_{\mathrm{max}}}^2\frac{1}{4}+\overline{\mu }^2`$, $`\overline{\mathrm{\Lambda }}_l^2=\overline{\omega }_{l_{\mathrm{max}}}^2\overline{\omega }_l^2`$, $`\overline{p}_l=(\stackrel{}{y},\overline{\omega }_l+i\overline{\mu })`$. $`(\mu ,T)`$ defines a $`(\mu ,T)`$-dependent “bag constant” and in this model
$$(0,0)=(0.102\eta )^4=(0.109\mathrm{GeV})^4,$$
(3.2.22)
which can be compared with the value $`(0.145\mathrm{GeV})^4`$ commonly used in bag-like models of hadrons. The positive value indicates that the confining phase of the model with DCSB is stable at $`\mu =0=T`$. The phase boundary in the $`(\mu ,T)`$-plane is defined by the line $`(\mu ,T)=0`$ , which is evident in Fig. 3.4. In this model the transition is first order except at $`\mu =0`$, which is clear from the nonzero derivative at the phase boundary, and the deconfinement and chiral symmetry restoring transitions are coincident, which is a consequence of the character of the solutions of the gap equation, Eqs. (3.2.143.2.18). The $`T=0`$ transitions are also coincident in the model obtained with $`𝒟_\mathrm{B}`$, Eq. (3.2.3. The chiral order parameters increase with increasing $`\mu `$, as they do in all models that preserve the momentum dependence of the dressed-quark self energies; e.g, Refs. , typically reaching a value $`<20`$% larger than that at $`\mu =0=T`$. This is an anticipated result of confinement: as long as $`\mu <\mu _c`$, each additional quark must be paired with an antiquark thereby increasing the density of condensate pairs. It cannot be observed in models that omit the momentum dependence of the dressed-quark self energy, such as those with an instantaneous interaction.
The criterion of maximal pressure was also employed in Ref. , wherein the model effective interaction saturates in the infrared and exhibits one-loop improved ultraviolet behaviour analogous to that in Eq. (2.2.35). The calculated $`\mu =0`$ critical temperature agrees with the estimate in Table 3.1 and the $`T=0`$ critical chemical potential: $`\mu _c=422`$MeV, is just $`13`$% larger than that estimated in Ref. . A similar interaction was employed in Ref. , wherein $`(\mu =0,T_c)`$ agrees with Table 3.1. However, the estimate of $`(\mu _c,T=0)`$ is $`60`$% larger than in Refs. and is likely too high. The usual range is $`\mu _c[300,400]`$MeV; e.g., see also Refs. .
As clear in Fig. 3.4, the chiral phase boundary traces out a curve in the $`(\mu ,T)`$-plane. In Ref. , beginning at $`(\mu _c,T=0)`$, this curve describes a line of first order critical points that ends with a tricritical point at $`(\mu _{\mathrm{tc}}0.5\mu _c,T_{\mathrm{tc}}0.8T_c)`$. For $`T>T_{\mathrm{tc}}`$ the transition is second order. In comparison with the $`𝒟_\mathrm{A}`$ model, this is a shift of the tricritical point away from the $`\mu =0`$ axis. We expect such a relocation to be a feature of all models with a better treatment of the ultraviolet behaviour of the effective interaction; e.g., there are preliminary indications that in the separable model of Ref. a tricritical point is located at $`(\mu _{\mathrm{tc}}0.4\mu _c,T_{\mathrm{tc}}0.9T_c)`$. The existence of a tricritical point has observable consequences . However, if it occurs at a value of $`\mu `$ as large as that estimated in these models it may be difficult to detect using the relativistic heavy ion collider \[RHIC\], which will focus on the domain of low baryon number density.
With $`(\mu ,T)>0`$, then, in the exploration of hadronic observables using the rainbow-ladder truncation, $`p_{\mathrm{\Sigma }_{\mathrm{NG}}}(\mu ,T)`$ describes the free energy of the “vacuum” or “ground state,” relative to which all excitations are measured $`(\mu ,T)𝒟`$; i.e., in the domain of confinement and DCSB. Therefore $`p_{\mathrm{\Sigma }_{\mathrm{NG}}}`$ is not directly observable and does not contribute to the thermodynamic pressure. Nevertheless, it does evolve with changes in $`(\mu ,T)`$, as evident in Fig. 3.4, and this evolution simply reflects the $`(\mu ,T)`$-dependence of the dressed-quark self energies, upon which it depends. \[NB. If it did not evolve there could never be a phase transition.\] The changes in the dressed-quark self energy represent the $`(\mu ,T)`$-dependent rearrangement of the vacuum and are indirectly observable; e.g., in the behaviour of hadron masses and decay constants. Hence, on the domain $`𝒟`$, they are evident in changes of the only true contributions to the pressure; i.e., those of the colour singlet hadronic correlations. At the phase boundary, denoted by $`𝒟`$: $`p_{\mathrm{\Sigma }_{\mathrm{NG}}}(\mu ,T)|_𝒟=p_{\mathrm{\Sigma }_\mathrm{W}}(\mu ,T)|_𝒟`$. This is the last vacuum reference point: hereafter dressed-quarks and -gluons are the active degrees of freedom.
The dressed-quarks’ contribution to the pressure is
$$P_q(\mu ,T)=\theta (𝒟)\left\{p_{\mathrm{\Sigma }_\mathrm{W}}(\mu ,T)p_{\mathrm{\Sigma }_\mathrm{W}}(\mu ,T)|_𝒟\right\},$$
(3.2.23)
where $`\theta (𝒟)`$ is a step function, equal to one for $`(T,\mu )𝒟`$. This pressure can be calculated numerically and is depicted in Fig. 3.5 for model $`𝒟_\mathrm{A}`$. As evident in Eq. (3.2.18), nonperturbative, momentum-dependent modifications of the dressed-quark’s vector self energy persist into the deconfined domain and this feature retards the approach of the pressure to the ultrarelativistic limit:
$$p_{\mathrm{UR}}(\mu ,T)=N_cN_f\frac{1}{12\pi ^2}\left(\mu ^4+2\pi ^2\mu ^2T^2+\frac{7}{15}\pi ^4T^4\right).$$
(3.2.24)
A qualitatively similar result is observed in numerical simulations of $`T0`$ lattice-QCD . Such a correspondence is impossible in models where the dressed-quark vector self energy is omitted.
The pressure in Eq. (3.2.23) is incomplete. It is zero in the confined domain where, as remarked above, the only contribution is that of colour singlet hadronic correlations. Those contributions can be estimated; e.g., using the hadronisation techniques of Ref. , but that has not yet been done. In the deconfined domain, one must add the contribution of the dressed-gluons. However, hitherto the model studies have employed “frozen” gluons; i.e., the structure of the effective interaction does not evolve with $`(\mu ,T)`$. $`T0`$ lattice-QCD simulations suggest that to be a good approximation, at least until very near the phase boundary , and Ref. indicates that changes in the interaction in this vicinity have little effect on the properties of the transitions. That is good because improving on the frozen approximation is difficult, requiring information about the effective interaction in a new domain. A first step is to allow the mass-scale characterising the infrared behaviour of the interaction to drop suddenly and significantly near the phase boundary, thereby mocking-up the $`T`$-dependence of the string tension. Such a study is presently lacking.
Not withstanding these weaknesses, the pressure depicted in Fig. 3.5 still provides a rudimentary model for the dressed-quark pressure that exhibits significant new qualitative features; e.g., the persistence into the deconfined domain of nonperturbative effects in the vector self energy yields an EOS much softer than that of a bag model. In this capacity it has been used to explore the structure and stability of nonstrange quark stars . References indicate that a transition to quark matter should occur at approximately three-times nuclear matter density, and the comparison between results obtained with the EOS depicted in Fig. 3.5 and a bag model EOS shows; e.g., that the prediction for the maximum stable radius of a pure quark star is model-independent: $`8`$$`10`$km, also in agreement with estimates based on a colour dielectric model . The maximum mass, however, is very sensitive to the EOS. The quark core described by Eq (3.2.23) must be surrounded by an hadronic crust and that will also affect the properties of the star. The problem rapidly becomes involved, and determining the composition of dense astrophysical objects and those signals that point to the realisation of quark matter in their core is an important focus of contemporary research; e.g., Refs. .
Superfluidity in Quark Matter In Sec. 2.3 we described aspects of diquark correlations and the manner in which they help to understand baryon properties. Another aspect follows from the meson-diquark auxiliary-field effective action : the steepest descent approximation to the vacuum pressure suggests the possibility of diquark condensation; i.e., quark-quark Cooper pairing, and that was first explored using a simple version of the Nambu–Jona-Lasinio model . A nonzero chemical potential promotes Cooper pairing in fermion systems and the possibility that it is exhibited in quark matter was considered using the rainbow-ladder truncation of the gap equation . Interest in that possibility has been renewed . A quark-quark Cooper pair is a composite boson with both electric and colour charge, and hence superfluidity in quark matter entails superconductivity and colour superconductivity. However, the last feature makes it difficult to identify an order parameter that can characterise a transition to the superfluid phase: the Cooper pair is gauge dependent and an order parameter is ideally describable by a gauge-invariant operator.
As we saw above, at $`\mu =0=T`$ QCD exhibits a nonzero quark-antiquark condensate but it is undermined by increasing $`\mu `$ and $`T`$, and there is a domain of the $`(\mu ,T)`$-plane, evident in Fig. 3.4, for which $`\overline{q}q=0`$. Increasing $`T`$ also opposes Cooper pairing. However, since increasing $`\mu `$ promotes it, there may be a (large-$`\mu `$,low-$`T`$)-subdomain in which quark matter exists in a superfluid phase. While that domain may not be accessible at RHIC, such conditions might exist in the core of dense astrophysical objects , which could undergo a transition to superfluid quark matter as they cool. Unambiguous signals of such a superfluid phase are actively being sought; e.g, Refs. .
We take Ref. as our exemplar, in which it was observed that a direct means of determining whether a SU$`{}_{c}{}^{}(N`$) gauge theory supports $`0^+`$ diquark condensation is to study the gap equation satisfied by
$$𝑫(p,\mu ):=𝒮(p,\mu )^1=\left(\begin{array}{cc}D(p,\mu )& \mathrm{\Delta }^i(p,\mu )\gamma _5\lambda _{}^i\\ \mathrm{\Delta }^i(p,\mu )\gamma _5\lambda _{}^i& CD(p,\mu )^\mathrm{t}C^{}\end{array}\right).$$
(3.2.25)
Here $`T=0`$, for illustrative simplicity and because temperature can only act to destabilise a Cooper pair, and, with $`\omega _{[\mu ]}=p_4+i\mu `$,
$$D(p,\mu )=i\stackrel{}{\gamma }\stackrel{}{p}A(\stackrel{}{p}^2,\omega _{[\mu ]}^2)+B(\stackrel{}{p}^2,\omega _{[\mu ]}^2)+i\gamma _4\omega _{[\mu ]}C(\stackrel{}{p}^2,\omega _{[\mu ]}^2);$$
(3.2.26)
i.e., the inverse of the dressed-quark $`2`$-point function, Eq. (3.1.11). In Eq. (3.2.25), $`\{\lambda _{}^i`$, $`i=1\mathrm{}n_c^{}`$, $`n_c^{}=N_c(N_c1)/2\}`$ are the antisymmetric generators of SU$`{}_{c}{}^{}(N_c)`$ and $`C`$ is the charge conjugation matrix, Eq. (2.3.8). Using such a gap equation to study superfluidity makes unnecessary a truncated bosonisation, which in all but the simplest models is a procedure difficult to improve systematically.
In addition to the usual colour, Dirac and isospin indices carried by the elements of $`𝑫(p,\mu )`$, the explicit matrix structure in Eq. (3.2.25) exhibits a quark bispinor index and is made with reference to
$$Q(x):=\left(\begin{array}{c}q(x)\\ \underset{¯}{q}(x):=\tau _f^2C\overline{q}^\mathrm{t}\end{array}\right),\overline{Q}(x):=\left(\begin{array}{cc}\overline{q}(x)\overline{\underset{¯}{q}}(x):=q^\mathrm{t}C\tau _f^2& \end{array}\right),$$
(3.2.27)
where $`\{\tau _f^i:i=1,2,3\}`$ are Pauli matrices, Eq. (2.1.7), that act on the isospin index. (Only two-flavour theories are considered in this exemplar. Additional possibilities open in three-flavour theories .)
As we have seen, a nonzero quark condensate: $`\overline{q}q0`$, is represented in the solution of the gap equation by $`B(\stackrel{}{p}^2,\omega _{[\mu ]}^2)0`$. The new but analogous feature here is that diquark condensation is characterised by $`\mathrm{\Delta }^i(p,\mu )0`$, for at least one $`i`$. That is clear if one considers the quark piece of the QCD Lagrangian density: $`L[\overline{q},q]`$. It is a scalar and hence $`L[\overline{q},q]^\mathrm{t}=L[\overline{q},q]`$. Therefore $`L[\overline{q},q]L[\overline{q},q]+L[\overline{q},q]^\mathrm{t}`$, and that sum, and hence the action, can be re-expressed in terms of a diagonal matrix using the bispinor fields in Eq. (3.2.27): $`\overline{Q}\mathrm{diag}[D,CD^\mathrm{t}C^{}]Q`$. It is plain now that a dynamically-generated lower-left element in $`𝑫(p,\mu )`$, the inverse of the dressed-bispinor propagator, corresponds to a $`\overline{\underset{¯}{q}}q`$ \[diquark\] correlation.
The bispinor gap equation can be written in the form
$$𝑫(p,\mu )=𝑫_0(p,\mu )+\left(\begin{array}{cc}\mathrm{\Sigma }_{11}(p,\mu )& \mathrm{\Sigma }_{12}(p,\mu )\\ \gamma _4\mathrm{\Sigma }_{12}(p,\mu )\gamma _4& C\mathrm{\Sigma }_{11}(p,\mu )^\mathrm{t}C^{}\end{array}\right),$$
(3.2.28)
where in the absence of a diquark source term
$$𝑫_0(p,\mu )=(i\gamma p+m)\tau _Q^0\mu \tau _Q^3,$$
(3.2.29)
with $`m`$ the current-quark mass, and the additional Pauli matrices: $`\{\tau _Q^\alpha ,\alpha =0,1,2,3\}`$, act on the bispinor indices. The structure of $`\mathrm{\Sigma }_{ij}(p,\mu )`$ specifies the theory and, in practice, also the approximation or truncation of it.
Two colour QCD \[QC<sub>2</sub>D\] provides an important and instructive example. In this case $`\mathrm{\Delta }^i\lambda _{}^i=\mathrm{\Delta }\tau _c^2`$ in Eq. (3.2.25), with $`\frac{1}{2}\stackrel{}{\tau _c}`$ the generators of $`SU_c(2)`$, and it is useful to employ a modified bispinor
$`Q_2(x)`$ $`:=`$ $`\left(\begin{array}{c}q(x)\\ \underset{¯}{q}_2:=\tau _c^2\underset{¯}{q}(x)\end{array}\right),`$ (3.2.32)
with $`\overline{Q}_2`$ the obvious analogue of $`\overline{Q}`$ in Eq. (3.2.27). Now the Lagrangian’s fermion-gauge-boson interaction term is simply
$$\overline{Q}_2(x)\frac{i}{2}g\gamma _\mu \tau _c^k\tau _Q^0Q_2(x)A_\mu ^k(x)$$
(3.2.33)
because SU$`{}_{c}{}^{}(2)`$ is pseudoreal; i.e., $`\tau _c^2\left(\stackrel{}{\tau }_c\right)^\mathrm{t}\tau _c^2=\stackrel{}{\tau }_c,`$ and the fundamental and conjugate representations are equivalent. (That the interaction term takes this form is easily seen using $`L[\overline{q},q]^\mathrm{t}=L[\overline{q},q]`$.)
Using the pseudoreality of SU$`{}_{c}{}^{}(2)`$ it can be shown that, for $`\mu =0`$ and in the chiral limit, the general solution of the bispinor gap equation is
$$𝑫(p)=i\gamma pA(p^2)+𝒱(𝝅)(p^2),𝒱(𝝅)=\mathrm{exp}\left\{i\gamma _5\underset{\mathrm{}=1}{\overset{5}{}}T^{\mathrm{}}\pi ^{\mathrm{}}\right\}=𝒱(𝝅)^1,$$
(3.2.34)
where $`\pi ^{\mathrm{}=1,\mathrm{},5}`$ are arbitrary constants and $`\{T^{1,2,3}=\tau _Q^3\stackrel{}{\tau _f},T^4=\tau _Q^1\tau _f^0,T^5=\tau _Q^2\tau _f^0\}`$, $`\{T^i,T^j\}=2\delta ^{ij}`$, so that
$$𝒮(p)=\frac{i\gamma pA(p^2)+𝒱(𝝅)(p^2)}{p^2A^2(p^2)+^2(p^2)}:=i\gamma p\sigma _V(p^2)+𝒱(𝝅)\sigma _S(p^2).$$
(3.2.35)
\[$`𝝅=(0,0,0,0,\frac{1}{4}\pi )`$ produces an inverse bispinor propagator with the simple form in Eq. (3.2.25).\] That the gap equation is satisfied for any constants $`\pi ^{\mathrm{}}`$ signals a vacuum degeneracy: if the interaction supports a mass gap, then that gap describes a five-parameter continuum of degenerate condensates:
$$\overline{Q}_2𝒱(𝝅)Q_20,$$
(3.2.36)
and there are 5 associated Goldstone bosons: 3 pions, a diquark and an anti-diquark. In this construction, Eq. (3.2.34), one has a simple elucidation of a necessary consequence of the Pauli-Gürsey symmetry of QC<sub>2</sub>D. For $`m0`$, the gap equation requires $`\mathrm{tr}_{FQ}\left[T^i𝒱\right]=0,`$ so that in this case only $`\overline{Q}_2Q_20`$ and the Goldstone bosons are now massive but remain degenerate.
For $`\mu 0`$ the general solution of the gap equation has the form
$$𝑫(p,\mu )=\left(\begin{array}{cc}D(p,\mu )& \gamma _5\mathrm{\Delta }(p,\mu )\\ \gamma _5\mathrm{\Delta }^{}(p,\mu )& CD(p,\mu )C^{}\end{array}\right),$$
(3.2.37)
and in the absence of a diquark condensate; i.e., for $`\mathrm{\Delta }0`$,
$$[U_B(\alpha ),𝑫(p,\mu )]=0,U_B(\alpha ):=\mathrm{e}^{i\alpha \tau _Q^3\tau _f^0},$$
(3.2.38)
which is a manifestation of baryon number conservation in QC<sub>2</sub>D. In this case the explicit form of the gap equation is complicated but its features and those of its solutions are easily illustrated using the model $`𝒟_\mathrm{A}`$ introduced above; e.g., it becomes clear that a chemical potential promotes fermion-pair condensation, at the expense of fermion-antifermion pairs, and $`\mathrm{\Delta }\text{ }\mathrm{R}`$, $`\mu `$. The rainbow-truncation gap equation was solved using this model in Ref. , and the relative stability of the quark- and diquark-condensed phases measured via the pressure difference
$$\delta p(\mu ):=p_{\mathrm{\Sigma }[B=0,\mathrm{\Delta }]}(\mu )p_{\mathrm{\Sigma }[B,\mathrm{\Delta }=0]}(\mu ),$$
(3.2.39)
where the pressure here is obtained as an obvious generalisation of Eq. (3.1.7). $`\delta p(\mu )`$ can be expressed in terms of $`\mu `$-dependent bag constants
$$_B(\mu ):=p_{\mathrm{\Sigma }[B,\mathrm{\Delta }=0]}(\mu )p_{\mathrm{\Sigma }[B=0,\mathrm{\Delta }=0]}(\mu ),_\mathrm{\Delta }(\mu ):=p_{\mathrm{\Sigma }[B=0,\mathrm{\Delta }]}(\mu )p_{\mathrm{\Sigma }[B=0,\mathrm{\Delta }=0]}(\mu ),$$
(3.2.40)
which measure the stability of a quark- or diquark-condensed vacuum relative to that with chiral symmetry realised in the Wigner-Weyl mode. \[NB. Improving on rainbow-ladder truncation may yield quantitative changes of $`<20`$% in the exemplary results that follow, however, the pseudoreality of QC<sub>2</sub>D and the equal dimension of the colour and bispinor spaces, which underly the theory’s Pauli-Gürsey symmetry, ensure that the entire discussion remains qualitatively unchanged.\]
The $`(m,\mu )=0`$ degeneracy of the quark and diquark condensates, Eq. (3.2.36), is manifest in
$$_B(0)=_\mathrm{\Delta }(0)=(0.13m_{J=1})^4,$$
(3.2.41)
where $`m_{J=1}`$ is the $`(m,\mu )=0`$ mass of the model’s vector meson, calculated in rainbow-ladder truncation. Increasing $`\mu `$ at $`m=0`$ and excluding diquark condensation, chiral symmetry is restored at
$$\mu _{2c}^{B,\mathrm{\Delta }=0}=0.40m_{J=1},$$
(3.2.42)
where $`_B(\mu )=0`$. However, $`\mu >0`$: $`\delta p(\mu )>0`$ and $`_\mathrm{\Delta }(\mu )>0,`$ with
$$_\mathrm{\Delta }(\mu _{2c}^{B,\mathrm{\Delta }=0})=(0.20m_{J=1})^4>_\mathrm{\Delta }(0).$$
(3.2.43)
Therefore the vacuum is unstable with respect to diquark condensation for all $`\mu >0`$, and this is always dynamically preferred over quark condensation.
$`(B=0,\mathrm{\Delta }0)`$ in Eq. (3.2.37) corresponds to $`𝝅=(0,0,0,0,\frac{1}{2}\pi )`$ in Eq. (3.2.36); i.e., $`\overline{Q}_2i\gamma _5\tau _Q^2Q_20`$. The usual chiral \[$`SU_A(2)`$\] transformations are realised via
$$𝑫(p,\mu )V(\stackrel{}{\pi })𝑫(p,\mu )V(\stackrel{}{\pi }),V(\stackrel{}{\pi }):=\mathrm{e}^{i\gamma _5\stackrel{}{\pi }\stackrel{}{T}},\stackrel{}{\pi }=(\pi ^1,\pi ^2,\pi ^3),$$
(3.2.44)
and therefore, since the anticommutator $`\{\stackrel{}{T},T^{4,5}\}=0`$, a diquark condensate does not break chiral symmetry. However, $`(B=0,\mathrm{\Delta }0)`$ does yield a dressed-bispinor propagator that violates reflection positivity and hence the model exhibits confinement to arbitrarily large densities. \[NB. Reference employs a frozen gluon approximation. In a more realistic analysis, the $`\mu `$-dependence of $`\eta `$, the mass-scale characterising the model, would be significant for $`\mu \mu _{2c}^{B,\mathrm{\Delta }=0}`$, and $`\eta 0`$ as $`\mu \mathrm{}`$, which would ensure deconfinement at large-$`\mu `$.\] Finally, although the $`\mu 0`$ theory is invariant under
$$Q_2U_B(\alpha )Q_2,\overline{Q}_2\overline{Q}_2U_B(\alpha ),$$
(3.2.45)
which is associated with baryon number conservation, the diquark condensate breaks this symmetry:
$$\overline{Q}_2i\gamma _5\tau _Q^2Q_2\mathrm{cos}(2\alpha )\overline{Q}_2i\gamma _5\tau _Q^2Q_2\mathrm{sin}(2\alpha )\overline{Q}_2i\gamma _5\tau _Q^1Q_2.$$
(3.2.46)
Hence, for $`(m=0,\mu 0)`$, one Goldstone mode remains. (These symmetry breaking patterns and the concomitant numbers of Goldstone modes in QC<sub>2</sub>D are also described in Ref. .)
For $`m0`$ and small values of $`\mu `$, the gap equation only admits a solution with $`\mathrm{\Delta }0`$; i.e., diquark condensation is blocked and this is because the current-quark mass is a source of quark condensation. However, with increasing $`\mu `$ a diquark condensate is generated and the $`𝒟_\mathrm{A}`$ model exhibits the following minimum chemical potentials for diquark condensation:
$$m=0.013m_{J=1}\mu ^{\mathrm{\Delta }0}=0.051m_{J=1},m=0.13m_{J=1}\mu ^{\mathrm{\Delta }0}=0.092m_{J=1}.$$
(3.2.47)
This retardation of diquark condensation by a nonzero current-quark mass can also be seen in Ref. .
The exploration of superfluidity in true QCD encounters two differences: the dimension of the colour space is greater than that of the bispinor space and the fundamental and conjugate representations of the gauge group are not equivalent. The latter is of obvious importance because it entails that the quark-quark and quark-antiquark scattering matrices are qualitatively different. As we saw above in connection with Eq. (2.3.42), these differences ensure that colour singlet meson bound states exist but \[necessarily coloured\] diquark bound states do not.
$`n_c^{}=3`$ in QCD and hence in canvassing superfluidity it is necessary to choose a direction for the condensate in colour space; e.g., $`\mathrm{\Delta }^i\lambda _{}^i=\mathrm{\Delta }\lambda ^2`$ in Eq. (3.2.25), so that
$$𝑫(p,\mu )=\left(\begin{array}{cc}D_{}(p,\mu )P_{}+D_{}(p,\mu )P_{}& \mathrm{\Delta }(p,\mu )\gamma _5\lambda ^2\\ & \\ \mathrm{\Delta }(p,\mu )\gamma _5\lambda ^2\text{}& CD_{}(p,\mu )C^{}P_{}+CD_{}(p,\mu )C^{}P_{}\end{array}\right),$$
(3.2.48)
where $`P_{}=(\lambda ^2)^2`$, $`P_{}+P_{}=\mathrm{diag}(1,1,1)`$, and $`D_{}`$, $`D_{}`$ are defined via obvious generalisations of Eqs. (3.2.25), (3.2.26). (NB. It is this selection of a direction in colour space that opens the possibility for colour-flavour locked diquark condensation in a theory with three effectively-massless quarks; i.e., current-quark masses $`\mu `$ .) In Eq. (3.2.48) the evident, demarcated block structure makes explicit the bispinor index. Here each block is a $`3\times 3`$ colour matrix and the subscripts: $``$, $``$, indicate whether or not the subspace is accessible via $`\lambda _2`$. The bispinors associated with this representation are given in Eqs. (3.2.27) and in this case the Lagrangian’s quark-gluon interaction term is
$`\overline{Q}(x)ig\mathrm{\Gamma }_\mu ^aQ(x)A_\mu ^a(x),\mathrm{\Gamma }_\mu ^a`$ $`=`$ $`\left(\begin{array}{cc}\frac{1}{2}\gamma _\mu \lambda ^a& 0\\ & \\ 0\text{}& \frac{1}{2}\gamma _\mu (\lambda ^a)^\mathrm{t}\end{array}\right).`$ (3.2.51)
It is straightforward to derive the gap equation at arbitrary order in the truncation scheme of Ref. and it is important to note that because
$$D_{}(p,\mu )P_{}+D_{}(p,\mu )P_{}=\lambda ^0\left\{\frac{2}{3}D_{}(p,\mu )+\frac{1}{3}D_{}(p,\mu )\right\}+\frac{1}{\sqrt{3}}\lambda ^8\left\{D_{}(p,\mu )D_{}(p,\mu )\right\}$$
(3.2.52)
the interaction: $`\mathrm{\Gamma }_\mu ^a𝒮(p,\mu )\mathrm{\Gamma }_\nu ^a`$, necessarily couples the $``$\- and $``$-components. That interplay is discarded in models that ignore the vector self energy of quarks, which, as we have repeatedly seen, is a qualitatively important feature of QCD .
Reference explored the possibility of diquark condensation in QCD via the gap equation using the model defined by $`𝒟_\mathrm{A}`$ in both the rainbow and a vertex-corrected truncation. The latter is illustrated in Fig. 3.6. For $`\mu =0`$ the rainbow-ladder truncation yields
$$\begin{array}{ccc}m_\omega ^2=m_\rho ^2=\frac{1}{2}\eta ^2,& \overline{q}q^0=(0.11\eta )^3,& _B(\mu =0)=(0.10\eta )^4,\end{array}$$
(3.2.53)
and momentum-dependent vector self energies, Eq. (3.2.14), which lead to an interaction between the $``$\- and $``$-components of $`𝑫`$ that blocks diquark condensation . This is in spite of the fact that $`\lambda ^a\lambda ^2(\lambda ^a)^\mathrm{t}=\frac{1}{2}\lambda ^a\lambda ^a,`$ which entails that the rainbow-truncation quark-quark scattering kernel is purely attractive and strong enough to produce diquark bound states . \[Remember that in the colour singlet meson channel the rainbow-ladder truncation gives the colour coupling $`\lambda ^a\lambda ^a`$.\] For $`\mu 0`$ and in the absence of diquark condensation, we saw in connection with Fig. (3.4) that the model exhibits coincident, first order chiral symmetry restoring and deconfining transitions at
$$\mu _{c,\mathrm{rainbow}}^{B,\mathrm{\Delta }=0}=0.28\eta =0.3\mathrm{GeV}.$$
(3.2.54)
For $`\mu 0`$, however, the rainbow-truncation gap equation admits a solution with $`\mathrm{\Delta }(p,\mu )0`$ and $`B(p,\mu )0`$. $`\delta p(\mu )`$ in Eq. (3.2.39) again determines whether the stable ground state is the quark-condensed or superfluid phase. With increasing $`\mu `$, $`_B(\mu )`$ decreases, very slowly at first, and $`_\mathrm{\Delta }(\mu )`$ increases rapidly from zero. That evolution continues until
$$\mu _{c,\mathrm{rainbow}}^{B=0,\mathrm{\Delta }}=0.25\eta =0.89\mu _{c,\mathrm{rainbow}}^{B,\mathrm{\Delta }=0},$$
(3.2.55)
where $`_\mathrm{\Delta }(\mu )`$ becomes greater-than $`_B(\mu )`$. This signals a first order transition to the superfluid ground state and at the boundary
$$\overline{Q}i\gamma _5\tau _Q^2\lambda ^2Q_{\mu =\mu _{c,\mathrm{rainbow}}^{B=0,\mathrm{\Delta }}}=(0.65)^3\overline{Q}Q_{\mu =0}.$$
(3.2.56)
The chemical potential at which the switch to the superfluid ground state occurs is consistent with other estimates made using models comparable to the rainbow-truncation class , as is the magnitude of the gap at this point .
A question that now arises is: how sensitive is this phenomenon to the nature of the quark-quark interaction? As we discussed in connection with Eq. (2.3.42), the inhomogeneous ladder BSE exhibits particle-like singularities in the $`0^+`$ diquark channels and such states do not exist in the strong interaction spectrum. Does diquark condensation persist when a truncation of the gap equation is employed that does not correspond to a BSE whose solutions exhibit diquark bound states? The vertex corrected gap equation depicted in Fig. 3.6 is just such a truncation and it was also studied in Ref. .
In this case there is a $`\mathrm{\Delta }0`$ solution even for $`\mu =0`$, and using $`𝒟_\mathrm{A}`$
$$\begin{array}{ccc}m_\rho ^2=(1.1)m_\rho ^{2\mathrm{ladder}},& \overline{Q}Q=(1.0)^3\overline{Q}Q^{\mathrm{rainbow}},& _B=(1.1)^4_B^{\mathrm{rainbow}},\end{array}$$
(3.2.57)
where the rainbow-ladder results are given in Eqs. (3.2.53), and
$$\overline{Q}i\gamma _5\tau _Q^2\lambda ^2Q=(0.48)^3\overline{Q}Q,_\mathrm{\Delta }=(0.42)^4_B.$$
(3.2.58)
The last result shows, unsurprisingly, that the quark-condensed phase is favoured at $`\mu =0`$. Precluding diquark condensation, the solution of the vertex-corrected gap equation exhibits coincident, first order chiral symmetry restoring and deconfinement transitions at
$$\mu _c^{B,\mathrm{\Delta }=0}=0.77\mu _{c,\mathrm{rainbow}}^{B,\mathrm{\Delta }=0}.$$
(3.2.59)
Admitting diquark condensation, however, the $`\mu `$-dependence of the bag constants again shows there is a transition to the superfluid phase, here at
$$\mu _c^{B=0,\mathrm{\Delta }}=0.63\mu _c^{B,\mathrm{\Delta }=0},\mathrm{with}\overline{Q}i\gamma _5\tau _Q^2\lambda ^2Q_{\mu =0.63\mu _c^{B,\mathrm{\Delta }=0}}=(0.51)^3\overline{Q}Q_{\mu =0}.$$
(3.2.60)
Thus the material step of eliminating diquark bound states leads only to small quantitative changes in the quantities characterising the still extant superfluid phase.
Solving the inhomogeneous BSE for the $`0^+`$ diquark vertex in the quark-condensed phase provides additional insight . At $`\mu =0`$ and zero total momentum: $`P=0`$, the additional \[confining\] contributions to the quark-quark scattering kernel generate an enhancement in the magnitude of the scalar functions in the Bethe-Salpeter amplitude. However, as $`P^2`$ evolves into the timelike region, the contributions become repulsive and block the formation of a diquark bound state. Conversely, increasing $`\mu `$ at any given timelike-$`P^2`$ yields an enhancement in the magnitude of the scalar functions, and as $`\mu \mu _c^{B,\mathrm{\Delta }=0}`$ that enhancement becomes large, which suggests the onset of an instability in the quark-condensed vacuum. This “robustness” of scalar diquark condensation is consistent with the observations in Ref. . However, the studies described herein do not obviate the question of whether the diquark condensed phase is stable with-respect-to dinucleon condensation , which requires further attention. \[NB. As remarked above, the inclusion of temperature undermines a putative diquark condensate and existing studies suggest that it will disappear for $`T>60`$$`100`$MeV. However, such temperatures are high relative to that anticipated inside dense astrophysical objects, which may therefore provide an environment for detecting quark matter superfluidity.\]
Finally, this discussion illustrates that, in some respects; such as the transition point and magnitude of the gap, the phase diagram of QC<sub>2</sub>D is quantitatively similar to that of QCD. That is a useful observation because the simplest superfluid order parameter is gauge invariant in QC<sub>2</sub>D, and the fermion determinant is real and positive, which makes tractable the exploration of superfluidity in QC<sub>2</sub>D using numerical simulations of the lattice theory . Hence, the results of those studies may provide an additional, reliable guide to the nature of quark matter superfluidity.
## 4 Density, Temperature and Hadrons
### 4.1 Masses and Widths
Temperature Hitherto we have canvassed the bulk thermodynamic properties of QCD at nonzero $`(\mu ,T)`$. However, the terrestrial formation of a QGP will be signalled by changes in the observable properties of those colour singlet mesons that reach detectors. In this connection, a primary feature of the QGP is chiral symmetry restoration and, since the properties of the pion (mass, decay constant, other vertex residues, etc.) are tied to the dynamical breaking of chiral symmetry, an elucidation of the $`T`$-dependence of these properties is important; particularly since a prodigious number of pions is produced in heavy ion collisions. Also important is understanding the $`T`$-dependence of the properties of the scalar analogues and chiral partners of the pion in the strong interaction spectrum. For example, should the mass of a putative light isoscalar-scalar meson fall below $`2m_\pi `$, the strong decay into a two pion final state can no longer provide its dominant decay mode. In this case electroweak processes will be the only open decay channels below $`T_c`$ and the state will appear as a narrow resonance . Analogous statements are true of isovector-scalar mesons.
The effective interaction denoted above as $`𝒟_\mathrm{C}`$ in Eq. (3.2.2), with $`\omega =1.2m_t`$, has been employed in an exploration of the $`T`$-dependence of scalar and pseudoscalar meson properties. As we saw in Sec. 2.3, mesons appear as simple poles in $`3`$-point vertices and these vertices alone provide information about the persistence of correlations away from the bound state pole, which can be useful in studying the $`T`$-evolution of a system with deconfinement. For $`T0`$ and two flavours, the ladder-truncation of the inhomogeneous BSE for the isovector $`0^+`$ vertex is
$$\mathrm{\Gamma }_{\mathrm{ps}}^i(p_{\omega _k};P_0;\zeta )=Z_4\frac{1}{2}\tau ^i\gamma _5_{l,q}^{\overline{\mathrm{\Lambda }}}\frac{4}{3}g^2D_{\mu \nu }(p_{\omega _k}q_{\omega _l})\gamma _\mu S(q_{\omega _l}^+)\mathrm{\Gamma }_{\mathrm{ps}}^i(q_{\omega _l};P_0;\zeta )S(q_{\omega _l}^{})\gamma _\nu ,$$
(4.1.1)
where $`q_{\omega _l}^\pm =q_{\omega _l}\pm P_0/2`$, and with $`P_0=(\stackrel{}{P},0)`$ this is the equation for the zeroth Matsubara mode. \[This is an extension of Eq. (2.3.4). Remember, $`m_R(\zeta )\mathrm{\Gamma }_{\mathrm{ps}}^i(p_{\omega _k};P_{\mathrm{\Omega }_n};\zeta )`$ is renormalisation point independent.\]
The solution of Eq. (4.1.1) has the form \[hereafter the $`\zeta `$-dependence is often implicit\]
$`\mathrm{\Gamma }_{\mathrm{ps}}^i(p_{\omega _k};\stackrel{}{P})=`$
$`\frac{1}{2}\tau ^i\gamma _5\left[iE_{\mathrm{ps}}(p_{\omega _k};\stackrel{}{P})+\stackrel{}{\gamma }\stackrel{}{P}F_{\mathrm{ps}}(p_{\omega _k};\stackrel{}{P})+\stackrel{}{\gamma }\stackrel{}{p}\stackrel{}{p}\stackrel{}{P}G_{\mathrm{ps}}^{}(p_{\omega _k};\stackrel{}{P})+\gamma _4\omega _k\stackrel{}{p}\stackrel{}{P}G_{\mathrm{ps}}^{}(p_{\omega _k};\stackrel{}{P})\right],`$
where the $`T0`$ analogues of the $`\sigma _{\mu \nu }`$-like contributions in Eq. (2.3.9) are omitted because they play a negligible role at $`T=0`$ . The breaking of $`O(4)`$ symmetry is responsible for making this $`T0`$ amplitude more complicated than its zero-temperature counterpart. For the higher Matsubara frequencies the form is still more complicated, with three additional terms. The scalar functions in Eq. (4.1) exhibit a simple pole at $`\stackrel{}{P}^2+m_\pi ^2=0`$ so that
$$\mathrm{\Gamma }_{\mathrm{ps}}^i(p_{\omega _k};\stackrel{}{P})=\frac{r_\pi (\zeta )}{\stackrel{}{P}^2+m_\pi ^2}\mathrm{\Gamma }_\pi ^i(p_{\omega _k};\stackrel{}{P})+\mathrm{regular},$$
(4.1.3)
where again regular means terms regular at this pole and $`\mathrm{\Gamma }_\pi ^i(p_{\omega _k};\stackrel{}{P})`$ is the bound state pion Bethe-Salpeter amplitude, canonically normalised via the obvious $`T0`$ extension of Eq. (2.3.6).
The pole position in Eq. (4.1.3) determines the spatial “screening-mass” of the pion’s zeroth Matsubara mode and its inverse describes the persistence length of that mode at equilibrium in the heat bath. There is a screening mass and amplitude for each mode, and each mode’s amplitude is canonically normalised. The full $`T0`$ bound state propagator can be calculated via any polarisation tensor that receives a contribution from the bound state but only once all the screening masses have been determined. (For the pion: the pseudoscalar or pseudovector polarisations will serve .) The propagator so obtained is defined only on a discrete set of points along what might be called the imaginary-energy axis and the “pole-mass;” i.e., the mass that yields the bound state energy pole for $`\stackrel{}{p}0`$, is obtained only after an analytic continuation of the propagator onto the real-energy axis. \[The loss of $`O(4)`$ invariance for $`T0`$ means that, in general, the pole mass and screening masses are unequal.\] That continuation is not unique but an unambiguous result is obtained by requiring that it yield a function that is bounded at complex-infinity and analytic off the real axis . From this description it is nonetheless clear that the screening masses completely specify the properties of $`T0`$ bound states. Furthermore, the often used calculational expedient of replacing the meson Matsubara frequencies by a continuous variable: $`\mathrm{\Omega }_ni\nu `$, and the identification of the energy scale thus obtained with a pole mass, is seen to be merely an artefice. However, since this prescription yields the correct result for free particle propagators, it might provide an illustrative guide.
The residue in Eq. (4.1.3) is
$$\delta ^{ij}ir_\pi =Z_4\mathrm{tr}_{l,q}^{\overline{\mathrm{\Lambda }}}\frac{1}{2}\tau ^i\gamma _5\chi _\pi ^j(q_{\omega _l};\stackrel{}{P}),$$
(4.1.4)
where $`\chi _\pi (q_{\omega _l};\stackrel{}{P}):=S(q_{\omega _l}^+)\mathrm{\Gamma }_\pi (q_{\omega _l};\stackrel{}{P})S(q_{\omega _l}^{})`$ is the unamputated Bethe-Salpeter wave function. \[Ref. employs the $`f_\pi =92`$MeV normalisation, cf. Eq. (2.3.22)\] Substituting Eq. (4.1.3) into Eq. (4.1.1) and equating pole residues yields the homogeneous pion BSE, which provides the simplest way to obtain the bound state amplitude. As already noted, the pion also appears as a pole in the axial-vector vertex and there the residue is the leptonic decay constant
$$\delta ^{ij}\stackrel{}{P}f_\pi =Z_2^A\mathrm{tr}_{l,q}^{\overline{\mathrm{\Lambda }}}\frac{1}{2}\tau ^i\gamma _5\stackrel{}{\gamma }\chi _\pi ^j(q_{\omega _l};\stackrel{}{P}).$$
(4.1.5)
At $`T=0`$ in this $`𝒟_\mathrm{C}`$ model the calculated chiral limit values are \[remember: $`f_\pi =92`$MeV normalisation\]
$$\begin{array}{ccc}f_\pi ^0=0.088\mathrm{GeV},& \overline{q}q_{1\mathrm{GeV}^2}^0=(0.235\mathrm{GeV})^3,& r_\pi ^0(1\mathrm{GeV}^2)=(0.457\mathrm{GeV})^2.\end{array}$$
(4.1.6)
The analogue of Eq. (4.1.1) for the $`0^{++}`$ vertex is
$$\mathrm{\Gamma }_\mathrm{s}^\alpha (p_{\omega _k};P_0;\zeta )=Z_4\frac{1}{2}\tau ^\alpha \mathrm{𝟏}_{l,q}^{\overline{\mathrm{\Lambda }}}\frac{4}{3}g^2D_{\mu \nu }(p_{\omega _k}q_{\omega _l})\gamma _\mu S(q_{\omega _l}^+)\mathrm{\Gamma }_\mathrm{s}^i(q_{\omega _l};P_{\mathrm{\Omega }_n};\zeta )S(q_{\omega _l}^{})\gamma _\nu ,$$
(4.1.7)
where $`\alpha =0,1,2,3`$. However, as discussed in connection with Eq. (2.3.23), the combination of rainbow and ladder truncations is not certain to provide a reliable approximation in the scalar sector. Nevertheless, in the absence of an improved, phenomenologically efficacious kernel, Ref. employed Eq. (4.1.7) in the expectation that it would provide some qualitatively reliable insight, an approach justified a posteriori.
The solution of Eq. (4.1.7) has the form
$`\mathrm{\Gamma }_\mathrm{s}^i(p_{\omega _k};\stackrel{}{P})=`$
$`\frac{1}{2}\tau ^\alpha \mathrm{𝟏}\left[E_\mathrm{s}(p_{\omega _k};\stackrel{}{P})+i\stackrel{}{\gamma }\stackrel{}{p}G_\mathrm{s}^{}(p_{\omega _k};\stackrel{}{P})+i\gamma _4\omega _kG_\mathrm{s}^{}(p_{\omega _k};\stackrel{}{P})+i\stackrel{}{\gamma }\stackrel{}{P}\stackrel{}{p}\stackrel{}{P}F_\mathrm{s}(p_{\omega _k};\stackrel{}{P})\right],`$
where here the requirement that the neutral mesons be charge conjugation eigenstates shifts the $`\stackrel{}{p}\stackrel{}{P}`$ term cf. the $`0^+`$ amplitude in Eq. (4.1). As already observed in connection with Eq. (2.3.69), the scalar functions in Eq. (4.1) exhibit a simple pole at $`\stackrel{}{P}^2+m_\sigma ^2=0`$ with residue
$$\delta ^{\alpha \beta }r_\sigma =Z_4\mathrm{tr}_{l,q}^{\overline{\mathrm{\Lambda }}}\frac{1}{2}\tau ^\alpha \chi _\sigma ^\beta (q_{\omega _l};\stackrel{}{P}).$$
(4.1.9)
However, since a $`VA`$ current cannot connect a $`0^{++}`$ state to the vacuum, the scalar meson does not appear as a pole in the vector vertex; i.e,
$$\delta ^{\alpha \beta }\stackrel{}{P}f_\sigma =Z_2^A\mathrm{tr}_{l,q}^{\overline{\mathrm{\Lambda }}}\frac{1}{2}\tau ^\alpha \stackrel{}{\gamma }\chi _\sigma ^\beta (q_{\omega _l};\stackrel{}{P})0.$$
(4.1.10)
The homogeneous equation for the scalar bound state amplitude is obtained, as usual, from Eq. (4.1.7) by equating pole residues.
The $`T`$-dependence of the pole positions in the solution of the inhomogeneous BSEs is illustrated in Fig. 4.1, from which it is evident that: 1) at the critical temperature, $`T_c=152`$MeV in Table 3.1, one has degenerate, massless pseudoscalar and scalar bound states; and 2) the bound states persist above $`T_c`$, becoming increasingly massive with increasing $`T`$. These features are also observed in numerical simulations of lattice-QCD . The bound state amplitudes are obtained from the homogeneous BSEs. Above $`T_c`$, all but the leading Dirac amplitudes: $`E_\pi `$, $`E_\sigma `$, vanish and the surviving amplitudes are pointwise identical. These results indicate that the chiral partners are locally identical above $`T_c`$, they do not just have the same mass.
The results are easily understood algebraically. The BSE is a set of coupled homogeneous equations for the Dirac amplitudes. Below $`T_c`$ each of the equations for the subleading Dirac amplitudes has an “inhomogeneity” whose magnitude is determined by $`B_0`$, the scalar piece of the quark self energy, which is dynamically generated in the chiral limit. $`B_0`$ vanishes above $`T_c`$ eliminating the inhomogeneity and allowing a trivial, identically zero solution for each of these amplitudes. Additionally, with $`B_00`$ the kernels in the equations for the dominant pseudoscalar and scalar amplitudes are identical, and hence so are the solutions. It follows from these observations that the Goldberger-Treiman-like relation
$$f_\pi ^0E_\pi (p_{\omega _k};0)=B_0(p_{\omega _k})$$
(4.1.11)
is satisfied for all $`T`$ only because both $`f_\pi ^0`$ and $`B_0(p_{\omega _k})`$ are equivalent order parameters for chiral symmetry restoration.
Using the calculated bound state amplitudes and dressed-quark propagators, the $`T`$-dependence of the matrix elements in Eqs. (4.1.4), (4.1.5), (4.1.9) follows. It is depicted in Fig. 4.2 for the chiral limit. Below $`T_c`$ the scalar meson residue in the scalar vertex, $`r_\sigma `$ in Eq. (4.1.9), is a little larger than the residue of the pseudoscalar meson in the pseudoscalar vertex, $`r_\pi `$ in Eq. (4.1.4). However, they are nonzero and equal above $`T_c`$, which is an algebraic consequence of $`B_00`$ and the vanishing of the subleading Dirac amplitudes. As a bona fide order parameter for chiral symmetry restoration
$$f_\pi ^0\sqrt{t},t=(T/T_c1)<0,$$
(4.1.12)
as anticipated from the discussion accompanying Table 3.1. The same result was obtained in Ref. , where it was observed too that
$$\frac{1}{(r_\pi ^0)_{\mathrm{em}}}f_\pi ^0;$$
(4.1.13)
i.e., the pion charge radius diverges at $`T_c`$ in the chiral limit. This is plausible but Ref. did not attempt its verification. Reference did confirm the behaviour of the ratio of pole residues observed in Ref. :
$$\frac{r_\pi ^0(\zeta )}{f_\pi ^0}\frac{1}{\sqrt{t}},t<0.$$
(4.1.14)
The concomitant results: $`f_\pi ^0=0`$ and $`r_\pi ^00`$ for $`t>0`$, demonstrate that the pion disappears as a pole in the axial-vector vertex but persists as a pole in the pseudoscalar vertex. Evident also in Fig. 4.2 is that
$$m_\sigma ^0\sqrt{t},t<0:$$
(4.1.15)
$`m_\sigma ^2`$ follows a linear trajectory for $`t<0`$. Such behaviour in the isoscalar-scalar channel might be anticipated because this channel has vacuum quantum numbers and hence the bound state is a strong interaction analogue of the electroweak Higgs boson .
$`m_{\mathrm{sym}}^2`$ in Fig. 4.2 is the mass obtained when the chirally symmetric solution of the quark DSE is used in the BSE. \[$`B_00`$ is always a solution in the chiral limit.\] For $`t>0`$, $`m_{\mathrm{sym}}^2(T)`$ is the unique meson mass-squared trajectory. However, for $`t<0`$, $`m_{\mathrm{sym}}^2<0`$; i.e., the solution of the BSE in the Wigner-Weyl phase exhibits a tachyonic solution. \[cf. The Nambu-Goldstone phase masses: $`m_\sigma ^2>m_\pi ^2=0`$.\] By analogy with the $`\sigma `$-model, this tachyonic mass indicates the instability of the Wigner-Weyl phase below $`T_c`$ and translates into the statement that the pressure is not maximal in this phase. Figure 4.3 depicts the evolution of this \[common\] meson mass at large $`T`$. As expected in a gas of weakly interacting quarks and gluons
$$\frac{m_{0^\pm \mathrm{meson}}}{2\omega _0}1^{},$$
(4.1.16)
where $`\omega _0=\pi T`$ is a quark’s zeroth Matsubara frequency and “screening mass.” Similar behaviour is observed for the $`\rho `$-meson mass in Ref. and can be demonstrated algebraically using the $`𝒟_\mathrm{A}`$ model for the effective interaction .
The results described hitherto were all calculated in the chiral limit. The extension to $`\widehat{m}0`$ is straightforward although calculations with the renormalisation group improved rainbow-ladder truncation become more time consuming. That is why simpler models, such as employed in Refs. , can be useful. For $`\widehat{m}0`$, chiral symmetry restoration with increasing $`T`$ is exhibited as a crossover rather than a phase transition. The solutions of the inhomogeneous BSEs again exhibit a pole for all $`T`$, with the bound state amplitudes obtained from the associated homogeneous equations. In this case the scalar and pseudoscalar bound states are locally identical for $`T>\frac{4}{3}T_c`$. Figure 4.4 is the $`\widehat{m}0`$ analogue of Fig. 4.2. An important result is that the axial-vector Ward-Takahashi identity is satisfied, both above and below the chiral transition temperature, which was demonstrated in Ref. via Eq. (2.3.21): the two sides remain equal $`T`$. The Gell-Mann–Oakes-Renner formula, however, which involves $`r_\pi ^0(\zeta )`$, fails for $`t>0.1`$, as observed too in the separable model study of Ref. .
As one can anticipate from Sec. 2.3, the calculated $`\sigma `$ and $`\pi `$ bound state amplitudes and dressed-quark propagator also make possible a study of two-body decays. For example, the impulse approximation to the isoscalar-scalar-$`\pi \pi `$ coupling is described by the matrix element
$`g_{\sigma \pi \pi }:=\pi (\stackrel{}{p_1})\pi (\stackrel{}{p_2})|\sigma (\stackrel{}{p})`$
$`=2N_c\mathrm{tr}_D{\displaystyle _{l,q}^{\overline{\mathrm{\Lambda }}}}\mathrm{\Gamma }_\sigma (k_{\omega _l};\stackrel{}{p})S_u(k_{++})i\mathrm{\Gamma }_\pi (k_{0+};\stackrel{}{p_1})S_u(k_+)i\mathrm{\Gamma }_\pi (k_0;\stackrel{}{p_2})S_u(k_{}),`$
$`k_{\alpha \beta }=k_{\omega _l}+(\alpha /2)\stackrel{}{p_1}+(\beta /2)\stackrel{}{p_2}`$, in terms of which the width is
$$\mathrm{\Gamma }_{\sigma (\pi \pi )}=\frac{3}{2}g_{\sigma \pi \pi }^2\frac{\sqrt{14m_\pi ^2/m_\sigma ^2}}{16\pi m_\sigma }.$$
(4.1.18)
The coupling and width obtained from Eq. (4.1) are depicted in Fig. 4.5, which indicates that both vanish at $`T_c`$ in the chiral limit. Again this can be traced to $`B_00`$. For $`\widehat{m}0`$, the coupling reflects the crossover but that is not observable because the width vanishes just below $`T_c`$ where the isoscalar-scalar meson mass falls below $`2m_\pi `$ and the phase space factor vanishes. \[See the right panel of Fig. 4.4.\] The evolution $`m_\sigma 2m_\pi `$ may, however, be observable via an enhancement in the $`\pi \pi \gamma \gamma `$ cross-section .
Additionally, the particular properties of the $`\pi ^0\gamma \gamma `$ decay, which is mediated by the “triangle anomaly,” \[see the paragraph after Eq. (2.3.41) on page 2.3\] make interesting the behaviour of this process at $`T0`$. At $`T=0`$, the anomalous contribution to the divergence of the axial-vector vertex is saturated by the pseudoscalar piece, $`E_\pi `$, of the pion Bethe-Salpeter amplitude
$$\widehat{T}_{\mu \nu }(k_1,k_2)=\mathrm{tr}_{l,q}^{\overline{\mathrm{\Lambda }}}S(q_1)\gamma _5\tau ^3iE_\pi (\widehat{q};P)S(q_2)iQ_e\mathrm{\Gamma }_\mu (q_2,q_{12})S(q_{12})iQ_e\mathrm{\Gamma }_\nu (q_{12},q_1),$$
(4.1.19)
where here, just to be specific, $`k_1=(\stackrel{}{k}_1,0)`$, $`k_2=(\stackrel{}{k}_2,0)`$, $`P=k_1+k_2`$, $`q_1=q_{\omega _l}k_1`$, $`q_2=q_{\omega _l}+k_2`$, $`\widehat{q}=\frac{1}{2}(q_1+q_2)`$, $`q_{12}=q_{\omega _l}k_1+k_2`$. Equation (4.1.19) involves the dressed-quark-photon vertex: $`\mathrm{\Gamma }_\mu `$, which also appeared in the calculation of $`F_\pi (Q^2)`$, described in Sec. 2.3. As we saw, quantitatively reliable numerical solutions of the $`T=0`$ vector vertex equation are now available . However, this anomalous coupling is insensitive to details and an accurate result requires only that the dressed vertex satisfy the vector Ward-Takahashi identity. For $`T=0`$ with real photons, Eq. (4.1.19) is expressed in terms of one scalar function:
$$\widehat{T}_{\mu \nu }(k_1,k_2)=\frac{\alpha _{\mathrm{em}}}{\pi }ϵ_{\mu \nu \rho \sigma }k_{1\rho }k_{2\sigma }𝒯(0).$$
(4.1.20)
Now, as long as $`\mathrm{\Gamma }_\mu `$ satisfies the Ward-Takahashi identity, Eq. (2.3.26), one finds algebraically in the chiral limit
$$f_\pi ^0𝒯(0):=g_{\pi ^0\gamma \gamma }=1/2,$$
(4.1.21)
completely independent of model details , as required since at $`T=0`$ the anomalies are a feature of the global aspects of DCSB . This value reproduces the experimental width. \[Remember, the normalisation in this subsection yields $`f_\pi =92`$MeV.\]
The $`T0`$ calculation requires only a valid extension of Eq. (2.3.31) and one such is
$`i\stackrel{}{\mathrm{\Gamma }}(q_{\omega _{l_1}},q_{\omega _{l_2}})`$ $`=`$ $`\mathrm{\Sigma }_A(q_{\omega _{l_1}}^2,q_{\omega _{l_2}}^2)i\stackrel{}{\gamma }+(\stackrel{}{q}_1+\stackrel{}{q}_2)[\frac{1}{2}iG(q_{\omega _{l_1}},q_{\omega _{l_2}})+\mathrm{\Delta }_B(q_{\omega _{l_1}}^2,q_{\omega _{l_2}}^2)],`$ (4.1.22)
$`i\mathrm{\Gamma }_4(q_{\omega _{l_1}},q_{\omega _{l_2}})`$ $`=`$ $`\mathrm{\Sigma }_C(q_{\omega _{l_1}}^2,q_{\omega _{l_2}}^2)i\gamma _4+(\omega _{l_1}+\omega _{l_2})[\frac{1}{2}iG(q_{\omega _{l_1}},q_{\omega _{l_2}})+\mathrm{\Delta }_B(q_{\omega _{l_1}}^2,q_{\omega _{l_2}}^2)],`$ (4.1.23)
$`G(q_{\omega _{l_1}},q_{\omega _{l_2}})`$ $`=`$ $`\stackrel{}{\gamma }(\stackrel{}{q}_1+\stackrel{}{q}_2)\mathrm{\Delta }_A(q_{\omega _{l_1}}^2,q_{\omega _{l_2}}^2)+\gamma _4(\omega _{l_1}+\omega _{l_2})\mathrm{\Delta }_C(q_{\omega _{l_1}}^2,q_{\omega _{l_2}}^2).`$ (4.1.24)
It is a particular case of the Ansatz in Ref. and satisfies the $`T0`$ vector Ward-Takahashi identity
$$(q_{\omega _{l_1}}q_{\omega _{l_2}})_\mu i\mathrm{\Gamma }_\mu (q_{\omega _{l_1}},q_{\omega _{l_2}})=S^1(q_{\omega _{l_1}})S^1(q_{\omega _{l_2}}).$$
(4.1.25)
For $`T0`$ the tensor structure of Eq. (4.1.20) survives to the extent that, with $`k_1`$, $`k_2`$ as defined, it ensures one of the photons is longitudinal (a plasmon) and the other transverse, with
$$\widehat{T}_{i4}(k_1,k_2)=\frac{\alpha _{\mathrm{em}}}{\pi }(\stackrel{}{k_1}\times \stackrel{}{k_2})_i𝒯(0).$$
(4.1.26)
The $`T`$-dependence of the anomalous coupling follows from that of $`𝒯(0)`$, and it and the $`T`$-dependence of the width are depicted in Fig. 4.6. In the chiral limit the interesting quantity is: $`𝒯(0)=g_{\pi ^0\gamma \gamma }^0/f_\pi ^0`$, and obvious in the figure is that it vanishes at $`T_c`$. (It vanishes with a mean field critical exponent .) Thus, in the chiral limit, the coupling to the dominant decay channel closes for both charged and neutral pions. These features were anticipated in Ref. . Further, the calculated $`𝒯(0)`$ is monotonically decreasing with $`T`$, supporting the perturbative O($`T^2/f_\pi ^2)`$ analysis in Ref. . For $`\widehat{m}0`$ both the coupling: $`g_{\pi ^0\gamma \gamma }/f_\pi `$, and the width exhibit the crossover, with a slight enhancement in the width as $`TT_c`$ due to the increase in $`m_\pi `$. This is similar to the results of Ref. , although the $`T`$-dependence depicted here is much weaker because the pion mass approaches twice the $`T0`$ free-quark screening-mass from below, never reaching it, Eq. (4.1.16); i.e., the continuum threshold is not crossed.
The $`T`$-dependence of meson properties illustrated here is robust: it agrees with the results obtained in lattice simulations when there is an overlap, and with the results of other models. The local equivalence exhibited by isovector chiral partners above $`T_c`$ might be expected as a general feature of the bound state spectrum in the Wigner-Weyl phase. However, an explicit demonstration is numerically challenging; e.g., in the $`\rho `$-$`a_1`$ complex the bound state amplitudes have eight independent functions even at $`T=0`$, compared with the four in the pseudoscalar and scalar amplitudes at $`T0`$.
Adding a third light quark introduces one qualitatively new aspect: the $`\eta `$-$`\eta ^{}`$ system and the restoration of $`U_A(1)`$ symmetry, which can affect the order of the chiral transition . As already observed, it is necessary to move beyond the rainbow-ladder truncation before that can be addressed using the DSEs. The question has been explored in lattice simulations but the results are not currently conclusive: the mass splittings used to characterise the symmetry breaking might become smaller near $`T_c`$ but strong, topological arguments can nevertheless be made in favour of the non-restoration of $`U_A(1)`$ symmetry . Much remains to be done and improved models can be useful.
Chemical Potential The relation between chemical potential and baryon number density can only be determined after the EOS is known; i.e., via Eq. (3.1.5). As described in connection with Eq. (3.2.23) on page 3.2.23, the dressed-quarks and -gluons contribute nothing to the EOS in the confined domain, even though they dominate it in the QGP, and the only true contributions to the pressure in the confined domain are those of colour singlet bound states. This physical requirement is overlooked in many model explorations of the density dependence of meson properties. For example, in applications of the Nambu–Jona-Lasinio model the EOS for a free fermion gas is used to define a baryon number density. While this is the EOS for that non-confining model it is not a good model for QCD’s EOS.
The application of DSEs in calculating the $`\mu `$-dependence of hadron properties is rudimentary. However, even that is significant given the problem is currently inaccessible in lattice simulations. The model obtained with $`𝒟_\mathrm{A}`$ in Eq. (3.2.2) again provides a useful, algebraic exemplar. As described after Eq. (3.2.22), the QGP transition is first order at $`(\mu ,T=0)`$ and the chiral order parameters increase with increasing $`\mu `$ when the dressed-quark self energy is momentum dependent. Mechanically, the latter is an obvious consequence of analyticity in the neighbourhood of the real axis: any function, $`O(4)`$ invariant at $`\mu =0=T`$, has the expansion
$$f(\stackrel{}{p}^{\mathrm{\hspace{0.17em}2}},\stackrel{~}{\omega }_k^2)\stackrel{T0\mu }{=}f(\stackrel{}{p}^{\mathrm{\hspace{0.17em}2}},0)+\stackrel{~}{\omega }_k^2\frac{f(\stackrel{}{p}^{\mathrm{\hspace{0.17em}2}},y)}{y}|_{y=\stackrel{~}{\omega }_k^2=0}+\mathrm{},$$
(4.1.27)
and since $`\text{ }\mathrm{R}[\stackrel{~}{\omega }_k^2]=\omega _k^2\mu ^2`$ then, if $`\text{ }\mathrm{R}[f(\stackrel{}{p}^{\mathrm{\hspace{0.17em}2}},\stackrel{~}{\omega }^2)]`$ decreases with $`T^2`$, the derivative is negative and $`\text{ }\mathrm{R}[f(\stackrel{}{p}^{\mathrm{\hspace{0.17em}2}},\stackrel{~}{\omega }^2)]`$ must increase with $`\mu ^2`$. \[Only the real-part is important because the imaginary-part of physical quantities vanishes after summing over the Matsubara frequencies. The derivative is zero in models with an instantaneous interaction.\]
Equation (4.1.27) is exemplified in the behaviour of the pion’s leptonic decay constant, which using the algebraic $`𝒟_\mathrm{A}`$ model is simply expressed:
$`f_\pi ^2`$ $`=`$ $`\eta ^2{\displaystyle \frac{16N_c}{\pi ^2}}\overline{T}{\displaystyle \underset{l=0}{\overset{l_{\mathrm{max}}}{}}}{\displaystyle \frac{\overline{\mathrm{\Lambda }}_l^3}{3}}\left(1+4\overline{\mu }^24\overline{\omega }_l^2\frac{8}{5}\overline{\mathrm{\Lambda }}_l^2\right),`$ (4.1.28)
where the notation is specified in connection with Eq. (3.2.21) on page 3.2.21. Its behaviour is depicted in Fig. 4.7, as is that of the pion’s mass. $`m_\pi `$ is almost insensitive to changes in $`\mu `$ and only increases slowly with $`T`$, until $`T`$ is very near the critical temperature, as already seen in Fig. 4.4, which was calculated with the renormalisation-group-improved $`𝒟_\mathrm{C}`$ model. The insensitivity to $`\mu `$ mirrors that to $`T`$ and is the result of compensating changes in $`r_\pi (\zeta )`$ and $`f_\pi `$; i.e., a consequence of the axial-vector Ward-Takahashi identity. All these features are also evident in Refs. .
A first step in exploring the properties of the $`\rho `$-meson is solving the the vector meson BSE. That too takes a particularly simple form in the $`𝒟_\mathrm{A}`$ model :
$$\mathrm{\Gamma }_M(p_{\omega _k};\stackrel{}{P})=\frac{\eta ^2}{4}\text{ }\mathrm{R}\left\{\gamma _\mu S(p_{\omega _k}+\frac{1}{2}\stackrel{}{P})\mathrm{\Gamma }_M(p_{\omega _k};\stackrel{ˇ}{P}_{\mathrm{}})S(p_{\omega _k}\frac{1}{2}\stackrel{}{P})\gamma _\mu \right\}.$$
(4.1.29)
There are two solutions: one longitudinal and one transverse with-respect-to $`\stackrel{}{P}`$:
$$\mathrm{\Gamma }_\rho =\{\begin{array}{c}\gamma _4\theta _{\rho +}\hfill \\ \left(\stackrel{}{\gamma }\frac{1}{\left|\stackrel{}{P}\right|^2}\stackrel{}{P}\stackrel{}{\gamma }\stackrel{}{P}\right)\theta _\rho \hfill \end{array},$$
(4.1.30)
where $`\theta _{\rho +}`$ labels the longitudinal and $`\theta _\rho `$ the transverse solution. Substituting, one finds that for $`\widehat{m}=0`$
$$m_\rho ^2=\frac{1}{2}\eta ^2,\text{independent}\text{ of }\mu \text{ and }T\text{.}$$
(4.1.31)
Even for nonzero current-quark mass, $`m_\rho `$ changes by less than 1% as $`\mu `$ and $`T`$ are increased from zero toward their critical values. This insensitivity is just what one would expect for the transverse mode: remember, there is no constant mass shift in the transverse polarisation tensor for a gauge-boson.
For the longitudinal component one obtains in the chiral limit:
$$m_{\rho +}^2=\frac{1}{2}\eta ^2+4(\pi ^2T^2\mu ^2),$$
(4.1.32)
where the anticorrelation between the response of $`m_{\rho +}`$ to $`T`$ and its response to $`\mu `$ is plain. From Eq. (4.1.32) and continuity at the second-order phase boundary it follows that
$$\frac{m_{\rho +}^2}{(2\pi T)^2}\stackrel{T\mathrm{}}{}\mathrm{\hspace{0.17em}1}^+,$$
(4.1.33)
which is analogous to Eq. (4.1.16). Here, however, the limit is approached from above because $`m_{\rho ^+}0`$ in the chiral limit and increases with $`T`$. (NB. It is only because this model exhibits confinement that such a result is possible. Studying the $`\rho `$-meson in non-confining dressed-quark-based models requires that some means be employed to suppress or eliminate the $`\rho \overline{q}q`$ threshold; e.g., Ref. . However, that artefice is merely an indigent expression of confinement. For the transverse component of the $`\rho `$, $`m_\rho ^{}^2/(2\omega _0)^21^{}`$ because of Eq. (4.1.31); e.g., Ref. .) The $`(\mu ,T)`$-dependence of the $`\rho `$-meson mass is depicted in Fig. 4.8 and; e.g., at $`T=0`$, $`M_{\rho +}(\mu _c)0.65M_{\rho +}(\mu =0)`$. As observed in the introduction, though, the connection between $`\mu `$-dependence and baryon-density-dependence cannot be determined until the EOS is calculated. Without it one can only observe that, in a two-flavour free-quark gas, the $`T=0`$ critical chemical potential corresponds to $`3\rho _0`$, see Fig. 3.4 on page 3.4. Therefore, at $`1`$$`2\rho _0`$ a mass reduction less-than this should be anticipated, plausibly no more than $`25`$.
The $`T`$-dependence described above was also observed in the confining model of Refs. , wherein too it was found that, because of the $`T`$-dependence of $`m_\pi `$, $`m_\rho ^{}`$, the dominant $`2\pi `$-decay mode of the $`\rho _{}`$ meson mode is phase-space blocked for $`T/T_c>1.2`$. \[cf. The $`2\pi `$ mode of the isoscalar-scalar discussed in connection with Eq. (4.1.18).\] The $`T`$-dependence of the $`\rho `$-meson’s dilepton decay width was also considered: it is suppressed by a factor of $`0.9`$ in the vicinity of $`(\mu =0,T_c)`$. However, its $`\mu `$-dependence is yet to be explored.
The anticorrelation, anticipated in Eq. (4.1.27), between the $`\mu `$ and $`T`$ dependence of mass-dimension-two observables; such as $`𝒳`$, $`f_\pi `$, $`m_{\rho ^+}`$, etc., is apparent. It entails that, in these cases at least, there is a trajectory in the $`(\mu ,T)`$-plane along which the observables are constant . It also means that observables calculated using the rainbow-ladder truncation do not exhibit a $`\mu `$-scaling law of the type conjectured for baryon-number-density in Ref. .
We anticipate that the mass of the $`a_1`$-meson \[the chiral partner of the $`\rho `$\] will decrease with increasing $`T`$ so that it can evolve to meet the increasing $`\rho `$-mass, with $`m_{\rho ^+}=m_{a_1^+}`$ at the phase boundary. However, since model $`𝒟_\mathrm{A}`$ is defective in not supporting an axial-vector bound state, that remains to be verified. (It fails to support a scalar bound state too .) As remarked in Sec. 2.3 and Ref. , much remains to be learnt about axial-vector mesons, in which connection confinement is an important element. Determining the $`\mu `$-dependence of $`m_{a_1}`$ is particularly interesting given the $`\mu `$ cf. $`T`$ anticorrelation exemplified above.
In the studies described here, and also in lattice simulations, $`\mu `$ is an intensive thermodynamic parameter whose presence modifies the propagation characteristics of dressed, confined particles, and this modification is transmitted to the observable hadrons they comprise. It is clear from the existence of a critical quark chemical potential, below which asymptotic quark states cannot be produced, that in the confined domain there is no simple proportionality between the quark chemical potential and the chemical potential associated with colour singlet baryons. In approaches based on elementary hadronic degrees of freedom; e.g., those reviewed in Ref. , this consideration is bypassed. Colour singlet baryon density is introduced directly via the expedient of in-medium elementary meson and nucleon propagators, which are then employed in calculating the myriad nuclear-matter many-body loop integrals that contribute to observable processes. The approach has a long history and yields a useful phenomenology. However, it ignores the dressed-quark level effects described above and also questions, such as, just what is represented by an off-shell hadron propagator? \[cf. The discussions associated with Eqs. (2.3.34)–(2.3.40), page 2.3.34, and (2.3.69)–(2.3.79), page 2.3.69.\] We judge that in understanding QCD at nonzero baryon density it is important to uncover the nature of the relationship between these approaches.
### 4.2 Hadronic Signatures of the Quark-Gluon Plasma
A QGP existed approximately one microsecond after the big bang, and primary goals of current generation experiments at CERN and Brookhaven are the terrestrial recreation of the plasma and an elucidation of its properties. A number of signals for QGP formation at high temperature and low baryon number density have been suggested and here we briefly review three of them. They and others are discussed more extensively in Refs. .
$`J/\mathrm{\Psi }`$ Suppression. Reference proposed “…that $`J/\mathrm{\Psi }`$ suppression in nuclear collisions should provide an unambiguous signature of quark-gluon plasma formation.” The reasoning is simple. A $`c\overline{c}`$-pair produced in a hard parton-parton interaction will evolve into a $`J/\mathrm{\Psi }`$-meson if the colour interaction is sufficiently strong to effect binding. That will always be the case unless the $`c\overline{c}`$-pair is produced in a heat-bath of deconfined, colour-carrying excitations that \[Debye-\]screen the $`c\overline{c}`$-attraction; i.e., unless the $`c\overline{c}`$-pair is produced in a QGP. QGP formation occurs only for temperatures greater than some critical value. Hence the $`J/\mathrm{\Psi }`$ production cross-section should evolve smoothly with controllable experimental parameters; such as, projectile and target mass numbers, and impact parameter, until a QGP is produced, when a dramatic suppression should follow.
Following this suggestion, the $`J/\mathrm{\Psi }`$ production cross-section has been systematically explored in relativistic heavy ion collisions at CERN using the Super Proton Synchrotron \[SpS\]. As illustrated in Fig. 4.9, all results are described by a simple scaling law:
$$(A_{\mathrm{projectile}}\times B_{\mathrm{target}})^\alpha ,\alpha =0.92\pm 0.01,$$
(4.2.1)
where $`A`$, $`B`$ are the mass numbers, except those for Pb-Pb collisions. The scaling relation is easily understood as a consequence of nuclear absorption and is termed “normal” $`J/\mathrm{\Psi }`$-suppression. While the Pb-Pb data agree with this normal pattern for peripheral collisions; i.e., $`L8`$fm, that is not the case for the most central collisions, which correspond to an energy density range of $`2`$$`3`$GeV$`/`$fm<sup>3</sup>. Hitherto, using standard in-medium hadronic tools, this “anomalous” suppression is inexplicable. However, an explanation can be found in the transition to a QGP; e.g., Refs. , and this signal has recently been claimed as “…evidence for the creation of a new state of matter in Pb-Pb collisions at the CERN SPS.”
Low-mass Dilepton Enhancement. Leptons produced in a relativistic heavy ion collision escape the interaction region without attenuation by strong interactions, which means they are a probe of phenomena extant in the early phase of the collision. Lepton pair production data collected in relativistic S-Au and Pb-Au collisions at the CERN SpS exhibits an excess with-respect-to proton-nucleus data in the “low-mass” region: $`0.25\mathrm{GeV}<M_{e^+e^{}}<0.70\mathrm{GeV}`$, which is illustrated in Fig. 4.10.
Since the clear excess appears above the $`\pi \pi `$ threshold and below the position of the $`\rho `$-meson peak in proton-nucleus collisions, the search for an explanation has focused on exploring the in-medium properties of the $`\rho `$-meson. Adequate explanations are found in: 1) collisional broadening; i.e., the $`\rho `$ has a shorter lifetime in an hadron-rich medium, and an increase in the $`\rho `$-meson’s width due to in-medium-modified hadron-loop contributions to the $`\rho `$-meson’s self-energy ; and 2) a simple reduction in the $`\rho `$-meson mass . The phenomenological models require large values of temperature \[$`T0.15`$GeV\] and baryon density \[$`\rho 1`$$`2\rho _0`$\] to describe the data. On this domain the $`\rho `$-meson screening masses, calculated in Refs. and described in connection with Fig. 4.8, are increased by $`<25`$% and, while $`\mathrm{\Gamma }_{\rho e^+e^{}}`$ is roughly unchanged, $`\mathrm{\Gamma }_{\rho \pi \pi }`$ is much reduced . These constraints and effects are ignored in contemporary analyses, even though such features can affect photon production rates .
Strangeness Enhancement. The lifetime of a terrestrial QGP can only be of the order: $`\tau 5`$$`10`$fm, which is much too short for weak interactions to be important. Hence strangeness, once produced, can only disappear through $`s`$-$`\overline{s}`$ annihilation. Such events are unlikely unless there is a super-abundance of strangeness. Therefore strangeness carrying reaction products in the debris are a good probe of the conditions created by a relativistic heavy ion collision.
The amount of strangeness produced in collisions can be quantified via the ratio
$$\lambda _s^{AA}\frac{2s+\overline{s}}{u+\overline{u}+d+\overline{d}},$$
(4.2.2)
where $`s+\overline{s}`$, etc., are the mean multiplicities of newly produced valence quark-antiquark pairs at primary hadron level before resonance decays. Experimentally, as illustrated in Fig. 4.11,
$$\lambda _s^{AA}2\lambda _s^{pp}.$$
(4.2.3)
In addition to this global enhancement, specific enhancements in the yields of $`K`$, $`\overline{K}`$, $`\mathrm{\Lambda }`$, $`\overline{\mathrm{\Lambda }}`$, etc., have been observed in CERN experiments . Detailed analyses indicate that quark degrees-of-freedom are necessary to describe the strangeness enhancement, with the suggestion that it may be indicative of QGP formation .
## 5 Toward a Kinetic Description
Hitherto we have described features of cold, sparse QCD: the nature of DCSB, confinement and bound state properties, and then the effect that the intensive thermodynamic parameters chemical potential and temperature have on these phenomena. In the latter we applied the methods of equilibrium statistical field theory and elucidated properties of the QGP phase. The terrestrial creation of this QGP, however, is expected to be effected via relativistic heavy ion collisions, which initially yield a quantum system far from equilibrium. This system must then evolve to form the plasma, and the study of that evolution and the signals that characterise the process are an important contemporary aspect of QGP research.
### 5.1 Preliminaries
The energy density in an ideal gas of eight gluons and two flavours of massless quarks is
$$ϵ_{g+u+d}=(2\times 8)\frac{\pi ^2T^4}{30}+(2\times 3\times 2)\frac{7\pi ^2T^4}{120}=\frac{37\pi ^2T^4}{30}.$$
(5.1.1)
As we saw in connection with Table 3.1, the critical temperature for QGP formation is $`T_c0.15`$GeV and at this temperature: $`ϵ_{g+u+d}=0.8`$GeV/fm<sup>3</sup>. Construction of RHIC at Brookhaven National Laboratory is complete and it will soon provide counter-circulating, colliding $`100`$A GeV $`^{197}`$Au beams to generate a total centre-of-mass energy of $`40`$TeV, which corresponds to an initial energy density: $`ϵ10`$$`100`$GeV/fm<sup>3</sup>. The Large Hadron Collider \[LHC\] project at CERN is scheduled for completion in 2005. Plans are for it to provide <sup>208</sup>Pb-<sup>208</sup>Pb collisions with $`\sqrt{s}>\mathrm{2\hspace{0.17em}000}`$TeV and a consequent initial energy density $`ϵ>\mathrm{1\hspace{0.17em}000}`$GeV/fm<sup>3</sup>. With these energy scales RHIC and LHC should certainly provide the conditions necessary for QGP formation.
Control over the conditions produced in a relativistic heavy ion collision can only be exerted via two experimental parameters: the beam/target properties and, to some extent, the impact parameter. They can be used together to analyse the debris collected in the detectors. It is in the behaviour of this debris, summed over many events, that signals of the evolution and formation of a QGP must be found. Predicting just what the signals are requires an understanding of these processes, including and perhaps especially their the non-equilibrium aspects.
Figure 5.1 illustrates the spacetime evolution of a relativistic heavy ion collision. All stages subsequent to the formation of a thermalised QGP are adequately addressed using hydrodynamical models . However, in this approach the plasma’s initial conditions \[energy density, temperature, etc.\] and EOS must be specified. These initial conditions can only be reliably determined by following the complete evolution of the system produced in the collisions; i.e., by understanding the pre-equilibrium stage. Furthermore, the very existence of a pre-equilibrium phase can lead to signals of QGP formation; e.g., plasma oscillations, a disoriented chiral condensate , and out-of-equilibrium photon and dilepton emission .
The QGP is a hot, equilibrated, many-parton agglomeration, and in recent years two main approaches have been employed in describing how such a system might be produced in a relativistic heavy ion collision. In the perturbative parton picture , the colliding nuclei are visualised as pre-formed clouds of quarks and gluons. The collision proceeds via rapid, multiple, short-range parton-parton interactions, which generate entropy and transverse energy in a cascade-like process. In the string picture , after passing through one another, the colliding nuclei are imagined to stretch a high energy-density flux tube between them, which decays via a nonperturbative particle-antiparticle production process. Each approach has its merits and limitations, and their simultaneous pursuit provides complimentary results. The analysis of data proceeds via one of the many Monte-Carlo event generators that have been developed for both pictures .
Flux Tube Model and Schwinger Mechanism. One intuitively appealing, semi-classical picture of confinement is provided by the notion that it is effected via the formation of a small-diameter colour flux tube between colour sources; and there is evidence in lattice simulations for the appearance of such flux tubes between heavy-quark-antiquark pairs; e.g., Ref. . This is a motivation for the string-like models just introduced, which have been used to study particle production in $`e^{}`$-$`e^+`$, $`p`$-$`p`$ and $`p`$-nucleus collisions . A flux tube yields a linearly-rising, confining quark-antiquark potential: $`V_{q\overline{q}}(r)=\sigma r`$, where the string tension: $`\sigma `$, can be estimated in static-quark lattice simulations. An overstretched flux tube can be viewed as a strong background field. As such it destabilises the QCD vacuum, which is corrected by particle-antiparticle production via a process akin to the Schwinger mechanism in QED . We use this particle production mechanism as the primary medium for our discourse.
The fermion production rate for a constant, homogeneous electric field $`E`$ in a flux tube is
$$S(p_{})=\frac{dN}{dtdVd^2p_{}}=|eE|\mathrm{ln}\left[1+\mathrm{exp}\left(\frac{2\pi (m^2+p_{}^2)}{|eE|}\right)\right],$$
(5.1.2)
where $`m`$ is the mass and $`e`$ the charge of the produced particles, and it is plain from this formula that the rate increases with increasing $`E`$ and is suppressed for large $`m`$ and/or $`p_{}`$. Particle-antiparticle production via this mechanism is analogous to a tunnelling process in quantum mechanics, as illustrated in Fig. 5.2. \[NB. The Schwinger formula follows immediately by integrating Eq. (5.1.2) over the transverse momentum: $`d^2p_{}`$.\]
What happens once the particles are produced? Naturally, they are accelerated in the tube by the field. This produces currents, which generate an electric field that works to screen the flux-tube \[background\] field. In the absence of other effects, the net field vanishes and particle production stops. However, the currents persist, now generating a field and renewed particle production that opposes their own existence. That continues until the net current vanishes. But at this point there is again a strong electric field … and the process repeats itself. This is the “back-reaction” phenomenon and the obvious consequence is time-dependent fields and currents; i.e., plasma oscillations characteristic of the system. In recent years its affects have been studied in both boson and fermion production and we exemplify the process in Sec. 5.3.
There are a number of shortcomings in extant applications of the flux-tube particle-production picture to QGP formation: 1) The non-Abelian nature of the chromoelectric field is often ignored because the QCD analogue of the Schwinger mechanism is poorly understood. Instead the flux tube is represented via a classical electromagnetic potential; 2) On the whole, finite volume effects are neglected, with the electric field assumed to be homogeneous in space. A more realistic description would account for the geometry of the interaction region. Some steps have been taken in this direction, with a consideration of effects produced by a cylindrical boundary and those in a finite-length flux tube with a confining transverse potential . Going further, the geometrical representation of a flux tube can be replaced, allowing the flux tube profile itself to be a dynamical quantity, whose presence and stability is affected by the charged particle currents ; 3) The time-evolution of the system is described using either mean-field theory, which retains quantum effects but makes problematic an exploration of the effects of collisions, or a transport equation , which neglects quantum effects. Below we will describe a partial reconciliation of these approaches; and 4) Little attention has been paid to particular non-Abelian features in transport equations . A study of Wong’s equation is one step in the direction of explicitly including colour algebra effects.
Quantum field theory can be applied directly to out-of-equilibrium plasma phenomena; for example, Refs. . However, kinetic equations provide an appealing alternative because of their intuitive character. This approach begins with a transport equation
$$p^\mu \frac{f}{q^\mu }Qp^\mu F_{\mu \nu }\frac{f}{p_\nu }=S(p,q)+C(p,q),$$
(5.1.3)
where: $`f`$ is the single particle distribution function, which gives the ensemble fraction of particles in a given phase-space cell \[$`p`$ is four-momentum and $`q`$ is a spacetime four-vector\]; $`F_{\mu \nu }`$ is the field strength tensor; and $`Q`$ is the particles’ charge. The source term: $`S(p,q)`$, describes the production of particle-antiparticle pairs and $`C(p,q)`$ is the collision term, whose origin is intuitively obvious but which is difficult to approximate accurately.
A Boltzmann equation of this form has been widely applied to particle production using a source term of the type in Eq. (5.1.2) both with and without a collision term, and in the following we exemplify this. However, to provide an introductory overview, Refs. employ a classical, Abelian electric field in the source term but completely ignore the effect of collisions, $`C0`$ in Eq. (5.1.3). The effect of collisions, represented via a “relaxation time approximation,” \[described in connection with Eq. (5.4.2)\] has been considered in an hydrodynamical approximation to Eq. (5.1.3), Refs , with the influence of back-reactions neglected in the former but explored in the latter. A comparison between the transport equation and mean-field theory approaches has also been made and the results are remarkably similar. That, however, is problematic since; e.g., the application of the Schwinger source term in the presence of a rapidly changing electric field is a priori unjustified and, although quantum field theory with its manifest microscopic time reversal invariance must underly the behaviour all quantum systems, experience confirms that systems far from equilibrium exhibit macroscopically irreversible behaviour that is amenable to treatment using \[inherently time-irreversible\] kinetic theory. The nature of the connection has recently been established .
At this point we re-emphasise that flux-tube models describe the nonperturbative production of soft partons. The production of hard and semi-hard partons is described by perturbative QCD and that mechanism is explored in Refs. . Simultaneously incorporating both types of particle production is challenging but Ref. is a step in that direction. Therein hard and semihard partons are produced via “minijet gluons” and provide the initial conditions necessary to solve the transport equations, and the subsequent evolution of the plasma is described by a classical but non-Abelian transport equation. Collisions are accounted for using a relaxation time approximation but quantum effects in the source term are neglected and, as will become apparent, they can be important in strong fields.
### 5.2 Quantum Vlasov Equation
The derivation of a quantum Vlasov equation in Refs. provides a connection between the quantum field theoretical and transport equation approaches to particle production and plasma evolution, and shows that the particle source term is intrinsically non-local in time; i.e., non-Markovian. Therefore calculating the plasma’s properties at any given instant requires a complete knowledge of the history of the formation process.
The derivation begins with the Dirac \[Klein-Gordon\] equation for fermions \[bosons\] in an external, time-dependent, spatially homogeneous vector potential: $`A_\mu `$, in Coulomb gauge: $`A_0=0`$, taken to define the $`z`$-axis: $`\stackrel{}{A}=(0,0,A(t))`$. The corresponding electric field is
$$E(t)=\dot{A}(t)=\frac{dA(t)}{dt},$$
(5.2.1)
also along the $`z`$-axis. The vacuum instability created by this external field is corrected via particle-antiparticle production, Fig. 5.2, which is a time-dependent process. The transition from the in-state to the instantaneous, quasi-particle state at time $`t`$ is achieved by a time-dependent Bogoliubov transformation, which effects the diagonalisation of the system’s Hamiltonian. By this means one obtains a kinetic equation for the single particle distribution function
$$f(\stackrel{}{P},t)=0|a_\stackrel{}{P}^{}(t)a_\stackrel{}{P}(t)|0,$$
(5.2.2)
which is defined as the vacuum expectation value, in the time-dependent basis, of creation and annihilation operators: $`a_\stackrel{}{P}^{}(t)`$, $`a_\stackrel{}{P}(t)`$, for particle states at time $`t`$ with three-momentum $`\stackrel{}{P}`$. That equation is
$$\frac{df_\pm (\stackrel{}{P},t)}{dt}=\frac{f_\pm (\stackrel{}{P},t)}{t}+eE(t)\frac{f_\pm (\stackrel{}{P},t)}{P_{}(t)}=\frac{1}{2}𝒲_\pm (t)_{\mathrm{}}^t𝑑t^{}𝒲_\pm (t^{})[1\pm 2f_\pm (\stackrel{}{P},t^{})]\mathrm{cos}[x(t^{},t)],$$
(5.2.3)
where the lower \[upper\] sign in Eq. (5.2.3) corresponds to fermion \[boson\] pair creation. The momentum is defined as $`\stackrel{}{P}=(p_1,p_2,P_{}(t))`$, with the longitudinal \[kinetic\] momentum $`P_{}(t)=p_{}eA(t)`$, $`p_{}=p_3`$. \[NB. $`eE(t)=\dot{P}_{}(t)`$, the particle velocity attained via acceleration by the field $`E(t)`$.\] For fermions and bosons the transition amplitudes are
$$𝒲_{}(t)=\frac{eE(t)\epsilon _{}}{\omega ^2(t)},𝒲_+(t)=\frac{eE(t)P_{}(t)}{\omega ^2(t)},$$
(5.2.4)
where the transverse energy $`\epsilon _{}=\sqrt{m^2+\stackrel{}{p}_{}^{\mathrm{\hspace{0.17em}2}}}`$, $`\stackrel{}{p}_{}=(p_1,p_2)`$, and $`\omega (t)=\sqrt{\epsilon _{}^2+P_{}^2(t)}`$ is the total energy. In Eq. (5.2.3),
$$x(t^{},t)=2[\mathrm{\Theta }(t)\mathrm{\Theta }(t^{})],\mathrm{\Theta }(t)=_{\mathrm{}}^t𝑑t^{}\omega (t^{}),$$
(5.2.5)
is the dynamical phase difference. Equation (5.2.3) is the precise analogue of directly solving QED with an external field in mean-field approximation, as done; e.g., in Ref. . The physical content is therefore equivalent and, in particular, the fundamental quantum character is preserved. The appeal of this kinetic equation, however, is that it simplifies: the identification of reliable approximations; widespread applications; and numerical studies.
This quantum Vlasov equation has three qualitatively important features: 1) The source is non-Markovian for two reasons – (i) the source term on the r.h.s. requires complete knowledge of the distribution function’s evolution from $`t_{\mathrm{}}t`$; and (ii) the integrand is a non-local function of time, which is apparent in the coherent phase oscillation term $`\mathrm{cos}x(t^{},t)`$ and reflects the quantum nature of the source term; 2) Particles are produced with a momentum distribution cf. the Schwinger source term, which produces particles with zero longitudinal momentum; and 3) the production rate is affected by the particles’ statistics, as evident in Eqs. (5.2.4) and also in the sign appearing in the statistical factor $`[1\pm 2f_\pm ]`$, which leads to different phase space occupation .
These features can have a material impact on the solution of the kinetic equation, and their importance depends on the field strength and the time-scales characterising the production process. The first time-scale is set by the Compton wavelength of the produced particles:
$$\tau _{\mathrm{qu}}\frac{1}{\epsilon _{}}.$$
(5.2.6)
Events taking place over times less-than this expose the negative-energy elements in particle wave-packets, a core quantum field theoretical feature. This is the time-scale of the rapid oscillations generated by the dynamical phase difference in $`\mathrm{cos}x(t^{},t)`$. The high frequencies involved mean that the main contributions arise when $`tt^{}`$, and therefore a local approximation can be justified for weak field plasmas . We identified a second time-scale in Fig. 5.2; i.e., the time taken for a particle to tunnel through the barrier. It can also be motivated by considering the transition amplitudes in Eqs. (5.2.4), which for a constant electric field can be written in the form; e.g.,
$$𝒲_{}(t)=\frac{\epsilon _{}/eE}{(tp_{}/eE)^2+(\epsilon _{}/eE)^2},$$
(5.2.7)
which is a Lorentzian characterised by the time-scale \[width\]
$$\tau _{\mathrm{tu}}\frac{\epsilon _{}}{eE}.$$
(5.2.8)
This is also the time required to accelerate a charged particle to the speed of light in an electric field.
The times scales in Eqs. (5.2.6) and (5.2.8) are comparable when $`eE\epsilon _{}^2`$. For weak fields, $`\tau _{\mathrm{qu}}\tau _{\mathrm{tu}}`$ and, in the integrand on the r.h.s. of Eq. (5.2.3), the $`\mathrm{cos}x(t^{},t)`$ oscillations occur on a time-scale so-much shorter than variations in the other elements that a stationary phase approximation is valid. That approximation yields a local source term and this important limit of Eq. (5.2.3) was discussed in Ref. .
$`\tau _{\mathrm{tu}}>\tau _{\mathrm{qu}}`$ for strong fields and hence the particle’s Compton wavelength extends across the gap. In this case wave-like \[quantum\] interference effects become important in the behaviour of $`f(p,q)`$, which changes at a rate comparable with $`\mathrm{cos}x(t^{},t)`$, so that a stationary phase approximation is not valid. Can this strong field scenario be relevant to contemporary relativistic heavy ion collisions? It is easy to make an estimate. The field strength in a flux tube is commensurate with the QCD string tension; i.e., $`eE\sigma 0.4`$GeV$`{}_{}{}^{2}(2.5\mathrm{\Lambda }_{\mathrm{QCD}})^2`$, and $`\epsilon _{}2`$$`3\mathrm{\Lambda }_{\mathrm{QCD}}`$ is achievable. Therefore the answer is yes: effects arising from the non-Markovian structure of the source term are likely to be exhibited in contemporary relativistic heavy ion collisions.
This was illustrated in Ref. , wherein a comparison was made between the complete solution and that obtained in the low density limit; i.e., retaining $`\mathrm{cos}x(t^{},t)`$ in the source term but making the replacement $`[1\pm 2f_\pm ]1`$, which is the assumption
$$f_\pm 0,\mathrm{almost}\mathrm{everywhere},$$
(5.2.9)
so that the source term becomes
$$S_\pm ^0(\stackrel{}{p},t)=\frac{1}{2}𝒲_\pm (t)_{\mathrm{}}^t𝑑t^{}𝒲_\pm (t^{})\mathrm{cos}x(t^{},t).$$
(5.2.10)
Using Eq. (5.2.10), the solution of Eq. (5.2.3) is
$$f_\pm ^0(\stackrel{}{p},t)=_{\mathrm{}}^t𝑑t^{}S_\pm ^0(\stackrel{}{p},t^{}),$$
(5.2.11)
and is depicted for fermions in Fig. 5.3.
Particle production begins at $`t=0`$ because, with $`p_{}=0`$, the transition amplitudes in Eq. (5.2.4) are maximal at this point \[the vector potential is zero\], and the distribution function rapidly approaches a Schwinger-like asymptotic value:
$$f^{\mathrm{full}}(t\mathrm{})=\mathrm{exp}(\pi \epsilon _{}^2/eE)=\begin{array}{ccc}e^{\pi /3}& & 0.35(\mathrm{strong})\\ e^{2\pi }& & 0.0019(\mathrm{weak})\end{array}.$$
(5.2.12)
The qualitative features of the results depicted in Fig. 5.3 are easily understood. A fermion, once produced with a certain momentum, retains it because there are no further interactions \[collisions are ignored\]. The number density is greater for strong fields because they produce more particles and this plainly means that the low-density limit will smoothly become invalid with increasing field strength. The low-density limit overestimates the fermion distribution function in strong fields because it eliminates Pauli blocking. The opposite effect is seen for bosons . For very strong fields: $`f_{}1`$, again because of Pauli blocking, but there is no such bound on the boson distribution function. Finally, since $`\sigma /\epsilon _{}^2>3`$ is achievable in contemporary relativistic heavy ion collisions, the low-density limit is quantitatively unreliable.
### 5.3 Back-reactions
The illustration above assumed a constant electric field. Allowing the more realistic case of a time-dependent field introduces another time-scale, namely that characterising the response of the field to the system’s evolution. This brings us to the phenomenon of back-reactions, which have been explored in connection with models in cosmology and recently much in connection with QGP evolution. In both cases the particles produced by the strong background field modify that field: in cosmology it is the time-dependent gravitational field, which couples via the masses, and in a QGP, it is the chromoelectric field affected by the partons’ colour charge.
The effect of feedback is incorporated by solving Maxwell’s equation, which for a spatially homogeneous but time dependent electric field is just
$$\dot{E}(t)=j(t).$$
(5.3.1)
The total field is a sum of two terms: an external field, $`E_{ex}(t)`$, excited by an external current, $`j_{ex}(t)`$, such as might represent a heavy ion collision \[that is a model input\]; and an internal field, $`E_{in}(t)`$, generated by the internal current, $`j_{in}(t)`$, which characterises the behaviour of the particles produced. The internal current has two components : continued spontaneous production of charged particle pairs creates a polarisation current, $`j_{pol}(t)`$, that depends on the particle production rate, $`S(\stackrel{}{p},t)`$; and the motion of the existing particles in the plasma generates a conduction current, $`j_{cond}(t)`$, that depends on their momentum distribution, $`f(\stackrel{}{p},t)`$. The internal field is therefore obtained from
$$\dot{E}_{in}(t)=j_{in}(t)=j_{cond}(t)+j_{pol}(t).$$
(5.3.2)
In mean field approximation the currents can be obtained directly from the constraint of local energy-density conservation: $`\dot{ϵ}=0`$, where
$$ϵ(t)=\frac{1}{2}E^2(t)+2\frac{d^3p}{\left(2\pi \right)^3}\omega (\stackrel{}{p},t)f(\stackrel{}{p},t),$$
(5.3.3)
and the factor of $`2`$ accounts for antiparticles. For bosons the constraint yields
$$\dot{E}(t)=2e\frac{d^3p}{\left(2\pi \right)^3}\frac{p_{}eA(t)}{\omega (\stackrel{}{p},t)}\left[f(\stackrel{}{p},t)+\frac{\omega (\stackrel{}{p},t)}{\dot{\omega }(\stackrel{}{p},t)}\frac{df(\stackrel{}{p},t)}{dt}\right],$$
(5.3.4)
and one easily identifies the currents
$$j_{cond}(t)=2e\frac{d^3p}{\left(2\pi \right)^3}\frac{p_{}eA(t)}{\omega (\stackrel{}{p},t)}f(\stackrel{}{p},t),j_{pol}(t)=\frac{2}{E(t)}\frac{d^3p}{\left(2\pi \right)^3}\omega (\stackrel{}{p},t)S(\stackrel{}{p},t).$$
(5.3.5)
Equation (5.3.2) now yields the internal field. This construction has been used extensively to study back-reactions; e.g., Refs. .
The expression for the polarisation current exhibits the usual short-range divergence associated with charge renormalisation, however, its regularisation and renormalisation is straightforward . That accomplished, Maxwell’s equation assumes the form
$`\ddot{A}^\pm (t)=\dot{E}^\pm (t)=j^{ex}(t)`$
$`g_\pm e{\displaystyle \frac{d^3P}{\left(2\pi \right)^3}\frac{P_{}(t)}{\omega (\stackrel{}{P},t)}\left[f_\pm (\stackrel{}{P},t)+\frac{1}{2}\left\{2\frac{S(\stackrel{}{P},t)}{𝒲_\pm (\stackrel{}{P},t)}\frac{e\dot{E}^\pm (t)P_{}(t)}{4\omega ^4(\stackrel{}{P},t)}\right\}\left(\frac{ϵ_{}}{P_{}(t)}\right)^{g_\pm 1}\right]},`$
where $`g_{}=2`$, $`g_+=1`$, and all fields and charges are understood to be fully renormalised.
The effect of back-reactions can now be illustrated by solving this equation in concert with Eq. (5.2.3). \[Collisions are still neglected: $`C0`$ in Eq. (5.1.3).\] A relativistic heavy ion collision can be mimicked by an impulse profile for the external field :
$$E_{ex}(t)=\frac{A_0}{b}\mathrm{sech}^2(t/b),$$
(5.3.7)
which “switches-on” at $`t2b`$ and off at $`t2b`$, with a maximum magnitude of $`E_{\mathrm{max}}=A_0/b`$ at $`t=0`$. Once this field has vanished only the induced internal field remains to create particles and affect their motion. The calculated field and current profiles are depicted in Fig. 5.4, and the qualitative features are easily understood. The impulse electric field is evident at $`t0`$. It produces particles and accelerates them, and their motion generates an internal current that opposes the impulse field. Shortly after the “collision” the current reaches a short-lived plateau, when the total field vanishes and particle production halts temporarily. At about this time the external, collision-mimicking field “turns-off.” Nevertheless, in its absence, the total field grows in magnitude but now acts in the opposite direction, decelerating the existing particles, causing new particles to be produced and accelerating them in the new direction. The effect of that is clear, after a time the total current must vanish. A pattern is now established and, in the absence of other influences such as collisions or radiation, it repeats itself, reaching a steady state once the wash from the collision-mimicking impulse configuration disappears completely. The oscillations characteristic of a plasma with field-current feedback have now set-in. The oscillation period is the new time scale:
$$\tau _{\mathrm{pl}},\mathrm{the}\mathrm{plasma}\mathrm{period}.$$
(5.3.8)
It decreases with increasing field strength.
One more feature of the results is the high frequency oscillations evident at the current’s peaks and troughs. They are not a numerical artefact and become more pronounced with increasing values of $`eE/m^2`$; i.e., when $`\tau _{\mathrm{qu}}\tau _{\mathrm{tu}}`$. This makes plain that they are a non-Markovian feature and result from interference between effects on the tunnelling and quantum time-scales. Of course, as illustrated via an analogous feature in Fig. 5.5, they disappear if a local approximation to the source term is used because the stationary phase approximation suppresses such interference effects.
One additional observation is important here. Our example employed an electric field whose magnitude, $`A_0`$, is large and hence the tunneling time, $`\tau _{tu}`$, is small, being inversely proportional to $`eE`$. The period of the plasma oscillations, $`\tau _{pl}`$, also decreases with increasing $`A_0`$ but nevertheless, as clear in Fig. 5.4, $`\tau _{tu}\tau _{pl}`$. Thus, in contrast to the effect it has on the production process , the temporal nonlocality of the non-Markovian source term is unimportant to the collective plasma oscillation . This is the reason why kinetic equations with a simple source term of the form in Eq. (5.1.2) are successful in describing plasma oscillations . Whether or not these oscillations are observable in relativistic heavy ion collisions depends on the effect of dissipative processes, which we now discuss.
### 5.4 Collisions and Evolution of the Quark-Gluon Plasma
Thus far we have illustrated the phenomena of quantum particle production in strong fields and field-current feedback. The stationary state is unrealistic because the dissipative processes: “collisions,” that lead to thermalisation have been neglected. \[Collisions can also lead to particle production and effect hadronisation but we defer that discussion.\]
In the presence of collisions the kinetic equation takes the form
$$\frac{df_\pm (\stackrel{}{p},t)}{dt}=S_\pm (\stackrel{}{p},t)+C_\pm (\stackrel{}{p},t).$$
(5.4.1)
A simple and widely-used model for the collision term is
$$C_\pm (\stackrel{}{p},t)=\frac{f_\pm ^{eq}(\stackrel{}{p},t)f_\pm (\stackrel{}{p},t)}{\tau _\mathrm{r}},$$
(5.4.2)
where $`\tau _\mathrm{r}`$ is the “relaxation time” and $`f^{eq}`$ is the thermal equilibrium distribution function
$$f_\pm ^{eq}(\stackrel{}{p},t)=\frac{1}{\mathrm{exp}[\omega (\stackrel{}{p},t)/T(t)]1},$$
(5.4.3)
with $`T(t)`$ a time-dependent temperature profile, discussed on page 5.6. The relaxation time is a fourth time-scale and it is plain that plasma oscillations can only be observable if $`\tau _{\mathrm{pl}}\tau _\mathrm{r}`$; i.e., if collisions take place much-less frequently than oscillations.
In many of the exploratory calculations hitherto undertaken, $`\tau _\mathrm{r}`$ is a constant. That might be argued to be inadequate because the collision time is supposed to characterise a system’s thermalisation, a process which is interdependent with time-evolving quantities such as density and temperature. A more realistic representation might therefore employ a time-dependent $`\tau _\mathrm{r}`$, which is calculated self-consistently as the plasma evolves. Reference , employing collisions in a gluon-minijet plasma, is a step in that direction. However, from another perspective, any sophisticated collision term should be derived from, and justified by, an underlying microscopic theory (see; e.g., Refs. .); and, furthermore, a relaxation time approximation of any sort can only be valid under conditions of quasi-equilibrium, which cannot be justified a priori in the presence of strong fields. A truly realistic collision term will introduce non-Markovian effects in addition to those already present in the source term, and exemplary studies exist in connection with: relativistic heavy ion collisions ; collective effects in nuclear matter ; nuclear fragmentation ; and the damping rates of giant dipole resonances . These studies make clear that even a binary collision approximation, as characteristic of a Boltzmann-Uehling-Uhlenbeck kinetic equation, can be inadequate under extreme conditions; e.g., in the presence of strong fields and/or when the particle density is high. These are precisely the conditions relevant to QGP formation. The patent complexity justifies the use in exploratory, illustrative studies of the simple $`\tau _\mathrm{r}=`$constant relaxation time approximation. Improvements are a contemporary challenge.
As intuition suggests, collisions effect a damping of the distribution functions’ time-dependent structure. That is illustrated in Fig. 5.6, with this particular calculation obtained using an Ansatz for the temperature profile
$$T(t)=T_{eq}+(T_mT_{eq})\mathrm{e}^{t^2/t_0^2},$$
(5.4.4)
where $`T_{eq}=\mathrm{\Lambda }_{\mathrm{QCD}}`$, $`T_m=2T_{eq}`$, $`t_0^2=10/\mathrm{\Lambda }_{\mathrm{QCD}}`$. \[The definition and meaning of $`T(t)`$ is discussed below.\] As evident in the figure, for large relaxation times the plasma oscillations are unaffected, as might be anticipated because this is the collisionless limit. However, for $`\tau _\mathrm{r}\tau _{\mathrm{pl}}`$, the collision term has a significant effect, with both the amplitude and frequency of the plasma oscillations being damped. Finally, there is a value of $`\tau _\mathrm{r}`$ below which oscillations are not possible, just as in the case of an overdamped harmonic oscillator, and the system evolves quickly and directly to thermal equilibrium. The time taken by the plasma to thermalise depends on the ratio $`\tau _{\mathrm{pl}}/\tau _\mathrm{r}`$, being longer for larger values.
Why use the temperature profile in Eq. (5.4.4) and, indeed, what is temperature in a system far from equilibrium? The notion of temperature is introduced via an assumption of quasi-equilibrium at each time $`t`$, which is only truly valid if the fields are not too strong; i.e., only as long as $`\tau _{\mathrm{tu}}\tau _{\mathrm{pl}}`$. This temporally local temperature can be calculated self-consistently with the evolution of the distribution functions, an approach which represents a slight improvement over employing Ansätze, such as Eq. (5.4.4), and, in fact, can provide an a posteriori justification for that expedient. The calculation of $`T(t)`$ can, however, proceed in a number of ways so that the resulting profile is not unique but the differences are only small and quantitative.
To illustrate the definition of a temperature profile we follow Ref. . The total energy density in the evolving plasma is given in Eq. (5.3.3), where the second term is the particle contribution: $`ϵ_f(t)`$, and the particle number density is
$$n(t)=2\frac{d^3p}{\left(2\pi \right)^3}f_{}(\stackrel{}{p},t).$$
(5.4.5)
An intuitive definition of the local temperature is to require that, at each $`t`$, the average energy-per-particle in the evolving plasma is that of a a quasi-equilibrium gas; i.e.,
$$\frac{ϵ_f(t)}{n(t)}=\frac{ϵ^{eq}(t)}{n^{eq}(t)},$$
(5.4.6)
where
$$ϵ^{eq}(t)=2\frac{d^3p}{\left(2\pi \right)^3}\omega (\stackrel{}{p},t)f_{}^{eq}(\stackrel{}{p},t),n^{eq}(t)=2\frac{d^3p}{\left(2\pi \right)^3}f_{}^{eq}(\stackrel{}{p},t).$$
(5.4.7)
With $`f_{}^{eq}`$ given in Eq. (5.4.3), Eq. (5.4.6) is an implicit equation for $`T(t)`$, which must be solved simultaneously with Eqs. (5.3) and (5.4.1.
In Ref. , this system of equations was solved with $`m=\mathrm{\Lambda }_{\mathrm{QCD}}`$, $`e=2`$, and $`\tau _\mathrm{r}=1/\mathrm{\Lambda }_{\mathrm{QCD}}1`$fm, and using two exemplary impulse field configurations: Eq. (5.3.7), with $`b=0.5/\mathrm{\Lambda }_{\mathrm{QCD}}0.5`$fm and either $`A_0^{RHIC}=4\mathrm{\Lambda }_{\mathrm{QCD}}`$ or $`A_0^{LHC}=20\mathrm{\Lambda }_{\mathrm{QCD}}`$. As evident in Fig. 5.7, $`A_0^{RHIC}`$ yields a calculated initial energy-density characteristic of RHIC conditions \[see Sec. 5.1\] while $`A_0^{LHC}`$ gives a very much greater value characteristic of the initial energy-density expected at LHC. For RHIC conditions the energy-density rises rapidly but, after reaching a maximum, decays monotonically. The low density approximation is valid here. For LHC conditions, with an initial energy-density twenty-times larger, the situation is different: the solution obtained in the low density limit is only a qualitative guide to the plasma’s evolution and plasma oscillations are evident on observable time-scales. These oscillations retard the equilibration of the plasma, roughly doubling the time taken cf. RHIC-like conditions. This application makes evident an expected terrestrial hierarchy of time-scales:
$$\tau _{\mathrm{qu}}\tau _{\mathrm{tu}}<\tau _{\mathrm{pl}}\tau _\mathrm{r}.$$
(5.4.8)
In a very intense field the produced particles travel almost at the speed of light and hence the plasma oscillation period can be estimated via the ultrarelativistic formula :
$$\frac{1}{2}\tau _{\mathrm{pl}}^{UR}\frac{\sqrt{n_{\mathrm{max}}m+E_{\mathrm{max}}^2}}{n_{\mathrm{max}}e}2,$$
(5.4.9)
using the input and LHC-like results of Ref. , which yield $`n_{\mathrm{max}}12/\mathrm{\Lambda }_{\mathrm{QCD}}^3`$ from Eq. (5.4.5). That this is a reliable guide for LHC-like conditions is evident in Fig. 5.7, remembering that $`ϵ(t)E^2(t)`$ so that the peaks exhibited in the lower panel are separated by $`\frac{1}{2}\tau _{\mathrm{pl}}`$.
The calculated temperature profile is depicted in Fig. 5.8. Under RHIC-like conditions, an initial temperature of $`T(t=0)0.5`$GeV is reached and the temperature decreases monotonically. The LHC-like source conditions yield an initial temperature twice as large: $`T^{\mathrm{LHC}}(t=0)0.9`$GeV, and the temperature fluctuates in tune with the energy density.
As described above, the exemplary impulse models in Ref. set initial energy density scales appropriate for RHIC and LHC. An improvement over such Ansätze is to calculate the initial conditions, as done; e.g., in Refs. for the case of a gluon minijet plasma. Even with this improvement, however, the results obtained for observable quantities; such as, $`ϵ(t)`$, $`n(t)`$, the temperature and equilibration time, are semi-quantitatively identical and all features are qualitatively unchanged.
It is a uniform prediction that plasma oscillations are present under LHC-like conditions and they may manifest themselves in the dilepton spectrum. The dilepton production rate is (see; e.g., Ref. )
$$\frac{dN_{l^+l^{}}}{dtd^3x}d^3p_1d^3p_2f(E_1)f(E_2)\sigma (M)v_{12},$$
(5.4.10)
where $`f`$ is the distribution function of the charged elements in the plasma, $`\sigma (M)`$ is the cross-section for particle-antiparticle annihilation into a dilepton pair with invariant mass $`M`$ and $`v_{12}`$ is the standard phase-space factor \[a “relative particle-antiparticle velocity”\]. The distribution function responsible for the features in Fig. 5.7 itself exhibits structure and, according to Eq. (5.4.10), that will be transmitted to the dilepton spectrum. The oscillations characterising the pre-equilibrium phase of the plasma could then be exhibited as repeated dilepton bursts. However, this signal will only be detectable if it is strong enough relative to other process, such as Drell-Yan. That quantitative question is currently unanswered.
### 5.5 Hadronisation
The penultimate stage of QGP evolution is hadronisation; i.e., the re-establishment of the conditions that prevail at zero density and temperature. That process in not instantaneous and hence here too a dynamical approach is necessary, one that provides the means of describing phenomena such as bound state formation, and the onset of DCSB and confinement. Consequently the development of a transport theory in which such effects are directly accessible is an important current focus.
Relativistic transport theory and hydrodynamics have long been useful tools in the study of non-equilibrium states of matter at nonzero density and temperature. They have been employed extensively in relativistic nuclear physics, both in the intermediate energy region , where it may be appropriate to neglect quark and gluon degrees of freedom, and in the high energy domain directly relevant to a non-equilibrium QGP . One approach to the derivation of Vlasov-like transport equations can be described as a relativistic generalisation of the Zubarev method . Another is the contour Green function technique , which has been applied formally; e.g., in Refs. , and also in quantitative model studies; e.g., Refs. . The latter is widely used because it admits a systematic definition and exploration of approximations, and will be our focus.
To be concrete, extant phenomenological applications begin with a model Lagrangian density; e.g., that of the Friedberg-Lee or Nambu–Jona-Lasinio model. Such models are chosen because of their efficacy in describing DCSB, and have been employed in studies of: the spacetime evolution of the dressed-quark mass , which characterises the onset of DCSB; and hadronisation via the models’ fermion-antifermion–bound-state interaction terms . In the former case, the onset of DCSB is explored by solving a quantum Vlasov equation coupled with a distribution-function-dependent gap equation, which describes the evolution of the particles’ mass and feeds this information back into the system. One important qualitative result is a parton mass that increases as the system moves toward equilibrium; e.g., Ref. , which results in a dynamical softening of the EOS. This softening can have observable consequences; e.g., slowing the \[hydrodynamical\] expansion of the plasma . These models, however, do not incorporate confinement and hence cannot describe the replacement of coloured degrees of freedom by hadronic matter through the transition.
An indication that the onset of confinement might significantly affect the evolution of the plasma was observed in Ref. , wherein a $`T`$-dependent dressed-quark mass was introduced by fitting the model’s quasi-particle energy density to that determined in lattice simulations. A self-consistent solution of the coupled Vlasov and gap equations then exhibits confinement; i.e., the partons cannot leave the QGP volume because their mass becomes infinite outside this region. The model does not include a parton-antiparton–hadron interaction and hence the confining effect also stops the hydrodynamical expansion. This study only admitted the importance of the scalar piece of the dressed-quark self energy. However, as we saw; e.g., in connection with Fig. 3.5, the vector self energy is more important in QCD’s deconfined domain. The effect it has on plasma evolution has not hitherto been explored.
The incorporation of such qualitatively robust features of DSE studies into the Vlasov equation is an obvious next step . It is complicated for \[at least\] two reasons: 1) The quasi-particle energy is a functional of the scalar and vector dressed-quark self energies, which are both nonzero and momentum dependent. The scalar self energy is small in the QGP. \[It vanishes in the chiral limit.\] However, the significant vector self energy remains ; 2) In the confined domain the dressed-parton $`2`$-point functions violate reflection positivity and hence a single particle distribution function for these excitations cannot be defined. That is as it should be because hadrons are the relevant degrees of freedom in this phase: a kinetic theory based on quarks and gluons is only appropriate in the deconfined QGP phase. As described in connection with Fig. 3.3, DSE models can adequately represent the transition from confined excitation to propagating quasiparticle , an attribute exploited in Ref. in calculating the thermodynamic functions describing a QGP. These observations suggest that a truly realistic description of the quark and gluon distribution functions in a QGP will see them vanish in the vicinity of the phase boundary.
We close this section with an illustration of the concepts introduced. In a QGP the evolution of the dressed-quark distribution function can be described using a Vlasov equation with the quasi-particle energy
$$E_{}(\stackrel{}{p},\stackrel{}{x},t)=p_0[1+\mathrm{\Sigma }_C(\stackrel{}{p},\stackrel{}{x},t)]=\sqrt{|\stackrel{}{p}_{}|^2+B(\stackrel{}{p},\stackrel{}{x},t)^2},$$
(5.5.1)
where $`\stackrel{}{p}_{}=\stackrel{}{p}[1+\mathrm{\Sigma }_A(\stackrel{}{p},\stackrel{}{x},t)]`$ is the rescaled three-momentum and $`B(\stackrel{}{p},\stackrel{}{x},t)=[m+\mathrm{\Sigma }_B(\stackrel{}{p},\stackrel{}{x},t)]`$. Using Eqs. (3.1.17), Eq. (5.5.1) can be rewritten
$$E_{}(\stackrel{}{p},\stackrel{}{x},t)=C(\stackrel{}{p},\stackrel{}{x},t)E(\stackrel{}{p},\stackrel{}{x},t)=C(\stackrel{}{p},\stackrel{}{x},t)\sqrt{|\stackrel{}{p}|^2(\stackrel{}{p},\stackrel{}{x},t)^2+(\stackrel{}{p},\stackrel{}{x},t)^2},$$
(5.5.2)
where $`=A/C`$ and $`=B/C`$. As our exemplar we follow Ref. and employ an instantaneous variant of the $`𝒟_\mathrm{A}`$ model in Sec. 3.2:
$$𝒟(p_{\mathrm{\Omega }_k})=3\pi ^2\eta \delta ^3(\stackrel{}{p}).$$
(5.5.3)
With this interaction and using rainbow truncation, which is akin to a mean-field approximation, the Matsubara sum can be evaluated algebraically because the self energies are $`p_0`$-independent, and the gap equation assumes the form
$$B(\stackrel{}{p},\stackrel{}{x},t)=m+\eta \frac{B(\stackrel{}{p},\stackrel{}{x},t)}{E_{}(\stackrel{}{p},\stackrel{}{x},t)}f_{}(\stackrel{}{p},\stackrel{}{x},t),C(\stackrel{}{p},\stackrel{}{x},t)=1+\frac{1}{2}\eta \frac{C(\stackrel{}{p},\stackrel{}{x},t)}{E_{}(\stackrel{}{p},\stackrel{}{x},t)}f_{}(\stackrel{}{p},\stackrel{}{x},t),$$
(5.5.4)
with $`A=C`$, so that $`=1`$, and $`f_{}=f/C`$. \[NB. $`f`$ appears explicitly here because, with a quasiparticle pole, the Matsurbara sum can be evaluated algebraically.\] Again, in spite of the model’s simplicity, the solution exhibits qualitative features that are characteristic of a realistic dressed-quark $`2`$-point function; e.g., momentum-dependent scalar and vector self energies, and the persistence of this aspect of the solutions in the deconfined domain . In this class of models, for a collisionless plasma, $`f_{}`$, satisfies
$`0=_tf_{}(\stackrel{}{p},\stackrel{}{x},t)`$
$`+{\displaystyle \frac{1}{E(\stackrel{}{p},\stackrel{}{x},t)}}\left\{\left[\stackrel{}{p}+(\stackrel{}{p},\stackrel{}{x},t)\stackrel{}{}_p(\stackrel{}{p},\stackrel{}{x},t)\right]\stackrel{}{}_xf_{}(\stackrel{}{p},\stackrel{}{x},t)(\stackrel{}{p},\stackrel{}{x},t)\stackrel{}{}_x(\stackrel{}{p},\stackrel{}{x},t)\stackrel{}{}_pf_{}(\stackrel{}{p},\stackrel{}{x},t)\right\}.`$
It is important to note that the dressed-quark mass function is momentum-dependent so that in general
$$\stackrel{}{}_p(\stackrel{}{p},\stackrel{}{x},t)\stackrel{}{}_xf_{}(\stackrel{}{p},\stackrel{}{x},t)0.$$
(5.5.6)
\[NB. The Vlasov equation in Ref. is obtained by neglecting this term, since $`\stackrel{}{}_pM0`$ when the interaction is momentum-independent, and setting $`C=1`$.\] Equations (5.5.4) and (5.5) provide a coupled system for the quark’s distribution function and self energies.
As remarked above, the derivation of a transport equation like Eq. (5.5) requires the existence of a quasi-particle mass shell. The instantaneous model ensures that, with $`p^0=E(\stackrel{}{p},\stackrel{}{x},t)`$. In this case $`\frac{f_{}}{p^0}=0`$ and $`\frac{}{p^0}=0`$, hence Eq. (5.5) can be written
$$p^\mu _\mu ^xf_{}+\left\{_\mu ^x_p^\mu f_{}_\mu ^p_x^\mu f_{}\right\}=0,$$
(5.5.7)
using the Minkowski space metric conventions of Ref. , which is directly comparable with the equation described; e.g., in Ref. . However, the momentum dependence of $``$, here driven explicitly via Eqs. (5.5.4), precludes a simple spherical space-volume scaling solution .
The coupled system, Eqs. (5.5.4), (5.5), illustrates some of the complexity to be anticipated in studying the re-emergence of DCSB and confinement in an expanding QGP. Even in this simple model, as in Ref. , proceeding further requires a numerical solution. Such studies are in their infancy but the qualitative feature is plain: the nontrivial propagation characteristics of the dressed-partons will significantly affect $`f`$, and hence QGP evolution, as $`TT_c`$ and the distribution of partons begins to resemble a heat bath.
## 6 Epilogue
Continuum strong QCD is a broad field and we have only presented a snapshot. Nevertheless, our discourse should serve to introduce many of the topics currently occupying practitioners. Since we employed the Dyson-Schwinger equations as our primary medium, it should also provide an update of the progress that has been made with this tool in the last decade and identify the current challenges. The text provides the details. However, highlights of the progress include: the application of a one-parameter model for the dressed-gluon $`2`$-point function in a description of the masses, decays and form factors of light pseudoscalar and vector mesons, and its simultaneous prediction of the critical temperature for quark gluon plasma formation and the properties of in-plasma correlations; significant progress in calculating the baryon spectrum, and leptonic and nonleptonic interactions involving baryons; a unified treatment of chemical potential and temperature, and their effect on the equation of state for a quark gluon plasma and hadron properties; and an incipient understanding of the evolution to equilibrium of a quark gluon plasma, and the dynamical influence of the re-appearance of confinement and dynamical chiral symmetry breaking on the plasma’s expansion and cooling. On the other side, a primary challenge is to comprehend the origin of the infrared enhancement in the kernel of the QCD gap equation that is necessary to ensure dynamical chiral symmetry breaking. In summarising we cannot improve on Ref. : “This programme has a long way to go, but \[we\] hope you are convinced it has come far.”
While the Dyson-Schwinger equations have provided the backbone for our discussion, we have made connections throughout with the results of other methods. Where there is agreement there can be little doubt that the phenomena described are real. That is the rationale underlying the simultaneous pursuit of complementary methods. Disagreement, in fact or interpretation, provides a challenge, which should be met, but also an opportunity for dialogue that should not be missed.
Acknowledgments
We are grateful to M.B. Hecht, P. Maris and P.C. Tandy for careful readings of the manuscript and useful comments. We note too that, while there is no dearth of references, fallibility ensures we have overlooked some contributions. Such omissions are inadvertent. This work was supported by the US Department of Energy, Nuclear Physics Division, under contract no. W-31-109-ENG-38. S.M.S. is grateful for financial support from the A.v. Humboldt foundation.
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# Giant Molecular Outflows Powered by Protostars in L1448
## 1 Introduction
It has long been an open question whether young stars could be the agents of dispersal of their parent molecular clouds through the combined effects of their outflows (Norman & Silk, 1979; Bertoldi & McKee, 1996). The answer to this question depends on whether the outflows have the requisite kinetic energy to overcome the gravitational binding energy of the cloud, as well as the efficiency with which outflows can transfer momentum, in both magnitude and direction, to the surrounding cloud.
For considerations of molecular cloud dispersal, addressing the question of the adequacy of outflow momenta has historically lagged behind determinations of outflow energetics. This is because evaluation of the available energy sources needed to account for the observed spectral linewidths in a cloud is adequate for quantitative estimates of outflow energies. However, in order to address whether the requisite momentum for cloud dispersal exists in a given case requires well-sampled, sensitive, large-scale mapping of sufficiently large areas to encompass entire molecular clouds. Such observing capability has been beyond reach until the last few years, with the implementation of “rapid” or “On-The-Fly” mapping capabilities at large-aperture millimeter telescopes.
The fact that many outflows powered by young stellar objects actually extend well beyond their parent molecular cloud boundaries has been recognized only recently, with the advent of large-scale, narrowband optical imaging surveys that have revealed shock-excited Herbig-Haro objects at parsec-scale separations from their exciting sources (Bally, Devine, & Reipurth, 1996a; Bally, Devine, & Alten, 1996b; Bally et al., 1997; Devine et al., 1997; Eislöffel & Mundt, 1997; Wilking et al., 1997; Gomez, Whitney, & Kenyon, 1997; Gomez, Whitney, & Wood, 1998; Reipurth, Devine, & Bally, 1998) and from equally large-area, sensitive, millimeter line maps that show parsec-scale molecular outflows (Dent, Matthews, & Walther, 1995; Lada & Fich, 1996; Bence, Richer, & Padman, 1996; Bence et al., 1998; O’Linger et al., 1999). The millimeter line maps of parsec-scale flows have been almost exclusively confined to instances of single, well-isolated cases, due to the tremendous confusion of multiple outflows in regions of clustered star formation, such as are found in NGC 1333 (Sandell & Knee 1998; Knee & Sandell 2000), $`\rho `$ Oph, Serpens (White, Casali, & Eiroa, 1995), or Circinus (Bally et al., 1999).
The L1448 dark cloud, with a mass of 100 M over its $``$ 1.3 pc $`\times `$ 0.7 pc extent as traced by C<sup>18</sup>O emission, (Bachiller & Cernicharo, 1986a), is part of the much more extensive (10 pc $`\times `$ 33 pc) Perseus molecular cloud complex, which contains $``$ 1.7 $`\times `$ 10<sup>4</sup> M, at a distance of 300 pc (Bachiller & Cernicharo, 1986b). The two dense ammonia cores within L1448 contain 50 M distributed over a 1 pc $`\times `$ 0.5 pc area (Bachiller & Cernicharo, 1986a; Anglada et al., 1989). The core at V<sub>LSR</sub> = 4.2 km s<sup>-1</sup> contains the Class 0 protostar L1448 IRS 2, while the other core, at V<sub>LSR</sub> = 4.7 km s<sup>-1</sup>, harbors four Class 0 protostars: L1448C, L1448N(A), L1448N(B), and L1448NW. (Barsony et al., 1998; O’Linger et al., 1999; Eislöffel, 2000). The Class I source, L1448 IRS 1, lies close to the western boundary of the cloud, just outside the lowest NH<sub>3</sub> contours in the maps of (Bachiller & Cernicharo, 1986b).
High-velocity molecular gas in L1448 was discovered a decade ago via CO J$`=`$2$``$1 and CO J$`=`$1$``$0 mapping of a $``$ 2 $`\times `$ 6 area centered on L1448C, acquired with 12<sup>′′</sup> and 20<sup>′′</sup> angular resolutions, respectively (Bachiller et al., 1990). Due to its brightness, high-velocity extent ($`\pm `$ 70 km s<sup>-1</sup>), and symmetrically spaced CO bullets, the L1448C molecular outflow has been the object of much study, unlike the flows from its neighbors, the 7<sup>′′</sup> (in projected separation) protobinary, L1448N(A) & (B), just 1.2 to the north, or L1448 IRS 2, 3.7 to the northwest (e.g., (Curiel et al., 1990; Guilloteau et al., 1992; Bally, Lada, & Lane, 1993; Bachiller et al., 1994; Davis et al., 1994; Bachiller et al., 1995; Dutrey, Guilloteau, & Bachiller, 1997)). Although outflow activity in the vicinity of the protobinary had been reported previously, the H<sub>2</sub> and CO flows, driven by L1448N(A) and L1448N(B), respectively (Bachiller et al., 1990; Davis & Smith, 1995), were not recognized as distinct until recently (Barsony et al., 1998). Identification of these flows was aided by noting the position angle of the low-excitation H<sub>2</sub> flow, centered on L1448N(A), to be distinct from the position angle of the CO flow from L1448N(B), defined by the direction of the line joining L1448N(B) with the newly discovered Herbig-Haro object, HH 196 (Bally et al., 1997).
Recent, wide-angle ($``$ 70 field-of-view), narrowband optical imaging of the entire extent of the L1448 cloud has resulted in the discovery of several systems of Herbig-Haro objects, some displaced several parsecs from any exciting source (Bally et al., 1997). In order to investigate the link between high-velocity molecular gas and the newly discovered Herbig-Haro objects, as well as to study the possibility of cloud dispersal via outflows, we acquired new, sensitive, large-scale CO J$`=`$1$``$0 maps of a substantial portion of the L1448 cloud. These new molecular line maps were acquired with the On-The-Fly (OTF) mapping technique as implemented at NRAO’s 12-meter millimeter telescope atop Kitt Peak, Arizona.
## 2 Observations and Data Reduction
The CO J$`=`$1$``$0 maps of L1448 presented in this paper were acquired using the spectral-line On-The-Fly (OTF) mapping mode of the NRAO’s<sup>1</sup><sup>1</sup>1The National Radio Astronomy Observatory is a facility of the National Science Foundation, operated under cooperative agreement by Associated Universities, Inc. 12-meter telescope on 23 June 1997, UT 13<sup>h</sup>53<sup>m</sup> $``$ UT 19<sup>h</sup>25<sup>m</sup>. We stress that the OTF technique allows the acquisition of large-area, high-sensitivity, spectral line maps with unprecedented speed and pointing accuracy. For comparison, it would have taken eight times the amount of telescope time, or nearly a week in practice, to acquire this same map using conventional, point-by-point mapping. Although OTF mapping is not a new concept, given the rigor of the position encoding that allows precise and accurate gridding of the data, the fast data recording rates that allow rapid scanning without beam smearing, and the analysis tools that are available, the 12-meter implementation is the most ambitious effort at OTF imaging yet.
To produce our CO maps of L1448, we observed a 47 $`\times `$ 7 field along a position angle P.A. $`=`$ 135, (measured East from North), centered on the coordinates of L1448 IRS 2 ($`\alpha _{1950}=`$ 3<sup>h</sup> 22<sup>m</sup> 17.9<sup>s</sup>, $`\delta _{1950}=`$ 30 34 41<sup>′′</sup>). The 115 GHz beamwidth was $``$ 55<sup>′′</sup>. We scanned a total of 33 rows at a rate of 50<sup>′′</sup>/s, along P.A. $`=`$ 135, with a row spacing of 12.7<sup>′′</sup>. (The row spacing is determined by the optimum spatial sampling and by the scanning position angle.) We calibrated and integrated on an absolute off position ($`\alpha _{1950}=`$ 3<sup>h</sup> 20<sup>m</sup> 00.0<sup>s</sup>, $`\delta _{1950}=`$ 31 00 00<sup>′′</sup>) at the start of every row. Each row took approximately one minute to scan. Each map coverage took 41 $``$ 56 minutes to complete. We performed six map coverages to attain an RMS in each OTF spectrum of T$`{}_{}{}^{}{}_{R}{}^{}`$ $``$ 0.11 K.
A dual-channel, single-sideband SIS receiver was used for all observations. The backend consisted of 250 kHz and 500 kHz resolution filterbanks, yielding velocity resolutions of 0.65 km s<sup>-1</sup> and 1.3 km s<sup>-1</sup>, respectively. The filterbanks were used in parallel mode, each of the two receiver polarization channels using 256 filterbank channels. The polarization channels were subsequently averaged together to improve signal-to-noise. Only the 250 kHz resolution data were used to produce the maps presented here.
Line temperatures at the 12-meter are on the T$`{}_{}{}^{}{}_{R}{}^{}`$ scale, and must be divided by the corrected main beam efficiency, $`\eta _m^{}`$, to convert to the main-beam brightness temperature scale. For our very extended source, $`\eta _m^{}`$ is approximately 1.0. Since the corrected main-beam efficiency is the fraction of the forward power in the main diffraction beam relative to the total forward power in the main beam plus error beam, contribution from the error beam can make $`\eta _m^{}`$ $`>`$ 1.0. At 115 GHz, the theoretical error beam width is $``$ 17, but the ratio of the error beam amplitude to the main beam amplitude is only 6 $`\times `$ 10<sup>-4</sup>, suggesting contribution from the error beam can be ignored. We used the NRAO standard source, B5 ($`\alpha _{1950}=`$ 3<sup>h</sup> 44<sup>m</sup> 52.4<sup>s</sup>, $`\delta _{1950}=`$ 32 44 28<sup>′′</sup>), to check absolute line temperatures.
The OTF data were reduced with the Astronomical Image Processing Software (AIPS), Version 15JUL95. AIPS tasks specific to OTF data are ‘OTFUV’, which converts a single 12-meter OTF map (in UniPOPS SDD format) to UV (single-dish) format, and ‘SDGRD’, which selects random position single-dish data in AIPS UV format in a specified field of view about a specified position and projects the coordinates onto the specified coordinate system. The data are then convolved onto a grid. OTF data maps were first combined, then gridded into a data cube and baseline-subtracted. Channel maps as well as individual spectra were inspected to ensure good baseline removal and to check for scanning artifacts. Only the first and last rows scanned contained corrupted spectra and were rejected.
## 3 Results
### 3.1 The Extent of High-Velocity CO in L1448
Figure 1a shows the extent of high-velocity blue- and redshifted CO J$`=`$1$``$0 emission found within our OTF map boundaries, outlined by the zig-zag lines. The red contours represent high-velocity redshifted emission integrated over the velocity interval 8.1 km s<sup>-1</sup> $``$ V<sub>LSR</sub> $``$ 17.8 km s<sup>-1</sup>, whereas the blue contours represent high-velocity blueshifted emission integrated over $``$12.1 km s<sup>-1</sup> $``$ V<sub>LSR</sub> $``$ $``$1 km s<sup>-1</sup>. The rest velocities of the two NH<sub>3</sub> cores are 4.2 km s<sup>-1</sup> and 4.7 km s<sup>-1</sup>. The mapped area comprises 47 $`\times `$ 7, corresponding to 4 pc $`\times `$ 0.6 pc at the source, at a distance of 300 pc. Stars indicate the positions of the known Class 0 sources in L1448, which include L1448C, L1448N (the unresolved, 7<sup>′′</sup> separation protobinary L1448N(A) & L1448N(B)), L1448NW (20<sup>′′</sup> northwest of the protobinary), and L1448 IRS 2, at map center. A filled square indicates the position of the Class I Young Stellar Object (YSO), L1448 IRS 1 (Cohen 1980; Eislöffel 2000). Figure 1b, the inset to Figure 1a, shows the scale of the previously mapped region of high-velocity CO, for comparison (from (Bachiller et al., 1990)). This earlier map includes only a small portion of the high-velocity gas associated with Class 0 protostars in L1448. Comparison of Figures 1a and 1b highlights the truly spectacular spatial extent of outflow activity in L1448.
Figure 2 indicates the positions of all the known Herbig-Haro objects (crosses) which are found within our map boundaries, superimposed on the CO map of Figure 1a. Several striking features are evident in Figure 2: (1) There is a $``$15 long blueshifted filament that is connected to both the L1448N/L1448C region and to L1448 IRS 2 in a wishbone-shaped structure which contains HH 197, HH 195A$``$D, HH 196, and culminates in HH 193; (2) There is extended redshifted outflow emission which suggests structure along three separate axes (at P.A.’s$``$129, 150, & 180) directly to the southeast of L1448N/L1448C; (3) In the immediate vicinity of IRS 2, the blueshifted gas peaks on HH 195E, whereas the center of the red peak is along a line drawn through IRS 2 and HH 195A$``$D; (4) There is redshifted emission that peaks $``$9 to the southeast of IRS 2, which lies on a line connecting IRS 2 and HH 193 at P.A.$``$152; (5) There is redshifted emission that peaks on HH 277, which appears to be oriented nearly perpendicular to the long axis of our map; (6) There is blueshifted CO emission associated with the HH 267 knots, which lie at the northwestern edge of our map.
In addition to the OTF map, we also acquired single-point spectra of the CO J$`=`$1$``$0 and <sup>13</sup>CO J$`=`$1$``$0 transitions at four positions, as indicated in Figure 3.
These spectra were used to help determine the appropriate velocity integration limits for the high-velocity CO emission shown in Figures 1a and 2. The mapped linewing emission velocities are indicated by the horizontal arrows for the the blue-shifted emission in Figures 3a & 3c, and for the red-shifted emission in Figures 3b & 3d. These spectra were also used to determine velocity-intervals free of line emission for baseline-subtraction of the gridded OTF data, and to check CO optical depths in the line wings. The spectra shown in Figure 3a were obtained just off the northwest corner of our map, near the HH 267 knots, whereas the spectra of Figures 3b, c, & d were obtained at positions of strong outflow emission within our map boundary. The spectra in Figure 3a show a separate velocity feature at 0 km s<sup>-1</sup> in both CO and <sup>13</sup>CO. The integration limits for the high-velocity gas were chosen conservatively, avoiding emission from the velocity feature at V$`{}_{LSR}{}^{}=`$0 km s<sup>-1</sup>. It is interesting to note that this velocity feature corresponds to the V<sub>LSR</sub> of three HH objects in this cloud: HH 196A, HH 196B, & HH 267B.
### 3.2 Background: Individual Outflow Sources
#### 3.2.1 The L1448C/L1448N Core
The CO outflow from L1448C has been studied extensively since its discovery, at which time it was recognized as a unique source due to its high collimation factor (approaching 10:1) and its extremely high velocities ($`\pm `$70 km s<sup>-1</sup> $``$ (Bachiller et al., 1990, 1995). Interferometric observations of the outflow, acquired with a 3<sup>′′</sup> $`\times `$ 2.5<sup>′′</sup> synthesized beam, were required to resolve the limb-brightened CO cavities at “low” velocities ($``$12 $``$ V<sub>LSR</sub> $``$ $`+`$16 km s<sup>-1</sup>) (Bachiller et al., 1995; Dutrey, Guilloteau, & Bachiller, 1997). By modelling the interferometric CO channel maps, the outflow inclination angle was found to be $`i=70^{}`$, implying actual jet velocities in excess of 200 km s<sup>-1</sup>. The initial conical outflow opening half-angle was found to be $`\varphi `$/2 $`=`$ 22.5 (Bachiller et al., 1995). The outflow cavity walls become parallel, however, $``$ 1 ($`=`$ 0.08 pc) downstream from the driving source, with a width of $``$ 20<sup>′′</sup>. Therefore, in our OTF map, the L1448C outflow remains unresolved along its width.
The redshifted lobe of the L1448C outflow is deflected by $``$ 20 from its initial direction, from an initial position angle of $`+`$160 to a final position angle of $`+`$180 (Dutrey, Guilloteau, & Bachiller, 1997; Eislöffel, 2000). This change in direction of the redshifted flow axis occurs abruptly near the position of the CO “bullet” known as R3 (Bachiller et al., 1990; Dutrey, Guilloteau, & Bachiller, 1997). A string of H<sub>2</sub> emission knots along P.A. $`=180^{}`$ extends for nearly two arcminutes, beginning about an arcminute southeast of L1448C (Eislöffel 2000). Similarly, the blueshifted CO outflow lobe powered by L1448C starts out at a position angle P.A. $`=`$ $`+`$159 (Bachiller et al., 1995), before being deflected through a total angle of $``$ 32 by the time it arrives an arcminute downstream (Davis & Smith 1995; Eislöffel 2000). This deflection is due to the collision of the blueshifted gas driven by L1448C with the ammonia core containing the protobinary L1448N(A) & (B) (Curiel et al., 1999). The strongly radiative shock emission at this interaction region is evident through various tracers: enhanced CO millimeter line emission at the site of the CO bullet “B3” of Bachiller et al. (1990), a far-infrared continuum emission peak (Barsony et al., 1998), high-excitation, shocked molecular hydrogen emission (Davis et al., 1994; Davis & Smith, 1995), and the optically visible shock-excited gas of HH 197 (Bally et al., 1997).
Extrapolating from the vicinity of HH 197 along P.A. $``$ 127, which is the axis of the blueshifted L1448C CO outflow after its deflection near HH 197 through knots I, R, and S (Davis & Smith 1995; Eislöffel 2000), leads directly to HH 267. The measured radial velocities of the HH 267 knots (HH 267A: V$`{}_{LSR}{}^{}=50`$ km s<sup>-1</sup>; HH 267B: V$`{}_{LSR}{}^{}=0`$ km s<sup>-1</sup>; HH 267C: V$`{}_{LSR}{}^{}=63`$ km s<sup>-1</sup> $``$ Bally et al. 1997) and the velocity extent of the blueshifted CO observed by Bachiller et al. (1990) agree well. The reported terminal velocity for the L1448C flow is 70 km s<sup>-1</sup> (Bachiller et al. 1990, 1995). These two facts led to the suggestion that HH 267 may be powered by L1448C (Barsony et al., 1998).
L1448N(A) and L1448N(B) form a close (7<sup>′′</sup> separation) protobinary (Terebey & Padgett 1997), which is unresolved in our OTF map. The redshifted portion of the L1448N(A) flow was first detected via its associated low-excitation molecular hydrogen emission (Davis & Smith, 1995). Its exciting source was first correctly identified by Barsony et al. (1998). The corresponding blueshifted lobe has been detected only in an extended, conical reflection nebulosity whose peak is estimated to lie $`1.^{\prime \prime }5`$ west and 6<sup>′′</sup> north of L1448N(A) (Bally et al. 1993). The position angle of $`150^{}`$ for this flow is determined by the symmetry axis of the U-shaped shocked molecular hydrogen emission (Barsony et al. 1998; see also Eislöffel 2000) and from the axis of the highly-collimated redshifted CO jet driven by L1448N(A) as seen in the V<sub>LSR</sub> $`=`$ $`+`$8 km s<sup>-1</sup> outflow channel map of Bachiller et al. (1995).
Redshifted gas associated with the L1448N(B) flow first appears in the paper reporting the discovery of the L1448C outflow (Bachiller et al. 1990). The corresponding blueshifted outflow lobe from L1448N(B) was partially mapped by Bontemps et al. (1996). Barsony et al. (1998) noted that HH 196, a series of blueshifted optical emission knots (HH 196A: V$`{}_{LSR}{}^{}=0`$ km s<sup>-1</sup>; HH 196B: V$`{}_{LSR}{}^{}=0`$ km s<sup>-1</sup>; HH 196C: V$`{}_{LSR}{}^{}=35`$ km s<sup>-1</sup>; HH 196D: V$`{}_{LSR}{}^{}=37`$ km s<sup>-1</sup> $``$ Bally et al. 1997) lie along the L1448N(B) outflow axis. The position angle of the L1448N(B) outflow is P.A. $``$129.
About 20<sup>′′</sup> northwest of L1448N(A) lies L1448NW. Recent observations suggest that L1448NW drives a small-scale H<sub>2</sub> outflow along an east-west direction (Eislöffel 2000). Although we do not detect a CO outflow associated with L1448NW with our spatial resolution, unpublished interferometric observations do indicate the presence of an E-W flow centered on L1448NW (Terebey 1998).
#### 3.2.2 The L1448 IRS 2 Core
L1448 IRS 2 was confirmed as a Class 0 protostar by O’Linger et al. (1999), who reported a CO outflow associated with this source. The outflow’s symmetry axis along P.A. $``$ 133 and full opening angle, $`\varphi =`$27, were initially inferred from (1) the locations of HH 195A$``$D, (2) a fan-shaped reflection nebulosity emanating from IRS 2 in the K images of Hodapp (1994), (3) the positions of CO “bullets”, detected at the 3$`\sigma `$ level along the outflow axis about 10 to the northwest of IRS 2, and (4) the apparent V-shaped morphology of the blueshifted CO emission. A more recent, H<sub>2</sub> image of this region has led to a refinement of the IRS 2 outflow axis determination to P.A. $``$ 138 (Eislöffel 2000).
### 3.3 Velocity Maps
In order to elucidate the velocity structure of the CO emission, we present Figures 4, 5, & 6. All three figures are presented in rotated coordinates, such that the major axis of our map along P.A. $`=`$ 135 now lies horizontally. Figure 4 is meant to be used as a key to identify the HH objects (crosses) and Young Stellar Objects (YSO’s–open stars for Class 0 protostars, filled square for the Class I protostar, L1448 IRS1) indicated by the same symbols in the CO veolocity channel maps of Figures 5 & 6. Our beamsize is indicated in the lower right-hand corner of each panel.
Figure 5 shows the contoured greyscale images of the blueshifted integrated CO intensities for five velocity intervals blueward of the ambient cloud velocity, proceeding top down from the highest velocities in the top panel, to the lowest, cloud-core velocities in the bottom panel. Figure 6 shows the same for the redshifted emission. In order to obtain good signal-to-noise, our highest velocity channel maps have been integrated over a 4 km s<sup>-1</sup> velocity interval. The other channel maps have been integrated over 2 km s<sup>-1</sup> intervals. These maps shed more light on the emission features enumerated in §3.1, and uncover additional information that is not obvious from the outflow map of Figure 2. It is immediately apparent from these maps that much of the emission within the line core delineates gas that has been entrained in the outflows. In particular, the $`15^{}`$ long blueshifted feature, which extends from IRS 2 to HH 193, is seen to some degree in all of the maps depicting emission blueward of the ambient cloud velocity (Figures 5a$``$e), as well as in the core gas redward of the ambient cloud velocity (Figures 6d & e). Redshifted, as well as blueshifted, CO emission surrounds HH 193. This is not surprising, since the radial velocities and linewidths of HH 193A, B, and C, are -18 & 40 km s<sup>-1</sup>, 10 & 70 km s<sup>-1</sup>, and -10 & 50 km s<sup>-1</sup>, respectively (Bally et al. 1997). This significant overlap of blueshifted and redshifted emission strongly suggests that the blueshifted feature is oriented close to the plane of the sky.
The blueshifted feature is part of a longer feature, which is bipolar about IRS 2 at P.A. $``$ 152. The blueshifted emission is visible as a well-defined structure from V$`{}_{LSR}{}^{}=`$-8 to 0 km s<sup>-1</sup>. Similarly, redshifted emission is visible as a well-defined structure from V$`{}_{LSR}{}^{}=`$8 to 16 km s<sup>-1</sup>. At the highest velocities (Figures 5a & 6a), the blue- and redshifted emission is highly collimated along P.A. $`152^{}`$, centered on L1448 IRS 2. However, the blueshifted emission intersects a second blueshifted feature emanating from the L1448N/L1448C region along P.A. $``$ 129 (Figure 5a). The intersection occurs $``$2 downstream of the HH 196 knots. A third blueshifted feature branches off to the west $``$3 downstream from the HH 196 knots (Figure 5b). The redshifted emission along P.A. $``$ 152 from IRS 2 extends $`9^{}`$ southeast of IRS 2 (Figures 6a$``$c), where it shows a prominent peak in Figure 6b, as well as in the CO total integrated intensity outflow map in Figure 2.
The blue- and redshifted peaks adjacent to L1448 IRS 2 are most prominent in the lowest-velocity outflow emission (Figures 5c & 6c). The blueshifted peak is spatially coincident with HH 195E. At higher velocities (Figures 5a & b), the blueshifted emission extends along a P.A. $``$ 125 from IRS 2, past IRS 1, and may continue past HH 194, which shows a local peak in blueshifted emission (Figures 5b & c). Curiously, the HH 194 knots are all redshifted, with very large linewidths (HH 194A: V$`{}_{LSR}{}^{}=`$66 km s<sup>-1</sup>, $`\mathrm{\Delta }`$V$`=`$110 km s<sup>-1</sup>; HH 194B: V$`{}_{LSR}{}^{}=`$66 km s<sup>-1</sup>, $`\mathrm{\Delta }`$V$`=`$150 km s<sup>-1</sup>; HH 194C: V$`{}_{LSR}{}^{}=`$110 km s<sup>-1</sup>, $`\mathrm{\Delta }`$V$`=`$60 km s<sup>-1</sup>). HH 194C has been associated with redshifted outflow emission from IRS 1 (Bally et al. 1997). Although there is a prominent local peak in the blueshifted CO emission at the position of HH 194, there is no evidence of redshifted CO emission (Figures 6a$``$c) at this position.
The extended redshifted emission directly to the southeast of L1448N/L1448C is also clearly present at core velocities, very strongly so in Figure 6d, and even in Figure 5e. The previously suggested structure along three separate axes is also seen in Figures 6a$``$c; mostly strikingly, in Figure 6b. These three axes, along P.A. $``$ 180, 150, & 129, correspond to the PAs of the redshifted lobes of the outflows associated with L1448C, L1448N(A), & L1448N(B), respectively. The redshifted feature that peaks on HH 277 is most prominent as a separate velocity feature in Figure 6c.
In the next section, we consider these features in connection with what is known about the individual outflows in order to interpret the outflow morphology in L1448.
## 4 DISCUSSION
### 4.1 Interpretation of Outflow Structure
Figures 7$``$9 present the various outflow extents and position angles. Figure 7 is relevant for the discussion of alternative interpretations of the outflow emission centered on the L1448 IRS 2 ammonia core. Figure 8 is used for the discussion of the outflows originating from sources embedded in the ammonia core associated with L1448C and L1448N. For both Figures 7 and 8, the outflow axes and extents are superposed on our integrated high-velocity CO linewing map of the L1448 cloud. Figure 9 shows the same outflow axes superposed on a much higher spatial-resolution ($``$ 1<sup>′′</sup> vs. $``$ 55<sup>′′</sup>) H<sub>2</sub> image of the L1448 region from Eislöffel (2000). (This is the only figure using J2000 coordinates.) In all cases (except for the two newly-identified outflow features seen in our CO maps associated with L1448 IRS 2 and the unidentified source outside our map boundaries), outflow position angles were determined from previously-published, arcsecond-scale outflow data. For all of the outflows, we find good agreement between large-scale CO features seen in our maps and flow axes that have been determined from the previous, higher-resolution observations. In Figures 7$``$9, solid, colored lines denote well-established flow position angles and extents, derived from our own CO data and the published literature, whereas dashed lines denote outflow position angles and extents that are consistent with our new CO data.
#### 4.1.1 Outflow Emission from the L1448 IRS 2 Ammonia Core
High-velocity CO outflow activity centered on L1448 IRS 2 was discovered from low-spatial resolution ($``$ 55<sup>′′</sup>) mapping (O’Linger et al. 1999). These authors suggested that L1448 IRS 2 was the source of a single outflow, with a constant opening angle, as depicted in Figure 7a. The previous outflow symmetry axis along P.A. $``$ 133 and opening angle, $`\varphi 27^{}`$, were derived from the high spatial resolution K images of Hodapp (1994). We derive a new outflow symmetry axis along P.A. $``$ 138 from the more recent H<sub>2</sub> images of Eislöffel (2000). Therefore, we have been able to determine more accurately the position angles for the proposed outflow cavity walls, which should lie along P.A. $``$ 152 and P.A. $``$ 125, if the IRS 2 outflow retains its initial opening angle out to large distances. This model explains the presence of high-velocity blue- and redshifted CO emission seen along these position angles at large distances from IRS 2, notably the V-shaped morphology of the blueshifted gas and the presence of several CO “bullets” located along the proposed outflow axis, well beyond HH 195A$``$D. In a constant opening angle outflow scenario, HH 193 lies along one arm of this V, at P.A. $``$ 152. The blueshifted gas along the other arm of the V (P.A. $``$ 125) would be confused with emission from the E-W outflow in this vicinity, but since the blueshifted emission extends past IRS 1 towards the redshifted HH 194 knots, at least part of this emission could be due to IRS 2. Although there is redshifted CO emission along P.A. $``$ 138 and along P.A. $``$ 152, there is little evidence for redshifted emission along P.A. $``$ 125. However, this may be due to confusion with the three redshifted lobes associated with L1448C, L1448N(A), & L1448N(B).
Figure 9 clearly shows an outflow associated with IRS 2 along a P.A. $``$ 138. The redshifted gas along the outflow axis is prominent in the H<sub>2</sub> emission, but is not clearly apparent in our CO maps beyond the redshifted peak about 1 southeast of IRS 2. Hints of more extended emission along P.A. $``$ 138 may be seen, however, in Figures 6b & c, and curiously, in blueshifted emission extended along this axis to the southeast of IRS 2 in Figure 5d. Such overlap of blueshifted gas along the redshifted outflow axis is expected for outflows oriented nearly in the plane of the sky. Indications of extended blueshifted emission along P.A. $``$ 138, at least a few arcminutes downstream of HH 195A$``$D, are seen in Figures 5b & c.
However, there are a few problems with the single outflow, constant opening angle model: (1) Figure 9 indicates that although the initial opening angle of the IRS 2 outflow is $`\varphi 27^{}`$, the two strands of H<sub>2</sub> defining this opening angle join in a bow shock structure in the vicinity of HH 195A-D; (2) The CO data show little evidence for emission along the cavity wall at P.A. $``$ 125, although we note that there is much confusion from high-velocity gas associated with other outflows along this position angle to both sides of IRS 2; (3) HH 193 lies precisely at the end of the outflow wall in this model, an unlikely location for a shock; (4) The velocity dispersion of the blueshifted feature along P.A. $``$ 152 is high, since it is prominent at both ambient cloud velocities and at highly blueshifted velocities (Figures 5, 6d&e); (5) The highest-velocity outflow emission should converge towards the outflow axis (P.A. $``$ 138) due to projection effects, but the highest-velocity outflow emission lies along P.A. $``$ 152, the proposed outflow wall.
Our maps are not sensitive enough to have picked up the highest-velocity outflow emission, however, which is severely diluted in our large ($``$1) beam. Nevertheless, the absence of CO emission along P.A. $``$ 138 in the highest velocities suggests that the feature along P.A. $``$ 152 defines a separate outflow axis. This has led to an alternate interpretation of the high-velocity CO associated with L1448 IRS 2 in which the presence of two outflows is required, as depicted in Figure 7b, with one outflow along P.A. $``$ 138, and a new, second outflow along P.A. $``$ 152, so prominent in the CO data. In this scenario, the new IRS 2 outflow would be responsible for exciting HH 193. Two outflows, along distinctly different position angles, would also suggest that IRS 2 is a binary system. Although there is currently no evidence to indicate this source to be binary from available continuum data obtained with the Submillimetre Common User Bolometer Array (SCUBA) at the JCMT on Mauna Kea, Hawaii (O’Linger et al. 1999), it is possible that IRS 2 is a compact binary on a scale smaller than 7<sup>′′</sup> (the resolution of the SCUBA 450 $`\mu `$m data). Recent work indicates a high incidence of binarity among young stellar systems (eg., Ghez, Neugebauer, & Matthews 1993; Looney, Mundy, & Welch 2000). Only arcsec/sub-arcsec imaging at either millimeter or centimeter wavelengths could test the binary hypothesis further.
Evidence of other outflow activity near L1448 IRS 2 is found by noting that the very confined outflow, whose lobes peak only $``$1 on either side of IRS 2, has its redshifted peak well-aligned along the P.A. $`=`$ 138 outflow that excites HH 195A-D, but the corresponding blue peak closest to IRS 2 is skewed at a somewhat shallower position angle closer to P.A. $``$ 125. No CO emission (Figure 2, Figures 6a$``$c) is seen along P.A. $``$ 125 on the opposite side from IRS 2, which would be expected if there were an outflow along this direction. The blue peak is spatially coincident with HH 195E, however, the only HH 195 knot which is off the P.A. $`=`$ 138 axis of the IRS 2 outflow. It has been argued that IRS 1 drives an east-west oriented outflow and is the most probable driving source of HH 195E (Bally et al. 1997; Eislöffel 2000). Although the presence of such an E-W oriented outflow in this region is undisputed, it is possible that an as yet undiscovered source, other than IRS 1, may be the responsible agent. Thus, the positioning of the blue peak so close to IRS 2 may be coincidental, and due primarily to local heating associated with HH 195 E, and overlapping outflows associated with blueshifted emission from IRS 2 and IRS 1. This picture is supported by the L1448 H<sub>2</sub> mosaic of Eislöffel (2000), shown in Figure 9, which shows the H<sub>2</sub> emission in the vicinity of HH 195E pointing toward IRS 1, not IRS 2.
#### 4.1.2 Outflow Emission from the L1448C/L1448N Ammonia Core
The position angles of the blue- and red-shifted outflowing gas powered by L1448C are indicated by the green lines in both Figures 8 & 9. Solid green lines indicate the L1448C outflow’s direction and extent as determined by previous workers (see the caption of Figure 8 for references). The dashed green line on the blue-shifted side represents the continuation of the L1448C outflow proposed by Barsony et al. (1998). The blueshifted L1448C outflow suffers a large deflection from its original direction, as can be seen clearly in Figure 9, which shows the blueshifted L1448C outflow axis passing directly through H<sub>2</sub> emission knots I & S, and about 20<sup>′′</sup> north of emission knot R. Emission knot R lies within more extended H<sub>2</sub> emission that appears to form a U-shaped or bow-shock structure which opens toward the southeast, bisected by the L1448C outflow axis. The sides of the U are separated by $``$ 40<sup>′′</sup>, with the brighter side (including knot R) lying to the south of the L1448C outflow axis.
This final, deflected, blue-shifted outflow axis lies along P.A. $``$ 127, where a long finger of high-velocity blueshifted CO emission is found. Blueshifted CO emission surrounds the HH 267 complex in a horseshoe shape about the proposed extension of the L1448C outflow, suggesting the outflow may be responsible for this emission. The dashed green line on the redshifted side in Figure 8 indicates the possible continuation of the L1448C outflow to the south.
The L1448N(A) molecular outflow axis and extent are depicted by the purple line in Figures 8 & 9. Only redshifted molecular gas associated with the L1448N(A) outflow, along P.A. $``$ 150, is detected in our CO maps, with a total length of at least 0.7 pc, as seen in Figure 8. The U-shaped molecular hydrogen emission that traces part of the redshifted outflow cavity wall from the L1448N(A) outflow, as seen in Figure 9, is unresolved in our CO maps. Knots of H<sub>2</sub> emission which trace the outflow wall are $``$30<sup>′′</sup> apart as far as 3 to the southeast of L1448N(A) (Barsony et al. 1998; Eislöffel 2000). Taken together with the observed length of the CO redshifted lobe, this suggests a collimation factor of at least 16:1. The opening half-angle of the L1448N(A) outflow was estimated to be $`\varphi /225^{}`$, from the morphology of the near-infrared reflection nebulosity to the north of L1448N(A), associated with what would be the blueshifted outflow lobe (Bally et al. 1993). The blueshifted flow powered by L1448N(A) is not apparent in our CO maps, judging by the drop in the blueshifted, high-velocity CO contour levels along the symmetry axis of its NIR reflection nebula. The lack of blueshifted CO emission from the L1448N(A) outflow is most likely accounted for by the likelihood that the cloud boundary has been reached in this direction, and that the flow has broken out of the molecular cloud.
The L1448N(B) molecular outflow axis and extent are depicted by the mustard-colored lines in Figures 8 & 9. The position angle, P.A. $``$ 129, of the L1448N(B) outflow was determined by the orientation of the redshifted CO outflow driven by L1448N(B) from Figure 1b and noting that this CO flow symmetry axis intersects HH 196 (Barsony et al. 1998). The true spatial extent of the L1448N(B) CO outflow, however, is demonstrated here for the first time. On the scale of our map, the L1448N(B) outflow remains unresolved along its width. The optical emission knots of HH 196 are, indeed, found to lie right along the L1448N(B) CO outflow axis, confirming the identification of L1448N(B) as their exciting source (Figure 8). Approximately 2 downstream of HH 196, the L1448N(B) flow becomes confused in projection with the P.A. $``$ 152 outflow from IRS 2. If the outflow continues along P.A. $``$ 129, it could account for the bulges in the blueshifted emission to the south and southwest of HH 193 (Figure 2, Figure 5c, Figure 8), and might even be responsible for the “C”–shaped blueshifted emission structure east of the HH 267 system, $``$2.5 parsecs from L1448N(B). The redshifted lobe associated with this source appears to terminate $``$2 northwest of HH 277, although this could be due to confusion with the high-velocity CO emission surrounding HH 277, which seems to be part of an outflow driven by an unidentified source outside of the boundaries of our map. Using the most conservative length for the L1448N(B) outflow, the major axis taken from the HH 196 knots through the end of the redshifted lobe ($``$12), the derived lower limit for the collimation factor is $``$12:1. Estimates of the initial opening angle and width of the L1448N(B) outflow await higher spatial resolution, interferometric imaging.
A dark blue dashed line in Figure 8 indicates a possible outflow axis passing through an otherwise unexplained, high-velocity redshifted feature associated with HH 277, in the southeast quadrant of our CO map. This redshifted CO velocity feature is most prominently seen in Figure 6c. The orientation of this structure is almost perpendicular to the general orientation of our map. Thus, HH 277 is probably driven by a source off the edge of our map.
Finally, although the origin of the HH 267 knots cannot definitively be resolved based on the CO data we present here, our observations do constrain the driving source. L1448N(A) & (B) can be ruled out as possible driving sources of HH 267, since the P.A.’s of their associated outflows are along completely different directions than the lines linking them with HH 267. Furthermore, blue-shifted molecular gas has yet to be detected from L1448N(A). Terminal velocities for the L1448 IRS 2 outflows can not be determined from our data, but the P.A.’s of both of these outflows also miss the HH 267 complex completely. However, the P.A. of the deflected blueshifted lobe of the L1448C flow goes right through HH 267, and the reported terminal velocity for this outflow (70 km s<sup>-1</sup>: Bachiller et al. 1990, 1995) is in good agreement with the measured HH 267 velocities (Bally et al. 1997), as suggested by Barsony et al. (1998).
### 4.2 Cloud Dispersal by Giant Protostellar Flows?
The most dramatic evidence for the direct effects of the outflows on the L1448 molecular cloud is seen in the distortions of the cloud contours at all velocities in Figures 5 & 6. Although we cannot estimate the masses and energetics of each individual outflow in our maps due to confusion in space and velocity, we can, nevertheless, estimate the total contribution of the outflows to the cloud’s energetics. The Local Thermodynamic Equilibrium (LTE) analysis used to estimate the combined mass of the outflows (M$`{}_{tot}{}^{}0.7`$ M) is discussed in O’Linger et al. (1999). Optically thin high-velocity CO emission was assumed, given the lack of observed high-velocity <sup>13</sup>CO emission in the velocity intervals outside the cloud core velocities (see Figure 3). Therefore, the resultant derived mass is a strict lower limit, since no attempt was made to correct for the considerable mass expected to be masked by the line core emission.
For highly inclined outflows ($`i>70^{}`$), the characteristic velocity which is used to calculate outflow energetic parameters is best chosen as the geometrical mean between the highest observed velocities, $`V_{CO}`$, and the inclination-corrected velocity, $`V_{CO}/cos(i\frac{\varphi }{2})`$, where $`\frac{\varphi }{2}`$ is the half-opening angle of the outflow (Cabrit & Bertout 1992). Assuming the outflow inclinations, 70 $``$ $`i`$ $``$ 90, V$`{}_{char}{}^{}=2234`$ km s<sup>-1</sup> for the L1448 outflows, the total momentum in all the flows is computed to be 16 M km s<sup>-1</sup> $``$ M<sub>tot</sub>V<sub>char</sub> $``$ 24 M km s<sup>-1</sup>. This range of values is nearly equivalent to the momentum content of the quiescent NH<sub>3</sub> cores, assuming 50 M total cores and $`v_{turb}`$ $``$ 0.5 km s<sup>-1</sup> (Bachiller & Cernicharo, 1986a). Even more striking, the total kinetic energy in all the flows, 2$`\times 10^{45}`$ ergs $``$ $`\frac{1}{2}`$M<sub>tot</sub>V$`{}_{}{}^{2}{}_{char}{}^{}`$ $``$ 8$`\times 10^{45}`$ ergs, exceeds the gravitational binding energy ($`GM^2/R`$ 5 $`\times `$ 10<sup>44</sup> ergs) of the NH<sub>3</sub> cores by an order of magnitude, and the gravitational binding energy of the 100 M C<sup>18</sup>O cloud (9 $`\times `$ 10<sup>44</sup> ergs), contained within a 1.3 pc $`\times `$ 0.7 pc region (Bachiller & Cernicharo, 1986b), by a factor of five.
The total outflow momentum quoted above is, in fact, a lower limit, since these outflows are still gaining momentum from the force provided by the central driving engines, and the total outflow mass may be grossly underestimated. The magnitude of both the total energy and momenta of the outflows suggests these outflows are capable of dispersing the NH<sub>3</sub> cores, with the caveat that it is unclear, both from the outflow and ambient cloud morphology, how the outflow momenta can be adequately transferred to the surrounding core. Possibly, this can be accomplished as the individual outflow opening angles increase with time.
## 5 Summary
$``$ Spectral-line “on-the-fly” mapping was used at the NRAO 12-meter millimeter telescope to produce a large-scale (47 $`\times `$ 7) CO (J$`=`$1$``$0) map of the L1448 dark cloud, sensitive enough (1 $`\sigma `$ $`=`$ 0.1K) to enable the detection of outflow activity on parsec-scales and the identification of six distinct molecular outflows. Large-scale, high-spatial resolution optical and near-infrared images of shocked gas emission regions associated with each outflow were crucial for identifying the CO counterparts of these outflows. Three of the outflows are associated with the Class 0 protostars, L1448C, L1448N(A), & L1448N(B). Two outflows are associated with the Class 0 protostar, L1448 IRS 2. A sixth outflow, apparently associated with HH 277, is probably driven by an unidentified source located outside of our map.
$``$ For all of the outflows which have previously been identified through high-resolution interferometric observations or small-scale shocked gas emission, we find good agreement between large-scale CO features and flow axes that have been determined from these higher-resolution observations.
$``$ We find evidence of two distinct outflows emanating from the recently-confirmed Class 0 protostar, L1448 IRS 2 (O’Linger et al. 1999), suggesting that IRS 2 is an unresolved binary system. One of these outflows lies along P.A. $``$ 138, and is apparent both in H<sub>2</sub> emission (Eislöffel 2000) and, to a lesser degree, in our CO data. The second outflow lies along P.A. $``$ 152 and is seen as a highly-collimated jet in our CO data, culminating, on the blueshifted side, in HH 193.
$``$ The ambient cloud emission contours are severely disturbed by the outflows, suggesting a large fraction of the ambient cloud in the mapped region has been entrained in, or stirred up by, the outflows.
$``$ The total outflow kinetic energy ($`>`$ 2 $`\times `$ 10<sup>45</sup> ergs) and combined outflow momenta ($`>`$ 16 M km s<sup>-1</sup>) indicate that the outflows are energetically, and probably dynamically, capable of dispersing the dense ammonia cores out of which the protostellar outflow driving sources of five of the six identified flows, L1448C, L1448N(A), L1448N(B), and L1448 IRS 2, are currently forming. However, it is unclear, from both the outflow and ambient cloud morphology, how the outflow momenta can be adequately transferred to the surrounding cores.
ACKNOWLEDGEMENTS: We thank Dr. Darrel Emerson, Dr. Eric Greisen, and Dr. Jeff Mangum of NRAO for the development, implementation, and improvement of the spectral-line On-The-Fly mapping capability of the 12-meter telescope. GWC, JO, and MB gratefully acknowledge financial support from NSF grant AST-0096087 while part of this work was carried out. Part of this work was performed while GWC held a President’s Fellowship from the University of California. MB’s NSF POWRE Visiting Professorship at Harvey Mudd College, NSF AST-9731797, provided the necessary time to bring this work to completion. JO acknowledges financial support by the NASA Grant to the Wide-Field Infrared Explorer Project at the Jet Propulsion Laboratory, California Institute of Technology. We would like to thank our referee, John Bally, for his many helpful suggestions which greatly improved this paper.
a. Blue contours indicate high-velocity blueshifted ($`12.1V_{LSR}1`$ km s<sup>-1</sup>) emission and red contours indicate high-velocity redshifted ($`+`$8.1 $`V_{LSR}`$+$`17.8`$ km s<sup>-1</sup>) CO J$`=`$1$``$0 emission in the 47$`{}_{}{}^{}\times 7^{}`$ region we mapped with the NRAO 12-meter telescope. Contour levels start at 2 K km s<sup>-1</sup> ($``$ 3 $`\sigma `$), and increase in 1.5 K km s<sup>-1</sup> intervals. For comparison, the dashed rectangle shows the approximate area previously mapped (in CO J$`=`$2$``$1), shown in the inset. Stars indicate the positions of L1448C, L1448N, L1448NW, and L1448 IRS 2. All of these are Class 0 sources; L1448N is a 7<sup>′′</sup> separation protobinary, consisting of L1448N(A) and L1448N(B). L1448NW lies $``$ 20<sup>′′</sup> to the northwest of the protobinary. The filled box indicates the position of the Class I source, L1448 IRS 1. The 55<sup>′′</sup> FWHM NRAO beamsize is indicated in the lower right-hand corner. b. The previous CO J$`=`$2$``$1 IRAM 30-meter map of high-velocity gas in L1448 (from (Bachiller et al., 1990)): Solid contours indicate high-velocity, blueshifted ($``$ 55 km s<sup>-1</sup> $``$ V<sub>LSR</sub> $``$ 0 km s<sup>-1</sup>) gas, and dotted contours indicate high-velocity, redshifted ($`+`$ 10 km s<sup>-1</sup> $``$ V<sub>LSR</sub> $``$ $`+`$ 65 km s<sup>-1</sup>) gas. First contour and contour intervals are at 10 K km s<sup>-1</sup>. In addition to the famous protostellar outflow powered by L1448C, the weaker outflow, powered by L1448N(B), is also detected . The 12<sup>′′</sup> FWHM IRAM beamsize is indicated in the lower right-hand corner.
Names and positions of all the Herbig-Haro objects (black crosses) are shown, as well as the five Class 0 sources (black stars) and the Class I source, L1448 IRS 1 (solid black box). Velocity intervals and contour levels are the same as in Figure 1a.
<sup>12</sup>CO (thin solid line) & <sup>13</sup>CO (thick solid line) J=1$``$0 spectra obtained at the four positions on our map whose B1950 coordinates are indicated: (a) CO spectra near the HH 267 knots (just off the northwest corner of our map) clearly show two separate velocity components: a brighter component at the same velocity as HH 267B, V$`{}_{LSR}{}^{}=`$0 km s<sup>-1</sup> (thin vertical dashed line), and a dimmer component at V$`{}_{LSR}{}^{}=`$4.25 km s<sup>-1</sup> (thick vertical dashed line). (b) $``$ (d) CO spectra show strong CO self-absorption at the ambient cloud velocity, V$`{}_{LSR}{}^{}=`$4.25 km s<sup>-1</sup>. Arrows indicate the velocity ranges used to determine integrated intensities, masses, and energetics for the outflow emission.
The positions of all the Herbig-Haro objects are indicated by tilted black crosses; the Class I source, L1448 IRS 1, is shown by a filled red box; the Class 0 protostars, L1448 IRS 2, L1448C, L1448N(A) & (B), and L1448NW are indicated by red stars. The separation between L1448C & L1448N is $``$ 80<sup>′′</sup>.
Contoured greyscale images of the integrated CO intensity over velocity intervals blueward of the ambient cloud emission. Velocity intervals, lowest contour levels, and contour intervals, are, respectively: (a) -8 to -4 km s<sup>-1</sup>, 1 K km s<sup>-1</sup>, 0.5 K km s<sup>-1</sup>; (b) -4 to -2 km s<sup>-1</sup>, 1 K km s<sup>-1</sup>, 0.5 K km s<sup>-1</sup>; (c) -2 to 0 km s<sup>-1</sup>, 1.5 K km s<sup>-1</sup>, 1 K km s<sup>-1</sup>; (d) 0 to 2 km s<sup>-1</sup>, 2 K km s<sup>-1</sup>, 2 K km s<sup>-1</sup>; and (e) 2 to 4 km s<sup>-1</sup>, 2 K km s<sup>-1</sup>, 2 K km s<sup>-1</sup>. Positions of Herbig-Haro objects (tilted yellow crosses), IRS 1 (filled red box), and Class 0 objects (red stars) are indicated.
Contoured greyscale images of the integrated CO intensity over velocity intervals redward of the ambient cloud emission. Velocity intervals, lowest contour levels, and contour intervals, are, respectively: (a) 12 to 16 km s<sup>-1</sup>, 1 K km s<sup>-1</sup>, 1 K km s<sup>-1</sup>; (b) 10 to 12 km s<sup>-1</sup>, 1 K km s<sup>-1</sup>, 1 K km s<sup>-1</sup>; (c) 8 to 10 km s<sup>-1</sup>, 2 K km s<sup>-1</sup>, 1 K km s<sup>-1</sup>; (d) 6 to 8 km s<sup>-1</sup>, 2 K km s<sup>-1</sup>, 2 K km s<sup>-1</sup>; and (e) 4 to 6 km s<sup>-1</sup>, 2 K km s<sup>-1</sup>, 2 K km s<sup>-1</sup>. Positions of Herbig-Haro objects (tilted yellow crosses), IRS 1 (filled red box), and Class 0 objects (red stars) are indicated.
a. One outflow constant-opening angle model (O’Linger et al. 1999). The outflow axis (P.A. $``$ 138) and corresponding extent of the H<sub>2</sub> emission are denoted by the solid black line. The dashed black line indicates possible extension of this outflow and is based on observed CO bullets along this axis. The dotted lines indicate the outflow walls, determined from a fan-shaped reflection nebulosity ($`\varphi 27^{}`$) seen in K emission (Hodapp 1994), as well as two strands of H<sub>2</sub> emission originating from IRS 2 which lie along these axes (Eislöffel 2000). b. Two outflows model. Solid and dashed black lines as in Figure 6a. The dashed orange line indicates the axis of the highly-collimated bipolar CO emission which is evident in our CO data. Positions of all the Herbig-Haro objects (black crosses) are shown, as well as the five Class 0 sources (black stars) and the Class I source, IRS 1 (solid black box).
Individual outflow position angles and extents superimposed on the CO outflow integrated intensity map. Positions of all the Herbig-Haro objects (black crosses) are shown, as well as the five Class 0 sources (black stars) and the Class I source, IRS 1 (solid black box). Position angles are indicated for the outflows associated with L1448C (green $``$ Bachiller et al. 1995; Davis & Smith 1995; Dutrey et al. 1997; Eislöffel 2000), L1448N(B) (mustard $``$ Bachiller et al. 1990; Bontemps et al. 1996; Barsony et al. 1998), L1448N(A) (purple $``$ Davis & Smith 1995; Barsony et al. 1998), L1448 IRS 2 (black $``$ O’Linger et al. 1999; orange $``$ this work), and high velocity gas of unknown origin associated with HH 277 (dark blue). The extents of the outflow lobes that are well-established from the literature and our data are indicated with solid lines. Features that are seen only in our CO data, and extrapolations that are consistent with our data, are shown with dashed lines.
Individual outflow position angles superimposed on an H<sub>2</sub> mosaic of L1448 (Eislöffel 2000). Note that the axes of this figure are in J2000 coordinates. The positions of Herbig-Haro objects (orange crosses), Class 0 sources (red stars), and IRS 1 (open red box), are indicated. Position angles of outflows are indicated as in Figure 8. Note that after its bend at emission knot I, the position angle of the L1448C outflow passes directly through emission knot S. Although emission knot R lies about 20<sup>′′</sup> south of the L1448C outflow axis, this knot appears to be part of extended emission that forms a U-shaped structure which opens toward the southeast, bisected by the outflow axis. The sides of the U are separated by $``$ 40<sup>′′</sup>, with the brighter side (including knot R) lying to the south of the L1448C outflow axis.
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# 𝐶𝑃 Asymmetries in Radiative B Decays with R–parity Violation
## I Introduction
$`CP`$ violation in radiative $`B`$ decays may provide a promising tool for probing new physics beyond the Standard Model (SM). The direct $`CP`$ asymmetry in the radiative $`bs\gamma `$ decay defined by
$$A_{CP}^{bs\gamma }(\delta )=\frac{\mathrm{\Gamma }(\overline{B}X_s\gamma )\mathrm{\Gamma }(BX_{\overline{s}}\gamma )}{\mathrm{\Gamma }(\overline{B}X_s\gamma )+\mathrm{\Gamma }(BX_{\overline{s}}\gamma )}|_{E_\gamma >(1\delta )E_\gamma ^{\mathrm{max}}}$$
(1)
is below 1 $`\%`$ in the SM . Therefore, the observation of a sizable $`CP`$ asymmetry would be a clean signal of new physics and may further discriminate various extensions of the SM. Recent model-independent analyses of $`CP`$ violating effects in the inclusive decays $`BX_s\gamma `$ in terms of the Wilson coefficients of the dipole moments operators have shown that models with enhanced chromo-magnetic dipole moments can naturally provide a large asymmetry and it can reach up to 30 $`\%`$ accommodating the observed branching ratios of $`bs\gamma `$ and $`bsg`$ decays . The specific predictions for the $`CP`$ asymmetry are of course model-dependent. The left-right model can yield $`A_{CP}`$ at the level of 1 $`\%`$ , and a two Higgs doublet model up to 10 $`\%`$ . In supersymmetric extensions of the Standard Model, there exists additional $`CP`$ phases which would give rise to large $`CP`$ violating effects. In the context of minimal supergravity models , the direct $`CP`$ asymmetry turns out to be less than 2 $`\%`$ due to strong constraints on supersymmetric $`CP`$ violating phases from the neutron and electron electric dipole moments . The asymmetry can be enlarged in non-minimal models up to 7 $`\%`$ with more freedom in $`CP`$ phases , or up to 15 $`\%`$ with generic sfermion mixing .
Another important $`CP`$ violating observable is the mixing-induced $`CP`$ asymmetry in exclusive radiative $`B`$ decays which can occur for radiative $`B_qM_q\gamma `$ decays, where $`M_{q=d,s}`$ is any hadronic self-conjugate state with $`CP`$ eigenvalue $`\xi =\pm 1`$. The $`CP`$ asymmetry in the time-dependent decay rates $`\mathrm{\Gamma }(t)`$ for $`B_qM_q\gamma `$ and $`\overline{\mathrm{\Gamma }}(t)`$ for $`\overline{B_q}M_q\gamma `$ is then
$$A(t)\frac{\mathrm{\Gamma }(t)\overline{\mathrm{\Gamma }}(t)}{\mathrm{\Gamma }(t)+\overline{\mathrm{\Gamma }}(t)}=\xi A_M\mathrm{sin}(\varphi _M\varphi _L\varphi _R)\mathrm{sin}(\mathrm{\Delta }mt),A_M\frac{2|C_{7L}C_{7R}|}{|C_{7L}|^2+|C_{7R}|^2}$$
(2)
where $`\varphi _M`$ and $`\mathrm{\Delta }m`$ are the phase and the mass difference of $`B_q`$$`\overline{B}_q`$ mixing, respectively, and $`C_{7L,7R}`$ are the effective coefficients of the left-handed and right-handed dipole moment operators for the $`bq\gamma `$ decays, and $`\varphi _{L,R}`$ are their phases, respectively. In the standard model, such asymmetries are expected to be small as one has $`C_{7R}/C_{7L}m_q/m_b`$ and thus $`A_M`$ of order 1 $`\%`$ and 10 $`\%`$ for the $`bd\gamma `$ and $`bs\gamma `$ decays, respectively. Unlike the direct $`CP`$ asymmetry, the left-right symmetric model may allow the mixing-induced asymmetry up to 50 $`\%`$ . The supersymmetric model with generic sfermion mixing can yield even larger asymmetry up to 90 $`\%`$ .
In this paper, we will analyze the effects of R–parity violation in the Minimal Supersymmetric Standard Model (MSSM) on the $`CP`$ asymmetries in radiative $`B`$ decays. R–parity violation in the MSSM introduces a large number of trilinear couplings which violate lepton and baryon number. These additional couplings can surely be sources of flavor and $`CP`$ violation, which might lead to a huge effect on the $`CP`$-odd observables in the B decays. $`CP`$ violating effects of R–parity violation in the hadronic $`B`$ decays have been considered previously in Refs., showing that significant modifications to the SM predictions can follow from R–parity violation. This work is devoted to the analysis of the $`CP`$ asymmetries in the radiative $`B`$ decay. As we will see, contrary to the cases with hadronic $`B`$ decays, various experimental constraints on the R–parity violating couplings coming particularly from rare $`B`$ decays strongly limit the amount of the direct $`CP`$ asymmetry in the $`bs\gamma `$ decay whereas a large mixing-induced $`CP`$ asymmetry can be induced from, in particular, baryon number violation.
This paper is organized as follows. In section II, we analyze the general effective Hamiltonian describing the $`bs\gamma `$ decay including the new operators induced by R–parity violation. We derive the anomalous dimensions and evolution matrix with the enlarged operator set. In section III, we discuss the $`CP`$ asymmetries arising from lepton number violating couplings with which we can have both the direct and mixing-induced $`CP`$ asymmetry. We deal with both cases in separated subsections. In section IV, we discuss the $`CP`$ asymmetries arising from baryon number violating couplings, in which case only mixing-induced $`CP`$ asymmetry can be obtained. We conclude in section V.
## II R–parity violation and effective Hamiltonian
Let us begin our discussion with defining our choice of the basis to describe the R–parity violating couplings. For comparison with experiments, it is convenient to work with R–parity violating couplings defined in the quark and lepton mass eigenbasis. In this prescription, we leave the neutrinos in the charged lepton mass eigenbasis since neutrinos can be taken to be massless for our purpose. The full superpotential of the MSSM fields including generic R–parity violating couplings is then given by
$`W`$ $`=`$ $`\mu H_1H_2+h_i^eH_1L_iE_i^c+h_i^d(H_1^0D_iD_i^cH_1^{}V_{ij}^{}U_jD_i^c)+h_i^u(H_2^0U_iU_i^cH_2^+V_{ij}D_jU_i^c)`$ (3)
$`+`$ $`\lambda _{ijk}(L_i^0E_jL_j^0E_i)E_k^c+\lambda _{ijk}^{}(L_i^0D_jD_k^cE_iV_{jl}^{}U_lD_k^c)+{\displaystyle \frac{1}{2}}\lambda _{ijk}^{\prime \prime }U_i^cD_j^cD_k^c,`$ (4)
where $`L_i=(L_i^0,E_i)`$ and $`Q_i=(U_i,D_i)`$ are the lepton and quark $`SU(2)`$ doublets, and $`E_i^c`$, $`U_i^c`$, $`D_i^c`$ are the $`SU(2)`$ singlet anti-lepton and anti-quark superfields. Here $`V_{ij}`$ is the Cabibbo–Kobayashi–Maskawa (CKM) matrix of quark fields. In order to ensure the longevity of proton, the products $`\lambda ^{}\lambda ^{\prime \prime }`$ have to be highly suppressed . For this reason, one usually assumes lepton or baryon number conservation to discard the couplings $`\lambda /\lambda ^{}`$ or $`\lambda ^{\prime \prime }`$, respectively. In this paper, we discuss both cases separately: one with R–parity and lepton number violation with nonvanishing $`\lambda ^{}`$, and the other with R–parity and baryon number violation with $`\lambda ^{\prime \prime }`$. The couplings $`\lambda `$ will be irrelevant for our discussion.
The presence of R–parity violating couplings in Eq. (3) give rise to new contributions to the radiative decay $`bs\gamma `$ through one-loop diagrams exchanging sleptons or squarks as depicted in FIG. 1. Note that R–parity violation may induce equally sizable dipole moments of the left-handed and right-handed type, respectively labeled by $`L`$ and $`R`$ as follows:
$$O_{7L}\overline{s_L^\alpha }\sigma ^{\mu \nu }b_R^\alpha F_{\mu \nu },O_{7R}\overline{s_R^\alpha }\sigma ^{\mu \nu }b_L^\alpha F_{\mu \nu }.$$
(5)
In the SM and many extensions of it, the coefficients of the second one is usually suppressed by the factor $`m_s/m_b`$. In our case with lepton or baryon number violation in the MSSM, the operator $`O_{7L}`$ and $`O_{7R}`$ are generated for nonvanishing combinations of couplings;
$$\lambda _{n3j}^{}\lambda _{n2j}^{}\mathrm{and}\lambda _{nj2}^{}\lambda _{nj3}^{}/\lambda _{n12}^{\prime \prime }\lambda _{n23}^{\prime \prime },$$
(6)
respectively. With the couplings in Eq. (6), there arise also other four-quark operators which should be taken into account in the complete effective Hamiltonian describing the $`bs\gamma `$ decay. The whole set of the effective operators arising from R–parity violation can be described by a simple generalization of the SM operator space by separating the standard operators $`O_{3,4,5,6}`$ for each quark flavor $`q=u,c,d,s,b`$. That is, we introduce the additional operators,
$`O_{3L}^q=\overline{s_L^\alpha }\gamma ^\mu b_L^\alpha \overline{q_L^\beta }\gamma _\mu q_L^\beta ,`$ $`O_{4L}^q=\overline{s_L^\beta }\gamma ^\mu b_L^\alpha \overline{q_L^\alpha }\gamma _\mu q_L^\beta ,`$ (7)
$`O_{5L}^q=\overline{s_L^\alpha }\gamma ^\mu b_L^\alpha \overline{q_R^\beta }\gamma _\mu q_R^\beta ,`$ $`O_{6L}^q=\overline{s_L^\beta }\gamma ^\mu b_L^\alpha \overline{q_R^\alpha }\gamma _\mu q_R^\beta ,`$ (8)
where $`\alpha ,\beta `$ are the color indices. The dipole moment operators for the $`bs\gamma `$ and $`bsg`$ are defined by
$`O_{7L}`$ $`=`$ $`{\displaystyle \frac{e}{16\pi ^2}}m_b\overline{s_L^\alpha }\sigma ^{\mu \nu }b_R^\alpha F_{\mu \nu },`$ (9)
$`O_{8L}`$ $`=`$ $`{\displaystyle \frac{g_s}{16\pi ^2}}m_b\overline{s_R^\alpha }\sigma ^{\mu \nu }T_{\alpha \beta }^ab_L^\beta G_{\mu \nu }^a.`$ (10)
We also have the right-handed counterpart of the operators which can be obtained by the exchange $`LR`$ in Eqs. (7) and (9).
The effective Hamiltonian at scale $`\mu 𝒪(m_W)`$ relevant for the $`bs\gamma `$ decay is now described in terms of the enlarged set of operators as follows:
$`_{\mathrm{eff}}`$ $`=`$ $`{\displaystyle \frac{4G_F}{\sqrt{2}}}\lambda _t[{\displaystyle \underset{i=1}{\overset{8}{}}}C_{iL}(\mu )O_{iL}(\mu )+{\displaystyle \underset{q}{}}{\displaystyle \underset{j=3}{\overset{6}{}}}C_{jL}^q(\mu )O_{jL}^q(\mu )+(LR)],`$ (11)
where $`\lambda _t=V_{ts}^{}V_{tb}`$, $`G_F`$ is the Fermi constant and $`C`$’s are the Wilson coefficients which will be determined later. The operators $`O_{iL}`$ with $`i=1,\mathrm{},8`$ are those considered in the SM . Note that $`O_{1L,2L}=O_{3L,4L}^c`$ and $`O_{iL}=_qO_{iL}^q`$ for $`i=3,\mathrm{}6`$.
Given the Wilson coefficients $`C_i,C_j^q`$ including the contributions from R–parity violation at the weak scale $`\mu =m_W`$, we need to calculate those at the scale $`\mu m_b`$ through the renormalization group (RG) evolution. Since the QCD running do not mix the left-handed and right-handed set of operators and its effect is identical, it is enough to calculate the RG equation at the one sector. At the leading order, it is rather straightforward to calculate the anomalous dimension matrix in the extended operator basis following the standard calculation . At this point, let us recall that it is convenient to use the so-called “effective coefficients” which are free from the regularization-scheme dependency in the mixing between the sets $`O_i^{(q)}`$ with $`i=1,\mathrm{},6`$ and $`O_{7,8}`$ resulting from two-loop diagrams. In terms of the effective coefficients,
$`C_7^{eff}(\mu )`$ $`=`$ $`C_7(\mu )+{\displaystyle \underset{I}{}}y_IC_I(\mu )`$ (12)
$`C_8^{eff}(\mu )`$ $`=`$ $`C_8(\mu )+{\displaystyle \underset{I}{}}z_IC_I(\mu ),`$ (13)
with the index $`I`$ running for 26 indices labeled by $`i`$ and $`jq`$ for the operators $`O_i`$ ($`i7,8`$) and $`O_j^q`$, the effective anomalous dimension matrix is given by
$$\gamma _{JI}^{eff}=\{\begin{array}{cc}\gamma _{J7}+_Ky_K\gamma _{JK}y_J\gamma _{77}z_J\gamma _{87},\text{for }I=7,J7,8\hfill & \\ \gamma _{J8}+_Kz_K\gamma _{JK}z_J\gamma _{88},\text{for }I=8,J7,8\hfill & \\ \gamma _{JI},\text{otherwise.}\hfill & \end{array}$$
(14)
The corresponding RG equations for the effective coefficients are then
$$\frac{d}{d\mathrm{ln}\mu }C_I^{eff}(\mu )=\frac{\alpha _s}{4\pi }\gamma _{JI}^{eff}C_J^{eff}(\mu ).$$
(15)
In the naive dimensional regularization scheme which we follow , the nonvanishing coefficients $`y_I,z_I`$ are
$`y_5=y_{5b}={\displaystyle \frac{1}{3}},y_6=y_{6b}=1`$ (16)
$`z_5=z_{5b}=1.`$ (17)
Now, keeping track of the flavor structure of the standard calculation of the matrix $`\gamma ^{eff}`$ , one can find the $`28\times 28`$ anomalous dimension matrix in a straightforward way. The results are presented in the Appendix. There, we also calculate the evolution matrix from which the Wilson coefficients at the scale $`\mu =m_b`$ can be obtained in terms of the coefficients determined at the weak scale. The contributions to these coefficients from R–parity violation will be discussed in detail in the following sections.
## III $`CP`$ asymmetries with lepton number violation
### A $`\lambda _{\mathrm{𝐧𝟑𝐣}}^{}\lambda _{\mathrm{𝐧𝟐𝐣}}^{}:`$ Direct $`CP`$ asymmetry
With the nonvanishing combination of the R–parity violating couplings $`\lambda _{n3j}^{}\lambda _{n2j}^{}`$, the effective Hamiltonian in Eq. (11) includes only the operators of the left-handed type: $`O_{6L}^{d_j}`$ induced by the tree-level diagram exchanging sneutrino $`\stackrel{~}{\nu }_n`$. In addition to these, there arises also additional effective semileptonic operators,
$$O_{9L}^{\nu _n}=\overline{s_L}\gamma ^\mu b_L\overline{\nu _{nL}}\gamma _\mu \nu _{nL},$$
(18)
through the tree diagram exchanging the right-handed down squark $`\stackrel{~}{d_j^c}`$. The corresponding Wilson coefficients are given by
$`C_{6L}^{d_j}(m_W)`$ $`=`$ $`{\displaystyle \frac{1}{4\sqrt{2}G_F\lambda _t}}{\displaystyle \underset{n=1}{\overset{3}{}}}{\displaystyle \frac{\lambda _{n3j}^{}\lambda _{n2j}^{}}{m_{\stackrel{~}{\nu }_n}^2}},`$ (19)
$`C_{9L}^{\nu _n}(m_W)`$ $`=`$ $`{\displaystyle \frac{1}{4\sqrt{2}G_F\lambda _t}}{\displaystyle \underset{j=1}{\overset{3}{}}}{\displaystyle \frac{\lambda _{n3j}^{}\lambda _{n2j}^{}}{m_{\stackrel{~}{d_j^c}}^2}}.`$ (20)
Note that $`C_{5L}^{d_j}(m_W)=0`$ and it remains vanishing at low energy scale under one-loop RG evolution as can be seen from Eq. (A9) in the Appendix. The R–parity violating contributions to the coefficients of the operators $`O_{7L}`$ and $`O_{8L}`$ come from one-loop diagrams exchanging sneutrinos and right-handed down squarks, and have been computed in Ref. . They can be conveniently re-written in terms of the coefficients $`C_{6L}^{d_j}`$ and $`C_{9L}^{\nu _n}`$ as follows:
$$C_{7L}^{R_p/}(m_W)=Q_dC_{8L}^{R_p/}(m_W)=\frac{Q_d}{6}\left[\underset{n}{}C_{9L}^{\nu _n}(m_W)+2\underset{j}{}C_{6L}^{d_j}(m_W)\right],$$
(21)
where $`Q_d=1/3`$. In deriving $`C_{7L,8L}^{R_p/}`$, we neglected the down-type squark mixing and the masses of down-type quarks compared to the mass of the sneutrino. As shown in Eq. (16), the effective coefficient $`C_{7L}^{eff}`$ at $`\mu =m_W`$ gets a nontrivial contribution from $`C_{6L}^{d_j}`$ with $`j=3`$ and thus we have
$`C_{7L}^{eff}(m_W)`$ $`=`$ $`C_{7L}(m_W)C_{6L}^b(m_W),`$ (22)
$`C_{8L}^{eff}(m_W)`$ $`=`$ $`C_{8L}(m_W).`$ (23)
Making use of the relation (21) and the formula (A22) in the Appendix, one can obtain the Wilson coefficients at the scale $`\mu =m_b`$ as follows;
$`C_{7L}^{eff}(m_b)`$ $`=`$ $`0.67C_{7L}^{\mathrm{SSM}}(m_W)+0.092C_{8L}^{\mathrm{SSM}}(m_W)0.17C_{2L}(m_W)`$ (24)
$``$ $`0.14[C_{6L}^d(m_W)+C_{6L}^s(m_W)]0.80C_{6L}^b(m_W)0.022C_{9L}^{\nu _n}(m_W),`$ (25)
$`C_{8L}^{eff}(m_b)`$ $`=`$ $`0.70C_{8L}^{\mathrm{SSM}}(m_W)0.080C_{2L}(m_W)`$ (26)
$`+`$ $`0.42[C_{6L}^d(m_W)+C_{6L}^s(m_W)]+0.60C_{6L}^b(m_W)+0.12C_{9L}^{\nu _n}(m_W),`$ (27)
where $`C_{7L,8L}^{\mathrm{SSM}}(m_W)`$ contain the contributions from the R–parity conserving MSSM sector as well as from the SM one. We assume that $`C_{2L}(m_W)`$ comes solely from the SM. To quantify $`C_{7L,8L}^{\mathrm{SSM}}(m_W)`$ for our purpose, we introduce the parameters $`\eta _{7,8}`$ which are defined by
$$C_{7L}^{\mathrm{SSM}}(m_W)\eta _7C_{7L}^{\mathrm{SM}}(m_W),C_{8L}^{\mathrm{SSM}}(m_W)\eta _8C_{8L}^{\mathrm{SM}}(m_W).$$
(28)
The parameters $`\eta _{7,8}`$ are complex in general.
Given the Wilson coefficients in Eq. (24), we are ready to analyze the $`CP`$ violating effects from R–parity violation in the radiative $`B`$ decay. Referring to the work by Kagan and Neubert for details, the direct $`CP`$ asymmetry $`A_{CP}`$ and the branching ratios for the decays $`bs\gamma `$ and $`bsg`$ are given by <sup>*</sup><sup>*</sup>*Here, we neglect the R–parity violating contributions to $`A_{CP}`$ through terms such as $`\mathrm{}[C_L^{R_p/}C_{7L}^{}]`$.
$`A_{CP}={\displaystyle \frac{1}{\left|C_{7L}\right|^2}}\{1.23\mathrm{}[C_{2L}C_{7L}^{}]9.52\mathrm{}[C_{8L}C_{7L}^{}]+0.10\mathrm{}[C_{2L}C_{8L}^{}]\}(\%),`$ (29)
$`\mathrm{B}(BX_s\gamma )2.57\times 10^3K_{\mathrm{NLO}}(\delta )\left({\displaystyle \frac{\mathrm{B}_{\mathrm{semi}}}{0.105}}\right),`$ (30)
$`\mathrm{B}(BX_{sg})0.96|C_{8L}|^2\mathrm{B}_{\mathrm{semi}},`$ (31)
where the coefficients $`C`$’s without arguments are understood to be the effective ones evaluated at the scale $`m_b`$. The quantity $`K_{\mathrm{NLO}}(\delta )=|C_7|^2+O(\alpha _s,1/m_b^2)`$ contains the corrections to the leading–order result and the specific forms of the corrections can be found in Ref. . We will take $`\delta =0.3`$ and $`\mathrm{B}_{\mathrm{semi}}=10.5`$ % for the branching ratio of the semileptonic decay $`BX_ce\overline{\nu }`$. In this work, we take the following values for the SM predictions for the $`C_{2L,7L,8L}^{\mathrm{SM}}`$ at the $`m_W`$ scale;
$$C_{2L}^{\mathrm{SM}}(m_W)1.0,C_{7L}^{\mathrm{SM}}(m_W)0.20,C_{8L}^{\mathrm{SM}}(m_W)0.10.$$
(32)
The above choice of parameters yields the values at the scale $`\mu =m_b=4.8`$ GeV;
$$C_{2L}^{\mathrm{SM}}(m_b)1.11,C_{7L}^{\mathrm{SM}}(m_b)0.32,C_{8L}^{\mathrm{SM}}(m_b)0.15.$$
(33)
Considering the above values of the SM coefficients, Eq. (24) suggests that a significant contribution from R–parity violation to the $`bs\gamma `$ decay can arise for $`|C_{9L}^{\nu _n}(m_W)|10`$, $`|C_{6L}^{d,s}(m_W)|2`$, or $`|C_{6L}^b(m_W)|0.3`$. Furthermore, as we will show, if a sizable $`|C_{6L}^b|0.3`$ is allowed, one can get the direct $`CP`$ asymmetry of the order $`|A_{CP}|10\%`$ satisfying the observed branching ratios of $`BX_s\gamma `$ and $`BX_{sg}`$ decay.
In order to figure out how large $`CP`$ asymmetry can come from R–parity violation, let us consider the experimental bounds on the new Wilson coefficients appeared in Eq. (24). First of all, those coefficients will be constrained by the experimental data for the branching ratios in Eq. (29). In this work, we use the CLEO data:
$`\mathrm{B}(BX_s\gamma )`$ $`=`$ $`(3.15\pm 0.35_{\mathrm{stat}}\pm 0.32_{\mathrm{stat}}\pm 0.26_{\mathrm{model}})\times 10^4,\text{[20]}`$ (34)
$`\mathrm{B}(BX_{sg})`$ $``$ $`6.8\%(90\%\mathrm{C}.\mathrm{L}.).\text{[21]}`$ (35)
More important constraints on the coefficients $`C_{6L}^{d_j}`$ and $`C_{9L}^{\nu _n}`$ come from experimental data on the various B meson decays. Let us discuss the relevant bound for each coefficient.
First, the coefficient $`C_{6L}^d`$ is constrained by the $`B`$ decay mode $`\overline{B^0}K^0\pi ^0`$, whose branching ratio is observed to be
$$\mathrm{B}(\overline{B^0}K^0\pi ^0)<4.1\times 10^5.$$
Following the similar method used in Ref. , we estimate the matrix element of the operator $`O_{6L}^d`$ as
$$<K^0\pi ^0|O_{6L}^d|\overline{B^0}>i\frac{m_K^2(m_B^2m_\pi ^2)}{2(m_s+m_d)(m_bm_d)}f_KF_1^{B\pi }(m_K^2).$$
(36)
Taking $`f_K=0.16`$ GeV, $`F_1^{B\pi }(m_K^2)=0.33`$, we obtain
$`\left|C_{6L}^d\right|<0.17.`$ (37)
Second, the coefficient $`C_{6L}^s`$ is constrained by considering the $`B`$ decay mode $`\overline{B^0}\varphi K^0`$. The experimental limits on the branching ratio of this decay mode is
$$\mathrm{B}(\overline{B^0}\varphi K^0)<3.1\times 10^5.$$
From our estimation of the matrix element of the operator $`O_{6L}^s`$;
$$<\varphi K^0|O_{6L}^s|\overline{B^0}>\frac{1}{4N}f_\varphi m_\varphi ϵ(P_B+P_K)F_1^{BK}(m_\varphi ^2),$$
(38)
where $`ϵ`$ is a polarization vector of $`\varphi `$ and we will take $`N=3`$. With $`f_\varphi =0.23`$ GeV, $`F_1^{BK^0}(m_\varphi ^2)=0.38`$, we obtain following limit :
$`\left|C_{6L}^s\right|<0.23.`$ (39)
The above bounds in Eqs. (37) and (39) tell us that the coefficients $`C_{6L}^{d,s}`$ cannot play any important role in the radiative B decays as can be seen from Eq. (24) The bounds in Eqs. (37) and (39) are at the scale $`m_b`$. Practically, the Wilson coefficients induced by the R–parity violating couplings at the scale $`m_b`$ are nearly the same as those at the scale $`m_W`$ \[see Eq. (A9)\]..
Now let us consider the constraint on $`C_{6L}^b`$. The most useful bound comes indirectly from the consideration of the decay mode $`\overline{B^0}X_s\nu _n\overline{\nu _n}`$ whose branching ratio can be calculated as
$$\mathrm{B}(\overline{B^0}X_s\nu _n\overline{\nu _n})=\left|\frac{\lambda _t}{V_{cb}}\right|^2\frac{\left|C_{9L}^{\nu _n}\right|^2}{f_{\mathrm{PS}}(m_c^2/m_b^2)}\mathrm{B}_{\mathrm{semi}}.$$
(40)
According to the analyses in Ref. , one obtains indirect experimental information on the above branching ratio:
$$\mathrm{B}(\overline{B^0}X_s\nu _n\overline{\nu _n})<3.9\times 10^4.$$
Taking this value and $`\left|\lambda _t/V_{cb}\right|=0.976`$, $`f_{\mathrm{PS}}(m_c^2/m_b^2)=0.5`$, we put the bound,
$$\left|C_{9L}^{\nu _n}\right|<0.044.$$
(41)
Thus the contribution of the coefficient $`C_{9L}^{\nu _n}`$ to Eq. (24) can also be neglected. Now, under the condition that the bound (41) is applied to each component of the coefficient \[see Eq. (19)\], we get
$$|C_{6L}^b|<0.044\left(\frac{m_{\stackrel{~}{d_3^c}}^2}{m_{\stackrel{~}{\nu _n}}^2}\right),$$
(42)
for each $`n=1,2,3`$. Here we remark that under the condition (41), the R–parity violating contribution to the semileptonic decay $`BX_cl_n\overline{\nu _n}`$ through the effective operator
$$_{eff}=\frac{\lambda _{n3j}^{}\lambda _{n2j}^{}}{2m_{\stackrel{~}{d}_j^c}^2}\overline{c_L}\gamma ^\mu b_L\overline{e}_{nL}\gamma _\mu \nu _{nL},$$
(43)
can be made small enough to satisfy the direct experimental bounds ;
$`|\lambda _{233}^{}\lambda _{223}^{}|<1.1\times 10^3\left({\displaystyle \frac{m_{\stackrel{~}{d_3^c}}}{100\mathrm{GeV}}}\right)^2,`$ (44)
$`|\lambda _{333}^{}\lambda _{323}^{}|<4.4\times 10^3\left({\displaystyle \frac{m_{\stackrel{~}{d_3^c}}}{100\mathrm{GeV}}}\right)^2.`$ (45)
Finally, we have to consider the neutrino mass coming from our choice of nonvanishing $`\lambda _{n33}^{}\lambda _{n23}^{}`$. With nonzero value of $`\lambda _{n33}^{}`$, the neutrino $`\nu _n`$ may get an undesirably large mass from one-loop diagrams with the exchange of bottom quark and squark . The one-loop contribution to the neutrino mass is given by
$$m_{\nu _n}\frac{3}{8\pi ^2}\frac{\lambda _{n33}^2Am_b^2}{\stackrel{~}{m}^2}$$
where $`A`$ denotes left-right sbottom mixing parameter and $`\stackrel{~}{m}`$ is an average sbottom mass. Taking $`\lambda _{n33}^2/\stackrel{~}{m}^2<10^7\mathrm{GeV}^2`$ from the consideration of Eq. (41), we obtain
$$m_{\nu _n}10\mathrm{keV}$$
(46)
with $`A=100`$ GeV. Thus it is well below the direct experimental limit for the muon (tau) neutrino which is 0.17 (18) MeV . Yet another indirect limit for the muon or tau neutrino mass comes from the observation of atmospheric muon neutrinos. The atmospheric neutrino data from the Super-Kamiokande are known to be nicely explained by the neutrino oscillation between largely mixed muon and tau neutrinos . A natural consequence of this would be that the muon and tau neutrinos are very light $`m_{\nu _\mu ,\nu _\tau }1`$ eV. If this is the case, there is no room at all for large $`CP`$ asymmetry from R–parity violation. However, having a large R–parity violation and thus a heavier muon or tau neutrino is not excluded completely as there exist some other viable options for the explanation of the Super-Kamiokande data, such as a neutrino decay .
From the above consideration, the only possibility for a significant enhancement of R–parity violating contribution to the $`bs\gamma `$ decay is to have a large coefficient $`C_{6L}^b`$ with a sfermion mass hierarchy $`m_{\stackrel{~}{d}_3^c}>m_{\stackrel{~}{\nu }_n}`$. For example, we need $`m_{\stackrel{~}{d}_3^c}5m_{\stackrel{~}{\nu }_n}`$ to get $`C_{6L}^b1`$. Taking into account all the above experimental limits on the Wilson coefficients, let us now analyze how large $`CP`$ asymmetry $`A_{CP}`$ can be obtained. In FIG. 2, we show $`A_{CP}`$ as a function of the branching ratio of the decay $`BX_s\gamma `$ varying $`|C_{6L}^b|`$ and $`\mathrm{Arg}(C_{6L}^b)`$ from 0 to $`2\pi `$. The other R–parity violating couplings are neglected. We take $`\eta _7=\eta _8=\eta `$ as a real number in FIG. 2 even though $`\eta `$ is a complex number generally. In this figure we consider additional experimental constraints coming from B($`BX_s\gamma `$), B($`BX_{sg}`$). If $`m_{\stackrel{~}{d}_3^c}^2`$ are larger than $`m_{\stackrel{~}{\nu }_n}^2`$, then sizable $`|C_{6L}^b|`$ is allowed evading the bounds Eq. (42). We find the $`CP`$ asymmetry can reach $`13\%`$ for $`m_{\stackrel{~}{d}_3^c}/m_{\stackrel{~}{\nu }_n}3.4`$ with vanishing R–parity conserving supersymmetric contributions ($`\eta =1`$) as shown in the left-upper frame of FIG. 2. In the case where the R–parity conserving supersymmetric contributions take the same sign as the SM values of $`C_{7L,8L}(m_W)`$, this $`CP`$ asymmetry can be larger as shown in FIG. 2 with $`\eta =2`$. On the other hand, the $`CP`$ asymmetry decreases when the R–parity conserving supersymmetric contributions take the opposite sign to the SM values of $`C_{7L,8L}(m_W)`$ as seen from FIG. 2 with $`\eta =0,1`$.
Before concluding this subsection, it is worthwhile to notice that the $`CP`$ asymmetry in the hadronic $`B`$ decays such as $`B_d\pi K_S`$ and $`B_d\varphi K_S`$ can be significantly affected by the R–parity violating couplings $`C_{6L}^{d,s}`$ even if the effects of the couplings on the direct $`CP`$ asymmetry in the radiative $`B`$ decay are negeligible .
### B $`\lambda _{\mathrm{𝐧𝐣𝟐}}^{}\lambda _{\mathrm{𝐧𝐣𝟑}}^{}:`$ Mixing-induced $`CP`$ asymmetry
Contrary to the previous case, the combination of R–parity violating couplings $`\lambda _{nj2}^{}\lambda _{nj3}^{}`$ leads only to the right-handed set of operators: $`O_{7R,8R}`$ and $`O_{5R,6R}^q`$ in the effective Hamiltonian (11). In addition, the coefficient of the operator $`O_{6R}^t`$ is also generated. Even though it does not contribute to the $`b`$ decay, it’s coefficient will be considered since it enters into $`C_{7R,8R}`$. As in the previous subsection, we have also the semileptonic four-Fermi operators as follows:
$`O_{9R}^{\nu _n}`$ $`=`$ $`\overline{s_R^\alpha }\gamma ^\mu b_R^\alpha \overline{\nu _{nL}}\gamma _\mu \nu _{nL},`$ (47)
$`O_{10R}^{e_n}`$ $`=`$ $`\overline{s_R^\alpha }\gamma ^\mu b_R^\alpha \overline{e_{nL}}\gamma _\mu e_{nL}.`$ (48)
The nonzero Wilson coefficients at $`m_W`$ induced from our R–parity violation are
$`C_{6R}^{q_j}(m_W)`$ $`=`$ $`{\displaystyle \frac{1}{4\sqrt{2}G_F\lambda _t}}{\displaystyle \underset{n=1}{\overset{3}{}}}{\displaystyle \frac{\lambda _{nj2}^{}\lambda _{nj3}^{}}{m_{\stackrel{~}{l}_n}^2}},`$ (49)
$`C_{9R}^{\nu _n}(m_W)`$ $`=`$ $`C_{10R}^{e_n}(m_W)={\displaystyle \frac{1}{4\sqrt{2}G_F\lambda _t}}{\displaystyle \underset{j=1}{\overset{3}{}}}{\displaystyle \frac{\lambda _{nj2}^{}\lambda _{nj3}^{}}{m_{\stackrel{~}{Q}_j}^2}},`$ (50)
where $`q_j`$ can be either $`u_j`$ or $`d_j`$, $`m_{\stackrel{~}{l}_n}`$ and $`m_{\stackrel{~}{Q}_j}`$ are the masses of the doublet slepton and squark, respectively. In deriving these coefficients, we do not consider the effect of the CKM mixing. Note that there can be also other four-Fermi operators involving two different flavors of quarks or leptons through the CKM mixing, which we neglect in our discussion as they give subleading contributions. Now, disregarding the contributions from the R–parity conserving supersymmetric sector, the Wilson coefficients for $`O_{7R,8R}`$ at the scale $`m_W`$ can be expressed in terms of the coefficients in Eq. (49) as follows
$`C_{7R}(m_W)=C_{7R}^{R_p/}(m_W)`$ $`=`$ $`{\displaystyle \frac{Q_d}{6}}\left[{\displaystyle \underset{n}{}}C_{9R}^{\nu _n}(m_W)+2{\displaystyle \underset{j}{}}C_{6R}^{q_j}(m_W)\right]`$ (52)
$`+{\displaystyle \frac{1}{18}}\left[8{\displaystyle \underset{n}{}}C_{10R}^{e_n}(m_W)+7{\displaystyle \underset{j}{}}C_{6R}^{q_j}(m_W)P_\gamma (x_j)\right],`$
$`C_{8R}(m_W)=C_{8R}^{R_p/}(m_W)`$ $`=`$ $`{\displaystyle \frac{1}{6}}[{\displaystyle \underset{n}{}}C_{9R}^{\nu _n}(m_W)+2{\displaystyle \underset{j}{}}C_{6R}^{q_j}(m_W)`$ (54)
$`+{\displaystyle \underset{n}{}}C_{10R}^{e_n}(m_W)+2{\displaystyle \underset{j}{}}C_{6R}^{q_j}(m_W)P_g(x_j)],`$
where
$$P_\gamma (x_j)=\frac{36}{7}[Q_uF_1(x_j)+F_2(x_j)],P_g(x_j)=6F_1(x_j)$$
with $`x_j=m_{u_j}^2/m_{\stackrel{~}{l}_n}^2`$. The functions $`F_{1,2}`$ are defined as
$`F_1(x)`$ $`=`$ $`{\displaystyle \frac{1}{12}}{\displaystyle \frac{(2+3x6x^2+x^3+6x\mathrm{ln}x)}{(1x)^4}},`$ (55)
$`F_2(x)`$ $`=`$ $`{\displaystyle \frac{1}{12}}{\displaystyle \frac{(16x+3x^2+2x^36x^2\mathrm{ln}x)}{(1x)^4}}.`$ (56)
Note that these functions are defined as $`P_\gamma (0)=P_g(0)=1`$ and their values for some selected $`x`$ are listed in TABLE II
In the case under consideration, there is no phase in $`C_{7R}C_{8R}^{}`$ contributing the direct $`CP`$ asymmetry in right-handed sector analogue of Eq. (29) since the same combination of the couplings generates both $`C_{7R,8R}`$. However, there can arise a sizable mixing-induced $`CP`$ asymmetry as defined in Eq. (2). That is, we may have $`|A_M|1`$ with $`|C_{7L}||C_{7R}|`$. Here, the Wilson coefficients are effective ones evaluated at the scale $`m_b`$. Combining our results in the Appendix (A26) and Eq. (52), we find
$`C_{7R}^{eff}(m_b)`$ $`=`$ $`0.40[C_{6R}^{q_1}(m_W)+C_{6R}^{q_2}(m_W)]+[0.80+0.26P_\gamma (x_t)+0.031P_g(x_t)]C_{6R}^{q_3}(m_W)`$ (58)
$`0.022C_{9R}^{\nu _n}(m_W)+0.31C_{10R}^{e_n}(m_W).`$
From this equation, it appears that we can easily obtain $`C_{7R}|C_{7L}|=|C_{7L}^{\mathrm{SSM}}|0.32|\eta _7|`$. To clarify this, we now consider the experimental constraints on the coefficients appearing in Eq. (58). Let us first note that the arguments in the previous subsection are applied here to get the bounds;
$$|C_{6R}^{q_1}|<0.17,|C_{6R}^{q_2}|<0.23,|C_{9R}^{\nu _n}|<0.044$$
(59)
as in Eqs. (37), (39) and (41), respectively. Thus, the contributions of these coefficients in Eq. (58) are not significant. Another important constraint on the relevant R–parity violating couplings comes from the decay $`BX_sl_n^+l_n^{}`$ induced by the second effective operator in Eq. (47). This consideration leads to
$$|C_{10R}^{e_n}|<0.017\left(\sqrt{\frac{\mathrm{B}(bX_sl_n^+l_n^{})^{\mathrm{expt}}}{5.7\times 10^5}}\right),$$
(60)
where $`\mathrm{B}(bX_sl_n^+l_n^{})^{\mathrm{expt}}`$ denotes the experimental upper limit on the branching ratio of the $`BX_sl_n^+l_n^{}`$ decay modes given by
$`\mathrm{B}(BX_se^+e^{})<5.7\times 10^5,`$ (61)
$`\mathrm{B}(BX_s\mu ^+\mu ^{})<5.8\times 10^5.`$ (62)
Even if there are no data for $`bX_s\tau ^+\tau ^{}`$, we expect it’s branching ratio is most probably less than that of $`BX_se^+e^{}`$. Then, similarly to Eq. (42), the bounds on $`C_{6R}^q`$ can be obtained indirectly from Eq. (60) as
$$|C_{6R}^{q_j}|<0.017\left(\sqrt{\frac{\mathrm{B}(BX_sl_n^+l_n^{})^{\mathrm{expt}}}{5.7\times 10^5}}\right)\left(\frac{m_{\stackrel{~}{Q}_j}}{m_{\stackrel{~}{l}_n}}\right)^2,$$
(63)
for each $`j`$ and $`n`$. Therefore, we may obtain a sizable $`C_{7R}^{eff}(m_b)`$ from the contribution of $`C_{6R}^{q_3}(m_W)`$ if there is again a hierarchy between sfermion masses. For example, taking $`m_{\stackrel{~}{Q}_3}/m_{\stackrel{~}{l}_n}5`$, $`|C_{6R}^{q_3}|=0.33`$ is allowed within the present experimental bound. Taking $`x_t=1`$ which gives $`P_\gamma (1)=5/14`$ and $`P_g(1)=1/4`$, one could obtain $`|C_{7R}^{eff}(m_b)|0.23`$. If $`|C_{7L}^{eff}(m_b)|=0.23`$, $`|A_M|1`$ is possible accommodating the measured $`\mathrm{B}(bs\gamma )`$. Finally, let us note that a large mixing-induced $`CP`$ asymmetry requires the (tau) neutrino to be heavy as discussed in the previous subsection.
Again we note that with the coupling $`\lambda _{n22}^{}\lambda _{n32}^{}`$ which does not affect the radiative $`B`$ decays, one can have important effects on the $`CP`$ asymmetries in the hadronic $`B`$ decays such as $`B_d\varphi K_S`$ and $`B^\pm \pi ^\pm K^0`$ .
## IV $`CP`$ asymmetries with baryon number violation
### $`\lambda _{\mathrm{𝐧𝟏𝟐}}^{\prime \prime }\lambda _{\mathrm{𝐧𝟏𝟑}}^{\prime \prime }:`$ Mixing-induced $`CP`$ asymmetry
Our final case is to have baryon number violation while lepton number is conserved. Then, the new operator set for the $`bs`$ transition contains again only right-handed ones with nonvanishing product of couplings $`\lambda _{n12}^{\prime \prime }\lambda _{n13}^{\prime \prime }`$: $`O_{3R,4R}^{u_n,d}`$. The Wilson coefficients of these operators calculated at the weak scale $`m_W`$ are
$`C_{3R}^{u_n}(m_W)`$ $`=`$ $`{\displaystyle \frac{1}{4\sqrt{2}G_F\lambda _t}}{\displaystyle \frac{\lambda _{n12}^{\prime \prime }\lambda _{n13}^{\prime \prime }}{m_{\stackrel{~}{d}_1^c}^2}}=C_{4R}^{u_n},`$ (64)
$`C_{3R}^d(m_W)`$ $`=`$ $`{\displaystyle \frac{1}{4\sqrt{2}G_F\lambda _t}}{\displaystyle \underset{n=1}{\overset{3}{}}}{\displaystyle \frac{\lambda _{n12}^{\prime \prime }\lambda _{n13}^{\prime \prime }}{m_{\stackrel{~}{u}_n^c}^2}}=C_{4R}^d.`$ (65)
Notice that the simultaneous presence of nonvanishing coefficients $`C_{3R,4R}`$ is due to color antisymmetry in the superpotential term, $`U^cD^cD^c`$. Following the similar steps as before, we get the relation
$`C_{7R}(m_W)=C_{7R}^{R_p/}(m_W)`$ $`=`$ $`{\displaystyle \frac{1}{9}}\left[4C_{4R}^d(m_W)5{\displaystyle \underset{n}{}}C_{4R}^{u_n}(m_W)P_\gamma ^{}(x_n)\right],`$ (66)
$`C_{8R}(m_W)=C_{8R}^{R_p/}(m_W)`$ $`=`$ $`{\displaystyle \frac{1}{6}}\left[C_{4R}^d(m_W)+{\displaystyle \underset{n}{}}C_{4R}^{u_n}(m_W)P_g^{}(x_n)\right],`$ (67)
where
$$P_\gamma ^{}(x_n)=\frac{36}{5}[Q_uF_1(x_n)Q_dF_2(x_n)],P_g^{}(x_n)=12[F_1(x_n)F_2(x_n)]$$
with $`x_n=m_{u_n}^2/m_{\stackrel{~}{d}^c}^2`$ and $`P_\gamma ^{}(0)=P_g^{}(0)=1`$. The Wilson coefficient $`C_{7R}^{eff}`$ for the $`bs\gamma `$ decay at $`m_b`$ is
$`C_{7R}^{eff}(m_b)`$ $`=`$ $`0.21C_{4R}^d(m_W)0.17[C_{4R}^u(m_W)+C_{4R}^c(m_W)]`$ (69)
$`[0.37P_\gamma ^{}(x_t)0.015P_g^{}(x_t)]C_{4R}^t(m_W).`$
Let us now consider the experimental limits for the various coefficients in Eq. (69). First of all, the R–parity violating couplings with $`n=1`$ are strongly constrained by the non-observation of nucleon-antinucleon oscillation and double nucleon decay; $`\lambda _{113}^{\prime \prime }<5\times 10^3`$ and $`\lambda _{112}^{\prime \prime }<10^6`$ with sfermion mass of 300 GeV . Thus $`C_{4R}^u`$ are negeligibly small. The constraint for $`C_{4R}^d(m_W)`$ comes again from the $`BK^0\pi ^0`$ decay. We estimate the matrix element of the operator $`O_{4R}^d`$ as
$$<K^0\pi ^0|O_{4R}^d|\overline{B^0}>i(m_B^2m_\pi ^2)f_KF_1^{B\pi }(m_K^2).$$
(70)
Using the similar values used in Eq. (37), we obtain
$$|C_{4R}^d|<0.15.$$
(71)
Note that this bound is also consistent with the data for the decay mode $`B^+\overline{K}^0\pi ^+`$ . Concerning the coefficient $`C_{4R}^c`$, the consideration of the $`BJ/\psi K_S`$ decay gives
$$|C_{4R}^c|<0.02.$$
(72)
Applying again the bound (71) to each component of $`C_{4R}^d`$, we get the bound on the $`C_{4R}^t`$;
$$|C_{4R}^t|<0.15\frac{m_{\stackrel{~}{t}^c}^2}{m_{\stackrel{~}{d}^c}^2}.$$
(73)
Thus, to get $`|C_{7R}^{eff}|>0.2`$ leading to a nearly maximal mixing-induced $`CP`$ asymmetry, we need a sizable $`C_{4R}^t`$ which can come about for $`m_{\stackrel{~}{t}^c}>3.5m_{\stackrel{~}{d}^c}`$ with $`x_t=1`$ The constraint on the relevant single baryon number violating coupling comes from $`\mathrm{\Gamma }(Zl\overline{l})/\mathrm{\Gamma }(Z\mathrm{hadrons})`$ , which gives $`|\lambda _{312,313,323}^{\prime \prime }|<0.5`$ for $`\stackrel{~}{m}=100`$ GeV. This constraint is so weak that $`|C_{3R,4R}^t|𝒪(1)`$ is easily allowed..
We conclude this section by pointing out that the direct $`CP`$ asymmetry in the decay $`B^\pm \pi ^\pm K^0`$ can arise maximally as is the case with the coupling $`\lambda _{n22}^{}\lambda _{n32}^{}`$ .
## V Conclusion
We have discussed the effects of R–parity violation on the $`CP`$ asymmetries in radiative $`B`$ decays. When we allow R–parity and lepton number violating couplings which generate at one-loop level the tau neutrino mass of order 10 keV, they can induce rather large $`CP`$ violating effects in the $`bs\gamma `$ decay. The direct $`CP`$ asymmetry can be as large as 17 % if the R–parity conserving supersymmetric contribution is comparable to the SM one. The mixing-induced $`CP`$ asymmetry can be almost maximal depending on the sfermion masses. For these sizable $`CP`$ violating effects, it is required to have moderate sfermion mass splittings by factor 4 or bigger. If the atmospheric neutrino data from the Super-Kamiokande are to be explained by the oscillation between the muon and tau neutrinos whose masses are very light $`m_{\nu _\mu ,\nu _\tau }1`$ eV, then the effects of R–parity violation on the radiative $`B`$ decays are negligible. A large mixing-induced $`CP`$ asymmetry is also possible with the R–parity and baryon number violating couplings with the similar order of sfermion mass splitting.
These results could be contrasted with the R–parity violating effects on the $`CP`$ asymmetries in the hadronic $`B`$ decays such as $`B_dJ/\psi K,\varphi K`$ or $`B^\pm \pi ^\pm K`$, which could be significant without sfermion mass splitting and are rather insensitive to the neutrino mass restriction.
Note added: While our work was being prepared, we encountered the paper which considers the R–parity violating effect on the radiative $`B`$ decay. It also deals with the RG running of the enlarged set of the Wilson coefficients which overlaps partly with our paper, and we find discrepancies in anomalous dimension matrix elements and the relation like (A10) and (A11).
## Acknowledgments
We thank G. Bhattacharyya for helpful discussions. The work of KH was supported by grant No. 1999-2-111-002-5 from the interdisciplinary Research program of the KOSEF and BK21 project of the Ministry of Education.
## A RG Equations with the extended operators
Here, we present the anomalous dimension matrix for the 28 operators including the standard set $`O_{1,\mathrm{}8}`$ and the extended set $`O_{3,\mathrm{},6}^q`$ with $`q=u,d,s,c,b`$. We drop the indices $`L,R`$ for the left-handed and right-handed set of operators as they have the identical anomalous dimension matrix. Omitting the usual $`8\times 8`$ matrix for the standard set of operators, we have the following nonvanishing block-diagonal elements of the whole $`28\times 28`$ matrix.
The $`2\times 2`$ submatrices mixing the operators $`O_{3,4}^q`$ and $`O_{5,6}^q`$ with themselves are
$$\left(\begin{array}{cc}2& 6\\ 6& 2\end{array}\right)\mathrm{and}\left(\begin{array}{cc}2& 6\\ 0& 16\end{array}\right),$$
(A1)
respectively. The $`2\times 8`$ submatrix mixing the operators $`O_{3,4}^q`$ with the standard set $`O_{1,\mathrm{},8}`$ is given by
$$\left(\begin{array}{cccccccc}0& 0& \frac{2}{9}\delta _{qb,s}& \frac{2}{3}\delta _{qb,s}& \frac{2}{9}\delta _{qb,s}& \frac{2}{3}\delta _{qb,s}& \frac{232}{81}\delta _{qb,s}& 3+\frac{70}{27}\delta _{qb,s}\\ 0& 0& \frac{2}{9}& \frac{2}{3}& \frac{2}{9}& \frac{2}{3}& \frac{416}{81}\delta _{qu}\frac{232}{81}\delta _{qd}& \frac{70}{27}+3\delta _{qb,s}\end{array}\right)$$
(A2)
where
$`\delta _{qb,s}=\{\begin{array}{cc}1\text{for }q=b\text{ or }s\hfill & \\ 0\text{otherwise,}\hfill & \end{array}`$ (A3)
$`\delta _{qu},\delta _{qd}=\{\begin{array}{cc}1\text{when }q\text{ is an up-type, or down-type quark}\hfill & \\ 0\text{otherwise.}\hfill & \end{array}`$ (A4)
Similarly, the $`2\times 8`$ submatrix mixing the operators $`O_{5,6}^q`$ with the 8 standard operators is
$$\left(\begin{array}{cccccccc}0& 0& 0& 0& 0& 0& \frac{32}{9}\delta _{qb}& 3\frac{14}{3}\delta _{qb}\\ 0& 0& \frac{2}{9}& \frac{2}{3}& \frac{2}{9}& \frac{2}{3}& \frac{448}{81}\delta _{qu}+\frac{200}{81}\delta _{qd}& \frac{119}{27}4\delta _{qb}\end{array}\right)$$
(A5)
where
$$\delta _{qb}=\{\begin{array}{cc}1\text{for }q=b\text{ }\hfill & \\ 0\text{otherwise.}\hfill & \end{array}$$
With this anomalous dimension matrix, the low energy Wilson coefficients are given by
$$\stackrel{}{C}^{eff}(\mu )=\widehat{U}^{eff}(\mu ,\mu _W)\stackrel{}{C}^{eff}(\mu _W)$$
(A6)
where
$$\widehat{U}^{eff}(\mu ,\mu _W)=\widehat{V}\left(\left[\frac{\alpha _s(\mu _W)}{\alpha _s(\mu )}\right]^{\frac{\stackrel{}{\gamma }^{eff}}{2\beta _0}}\right)_D\widehat{V}^1$$
(A7)
where $`\widehat{V}`$ diagonalizes $`\widehat{\gamma }^{effT}`$
$$\widehat{\gamma }_D^{eff}=\widehat{V}^1\widehat{\gamma }^{effT}\widehat{V}$$
(A8)
and $`\stackrel{}{\gamma }^{eff}`$ is the vector containing the diagonal elements of the diagonal matrix of eigenvalues of $`\widehat{\gamma }^{eff}`$.
By solving Eqs. (A4)–(A6), we find the following low energy Wilson coefficients at $`\mu _b`$ in terms of nonzero $`C_6^{d_i}(\mu _W),C_6^{u,c}(\mu _W)`$ in the L or R sector induced at $`\mu _W`$ by $`R`$ parity violating coupling $`\lambda _{n2j}^{}\lambda _{n3j}^{}`$ or $`\lambda _{nj2}^{}\lambda _{nj3}^{}`$, respectively, and $`C_{3,4}^{u,d,c}(\mu _W)`$ in the R sector by $`\lambda _{n12}^{\prime \prime }\lambda _{n13}^{\prime \prime }`$, and the SM contribution $`C_{2L}(\mu _W)`$ in L sector:
$`C_7^{eff}(\mu _b)`$ $`=`$ $`\eta ^{16/23}C_7^{eff}(\mu _W)+{\displaystyle \frac{8}{3}}\left(\eta ^{14/23}\eta ^{16/23}\right)C_8^{eff}(\mu _W)+C_2(\mu _W){\displaystyle \underset{i=1}{\overset{10}{}}}h_i\eta ^{a_i}`$ (A9)
$`+`$ $`\left(C_6^d(\mu _W)+C_6^s(\mu _W)\right){\displaystyle \underset{i=1}{\overset{10}{}}}r_i^1\eta ^{a_i}+C_6^b(\mu _W){\displaystyle \underset{i=1}{\overset{10}{}}}r_i^2\eta ^{a_i}+\left(C_6^u(\mu _W)+C_6^c(\mu _W)\right){\displaystyle \underset{i=1}{\overset{10}{}}}r_i^3\eta ^{a_i}`$ (A10)
$`+`$ $`C_3^d(\mu _W){\displaystyle \underset{i=1}{\overset{10}{}}}r_i^4\eta ^{a_i}+C_4^d(\mu _W){\displaystyle \underset{i=1}{\overset{10}{}}}r_i^5\eta ^{a_i}`$ (A11)
$`+`$ $`\left(C_3^u(\mu _W)+C_3^c(\mu _W)\right){\displaystyle \underset{i=1}{\overset{10}{}}}r_i^6\eta ^{a_i}+\left(C_4^u(\mu _W)+C_4^c(\mu _W)\right){\displaystyle \underset{i=1}{\overset{10}{}}}r_i^7\eta ^{a_i}`$ (A12)
$`C_8^{eff}(\mu _b)`$ $`=`$ $`\eta ^{14/23}C_8^{eff}(\mu _W)+C_2(\mu _W){\displaystyle \underset{i=1}{\overset{10}{}}}\overline{h}_i\eta ^{a_i}`$ (A13)
$`+`$ $`\left(C_6^d(\mu _W)+C_6^s(\mu _W)\right){\displaystyle \underset{i=1}{\overset{10}{}}}\overline{r}_i^1\eta ^{a_i}+C_6^b(\mu _W){\displaystyle \underset{i=1}{\overset{10}{}}}\overline{r}_i^2\eta ^{a_i}+\left(C_6^u(\mu _W)+C_6^c(\mu _W)\right){\displaystyle \underset{i=1}{\overset{10}{}}}\overline{r}_i^3\eta ^{a_i}`$ (A14)
$`+`$ $`C_3^d(\mu _W){\displaystyle \underset{i=1}{\overset{10}{}}}\overline{r}_i^4\eta ^{a_i}+C_4^d(\mu _W){\displaystyle \underset{i=1}{\overset{10}{}}}\overline{r}_i^5\eta ^{a_i}`$ (A15)
$`+`$ $`\left(C_3^u(\mu _W)+C_3^c(\mu _W)\right){\displaystyle \underset{i=1}{\overset{10}{}}}\overline{r}_i^6\eta ^{a_i}+\left(C_4^u(\mu _W)+C_4^c(\mu _W)\right){\displaystyle \underset{i=1}{\overset{10}{}}}\overline{r}_i^7\eta ^{a_i}`$ (A16)
$`C_2(\mu _b)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\eta ^{6/23}+\eta ^{12/23}\right)C_2(\mu _W)`$ (A17)
$`C_5^{d_i,u,c}(\mu _b)`$ $`=`$ $`0`$ (A18)
$`C_6^{d_i,u,c}(\mu _b)`$ $`=`$ $`C_6^{d_i,u,c}(\mu _W)`$ (A19)
$`C_3^{u,d,c}(\mu _b)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\eta ^{6/23}+\eta ^{12/23}\right)C_3^{u,d,c}(\mu _W)+{\displaystyle \frac{1}{2}}\left(\eta ^{6/23}\eta ^{12/23}\right)C_4^{u,d,c}(\mu _W)`$ (A20)
$`C_4^{u,d,c}(\mu _b)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\eta ^{6/23}+\eta ^{12/23}\right)C_4^{u,d,c}(\mu _W)+{\displaystyle \frac{1}{2}}\left(\eta ^{6/23}\eta ^{12/23}\right)C_3^{u,d,c}(\mu _W)`$ (A21)
where the quantities $`a_i,h_i,\overline{h}_i,r_i^n,\overline{r}_i^n`$ are shown in the TABLE I.
The explicit numerical expressions with the choice $`\mu _W=m_W`$, $`\mu _b=4.8`$ GeV, $`\alpha _s(\mu _W)=0.120`$ and $`\alpha _s(\mu _b)=0.214`$ are
$`C_{7L}^{eff}(\mu _b)`$ $`=`$ $`0.6687C_{7L}(\mu _W)+0.0920C_{8L}(\mu _W)0.1732C_{2L}(\mu _W)`$ (A22)
$``$ $`0.0974\left(C_{6L}^d(\mu _W)+C_{6L}^s(\mu _W)\right)(0.6687+0.0875)C_{6L}^b(\mu _W)`$ (A23)
$`C_{8L}^{eff}(\mu _b)`$ $`=`$ $`0.7032C_{8L}(\mu _W)0.0801C_{2L}(\mu _W)`$ (A24)
$`+`$ $`0.1893\left(C_{6L}^d(\mu _W)+C_{6L}^s(\mu _W)\right)+0.3670C_{6L}^b(\mu _W)`$ (A25)
$`C_{7R}^{eff}(\mu _b)`$ $`=`$ $`0.6687C_{7R}(\mu _W)+0.0920C_{8R}(\mu _W)`$ (A26)
$``$ $`0.0974\left(C_{6R}^d(\mu _W)+C_{6R}^s(\mu _W)\right)(0.6687+0.0875)C_{6R}^b(\mu _W)`$ (A27)
$`+`$ $`0.2506\left(C_{6R}^u(\mu _W)+C_{6R}^c(\mu _W)\right)0.0170C_{3R}^d(\mu _W)+0.0880C_{4R}^d(\mu _W)`$ (A28)
$`+`$ $`0.0147\left(C_{3R}^u(\mu _W)+C_{3R}^c(\mu _W)\right)0.1732\left(C_{4R}^u(\mu _W)+C_{4R}^c(\mu _W)\right)`$ (A29)
$`C_{8R}^{eff}(\mu _b)`$ $`=`$ $`0.7032C_{8R}(\mu _W)`$ (A30)
$`+`$ $`0.1893\left(C_{6R}^d(\mu _W)+C_{6R}^s(\mu _W)\right)+0.3670C_{6R}^b(\mu _W)`$ (A31)
$`+`$ $`0.1893\left(C_{6R}^u(\mu _W)+C_{6R}^c(\mu _W)\right)`$ (A32)
$``$ $`0.0894\left(C_{3R}^d(\mu _W)+C_{3R}^u(\mu _W)+C_{3R}^c(\mu _W)\right)`$ (A33)
$``$ $`0.0801\left(C_{4R}^d(\mu _W)+C_{4R}^u(\mu _W)+C_{4R}^c(\mu _W)\right)`$ (A34)
where we have used $`C_7^{eff}(\mu _W)=C_7(\mu _W)C_6^b(\mu _W)`$ and $`C_8^{eff}(\mu _W)=C_8(\mu _W)`$ on the right-hand sides of the above equations.
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# Quasiparticle Localization in Disordered 𝑑-Wave Superconductors
\[
## Abstract
An extensive numerical study is reported on disorder effect in two-dimensional $`d`$-wave superconductors with random impurities in the unitary limit. It is found that a sharp resonant peak shows up in the density of states at zero energy and correspondingly the finite-size spin conductance is strongly enhanced which results in a non-universal feature in one-parameter scaling. However, all quasiparticle states remain localized, indicating that the resonant density peak alone is not sufficient to induce delocalization. In the weak disorder limit, the localization length is so long that the spin conductance at small sample size is close to the universal value predicted by Lee (Phys. Rev. Lett. 71, 1887 (1993)).
\]
Since the discovery of the $`d`$-wave pairing symmetry in high-$`T_c`$ cuprates, there has been increased interest in low-energy quasiparticle properties in unconventional superconductors. In such a $`d`$-wave superconductor, quasiparticles are gapless along the four nodal directions on the essentially cylindrical Fermi surface, in contrast to the conventional $`s`$-wave superconductors, where quasiparticles are gapped everywhere. The issues of how the disorder affects the low energy quasiparticle excitations and whether these quasiparticles are localized remain unresolved. Some perturbative self-consistent $`T`$-matrix (SCTM) calculations and a non-perturbative one predicted a nonzero constant density of states (DOS) in low energy region in the presence of weak disorder. However, most of non-perturbative calculations and one numerical study showed that the DOS at zero energy vanishes. With this constant or even vanishing DOS, some groups suggested that all quasiparticle states are localized. On the other hand, it has been shown that a single unitary scattering impurity produces a zero-energy quasiparticle resonant state while the long-range overlap between these impurity states may lead to extended quasiparticle band near zero energy. More recently, a singularity in the DOS at zero energy was obtained by non-perturbative $`T`$-matrix method for the random distributed unitary impurities. It is noteworthy that in one dimension, there is a direct relation between the localization length and the DOS ; thus a singularity in zero-energy DOS signals the delocalization in the system. However, in two dimensions, this theorem does not hold generally . For example, in two dimensional integer quantum Hall systems, it has been shown that the delocalization property at the quantum critical point is not changed by the changing of the DOS due to strong electron-electron interaction . Therefore, the localization of quasiparticles in a $`d`$-wave superconductor in the presence of nonmagnetic unitary impurities is still an open question.
In this Letter, we numerically examine the disorder effect in $`d`$-wave superconductors with nonmagnetic impurities in the unitary limit. The quasiparticle DOS is calculated by exact diagonalization and the spin conductance is computed by the transfer matrix method. It is found that, depending on the particle-hole symmetry of the Hamiltonian, a sharp DOS peak can occur at zero energy and correspondingly the spin conductance is strongly enhanced at finite sample size. However, using one parameter scaling analysis we show that all the quasiparticle states are always localized regardless of the existence of the zero-energy peak in the DOS. In weak disorder limit, the localization length is so long that the spin conductance at small sample size remains close to the universal value $`2\xi _0/a`$ ($`\xi _0`$ is the coherence length of the superconductor and $`a`$ is the lattice constant) in agreement with the theoretical prediction by Lee . A non-universal feature in one-parameter scaling of conductance related to the resonant DOS peak is also discussed.
We begin with a lattice Hamiltonian for the $`d`$-wave superconductor
$`H`$ $`=`$ $`{\displaystyle \underset{ij,\sigma }{}}c_{i\sigma }^{}c_{j\sigma }+{\displaystyle \underset{i,\sigma }{}}(U_i\mu )c_{i\sigma }^{}c_{i\sigma }`$ (2)
$`+{\displaystyle \underset{ij}{}}[\mathrm{\Delta }_{ij}c_i^{}c_j^{}+\text{h.c.}],`$
where $`ij`$ refers to two nearest neighboring sites with the hopping integral taken as the unit, $`\mu `$ is the chemical potential and $`U_i`$ is the impurity potential. We mainly consider the unitary limit where $`U_i`$ takes a nonzero value $`U_0`$ only at a fraction $`n_i`$ of the sites which are randomly distributed in space. The $`d`$-wave symmetry is imposed by choosing order parameters: $`\mathrm{\Delta }_{i,i\pm \widehat{x}}=\mathrm{\Delta }_{i,i\pm \widehat{y}}=\mathrm{\Delta }_d`$, which yields the excitation spectrum $`E_k=\sqrt{ϵ_k^2+\mathrm{\Delta }_k^2}`$, with $`ϵ_k=2(\mathrm{cos}k_x+\mathrm{cos}k_y)\mu `$ and $`\mathrm{\Delta }_k=2\mathrm{\Delta }_d(\mathrm{cos}k_x\mathrm{cos}k_y)`$. Therefore, gapless quasiparticle states exist along the direction $`|k_x|=|k_y|`$ in the momentum space. Unless otherwise stated, $`\mu =0`$, $`\mathrm{\Delta }_d=0.1`$, and $`U_0=100`$ are taken throughout the work. In the presence of a single impurity in the unitary scattering limit, earlier study has shown that the order parameter is strongly suppressed near the impurity site on a scale of a few lattice constant. To take into account this effect, the order parameters on bonds connecting with the strong impurity sites are taken as zero. By exactly diagonalizing the Hamiltonian (1), one can calculate the quasiparticle DOS, which is defined as
$$\rho (E)=\frac{1}{N_L}\underset{n}{}\delta (EE_n),$$
(3)
where $`N_L=L\times L`$ with $`L`$ the linear dimension of the system in units of lattice constant ($`a=1`$). In Fig. 1(a), the DOS is plotted as a function of quasiparticle energy $`E`$ with impurity density $`n_i=0.1`$ ($`\mu =0`$ and $`\mathrm{\Delta }_d=0.1`$) at $`N_L=90\times 90`$. In the calculation the periodic boundary condition is used and it has been checked that the results do not depend on the boundary condition for large $`L`$ considered here. As shown in Fig. 1(a), we find that a sharp zero-energy peak shows up in the DOS, which can be fitted by the analytical form $`c_0n_i/2|E|(\mathrm{ln}^2|E/E_g|+\pi ^2/4)`$ from Ref. with $`c_0=0.66`$ (the superconducting gap $`E_g4\times \mathrm{\Delta }_d`$) as shown in the inset of Fig. 1(a). The fitting breaks down at an energy scale close to $`1/N_L`$ . The strength of the zero-energy peak is reduced when $`n_i`$ is changed to $`0.04`$ as shown in $`E<0`$ part of Fig. 1(b). The overall shape of the peak is sample size independent (from $`N_L=25\times 25`$ to $`120\times 120`$ as well as a strip system $`30\times 300`$) with the DOS value at the peak position increasing with $`N_L`$ very slowly. Note that in the presence of disorder, the order parameter $`\mathrm{\Delta }_{ij}`$ is in principle subject to the self-consistency condition: $`\mathrm{\Delta }_{ij}=g_{ij}c_jc_i`$ where $`g_{ij}=g_0`$ is the attracting interaction for $`d`$-wave pairing. We have also calculated the DOS by diagonalizing the Hamiltonian Eq. (2) self-consistently for each disorder configuration with $`N_L=25\times 25`$. The obtained DOS is also shown in Fig.1(b) ($`E>0`$ part), which is averaged over 50 impurity configurations and $`8\times 8`$ wavevectors in the supercell Brillouin zone. As can be seen, the DOS value at $`E_g`$ is reduced due to the suppression of $`\mathrm{\Delta }_{ij}`$ around each impurity site. However, all other features remain essentially unchanged (compared to the $`E<0`$ part of Fig. 1(b)).
The presence of zero-energy peak in the DOS crucially depends on the symmetry of the Hamiltonian. Since a repulsive or attractive impurity center with infinite strength (i.e., unitary limit) is equivalent to the exclusion of a lattice site, both the local and global particle-hole symmetry remain if the chemical potential $`\mu =0`$, which then produces a resonant peak at $`E=0`$ . This feature was not exhibited in earlier works because of either the breaking of local particle-hole symmetry by the soft impurity scattering or the breaking of the band particle-hole symmetry by considering $`\mu 0`$ , which indicates the importance of the realization of disorder model. In Fig. 1(c), the DOS is presented for random on-site disorders with $`U_i`$ uniformly distributed between \[-1,1\] with a homogeneous order parameter at all sites. Due to the absence of the local particle-hole symmetry in this case, $`\rho (E)`$ has a finite value at low energy region down to a mesoscopic scale $`E1/N_L`$, below which it shows a zero-energy dip, in agreement with those of the non-perturbative calculations . We have also relaxed the unitary scattering limit by taking $`U_0=10`$ (comparable to the band width) or (and) broken the band particle-hole symmetry by taking $`\mu =1`$, and found that the DOS at $`E=0`$ is always strongly suppressed, which is similar to the results obtained in Ref. . There is a smooth crossover from the zero-energy peak to dip in the DOS with the varying of model parameters as long as the Hamiltonian is driven away from the perfect particle-hole symmetry. Therefore, our numerical result indicates that different realizations of disorder give rise to different profiles of the DOS as the energy approaches to the Fermi level.
In the absence of the zero-energy resonant peak in DOS due to the breaking of band particle-hole symmetry by $`\mu 0`$, Franz et al. have studied the similar problem by examining the sensitivity of the wave function to the boundary conditions and by analyzing the finite-size dependence of inverse participation ratios, and presented a strong evidence for the localization of low energy quasiparticles. Here we are concerned with the question of quasiparticle localization or delocalization in the disordered $`d`$-wave superconductor when the zero-energy peak appears in the DOS. We employ the transfer-matrix method to calculate the finite-size localization length and the longitudinal conductance. We consider the quasi-one-dimensional strip sample with the length $`L10^5`$ and width $`M`$. The quasiparticle wavefunction amplitudes in the $`ix`$-th and $`(ix+1)`$-th slices satisfy the following equation:
$$\left(\begin{array}{c}\widehat{\varphi }_{ix+1}\\ \widehat{\varphi }_{ix}\end{array}\right)=T_{ix}\left(\begin{array}{c}\widehat{\varphi }_{ix}\\ \widehat{\varphi }_{ix1}\end{array}\right),$$
(4)
where $`\widehat{\varphi }_{ix}`$ is a $`2M`$-component vector of the Bogoliubov amplitudes for quasiparticle states, and $`T_{ix}`$ is a $`4M\times 4M`$ transfer matrix. The transfer matrix through the whole system, $`P_L=_{ix=1}^LT_{ix}`$, has a set of $`2M`$ pairs of Lyapunov exponents, which determine the inverse of length $`\lambda _i`$ ($`i=1,2,\mathrm{},2M`$). The usual orthonormalization procedure is taken in our calculation. Correspondingly, the longitudinal conductance $`g_s`$ extrapolated for square sample with width $`M`$ is given by :
$$g_s(M,n_i)=\underset{j=1}{\overset{2M}{}}\mathrm{cosh}^2\mathrm{\Lambda }_j,$$
(5)
where $`\mathrm{\Lambda }_j=\lambda _j/M`$. Note that $`g_s`$ corresponds to the spin conductance as the spin carried by quasiparticle is conserved . As shown in Fig. 2(a), $`g_s`$ as a function of $`M`$ monotonically decreases with the increase of $`M`$ at $`E=0`$ for each selected impurity density from $`n_i=0.01`$ to $`0.16`$, consistent with localization in the large $`M`$ limit. All the data between $`M=32`$ and $`M=120`$ at different $`n_i`$ can be collapsed onto a single curve:
$$g_s(M,n_i)=f\left(\frac{\xi (n_i)}{M}\right),$$
(6)
as shown in Fig. 2(b) in accordance with the one-parameter scaling law. Here $`\xi (n_i)`$ is the thermodynamic localization length which only depends on $`n_i`$ as shown in the inset of Fig. 2(b), and it remains finite for all the disorder density $`n_i`$, suggesting that all the states are localized even in the unitary limit with the presence of the zero energy resonant peak. In addition, we display in Fig. 3 the conductance $`g_s`$ as a function of quasiparticle energy $`E`$ at sample width $`M=48`$ and $`96`$ with $`n_i=0.1`$. It is found that $`g_s`$ is strongly enhanced as $`E0`$ in the region $`E0.01`$ corresponding to the width of the DOS peak while away from this region $`g_s`$ generally increases with the increase of $`E`$. For quasiparticle at low energy with $`E<0.1`$, it has been found that $`g_s`$ always decreases with the increase of $`M`$ and all the states are localized in large $`M`$ limit. However, the scaling curve $`g_s(M,n_i)=f(\frac{\xi (n_i)}{M})`$ found at $`E`$ away from the DOS peak ($`E>0.01`$) is different from that for $`E=0`$ at the peak of the DOS, indicating the breaking-down of the universal one-parameter scaling law due to the presence of the resonant peak in the DOS. In addition, for all values of $`\mathrm{\Delta }_d`$, $`g_s`$ decreases with $`M`$, implying localization in all the parameter region. At fixed $`M`$, $`g_s`$ always decreases monotonically with the increasing $`\mathrm{\Delta }_d`$, in agreement with the general argument that quasiparticle states in a superconducting phase are always more localized than the corresponding normal state. Given the localization of the quasiparticle states, the spin conductance at the longest length scales must vanish. However, in the weak disorder (the impurity density $`n_i4\%`$), i.e., the Born limit, we found that the localization length is so long that at small sample size $`M32`$, the spin conductance $`g_s^0`$ follows the universal form $`g_s^0=2\xi _0/a`$ as the coherence length $`\xi _0`$ is changed from $`1.6`$ to $`6.4`$ by changing $`\mathrm{\Delta }_d`$ between $`0.05`$ and $`0.20`$. For the strong disorder, the conductance is generally smaller than the Born limit value due to the onset of localization effect.
In conclusion, we have studied the quasiparticle states in 2D $`d`$-wave superconductors with randomly distributed strong impurities in the unitary scattering limit. As the particle-hole symmetry holds, a very sharp DOS peak is obtained at zero energy. Such a DOS peak enhances the finite-size conductance. However, using one-parameter scaling analysis we have shown that all the quasiparticle states at low energy are still localized and the localization effect is generally enhanced with the increase of superconducting order parameter.
Note added: After the submission of this paper, we received a preprint from Atkinson et al. where similar results for the DOS in the unitary limit of the symmetric band were obtained.
Acknowledgments \- The authors would like to acknowledge A. V. Balatsky, M. P. A. Fisher, P. A. Lee, X.-G. Wen, and Z. Y. Weng for helpful and stimulating discussions. JXZ also thanks W. A. Atkinson for patiently explaining their work to him. This work is supported by the State of Texas through ARP Grant No. 3652707, the Texas Center for Superconductivity at University of Houston, and the Robert A. Welch Foundation.
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# On the fine structure of stationary measures in systems which contract-on-average
## 1. Introduction
Suppose $`\{f_1,\mathrm{},f_m\}`$ is a set of Lipschitz maps of $`^d`$. We may form an iterated function system (IFS) by independently choosing the maps so that the map $`f_i`$ is chosen with probability $`p_i`$ ($`_{i=1}^mp_i=1`$). We denote the the probability vector by $`\overline{p}:=(p_1,\mathrm{},p_m)`$, and the IFS itself will be denoted by $`\mathrm{\Phi }`$.
More precisely, let $`\mathrm{\Omega }=_0^{\mathrm{}}\{1,\mathrm{},m\}`$ and equip $`\mathrm{\Omega }`$ with the product probability measure $`\nu `$ induced in the standard way on cylinder sets by the probability vector $`\overline{p}`$. Let $`x_0^d`$. For any $`\omega \mathrm{\Omega }`$ and any $`n`$ we define the point
$$x_n(\omega ):=f_{\omega _0}\mathrm{}f_{\omega _{n1}}(x_0).$$
If $`lim_n\mathrm{}x_n(\omega )`$ exists, then we define
(1.1)
$$\varphi (\omega )=\underset{n\mathrm{}}{lim}x_n(\omega ).$$
We are going to formulate the hypothesis on $`\mathrm{\Phi }`$ such that the mapping $`\varphi :\mathrm{\Omega }^d`$ will be defined $`\nu `$-a.e. Define
$$h(\overline{p}):=\underset{i=1}{\overset{m}{}}p_i\mathrm{log}p_i.$$
Note that $`h(\overline{p})`$ is the measure-theoretic entropy of the Bernoulli shift $`\sigma :\mathrm{\Omega }\mathrm{\Omega }`$ with the probabilities $`(p_1,\mathrm{},p_m)`$.
For any Lipschitz map $`g`$ of $`^d`$ we let $`g`$ denote the Lipschitz constant of $`g`$. We assume a contraction on average (sometimes called logarithmic average contractivity) condition to hold: for $`\nu `$-a.e. $`\omega \mathrm{\Omega }`$,
(1.2)
$$\underset{n\mathrm{}}{lim}\frac{1}{n}\mathrm{log}^+f_{\omega _0}\mathrm{}f_{\omega _{n1}}=\chi (\mathrm{\Phi })<0.$$
We will call $`\chi (\mathrm{\Phi })`$ the Lyapunov exponent of the system.
Note that the condition (1.2) is implied by the condition
$$\underset{i=1}{\overset{m}{}}p_i\mathrm{log}f_i<0.$$
Contraction on average implies that $`\varphi `$ is a well defined $`\nu `$-measurable function defined by the formula (1.1), it is independent of the choice of initial point $`x_0`$ and that there exists an invariant attracting set $`A\mathrm{\Omega }\times ^d`$ which is the graph of $`\varphi `$. This result is standard and can be found for instance in \[10, Theorem 3\], \[27, Theorem 1.4\]\[5, Theorem 5\], \[2, Theorem 4\] and \[4, Proposition 2.3\].
It is well known (see P. Diaconis and D. Freedman ) that for any such IFS there exists a unique stationary measure $`\mu `$ on $`^d`$ independent of the choice of initial point, i.e, such that $`L^{}\mu =\mu `$, where $`L`$ is the Perron-Frobenius operator for the IFS $`\mathrm{\Phi }`$:
$$L\psi (x)=\underset{i=1}{\overset{m}{}}p_i\psi (f_ix).$$
In fact the measurable mapping $`\varphi `$ induces $`\mu `$ on Borel sets of $`^d`$ by $`\mu (B)=\nu \varphi ^1(B)`$. Sometimes we will also call $`\mu `$ the invariant measure.
By results of L. Dubins and D. Freedman (see also M. Barnsley and J. Elton \[3, Proposition 1\]) on Markov operators, $`\mu `$ must be of pure type, i.e., either absolutely continuous or purely singular with respect to Lebesgue measure on $`^d`$. There are also results due to M. Barnsley and J. Elton \[3, Theorem 3\] about the structure of the support of $`\mu `$ when $`d=1`$ and the maps $`\{f_i\}`$ are affine (see Section 3).
An important classical example of an IFS is the one-parameter family
$`f_0(x)`$ $`=\lambda ^1x,`$
$`f_1(x)`$ $`=\lambda ^1x+1\lambda ^1`$
with $`p_1(0,1)`$. It has been extensively studied since the 1930’s. In recent work by B. Solomyak it was shown that if $`p_1=p_2=\frac{1}{2}`$, then a.e. $`\lambda (1,2)`$ induces an absolutely continuous measure $`\mu `$ on the interval $`[0,1]`$. A similar result was later obtained by the same authors for $`p_1[1/3,2/3]`$ (see Section 3). However the problem of whether the invariant measure (usually called the Bernoulli convolutions or the Erdös measure) for this system is absolutely continuous or singular for a given value of $`\lambda `$ (known as the Erdös Problem), is very hard and only few concrete results are known (see for a nice review and collection of references).
The purpose of this paper is to investigate conditions on IFS which contract on average under which their invariant measure is known to be singular or absolutely continuous. The structure of the paper is as follows: in Section 2 we present an upper bound for the Hausdorff dimension of the invariant measure $`\mu `$ and describe sufficient conditions for $`\mu `$ to be singular in terms of $`\chi (\mathrm{\Phi }),h(\overline{p})`$ and the expansion rate of the semigroup generated by $`\{f_i\}`$. In Section 3 we present several examples showing how to apply the main theorem. Thus, we have reduced the problem of estimating $`dim_H(\mu )`$ to certain combinatorial and algebraic issues concerning the semigroup in question.
We would like to emphasize that although our results apply to a general IFS which contracts-on-average, the most interesting case for us will be the systems in which not all of $`f_i`$ are uniformly contracting, i.e., such that the support of $`\mu `$ is unbounded. One of the reasons for doing so is that there are some indications that $`\text{supp}(\mu )`$ in this case will be “less fractal” than for uniformly contracting systems (see examples below).
## 2. sufficient conditions for singularity of the invariant measure
Let $`G^+`$ denote the semigroup generated by the maps $`\{f_1,\mathrm{},f_m\}`$. Its elements are all compositions $`f_{\omega _0}\mathrm{}f_{\omega _{n1}}`$ for any $`n`$ and $`\omega _k\{1,\mathrm{},m\}`$. It is clear that $`G^+`$ can be either the free semigroup $`_m^+`$ (if all such compositions are different) or a proper subsemigroup of $`_m^+`$. Both possibilities can occur (see Section 3).
Let $`D_n`$ denote the set of all words of length $`n`$ in $`G^+`$. In other words, $`D_n`$ is the set of equivalence classes in $`_0^{n1}\{1,\mathrm{},m\}`$, namely:
$$(\omega _0^{},\mathrm{},\omega _{n1}^{})(\omega _0^{},\mathrm{},\omega _{n1}^{})\text{ if }f_{\omega _0^{}}\mathrm{}f_{\omega _{n1}^{}}=f_{\omega _0^{}}\mathrm{}f_{\omega _{n1}^{}}.$$
From general considerations the growth of $`G^+`$ is exponential, i.e., there exists $`\theta [1,m]`$ such that
(2.3)
$$\theta =\underset{n+\mathrm{}}{lim}\sqrt[n]{\mathrm{\#}D_n}.$$
Indeed, let $`d_n=\mathrm{\#}D_n`$; then $`d_{n+k}d_nd_k`$, because $`(\omega _0,\mathrm{},\omega _{n1})(\omega _0^{},\mathrm{},\omega _{n1}^{})`$ and $`(\omega _n,\mathrm{},\omega _{n+k1})(\omega _n^{},\mathrm{},\omega _{n+k1}^{})`$ imply $`(\omega _0,\mathrm{},\omega _{n+k1})(\omega _0^{},\mathrm{},\omega _{n+k1}^{})`$. Hence, with a little work, (2.3) follows. Obviously, $`\theta 1`$ and if $`G^+`$ is abelian, then $`\theta =1`$.
Let $`H_\mu `$ denote the entropy of the random walk on the semigroup $`G^+`$ with probabilities $`\{p_1,\mathrm{},p_m\}`$. It is defined as follows: let $`\mu _n`$ be the $`n`$’th convolution of $`(p_1,\mathrm{},p_m)`$ on $`D_n`$, i.e.,
$$\mu _n([\omega _0^{},\mathrm{},\omega _{n1}^{}])=\underset{(\omega _0^{},\mathrm{},\omega _{n1}^{})(\omega _0^{},\mathrm{},\omega _{n1}^{})}{}\nu (\omega _0=\omega _0^{},\mathrm{},\omega _{n1}=\omega _{n1}^{}),$$
where $`[]`$ denotes the equivalence class.
We define
$$H_n:=\underset{[y]D_n}{}\mu _n([y])\mathrm{log}\mu _n([y]),$$
and finally, for $`\theta >1`$,
(2.4)
$$H_\mu :=\underset{n\mathrm{}}{lim}\frac{H_n}{n\mathrm{log}\theta }$$
(it is a standard argument that such a limit exists and equals the infimum of the corresponding sequence). For $`\theta =1`$ we set $`H_\mu :=0`$; it is natural, because $`H_n\mathrm{log}\mathrm{\#}D_n`$, whence $`lim_nH_n/n=0`$ in this case. By the definition of $`H_\mu `$ we have
(2.5)
$$0H_\mu \frac{h(\overline{p})}{\mathrm{log}m}1.$$
We will need a version of Shannon’s Theorem for random walks. In the case of discrete groups it was proved independently by Y. Derriennic and V. Kaimanovich and A. Vershik . We will adapt the proof from to our “semigroup” context (see also \[14, Theorem 1.6.4\]).
###### Lemma 2.1.
Let $`\omega \mathrm{\Omega }`$ and
$$_n(\omega )=\{\omega ^{}\mathrm{\Omega }(\omega _0,\mathrm{},\omega _{n1})(\omega _0^{},\mathrm{},\omega _{n1}^{})\}.$$
Then for $`\nu `$-a.e. $`\omega `$,
$$\underset{n+\mathrm{}}{lim}\frac{\mathrm{log}\nu _n(\omega )}{n}=H_\mu \mathrm{log}\theta .$$
Proof: Let $`[\omega ]_n`$ denote the set of all words of length $`n`$ equivalent to $`(\omega _0,\mathrm{},\omega _{n1})`$. We will identify $`[\omega ]_n`$ with $`_n(\omega )`$. Hence $`\nu (_n(\omega ))=\mu _n([\omega ]_n)`$.
Let $`f_n(\omega ):=\mathrm{log}\mu _n([\omega ]_n)`$. By the same reason as in the proof of formula (2.3),
$$\mu _{n+k}([\omega ]_{n+k})\mu _n([\omega ]_n)\mu _k([\omega _n,\mathrm{},\omega _{n+k1}])$$
for any $`\omega \mathrm{\Omega },n1,k1`$. Hence
$$f_{n+k}(\omega )f_n(\omega )+f_k(\sigma ^n\omega ),n,k1,$$
and by Kingman’s Subadditive Ergodic Theorem, there exists the limit $`\overline{f}(\omega )=lim_n\frac{1}{n}f_n(\omega )`$ for $`\nu `$-a.e. $`\omega `$ and
$$\frac{1}{n}_{D_n}f_n𝑑\mu _n\overline{f},n+\mathrm{}$$
as well. It suffices to note that $`_{D_n}f_n𝑑\mu _n=_{[y]D_n}\mu _n([y])\mathrm{log}\mu _n([y])=H_n`$ and apply (2.4).∎
Now we are ready to formulate the main result of this paper.
###### Theorem 2.2.
Suppose $`\mathrm{\Phi }`$ is an IFS on $`^d`$ which contracts on average. Then
(2.6)
$$dim_H(\mu )\frac{H_\mu \mathrm{log}\theta }{\chi (\mathrm{\Phi })}.$$
This has two immediate corollaries:
###### Corollary 2.3.
If
$$\frac{h(\overline{p})\mathrm{log}\theta }{\chi (\mathrm{\Phi })}<d\mathrm{log}m,$$
then $`\mu `$ is singular, and
$$dim_H(\mu )\frac{h(\overline{p})\mathrm{log}\theta }{\chi (\mathrm{\Phi })\mathrm{log}m}<d.$$
In particular, if $`p_1=\mathrm{}=p_m=\frac{1}{m}`$, then
(2.7)
$$dim_H(\mu )\frac{\mathrm{log}\theta }{|\chi (\mathrm{\Phi })|}.$$
Proof: follows from (2.5) and the fact that if $`dim_H(\mu )<d`$, then $`\mu `$ is singular.∎
###### Corollary 2.4.
The measure $`\mu `$ is singular for any $`\mathrm{\Phi }`$ such that
(2.8)
$$d|\chi (\mathrm{\Phi })|>|h(\overline{p})|.$$
Besides, if (2.8) is satisfied, then
$$dim_H(\mu )\frac{h(\overline{p})}{\chi (\mathrm{\Phi })}<d.$$
###### Remark 2.5.
As far as we know, there have been no analogs of Theorem 2.2 in such a general framework. However, F. Przytycki and M. Urbański proved the inequality (2.6) for the case of the Erdös measure $`\mu `$ and Pisot number $`\lambda `$ (and the equality in (2.6) was shown by S. Lalley – see Example 3.1 below).
V. Kaimanovich obtained a similar result for the Hausdorff dimension of the harmonic measure on trees with applications to certain classes of random walks. R. Lyons has a result analogous to Corollary 2.4 in the context of random continued fractions. S. Pincus has related results in the context of mappings on the line and $`2\times 2`$ matrices in the plane. K. Simon, B. Solomyak and M. Urbański have a theorem similar to Corollary 2.4 in the context of parabolic iterated function systems on the real line. Moreover, they were able to establish certain parameter values of their system for which the measure $`\mu `$ is absolutely continuous a.e.
###### Example 2.6.
Let us give a simple example. Suppose $`f_1(x)=2x+1,f_2(x)=\frac{1}{16}x+1`$ chosen with probabilities $`p_1=p_2=\frac{1}{2}`$. Then $`\chi (\mathrm{\Phi })=\frac{3}{2}\mathrm{log}2<h(\overline{p})=\mathrm{log}2`$ and hence by Corollary 2.4, the invariant measure $`\mu `$ is singular with respect to Lebesgue measure, and $`dim_H(\mu )\frac{2}{3}`$. However it is easy to show that the support of the invariant measure is the interval $`[1,\mathrm{})`$. Note that a more detailed analysis shows that since $`f_1f_2f_1^3f_2f_1=f_2f_1^5f_2`$, we have $`\theta <1.9836`$, whence by the formula (2.7), $`dim_H(\mu )\frac{2}{3}\mathrm{log}_2\theta <0.6588`$. For more examples see Section 3.
Proof of Theorem 2.2:
We let $`B(x,r)`$ denote the ball of radius $`r`$ about the point $`x^d`$. It is known (see \[12, Page 171\]) that for any Borel measure $`\mu `$ on $`^d`$,
(2.9)
$$dim_H(\mu )=\mu \text{-ess}sup\left\{\underset{r0}{lim\; inf}\frac{\mathrm{log}\mu B(x,r)}{\mathrm{log}r}\right\}.$$
Let $`B_\omega :=B(\varphi (\omega ),1)`$ and $`f_\omega ^{(n)}:=f_{\omega _0}\mathrm{}f_{\omega _{n1}}`$. Our goal is to show that
(2.10)
$$\mu f_\omega ^{}^{(n)}(B_\omega ^{})\mu (B_\omega ^{})\nu _n(\omega ^{}).$$
We have $`\mu f_\omega ^{}^{(n)}(B_\omega ^{})=\nu (\varphi ^1f_\omega ^{}^{(n)}(B_\omega ^{}))`$ and
$`\varphi ^1f_\omega ^{}^{(n)}(B^{})=\{`$ $`\omega :\underset{k\mathrm{}}{lim}f_{\omega _0}\mathrm{}f_{\omega _n}\mathrm{}f_{\omega _{n+k1}}(x_0)f_\omega ^{}^{(n)}(B^{})\}`$
$`\{`$ $`\omega :(\omega _0,\mathrm{},\omega _{n1})(\omega _0^{},\mathrm{},\omega _{n1}^{}),`$
$`\underset{k\mathrm{}}{lim}f_{\omega _n}\mathrm{}f_{\omega _{n+k1}}(x_0)B_\omega ^{}\},`$
whence by the fact that $`\nu `$ is a product measure (2.10) follows. Hence by Lemma 2.1 for any fixed $`\delta >0`$ for $`\nu `$-a.e. $`\omega ^{}`$ for all sufficiently large $`n`$,
$$\mu f_\omega ^{}^{(n)}(B_\omega ^{})\mu (B_\omega ^{})\theta ^{n(H_\mu +\delta )}.$$
We define $`\gamma =\gamma (\mathrm{\Phi }):=\mathrm{exp}\chi (\mathrm{\Phi })`$, fix $`\delta >0`$ sufficiently small that $`0<\gamma \delta <\gamma +\delta <1`$ and define the sets $`G_N^1`$ and $`G_N^2`$ as follows:
$`G_N^1`$ $`:=\{\omega \mathrm{\Omega }:n>N,(\gamma \delta )^n<f_{\omega _0}\mathrm{}f_{\omega _{n1}}<(\gamma +\delta )^n\},`$
$`G_N^2`$ $`:=\{\omega \mathrm{\Omega }:n>N,\mu f_\omega ^{}^{(n)}(B^{})\mu (B_\omega ^{})\theta ^{n(H_\mu +\delta )}\}.`$
Let $`G_N=G_N^1G_N^2`$. We may choose $`N`$ sufficiently large that $`\nu (G_N)>\frac{3}{4}`$.
Note that $`\nu `$-a.e. $`\omega \mathrm{\Omega }`$ has the property that $`\varphi (\omega )\text{ supp}(\mu )`$. Let $`\omega `$ be such a sequence and define $`A(N):=\{\omega :\mu (B_\omega )>\frac{1}{N}\}`$. Since $`\mu (B_\omega )>0`$, we have $`\nu \left(_{N=1}^{\mathrm{}}A(N)\right)=1`$. As $`A(N+1)A(N)`$, we have $`lim_N\nu (A(N))=1`$ as well.
Hence we may fix $`\alpha >0`$ sufficiently small such that $`\nu \{\omega :\mu (B_\omega )>\alpha \}>\frac{3}{4}`$. Define
$$B:=\{\omega G_N\mu (B_\omega )>\alpha \}.$$
Then $`\nu (B)>\frac{1}{2}`$ and by the fact that $`\sigma `$ preserves the measure $`\nu `$, we have $`\nu (\sigma ^n(B))>\frac{1}{2}`$ for any $`n>0`$.
We claim that for any $`n>N`$ and any $`x\varphi (\sigma ^n(B))`$,
(2.11)
$$\frac{\mu B(x,r)}{_dB(x,r)}C^{}(d)\alpha \theta ^{n(H_\mu +\delta )}(\gamma +\delta )^{dn}$$
for some $`r>0`$ (here $`C^{}(d)`$ is a constant which depends upon the dimension $`d`$ and $`_d`$ is $`d`$-dimensional Lebesgue measure). To prove (2.11), for $`n>N`$ and $`\omega ^{}B`$ we let $`x=\varphi (\sigma ^n\omega ^{})`$ and $`r=f_\omega ^{}^{(n)}`$. Since $`f_\omega ^{}^{(n)}(B^{})B(x,r)`$, we have
$$\mu B(x,r)\alpha \theta ^{n(H_\mu +\delta )}.$$
To estimate $`_dB(x,r)`$, we note that since $`\omega ^{}G_N`$, we have $`r<(\gamma +\delta )^n`$, whence
$$_dB(x,r)C_d(\gamma +\delta )^{dn},$$
where $`C_d`$ is the volume of the unit ball in $`^d`$. This proves (2.11) with $`C^{}(d)=1/C_d`$.
Since $`(\mathrm{\Omega },\sigma ,\nu )`$ is ergodic, for $`\nu `$-a.e. $`\omega `$ we have $`\omega \sigma ^n(B)`$ for infinitely many integers $`n`$. Hence for $`\mu `$-a.e. $`x\text{supp}(\mu )`$ we have $`x\varphi (\sigma ^nB)`$ infinitely often.
This establishes the fact that for a $`\mu `$-generic $`x\text{supp}(\mu )`$ there exists a subsequence $`r_n0`$ such that,
$$\frac{\mu B(x,r_n)}{_dB(x,r_n)}C^{}(d)\alpha \theta ^{n(H_\mu +\delta )}(\gamma +\delta )^{dn},$$
which is equivalent to
$$\frac{\mu B(x,r_n)}{r_n^d}\alpha \theta ^{n(H_\mu +\delta )}(\gamma +\delta )^{dn},$$
since $`C^{}(d)=1/C_d`$ and $`_dB(x,r_n)=C_dr_n^d`$. Taking logarithms and dividing by $`\mathrm{log}r_n`$, we have
$$\frac{\mathrm{log}\mu B(x,r_n)}{\mathrm{log}r_n}d\frac{\mathrm{log}\alpha }{\mathrm{log}r_n}\frac{n}{\mathrm{log}r_n}\left((H_\mu +\delta )\mathrm{log}\theta +d\mathrm{log}(\gamma +\delta )\right).$$
Since $`x=\varphi (\sigma ^n\omega ^{})`$, where $`\omega ^{}G_N`$, we have $`(\gamma \delta )^nr_n(\gamma +\delta )^n`$, whence it follows that for $`\mu `$-a.e. $`x`$,
$`\underset{r0}{lim\; inf}{\displaystyle \frac{\mathrm{log}\mu B(x,r)}{\mathrm{log}r}}`$ $`\underset{r_n0}{lim\; inf}{\displaystyle \frac{\mathrm{log}\mu B(x,r_n)}{\mathrm{log}r_n}}`$
$`d{\displaystyle \frac{(H_\mu +\delta )\mathrm{log}\theta }{\mathrm{log}\gamma }}d{\displaystyle \frac{\mathrm{log}(\gamma +\delta )}{\mathrm{log}\gamma }}.`$
Since $`\delta >0`$ may be taken arbitrarily small and $`\mathrm{log}\gamma =\chi (\mathrm{\Phi })`$, we finally obtain
$$\underset{r0}{lim\; inf}\frac{\mathrm{log}\mu B(x,r)}{\mathrm{log}r}\frac{H_\mu \mathrm{log}\theta }{\chi (\mathrm{\Phi })},$$
and by (2.9) inequality (2.6) holds, which completes the proof.∎
## 3. Examples
We are going to consider several examples, all of which are affine IFS.
###### Example 3.1.
(Bernoulli convolutions). Put $`\mathrm{\Omega }:=_0^{\mathrm{}}\{0,1\}`$ and let $`\lambda >1,f_0(x)=\lambda ^1x,f_1(x)=\lambda ^1x+1\lambda ^1,p_1=p_2=\frac{1}{2}`$ (see Introduction). In this case $`\chi (\mathrm{\Phi })=\mathrm{log}\lambda `$, and
$$f_{\omega _0}\mathrm{}f_{\omega _{n1}}(0)=(1\lambda ^1)\underset{k=0}{\overset{n1}{}}\omega _k\lambda ^k,$$
thus,
$$\varphi (\omega )=(1\lambda ^1)\underset{k=0}{\overset{\mathrm{}}{}}\omega _k\lambda ^k.$$
Hence $`(\omega _0,\mathrm{},\omega _{n1})(\omega _0^{},\mathrm{},\omega _{n1}^{})`$ iff $`_0^{n1}\omega _k\lambda ^k=_0^{n1}\omega _k^{}\lambda ^k`$. Assume first $`\lambda >2`$. Then $`\text{supp}(\mu )`$ is known to be a Cantor set, the semigroup $`G^+`$ is obviously free (as $`f_0([0,1])f_1([0,1])=\mathrm{}`$), and from Corollary 2.4 it follows
(3.12)
$$dim_H(\mu )\frac{\mathrm{log}2}{\mathrm{log}\lambda }<1.$$
If $`1<\lambda <2`$, then $`\text{supp}(\mu )=[0,1]`$; if $`\lambda `$ is transcendental, it is easy to see that $`G^+=_2^+`$, whence $`\theta =2`$ and $`H_\mu =1`$. Hence Corollary 2.4 again gives us the estimate (3.12), which is unfortunately useless, as $`\lambda <2`$. However, in some cases of algebraic $`\lambda `$ Theorem 2.2 can be used more efficiently.
More specifically, assume $`\lambda `$ to be a Pisot number, i.e., an algebraic integer greater than 1 whose conjugates have moduli less than 1. The famous Separation Lemma due to A. Garsia states that there exists a constant $`C=C(\lambda )>0`$ such that if $`_0^n\omega _k\lambda ^k_0^n\omega _k^{}\lambda ^k`$, then $`\left|_0^n(\omega _k\omega _k^{})\lambda ^k\right|C\lambda ^n`$. Hence it is easy to see that $`\theta =\lambda `$, and from (2.6) follows $`dim_H(\mu )H_\mu `$.
In work by S. Lalley the Separation Lemma was used to show that in fact
$$dim_H(\mu )=H_\mu <1.$$
The most transparent subcase is $`\lambda =\frac{1+\sqrt{5}}{2}`$. It was studied in several papers (see references in ); in particular, for this $`\lambda `$ we have $`G^+=a,bab^2=ba^2`$ and $`dim_H(\mu )=H_\mu =0.995713\mathrm{}`$ (this numerical result is due to J. C. Alexander and D. Zagier ). Besides, the measure $`\mu `$ was shown to be quasi-invariant under the $`\beta `$-shift (for $`\beta =\lambda `$) $`\tau _\lambda :[0,1)[0,1)`$ defined by the formula $`\tau _\lambda x=\{\lambda x\}`$ and the corresponding density is also known (see ).
Similar results hold for a more general Bernoulli convolution $`\mu =B_\lambda (p,1p)`$, i.e., the one for which the probability of taking $`f_0`$ is $`p(0,1)`$.
We believe the techniques of \[17, Proposition 4\] can be used to show that the equality holds in a more general situation. Let us formulate the corresponding conjecture; put as above, $`x_n(\omega ):=f_{\omega _0}\mathrm{}f_{\omega _{n1}}(x_0)`$ and $`\gamma =\mathrm{exp}\chi (\mathrm{\Phi })`$.
Conjecture. Suppose we have an affine IFS on the real line (i.e., $`f_i(x)=\lambda _ix+b_i`$) and $`|\lambda _i|1`$ for at least one $`i\{1,\mathrm{},m\}`$; assume the following Weak Separation Condition to be satisfied (we borrow this term from ): for $`\nu ^2`$-a.e. pair $`(\omega ,\omega ^{})\mathrm{\Omega }^2`$ and arbitrary $`\delta >0`$,
(3.13)
$$|x_n(\omega )x_n(\omega ^{})|\text{const}(\gamma \delta )^n$$
whenever $`x_n(\omega )x_n(\omega ^{})`$. We conjecture that the inequality (2.6) in this context is actually an equality. We state without proof that (3.13) does hold in the framework of Example 3.2 with $`\lambda `$ being a Pisot number (see below).
Suppose $`\mathrm{\Phi }`$ is an affine IFS in $``$. If $`\mathrm{\Phi }`$ is not uniformly contracting (i.e., there exists $`i`$ such that $`|\lambda _i|1`$), then by the result from mentioned above, the support of the invariant measure in this case is either a single point or $``$ or $`[d,+\mathrm{})`$ or finally $`(\mathrm{},d]`$ for some $`d`$. We may rule out the first case. The fact that $`\text{supp}(\mu )`$ is connected makes the problem about the fine structure of $`\mu `$ nontrivial.
###### Example 3.2.
Let $`\lambda >1`$ and
(3.14)
$$f_1(x)=\lambda ^1x,f_2(x)=x+1\text{ and }p_1=p_2=\frac{1}{2}.$$
The support of $`\mu =\mu (\lambda )`$ is $`[0,+\mathrm{})`$, and $`\chi (\mathrm{\Phi })=\frac{1}{2}\mathrm{log}\lambda `$. Hence by Corollary 2.4, for $`\lambda >4`$ the measure $`\mu `$ is singular, and $`dim_H(\mu )\frac{2\mathrm{log}2}{\mathrm{log}\lambda }<1`$.
We claim that for any transcendental $`\lambda `$ the semigroup $`G^+`$ is free. A trivial induction argument shows that
$$f_1^{n_1}f_2^{k_1}\mathrm{}f_1^{n_s}f_2^{k_s}(x)=\lambda ^{_1^sn_j}x+\underset{j=1}{\overset{s}{}}k_j\lambda ^{_{i=1}^jn_i},$$
whence if $`\lambda `$ is not algebraic,
$$f_1^{n_1}f_2^{k_1}\mathrm{}f_1^{n_s}f_2^{k_s}=f_1^{n_1^{}}f_2^{k_1^{}}\mathrm{}f_1^{n_s^{}}f_2^{k_s^{}}$$
implies $`n_jn_j^{},k_jk_j^{},j=1,\mathrm{},s`$. Hence $`\mathrm{\#}D_n=2^n,\theta =2`$ and $`G^+=_2^+`$.
It is worth noting that for certain particular values of $`\lambda (1,4)`$ the measure $`\mu `$ is nonetheless singular (similarly to the Bernoulli convolutions). Namely, since $`\mu `$ is invariant under the IFS $`\mathrm{\Phi }`$, we have the following self-similar relation for its Fourier transform:
$$\widehat{\mu }(x)=\frac{1}{2}\widehat{\mu }(\lambda ^1x)+\frac{1}{2}e^{ix}\widehat{\mu }(x),$$
whence
(3.15)
$$\widehat{\mu }(x)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{2\mathrm{exp}(i\lambda ^nx)}.$$
Assume again $`\lambda `$ to be a Pisot number. Then as is well known, the distance to the nearest integer for $`\lambda ^n`$ tends to 0 at exponential rate. Following the line of the proof of the classical work (see also for the case $`\lambda =2`$), we can consider the subsequence $`x_n=2\pi \lambda ^n`$ and show that $`\widehat{\mu }(x_n)\to ̸0`$ as $`n+\mathrm{}`$, which implies that the Riemann-Lebesgue Lemma is not satisfied, whence $`\mu `$ cannot be absolutely continuous. Therefore, it is singular by the Law of Pure Types.
Thus, we proved that $`\mu `$ is singular for $`\lambda 4`$ (as 4 is a Pisot number); at the same time it is singular for an infinite number of parameters $`\lambda (1,4)`$ as well. It is an open question whether its Hausdorff dimension is less than 1 for a Pisot number $`\lambda `$ (for the Bernoulli convolutions it is true, see above).
It is worth mentioning that in this example the stationary measure has an “arithmetic” interpretation as well. Namely, let $`\mathrm{\Sigma }=_1^{\mathrm{}}_+`$ and $`\xi `$ denote the stationary product measure on $`\mathrm{\Sigma }`$ with the following geometric distribution: $`\xi (\epsilon _n=k)=2^{k1},k=0,1,\mathrm{}`$ Let $`L_\lambda :\mathrm{\Sigma }_+`$ be defined as follows:
$$L_\lambda (\epsilon _1,\epsilon _2,\mathrm{})=\underset{n=1}{\overset{\mathrm{}}{}}\epsilon _n\lambda ^n$$
(it is obvious that $`L_\lambda `$ is well defined for $`\xi `$-a.e. $`\epsilon \mathrm{\Sigma }`$). Then from (3.15) it follows
$$\mu =\xi L_\lambda ^1,$$
as $`\widehat{\xi }(x)=\frac{1}{2}+\frac{1}{4}e^{ix}+\frac{1}{8}e^{2ix}+\mathrm{}=\frac{1}{2e^{ix}}`$. Thus, the essential difference with the case of Bernoulli convolutions is that the set of “digits” here is infinite.
When this paper was in preparation, Y. Peres suggested the following claim.
###### Lemma 3.3.
For $`_1`$-a.e. $`\lambda (1,32^{2/3})`$ the measure $`\mu `$ is absolutely continuous.
Proof. Since
$$\frac{1}{2e^{ix}}=\frac{2+e^{ix}}{4e^{2ix}}=\frac{\frac{2}{3}+\frac{1}{3}e^{ix}}{\frac{4}{3}\frac{1}{3}e^{2ix}},$$
by (3.15) the measure $`\mu `$ is the convolution of the Bernoulli measure $`B_\lambda (2/3,1/3)`$ and the stationary measure for the IFS
$`g_1(x)`$ $`=\lambda ^2x,`$
$`g_2(x)`$ $`=x+1`$
with $`p_1=\frac{1}{4},p_2=\frac{3}{4}`$. In it was shown that for any $`p[1/3,2/3]`$ the Bernoulli measure is absolutely continuous for a.e. $`\lambda (1,p^p(1p)^{1+p})`$. Hence the measure $`B_\lambda (2/3,1/3)`$ will be absolutely continuous for a.e. $`\lambda (1,32^{2/3})`$ and so will be $`\mu `$.∎
###### Remark 3.4.
Since
$$\frac{1}{1z}=\underset{n=0}{\overset{\mathrm{}}{}}\left(1+z^{2^n}\right),|z|<1,$$
it is easy to deduce that
$$\frac{1}{2e^{ix}}=\underset{n=0}{\overset{\mathrm{}}{}}\left(\frac{2^{2^n}}{2^{2^n}+1}+\frac{1}{2^{2^n}+1}\mathrm{exp}\left(2^nix\right)\right),$$
and by (3.15)
(3.16)
$$\mu =B_\lambda (\frac{2}{3},\frac{1}{3})B_{\lambda ,2}(\frac{4}{5},\frac{1}{5})\mathrm{}B_{\lambda ,2^n}(\frac{2^{2^n}}{2^{2^n}+1},\frac{1}{2^{2^n}+1})\mathrm{},$$
where $`B_{\lambda ,t}(p,1p)`$ is the infinite convolution of two-point discrete measures whose “basic” distribution is supported on the points 0 and $`t`$ with probabilities $`p`$ and $`1p`$ respectively (hence $`B_\lambda =B_{\lambda ,1}`$).
Summing up, we have the following
###### Proposition 3.5.
For the IFS (3.14) the following properties are satisfied:
1. For $`\lambda 4`$ the measure $`\mu `$ is singular, and $`dim_H(\mu )\frac{2\mathrm{log}2}{\mathrm{log}\lambda }`$;
2. for $`_1`$-a.e. $`\lambda (1,32^{2/3})`$ it is absolutely continuous;
3. for any Pisot $`\lambda `$ it is singular.
A natural question to ask is what happens between $`32^{2/3}1.8899`$ and 4. Note that if $`\lambda >32^{2/3}`$, then all convolutions in (3.16) are singular. Nonetheless, we conjecture that for $`_1`$-a.e. $`\lambda (32^{2/3},4)`$ the measure $`\mu `$ will be absolutely continuous as well.
###### Example 3.6.
Let $`\lambda >1,f_1(x)=\lambda x+1,f_2(x)=\lambda ^2x`$, and $`p_1=p_2=\frac{1}{2}`$. Here $`\chi (\mathrm{\Phi })=\frac{1}{2}\mathrm{log}\lambda `$. Let for the simplicity of notation $`a=f_1,b=f_2`$. Since
(3.17)
$$ababa=ba^3b,$$
we have $`d_{n+1}2d_nd_{n4}`$, whence $`\theta <1.9277`$. Thus, from the estimate (2.7) it follows that the measure $`\mu `$ is singular at least for $`\lambda >1.9277^23.716`$. However, it is clear that the actual estimate must be even sharper, because in fact there are infinitely many relations between $`a`$ and $`b`$.
Namely, from general considerations it follows that $`G^+`$ will be one and the same for any transcendental $`\lambda `$. We do not know whether $`G^+`$ is finitely presented but at least written in the generators $`a,b`$ it is not. For example, for any $`n=2,3,\mathrm{}`$ and any $`k=1,2,\mathrm{},n1`$ we have in addition to the relation (3.17),
$$a^{2k}b^na^{2(nk)}=b^{nk}a^{2n}b^k$$
(these are a direct consequence of the fact that $`a^{2k}b^k`$ and $`a^{2n}b^n`$ always commute). These relations are independent, i.e., none of them is a consequence of any other one. At the same time, there are relations that cannot be deduced from any of those described above; for instance, $`abab^2a^3=b^2a^5b`$. There are indications that actually $`\theta `$ is at least less than 1.7.
As far as we are concerned, there are no general results on the structure of $`\text{supp}(\mu )`$ in the case of higher dimensions. We are going to present an example of a family of IFS for $`d=2`$ such that $`\text{supp}(\mu )=^2`$, whereas $`\mu `$ is singular; at the same time this system does not “split” into one-dimensional actions. In a certain sense the following example is a two-dimensional generalization of Example 3.2.
###### Example 3.7.
Let $`\alpha `$ be a real number such that $`\alpha /\pi `$ is irrational and let $`R_\alpha `$ denote the rotation of $`^2`$ by the angle $`\alpha `$, i.e.,
$$R_\alpha =\left(\begin{array}{cc}\mathrm{cos}\alpha & \mathrm{sin}\alpha \\ \mathrm{sin}\alpha & \mathrm{cos}\alpha \end{array}\right).$$
Let $`\lambda >1`$ and the one-parameter family of IFS $`\mathrm{\Phi }_\lambda `$ be defined as follows:
$`f_1(x)`$ $`=A_\lambda ^1x,`$
$`f_2(x)`$ $`=x+\overline{e},`$
where $`A_\lambda =\lambda R_\alpha `$ and $`\overline{e}S^1`$ is fixed. As above, we assume $`p_1=p_2=\frac{1}{2}`$.
###### Proposition 3.8.
1. For any $`\lambda >1`$ the IFS $`\mathrm{\Phi }_\lambda `$ contracts on average and the support of the invariant measure $`\mu _\lambda `$ is full, i.e.,
(3.18)
$$\text{supp}(\mu _\lambda )=^2.$$
2. For $`\lambda >2`$ the measure $`\mu _\lambda `$ is singular, and
(3.19)
$$dim_H(\mu _\lambda )<\frac{2\mathrm{log}2}{\mathrm{log}\lambda }<2.$$
Proof: By the same reason as in the previous examples, we have $`\chi (\mathrm{\Phi }_\lambda )=\frac{1}{2}\mathrm{log}\lambda `$, whence for any $`\lambda >1`$ the system contracts on average. From Corollary 2.4 it follows that $`\lambda >2`$ implies the singularity of $`\mu _\lambda `$ together with (3.19).
The most delicate part of the proposition is the relation (3.18). Let us prove it. Assume $`:=\text{supp}(\mu _\lambda )^2`$; then there exists a disc $`B(x,\delta )`$ whose intersection with $``$ is empty. Hence by definition, $`f_1^nB(x,\delta )=\mathrm{}`$ for any $`n1`$. We have
$$f_1^nB(x,\delta )=B(y_n,\lambda ^n\delta ),$$
where $`y_n=f_1^n(x)=\lambda ^nR_\alpha ^n(x)`$. Since $`\alpha /2\pi `$ is irrational, the rotation $`R_\alpha :S^1S^1`$ is minimal, i.e., the orbit of every point is dense in $`S^1`$ (see, e.g., ). We apply this claim to the circle of radius $`x`$. Thus, for any $`\epsilon >0`$ there exists $`n`$ such that
(3.20)
$$R_\alpha ^n(x)x\overline{e}<\epsilon .$$
Fix $`r>1,\epsilon =\delta /2`$ and $`n`$ large enough to satisfy $`\lambda ^n2r/\delta `$ together with (3.20). Let $`z=\lambda ^nx\overline{e}`$; we claim that
$$B(z,r)B(y_n,\lambda ^n\delta ).$$
Indeed, let $`yB(z,r)`$, i.e., $`yzr`$. Hence
$`yy_n`$ $`yz+zy_n`$
$`r+\lambda ^n\epsilon =r+{\displaystyle \frac{1}{2}}\lambda ^n\delta <\lambda ^n\delta .`$
Hence $`B(z,r)=\mathrm{}`$, and $`f_2^kB(z,r)=\mathrm{}`$ as well for any $`k0`$. Since $`\overline{e}=1`$ and $`z`$ belongs to the half-line $`\{t\overline{e},t0\}`$, there exists $`k0`$ such that $`f_2^kB(z,r)B(0,r1)`$. Hence for any $`r>1,B(0,r1)=\mathrm{}`$, which means $`=\mathrm{}`$. The proposition is proven.∎
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# Naturalness Reach of the Large Hadron Collider in Minimal Supergravity
## 1 Introduction
A possibility for new physics beyond the standard model is supersymmetry (SUSY). If fermionic generators are added to the bosonic generators of the Lorentz group, the new space-time symmetry is supersymmetry. As a result of exact supersymmetry, all particles have a partner of equal mass but opposite spin-statistics. Cancellations between bosonic and fermionic loops prevent radiative corrections from driving scalar masses up to the highest scale present, assumed to be the GUT or Planck scale, $`10^{16}`$ to $`10^{19}`$ GeV, solving the naturalness problem of the standard model. In addition, the renormalised electromagnetic, weak, and strong couplings can be made to converge to an approximately common value at the grand unification scale.
Since supersymmetry is not observed amongst the already discovered particles, it must be a broken symmetry. The scale at which supersymmetry is broken, $`M_{\text{SUSY}}`$, would be the typical mass of the as yet undiscovered superpartners of the standard model particles, and represents the scale at which this new physics becomes relevant. Considerations of general new physics beyond the standard model can result in upper bounds on new physics scales of order a few TeV. However, in supersymmetric models , $`M_{\text{SUSY}}`$ is expected to be at most around 1 TeV. If the SUSY breaking scale is too large, then, unless there are large cancellations between SUSY breaking parameters, electroweak symmetry breaking also will be of order $`M_{\text{SUSY}}`$, and the $`W`$ and $`Z`$ bosons would have masses inconsistent with their measured values.
If $`M_{\text{SUSY}}`$ is at or below the TeV scale, then supersymmetric particles will almost certainly be discovered at the Large Hadron Collider (LHC), being built at CERN. Indeed, detailed studies of how the SUSY parameters would be measured by the LHC general-purpose experiments ATLAS and CMS have been made .
The simplest possible SUSY extension of the standard model, with a superpartner for each standard model particle, and the addition of a second Higgs scalar doublet, is called the minimal supersymmetric standard model, or MSSM. The most studied sub-category of the MSSM is minimal supergravity, mSUGRA. Supergravity, where supersymmetry is a local, rather than a global symmetry, at one time motivated unification assumptions amongst the MSSM SUSY breaking parameters, reducing the number of parameters, from the hundred or so of the MSSM, to just four, plus one sign. Currently, the suppression of flavour changing neutral currents, not supergravity, is the main motivator for these assumptions. The theory is fully specified by these parameters together with those of the standard model.
As mentioned above, it is possible to avoid the problem of large electroweak boson masses if there are extra cancellations amongst the SUSY masses. However, such situations, where the parameters of a theory are carefully tuned to avoid unphysical results, are often thought to be unsatisfactory. Fundamental parameters, it is argued, should be independent, uncorrelated inputs.
These ‘naturalness’ arguments are often quantified in terms of ‘fine tuning’ . There are a few different fine-tuning measures , and all are intended to be measures of the degree of cancellation required between fundamental parameters. A value of fine-tuning above which a theory becomes unacceptable is often advanced, and used to support the argument that $`M_{\text{SUSY}}`$ should be small.
The authors of discuss a more sophisticated measure of fine-tuning and use it to assess the status of the MSSM if superparticles are not found at the LHC. In fine-tuning motivated upper bounds on MSSM masses are obtained using a complete one-loop effective potential. discusses the use of fine-tuning to compare different high-energy supergravity scenarios, while uses fine-tuning to compare non-minimal supersymmetric models.
Since its status as a possible solution of the naturalness problem of the standard model is one of the main reasons for investigating supersymmetry, naturalness arguments have increased relevance to studies of supersymmetry.
We therefore contend that fine-tuning is a relevant way to compare SUSY models, experiments, and search channels, and is a useful measure of experimental discovery reach. As experiments push the lower bounds on SUSY parameters upwards, the minimum fine-tuning which SUSY can have increases, and our confidence that the universe is supersymmetric at low energies falls. However, we do not believe that a high fine-tuning in itself can be used to rule a theory out.
Within mSUGRA, the fundamental SUSY breaking parameters are boundary conditions on the running SUSY breaking masses and couplings imposed at a high scale, usually taken to be $`10^{16}`$ GeV. Physical masses of superparticles are obtained by evolving the MSSM parameters to the weak scale using the renormalisation group equations (RGEs). The RGE evolution of minimal supergravity shows a ‘focus point’ behaviour . A relatively large region of GUT scale parameters exists for which the RGE trajectories converge towards a small range of measurable properties. Specifically, the renormalisation group trajectories of the mass squared of a Higgs doublet ($`m_{H_2}^2`$) cross close to the electroweak scale.
As a result, the electroweak symmetry breaking is insensitive to the GUT scale SUSY breaking parameters , and fine-tuning is smaller than expected. This focus point corresponds to a region in which the scalar SUSY breaking masses, governed by the mSUGRA parameter $`m_0`$, may be large.
Previous predictions of the discovery reach of the LHC in mSUGRA parameter space, using ISAJET , went only as far as $`m_0<2`$ TeV . The purpose of the present investigation is to extend this reach to higher $`m_0`$ and to present it in terms of a naturalness measure. We seek to determine how the standard SUSY search channels perform in this region, where charginos and neutralinos would be the dominant SUSY particles. For large $`m_0`$, squarks, sleptons, and the heavy Higgs particle could avoid detection at the LHC, and the determination of the $`m_0`$ reach of searches for these particles would be an interesting further study.
In section 2, we discuss our fine-tuning measure, and which parameters to include in its definition. We discuss in section 3 the matter of the electroweak symmetry breaking excluded region. In section 4 we discuss SUSY search channels considered for use at the LHC, and our simulation of the discovery reach using the HERWIG event generator. We present in section 5 the fine-tuning reach in each channel, and the overall fine-tuning reach of the LHC.
## 2 Fine tuning
At tree-level, in the MSSM, the $`Z`$ boson mass is determined to be :
$$\frac{1}{2}M_Z^2=\frac{m_{H_1}^2m_{H_2}^2\mathrm{tan}^2\beta }{\mathrm{tan}^2\beta 1}\mu ^2$$
(1)
by minimising the Higgs potential. $`\mathrm{tan}\beta `$ is the ratio of Higgs vacuum expectation values (VEVs) $`v_1/v_2`$ and $`\mu `$ is the Higgs mass parameter in the MSSM superpotential. In mSUGRA, $`m_{H_2}`$ has the same origin as the super-partner masses ($`m_0`$). Thus as search limits put lower bounds upon super-partners’ masses, the lower bound upon $`m_0`$ rises, and consequently so does $`|m_{H_2}|`$. A cancellation is then required between the terms of equation 1 in order to provide the measured value of $`M_Z|m_{H_2}|`$. Various measures have been proposed in order to quantify this cancellation .
The definition of naturalness $`c_a`$ of a ‘fundamental’ parameter $`a`$ employed in reference is
$$c_a\left|\frac{\mathrm{ln}M_Z^2}{\mathrm{ln}a}\right|.$$
(2)
From a choice of a set of fundamental parameters $`\{a_i\}`$, the fine-tuning of a particular model is defined to be $`c=\text{max}(c_a)`$. Our initial choice of free, continuously valued, independent and fundamental mSUGRA parameters also follows ref. :
$$\{a_i\}=\{m_0,M_{1/2},\mu (M_{\text{GUT}}),A_0,B(M_{\text{GUT}})\}$$
(3)
where $`M_{\text{GUT}}10^{16}`$ GeV is the GUT scale. It is this selection which gives rise to low fine-tuning for large $`m_0`$.
We have calculated $`c`$ numerically to one-loop accuracy in soft masses, with two-loop accuracy in supersymmetric parameters. Dominant one-loop top/stop corrections were added to the Higgs potential and used to correct eq. (1). The Higgs potential was minimised at $`Q=\sqrt{(m_{\stackrel{~}{t}_1}m_{\stackrel{~}{t}_2})}`$, where its scale dependence is small.
Full one-loop sparticle and QCD corrections were used to determine $`m_t(m_t)^{\overline{\text{DR}}}`$ and the running of $`\mathrm{tan}\beta `$ was taken into account in order to calculate the Yukawa couplings from fermion running masses. Fermion running masses were determined at $`M_Z`$ by evolving them with 3-loop QCD $``$ 1-loop QED. $`m_b(M_Z)^{\overline{\text{DR}}}`$ and $`m_\tau (M_Z)^{\overline{\text{DR}}}`$ were determined by including one-loop SUSY QCD and third family corrections . The $`\overline{\text{DR}}`$ Higgs vacuum expectation value $`v\sqrt{v_1^2+v_2^2}`$ was determined by the approximate formula
$$v(Q)=248.6+0.9\mathrm{ln}\left(\frac{m_{\stackrel{~}{u}_L}}{Q}\right)\text{GeV}.$$
(4)
One-loop top, gluino and squark corrections were used in order to deduce $`\alpha _S^{\overline{\text{DR}}}(M_Z)`$ from $`\alpha _S^{\overline{\text{MS}}}(M_Z)=0.119`$. We will examine the implications of relaxing the accuracy of the calculation below.
Note that our code does not yet include 2-loop soft terms in the RGEs, finite corrections to the electroweak gauge couplings, a full one-loop calculation of the Higgs vacuum expectation value, or tadpole contributions to the Higgs potential. It is well known that improving the accuracy of the calculation of $`c`$ at this level can make a significant difference to its numerical value. Since its exact quantitative interpretation is obscure anyway, this fact is not in conflict with the proposed comparison of naturalness reaches in different channels.
We also consider the case where the top Yukawa coupling $`h_t(M_{\text{GUT}})`$ is added to the list of fundamental parameters in eq.(3). In refs. , it is shown that the focus point scenario with heavy scalars has a large fine-tuning if $`h_t`$ is included. The inclusion of $`h_t(M_{\text{GUT}})`$ into the definition of fine tuning thus increases the naturalness reach of the LHC, as our results in section 5 show.
## 3 Radiative electroweak symmetry breaking
Most plots of search reach in the $`m_0`$,$`M_{1/2}`$ plane found in the literature show a large excluded triangular region, for large $`m_0`$ and small $`M_{1/2}`$. In this region, mSUGRA does not have the required properties for radiative electroweak symmetry breaking (REWSB). However, the REWSB constraint, unlike most of the SUSY spectrum, is very sensitive to details of the calculation and input parameters, particularly the value of the running top mass $`m_t(m_t)`$, which must be calculated from the pole top mass in the renormalisation scheme being used. The precise relationship between the running and pole masses depends on the masses of the superparticles , and hence on the point in the SUSY parameter space being investigated.
In version 7.14 of ISASUGRA, which was used in to generate SUSY masses and mixings and to obtain the REWSB excluded region, the region above $`m_02M_{1/2}+1`$ TeV was excluded for an input pole top mass of 174 GeV. Here one-loop RGEs were used and minimisation of the scalar potential was performed at scale $`M_Z`$ .
Between 7.14 and more recent versions such as 7.42, the excluded region shifted to very high $`m_0`$, around 6 TeV for an input pole top mass of 174 GeV, as shown in figure 1. The REWSB constraint therefore no longer provides a useful limit on the values of $`m_0`$ which should be investigated. Here and in later versions, the scalar potential minimisation takes place at $`Q=\sqrt{m_{\stackrel{~}{t}_L}m_{\stackrel{~}{t}_R}}`$ . In the most recent version of ISASUGRA, 7.51, the region has returned once again to low values of $`m_0`$, as shown also in in figure 1. Here, two-loop corrections to the RGEs and one-loop sparticle mass corrections to the inputs have been added .
The REWSB exclusion region was also found using our fine-tuning calculation code, where the physics included was as described in the previous section. The dependence of the exclusion region on various aspects of the physics is shown in figure 2. We note again a strong dependence on the top mass, and on the treatment of the running mass and strong coupling.
Since the REWSB constraint is so heavily dependent upon the precise values of input parameters and details of the calculation, there is no reason to suppose that the exclusion region will remain stable under further refinements. Hence we believe it should not be taken as a limit upon the regions of mSUGRA parameter space to be investigated experimentally. Both our own code and ISASUGRA could be subject to large corrections. We therefore use ISASUGRA7.42 to compute the discovery reach, since the allowed region then extends up to values of $`m_0=10000`$ GeV (see figure 1). As we shall show, the discovery contour in the allowed region is not so sensitive to the version number or top mass. To produce the REWSB excluded regions displayed in the plots of our results, we used our fine-tuning code with the pole top masses specified in the captions.
## 4 Monte-Carlo simulation of LHC discovery reach
We now turn to the discussion of the LHC mSUGRA search.
In this study, as in , the results of which were used in , the SUSY discovery reach in the $`m_0`$, $`M_{1/2}`$ plane was found through a variety of signals involving hard isolated leptons, jets and missing transverse momentum. These cuts can be applied with a range of cut energies $`E_{cut}`$, and are listed below. The cuts represent typical SUSY discovery cuts that might be used at a general-purpose LHC experiment, such as ATLAS or CMS .
* Missing transverse momentum $`/p_T>E_{cut}`$
* Either: At least 2 jets, with pseudo-rapidity $`|\eta |<3.2`$, and $`p_T>E_{cut}`$, using a cone algorithm with cone-size 0.7 units of $`\eta ,\varphi `$. Or: No jets with $`p_T>`$ 25 GeV, using the same jet-finding algorithm.
* Any final state ($`\mu `$ or $`e`$) leptons, with $`|\eta |<2.0`$ and $`p_T>20`$ GeV, and lying further in $`\eta ,\varphi `$ space than 0.4 units from the centre of any 15 GeV jet with cone-size 0.4, and with less than 5 GeV of energy within 0.3 units of the lepton.
Channels with no leptons, one lepton, two leptons of opposite or the same sign, and three leptons were investigated.
* In the one-lepton channel, an extra cut was imposed to reduce the background due to standard model $`W`$ decay, involving the transverse mass:
$$M_T=\sqrt{2(|𝐩_𝐥||/𝐩_𝐓|𝐩_𝐥/𝐩_𝐓)}>100\text{ GeV},$$
(5)
where $`𝐩_𝐥`$, $`/𝐩_𝐓`$ are transverse two-component momenta, for the lepton and missing $`p_T`$ respectively.
* $$S_T=\frac{2\lambda _2}{\lambda _1+\lambda _2}>0.2,$$
(6)
where $`\lambda _i`$ are the eigenvalues of the matrix $`S_{ij}=\mathrm{\Sigma }_{ij}p_ip_j`$, the sum being taken over all detectable final state particles, and $`p_i`$ being the two-component transverse momentum of the particle. $`S_T`$ is often called the transverse circularity or transverse sphericity.
The search channels used are shown in table 1.
The mSUGRA events were simulated by employing the ISASUSY part of the ISAJET7.42 package to calculate sparticle masses and branching ratios, and HERWIG6.1 to simulate the events themselves. The expected SUSY signal was generated, together with the backgrounds due to standard model top anti-top and $`W`$ plus jet events. The $`W`$ background was generated using events with one of the jets produced in the hard subprocess, and the rest in QCD parton cascades. This produces an underestimate of the $`W`$ background, as discussed in , and the results were rescaled accordingly.
The discovery limit is set at values of $`m_0`$ and $`M_{1/2}`$ where $`S/\sqrt{B}>5`$ and $`S>10`$ where S and B are the expected number of events in the the SUSY signal and total background. The former constraint is an approximation to the requirement that the total number of observed events in some experiment will be significantly above the background at the 5 $`\sigma `$ level.
For each value of $`E_{cut}`$, the minimum final state transverse momentum in the hard subprocess was selected to obtain some events passing the cuts, within realistic computer time constraints. Where the minimum transverse momentum had to be increased above $`E_{cut}`$, results from the j0 channel were used to obtain linear correction factors to compensate for the resulting underestimate of the background.
In order to find the hard leptons, and to determine $`/𝐩_𝐓`$ and $`S_T`$, the Monte Carlo output was examined directly, particles with $`|\eta |>5.0`$ being ignored. For the jet count and lepton isolation check, a calorimeter simulation was used, with hadronic resolution of $`70\%/\sqrt{E}`$, electromagnetic resolution of $`10\%/\sqrt{E}`$ and extending to pseudo-rapidities up to 5.0. These parameters have been selected to simulate a general-purpose experiment at the LHC.
Calculations were made for mSUGRA with $`A_0=0`$, $`\mathrm{tan}\beta =10`$, pole top mass 174 GeV, $`\mu >0`$ and $`=10`$ fb<sup>-1</sup> of integrated luminosity. The integrated luminosity chosen is equivalent to one year of running in the low luminosity mode, $`10^{33}`$ cm<sup>-2</sup> s<sup>-1</sup>. We note that the LHC experiments expect to collect around 300 fb<sup>-1</sup> of data.
A low-statistics calculation of the discovery reach with a top mass of 179 GeV was also made, to determine whether the reach is independent of the top mass. There is a small decrease in the production cross section for $`t\overline{t}`$ background with increasing top mass, and a compensating increase in the energy of the resulting leptons and jets, which causes a greater proportion to pass the selection cuts. Thus for this study the background could safely be assumed to be independent of the top mass.
## 5 Results
Preliminary results on the naturalness reach of the LHC for the j1 channel, with $`E_{cut}=400`$ GeV, were presented in . The top Yukawa coupling was not included in the definition of fine-tuning. Here we obtain the discovery limit in all the channels in table 1, for several values of $`E_{cut}`$, and with both definitions of fine-tuning (with/without the top Yukawa coupling).
Table 2 shows the background for each generated channel, the total background, and the resulting expected number of supersymmetry events required to fulfill the discovery criteria detailed above. No discernible background in the jet veto channels (v2,v3) was obtained.
For each of the channels, and each value of $`E_{cut}`$, a plot has been used to determine the fine-tuning reach with both definitions of fine-tuning. The fine-tuning reach is the largest fine tuning where the fine-tuning contour is completely within the discovery limit, for $`m_0<4000`$ GeV. The values of $`E_{cut}`$ giving the best fine-tuning reach, and the corresponding reach, are shown in table 3. The results using $`m_t=174`$ GeV and $`m_t=179`$ GeV to calculate the fine-tuning and REWSB exclusion are given. $`E_{cut}=`$100, 200, 300, 400 and 500 GeV were investigated. No results can be presented for the jet veto channels, as the number of events in these channels is too small. The reach in these channels could be obtained by considering specific SUSY processes, but is expected to be very limited.
Figures 3 and 4 show the best obtainable discovery limit using the channel providing the best reach, the $`E_{cut}`$ = 400 GeV j1 channel, as a dashed line in the $`(m_0,M_{1/2})`$ plane, and the naturalness density without and with the top Yukawa coupling included, respectively. The resulting fine-tuning limits, as defined above, are 210 and 500 units. Excluded regions due to experimental limits on the chargino mass, and due to the cosmological requirement that the LSP be neutral, are included in the figures. The current limit on the light Higgs mass does not exclude any of the region of mSUGRA parameter space illustrated, due to the large value of $`\mathrm{tan}\beta `$ used here.
Figure 5 shows consistency of our results for $`m_0<2`$ TeV in the j1 channel with those obtained using ISAJET 7.14 in . The same figure shows the discovery reach calculated with ISAJET7.51 (and lower statistics). We see that that the differences in approximation between the various ISAJET versions which shifted the REWSB exclusion, as discussed in section 3, do not appreciably alter the discovery reach.
Also shown in figure 5 is the discovery reach for a 1$`\sigma `$ increase in the top mass, $`m_t=179`$ GeV. Unlike the REWSB exclusion region, the discovery contours are almost identical. We may therefore use our high-statistics $`m_t=174`$ GeV ISAJET 7.42 discovery contour to calculate the fine-tuning reach for $`m_t=179`$. As shown in table 3 and figure 6, the reach extends to a fine-tuning value of 260.
For $`m_0>4000`$ GeV the discovery reach in $`M_{1/2}`$ becomes roughly independent of $`m_0`$. Here, the SUSY processes involved are dominated by the gauginos, the squark masses being very high (larger than 3 TeV), as illustrated in figure 7. This figures also shows that scalar production, as a fraction of total SUSY processes, decreases at high $`m_0`$, being typically a percent around $`m_03`$ TeV. However, it is clear that at values of $`M_{1/2}`$ lower than the SUSY discovery contour but at $`m_0>2`$ TeV, it is possible to produce scalar SUSY particles at the LHC. A discussion of how to attempt to obtain a limit on the region of mSUGRA parameter space where the squarks in particular (rather than SUSY in general) can be discovered is beyond the scope of this paper.
We terminate the calculation of the fine-tuning and search-reach at $`m_0=`$ 4000 GeV, since for higher values of $`m_0`$ the $`M_{1/2}`$ discovery reach gives an adequate representation of the overall SUSY discovery power, through the gauginos, in each channel. This expression of the reach is also shown in table 3.
## 6 Summary
In this paper, we have obtained the discovery reach of the LHC into mSUGRA parameter space, using the new supersymmetry routines in HERWIG. Where our investigation repeats the calculation of , this provides a useful check on the consistency of the two Monte-Carlos. In addition, our use of the latest software for calculating the mSUGRA spectrum updates the old results, and allows us to move into the region of high scalar masses $`m_0`$. It has been suggested that a focus point gives this region increased naturalness, and the extent to which radiative electroweak symmetry breaking excludes high $`m_0`$ remains uncertain, so this region should be explored thoroughly. We demonstrate that even for arbitrarily high $`m_0`$, the standard SUSY searches at the LHC can discover supersymmetry, through events involving gauginos, provided they are not too heavy ($`M_{1/2}<460`$).
We have introduced the possibility of using fine-tuning as a quantitative way to compare the discovery reach of various channels. Fine-tuning can provide physicists with a quantitative measure of discomfort with a theory, which increases as the experimental bounds are improved. Since it is this disquiet which leads, in the end, to the abandonment of a theory such as supersymmetry, the fine-tuning reach represents the potential of these discovery channels for removing mSUGRA from the list of candidate theories for physics beyond the standard model. The work could be repeated using some different high-energy unification assumptions, instead of mSUGRA. The mSUGRA reach into $`A_0`$ and $`\mathrm{tan}\beta `$ could also be investigated.
The best fine tuning reach is found in a mono-leptonic channel, where for $`\mu >0,A_0=0`$,$`\mathrm{tan}\beta =10`$ (within the focus point region), and $`m_t=174`$ GeV, all points in the $`m_0`$,$`M_{1/2}`$ plane with $`m_0<4000`$ GeV and a fine tuning measure up to 210 (500) are covered by the search, where the definition of fine-tuning excludes (includes) the contribution from the top Yukawa coupling. For $`m_0>4000`$, all values of $`M_{1/2}<460`$ GeV are covered by the search.
###### Acknowledgments.
Part of this work was produced using the Cambridge University High Performance Computing Facility. BCA would like to thank K. Matchev for valuable discussions on checks of the numerical results. JPJH would like to thank C.G. Lester for useful discussions and CERN Theory Division for hospitality. The authors also thank G.G. Ross for helping to motivate this study. This work was funded by the U.K. Particle Physics and Astronomy Research Council.
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# Retarded versus time-nonlocal quantum kinetic equations
## I Introduction
The generalization of the Boltzmann equation (BE) towards dense interacting quantum systems is a still demanding and unsolved task. In dense systems, the mean time between successive collisions becomes comparable to the collision duration. In this case all particles cannot be described as a system of weakly interacting quasiparticles but a set of particles bound in two - particle correlations has to be accounted for in a distinct manner. In the spirit of chemical reactions, one can introduce the quasiparticle density, $`n_f`$ (a density of single-atom molecules), and the density of particles in the correlated states, $`n_c`$ (twice the density of bi-atomic molecules). The total density $`n`$ is their sum
$$n=n_f+n_c.$$
(1)
In parallel with the partial pressure of molecules, $`𝒫=k_BT(n_f+\frac{1}{2}n_c)`$, and the Guldberg-Waage law of mass action law, $`n_c=Kn_f^2`$, even a small fraction of the correlated density, $`n_cn_f`$, contributes to the second virial coefficient, $`𝒫k_BT\left(n\frac{1}{2}Kn_f^2\right)`$. The correlated density thus plays an important role in the thermodynamic behavior of the system.
While the equation of state of the ideal gas, $`𝒫=k_BTn`$, does not reflect any microscopic properties of particles in the system, the second virial coefficient directly follows from the particle-particle interaction. In the early time of statistical mechanics, the experimental second virial coefficient has been used to deduce interaction potentials. Surprisingly, the firmly established concept of the equilibrium virial expansion has never found a corresponding position in the theory of non-equilibrium systems although a number of attempts have been made to modify the BE so that its equilibrium limit will cover at least the second virial coefficient . The achieved corrections to the BE have a form of gradient or nonlocal contributions to the scattering integral. For a hard-sphere gas the non-local correction is trivial, therefore the main theoretical focus was on the statistical correlations . It turned out that the treatment of higher order contributions is far from trivial as the dynamical statistical correlations result in divergences that are cured only after a re-summation of an infinite set of contributions. Naturally, this re-summation leads to non-analytic density corrections to the scattering integral . The moral of these hard-sphere studies is that beyond the non-local corrections one must sum up an infinite set of contributions; the plain expansion leads to incorrect results. From this point of view, the approximate statistical virial corrections implemented by , although possibly reasonable, are not sufficiently justified by the theory.
Real particles do not interact like hard spheres, in particular when their de Broglie wave lengths are comparable with the potential range. An effort to describe the virial corrections for more realistic systems resulted in various generalizations of Enskog’s equation . By closer inspection one finds that all tractable quantum theories deal exclusively with the non-local corrections. The statistical correlations in quantum systems would require an adequate solution of three-particle collisions (say from Fadeev equations) which are now intensively being studied . A systematic incorporation of three-particle collisions into the kinetic equation, however, is not yet fully understood, therefore we discuss only binary processes.
An alternative way to describe the correlation has been developed for Fermi systems at very low temperatures. Since the Pauli principle excludes all but the zero-angle scattering channels, the correlation can be re-cast into a renormalization of the single-particle energies known as the quasiparticle energies. While quasiparticles behave in many aspects like free particles, the density dependence of the quasiparticle energies results in non-trivial thermodynamic properties. In accordance with the focus either on zero-angle or finite-angle scattering, in the literature one can find two distinct classes of quantum kinetic equations.
The zero-angle class are equations determining the time evolution of the momentum and space dependent quasiparticle distribution function $`f(k,r,t)`$. This leads to kinetic equations of the Landau-Silin type
$$\frac{f}{t}+\frac{\epsilon }{k}\frac{f}{r}\frac{\epsilon }{r}\frac{f}{k}=z((1f)\sigma _\epsilon ^<f\sigma _\epsilon ^>).$$
(2)
Here the drift term is characterized by the quasiparticle energies $`\epsilon (k,r,t)`$ and the collision integral consists of scattering-in and -out terms which are summarized in the selfenergy as functionals of the quasiparticle distribution function, $`\sigma ^{}[f]`$. The collision rates are reduced by the wave function renormalization, which is given by the frequency derivative of the real part of the selfenergy $`\sigma `$ at the pole value $`z^1=1\frac{\sigma }{\omega }|_{\omega =\epsilon }`$. This kinetic equation is local in time and space, i.e., it describes the collisions as instantaneous point–like events. Similarly to the BE, the local collisions yield no virial corrections. The virial corrections are thus covered exclusively by the quasiparticle energy while no correlated density appears.
The finite-angle class are equations which include the finite duration of collisions, i.e., they have non-Markovian scattering integrals. These kinetic equations are usually developed for the reduced density matrix or the Wigner function $`\rho (k,r,t)`$. The Wigner distribution obeys a Levinson-type of kinetic equation
$`{\displaystyle \frac{\rho _t}{t}}`$ $`+`$ $`{\displaystyle \frac{ϵ^{\mathrm{hf}}}{k}}{\displaystyle \frac{\rho _t}{r}}{\displaystyle \frac{ϵ^{\mathrm{hf}}}{r}}{\displaystyle \frac{\rho _t}{k}}`$ (3)
$`=`$ $`2\mathrm{Im}{\displaystyle \underset{\mathrm{}}{\overset{t}{}}}𝑑\overline{t}G_{t,\overline{t}}^R\left[(1\rho _{\overline{t}})\mathrm{\Sigma }_{\overline{t},t}^<\rho _{\overline{t}}\mathrm{\Sigma }_{\overline{t},t}^>\right].`$ (4)
Here the drift term is essentially determined by the Hartree-Fock mean field, $`ϵ^{\mathrm{hf}}=k^2/2m+\sigma _{\mathrm{hf}}`$, and the double-time selfenergies $`\mathrm{\Sigma }_{\overline{t},t}^{}`$ are given as functionals of the Wigner distribution. Explicit time arguments of the collision integral, $`\overline{t}<t`$, show the retardation which requires to include the propagation, $`G_{t,\overline{t}}^R`$, from the retarded time $`\overline{t}`$ to the time $`t`$ of balancing changes. Having the non-Markovian form, this kinetic equation is capable to describe finite duration effects and corresponding virial corrections.
The Landau-Silin and Levinson equations have common features with the original BE. In particular, they express the balance between the drift given by the gradient terms on the left hand side and the dissipation given by the scattering integral on the right hand side. On the other hand, there is a remarkable difference in the actual physical contents of these common ingredients, namely $`\rho f`$, $`\epsilon ϵ^{\mathrm{hf}}`$, and the right hand sides of (2) and (4) are rather different.
In the present paper we argue that the differences between the phenomenological quasiparticle picture, Eq. (2), and the BBGKY-type<sup>*</sup><sup>*</sup>*Equations derived from the Born-Bogoliubov-Green-Kirkwood-Yvon (BBGKY) hierarchy of reduced densities. equation (4) can be characterized in terms of the retardation of the collision integral of (4). We will show that there are various contributions to the total retardation, each responsible for different features seen in the quasiparticle picture extended by non-local corrections.
If both approaches are equivalent, the Levinson equation (4) must contain terms on the collisional side which should be possible to rearrange into gradients such that the quasiparticle energies of the drift side in the Landau-Silin equation (2) are obtained. The mutual relation between the Wigner distribution $`\rho `$ and the quasiparticle distribution $`f`$ is of essential importance for such a rearrangement.
In Sec. II we show that a part of the retardation of the Levinson equation (4) describes the off-shell motion of particles. This off-shell motion can be eliminated from the kinetic equation which requires to introduce an effective distribution (the quasiparticle distribution $`f`$) from which the Wigner distribution $`\rho `$ can be constructed (72)
$$\rho =f+\frac{d\omega }{2\pi }\frac{\mathrm{}}{\omega \epsilon }\frac{}{\omega }\left((1f)\sigma _\omega ^<f\sigma _\omega ^>\right).$$
(5)
This relation is the extended quasiparticle picture derived for small scattering rates. The limit of small scattering rates has been first introduced by Craig . An inverse functional $`f[\rho ]`$ has been constructed in . For equilibrium non-ideal plasmas this approximation has been employed by and under the name of the generalized Beth-Uhlenbeck approach has been used by in nuclear matter for studies of the correlated density. The authors in have used this approximation with the name extended quasiparticle approximation for the study of the mean removal energy and high-momenta tails of Wigner’s distribution. The non-equilibrium form has been derived finally as the modified Kadanoff and Baym ansatz . We will call it extended quasiparticle approximation.
We show that the retardation is responsible for gradient terms by which the mean-field drift of the Levinson equation (4) differs from the quasiparticle drift of the Landau-Silin equation (2). In Sec. III we discuss the remaining parts of the retardation, a piece which is compensated by the decay of internal propagators, and a piece which describes the collision delay. Special emphasize is put on internal double counts of correlations in the Levinson equation. Related to this question we present in appendix C a discussion of a common met pitfall in literature. In Sec. IV we complete the kinetic equation with the finite collision delay and provide related non-local corrections in space. Our main result is the nonlocal kinetic equation derived here in an alternative way to ,
$`{\displaystyle \frac{f_a}{t}}+{\displaystyle \frac{\epsilon _a}{k}}{\displaystyle \frac{f_a}{r}}{\displaystyle \frac{\epsilon _a}{r}}{\displaystyle \frac{f_a}{k}}`$ (6)
$`={\displaystyle \underset{b}{}}{\displaystyle 𝒫\left\{f_a^{}f_b^{}(1f_a)(1f_b)f_af_b(1f_a^{})(1f_b^{})\right\}}`$ (7)
(8)
with $`f_a^{}=f_a(kq\mathrm{\Delta }_K,r\mathrm{\Delta }_3,t\mathrm{\Delta }_t)`$, $`f_b^{}=f_b(p+q\mathrm{\Delta }_K,r\mathrm{\Delta }_4,t\mathrm{\Delta }_t)`$, $`f_a=f_a(k,r,t)`$ and $`f_b=f_b(p,r\mathrm{\Delta }_2,t)`$. The differential cross section $`𝒫|t^R|^2`$ is proportional to the square of the amplitude of the T-matrix. All non-local corrections are given by derivatives of the scattering phase shift $`\varphi =\mathrm{Im}\mathrm{ln}t^R(\omega ,k,p,q,t,r)`$ ,
$`\mathrm{\Delta }_t`$ $`=`$ $`{\displaystyle \frac{\varphi }{\omega }},\mathrm{\Delta }_E={\displaystyle \frac{1}{2}}{\displaystyle \frac{\varphi }{t}},\mathrm{\Delta }_K={\displaystyle \frac{1}{2}}{\displaystyle \frac{\varphi }{r}},\mathrm{\Delta }_3={\displaystyle \frac{\varphi }{k}},`$ (9)
$`\mathrm{\Delta }_2`$ $`=`$ $`{\displaystyle \frac{\varphi }{p}}{\displaystyle \frac{\varphi }{q}}{\displaystyle \frac{\varphi }{k}},\mathrm{\Delta }_4={\displaystyle \frac{\varphi }{k}}{\displaystyle \frac{\varphi }{q}}.`$ (10)
We discuss in Sec. IV the space-time symmetry of this collision integral and link the correlated density with the collision delay. In Sec. V we summarize and refer to numerical solutions of the derived nonlocal kinetic equation. The most important formulae for the T-matrix are presented in appendix A and the technique how to derive nonlocal corrections from gradient expansion is summarized in appendix B.
## II Kinetic equations
First we want to show that the Levinson equation can be transformed into the kinetic equation of Landau-Silin type and vice versa. Therefore we first derive the kinetic equation for the Wigner distribution from the real-time Green function technique. For a homogeneous system similar investigations have been presented by Bornath at al . Here we extend their approach to inhomogeneous systems and want to focus specially on the physical contents of both equations. This approach will furnish us with a set of Kramers-Kronig relations needed to link the Wigner and the quasiparticle distributions.
We consider the two independent correlation functions for Fermionic operators $`G^>(1,2)=a(1)a^+(2)`$ and $`G^<(1,2)=a^+(2)a(1)`$, where cumulative variables mean time and space, $`1t,x`$. The time diagonal part of $`G^<`$ yields the Wigner distribution, $`\rho (x_1,x_2,t)=G^<(1,2)_{t_{1,2}=t}`$. For the correlation function $`G^<`$ the Kadanoff and Baym equation of motion reads
$`i\left(G_0^1G^<G^<G_0^1\right)=`$ $``$ $`i\left(\mathrm{\Sigma }^RG^<G^<\mathrm{\Sigma }^A\right)`$ (11)
$`+`$ $`i\left(G^R\mathrm{\Sigma }^<\mathrm{\Sigma }^<G^A\right).`$ (12)
The retarded and advanced functions are introduced as $`B^R(1,2)=i\theta (t_1t_2)(B^>+B^<)`$ and $`B^A(1,2)=i\theta (t_2t_1)(B^>+B^<)`$. The products are understood as integrations over intermediate variables (time and space). The Hartree-Fock drift term has the form
$`G_0^1(1,2)`$ $`=`$ $`\left(i{\displaystyle \frac{}{t_1}}+{\displaystyle \frac{_1^2}{2m}}\right)\delta (12)`$ (13)
$``$ $`\mathrm{\Sigma }^{\mathrm{hf}}(x_1,x_2,t_1)\delta (t_1t_2).`$ (14)
The Kadanoff and Baym equation (12) is more general than any of the above kinetic equations which can be derived from it with the help of specific approximations. One possibility is to rearrange the terms such that they yield the Landau-Silin kinetic equation for the quasiparticle distribution function (72) when the gradient expansion is applied . The other possibility is to collect the terms on the right hand side of (12) into a non-Markovian collision integral which results in the Levinson equation (4) for the Wigner distribution . The connection between both equations, however, has remained dubious because in the Levinson equation the drift term is controlled exclusively by the Hartree-Fock selfenergy while in the Landau-Silin equation the drift term is represented by the full quasiparticle energy. In the following we demonstrate how the Levinson equation can be converted into the Landau-Silin equation. To this end we first derive the Levinson equation from (12).
### A Precursor of the Levinson equation
On the time diagonal, $`t_{1,2}=t`$, the Kadanoff and Baym equation (12) yields an identity
$`i(G_0^1G^<`$ $`G^<G_0^1)(t,x_1,t,x_2)`$ (15)
$`={\displaystyle \underset{\mathrm{}}{\overset{t}{}}}𝑑t^{}{\displaystyle 𝑑x^{}}`$ $`(G^>(t,x_1,t^{},x^{})\mathrm{\Sigma }^<(t^{},x^{},t,x_2)`$ (17)
$`+\mathrm{\Sigma }^<(t,x_1,t^{},x^{})G^>(t^{},x^{},t,x_2))`$
$`{\displaystyle \underset{\mathrm{}}{\overset{t}{}}}𝑑t^{}{\displaystyle 𝑑x^{}}`$ $`(G^<(t,x_1,t^{},x^{})\mathrm{\Sigma }^>(t^{},x^{},t,x_2)`$ (19)
$`+\mathrm{\Sigma }^>(t,x_1,t^{},x^{})G^<(t^{},x^{},t,x_2)).`$
On the left hand side there is the Hartree-Fock drift and the right hand side contains a non-Markovian collision integral. Note that correlations beyond the Hartree-Fock field are exclusively in the collision integral.
#### 1 The Wigner representation
Now we concentrate on the gradient expansion of the collision integral. The quasi-classical approximation is achieved using the Wigner mixed representation for space arguments
$`\widehat{\sigma }(k,{\displaystyle \frac{x_1+x_2}{2}},{\displaystyle \frac{t_1+t_2}{2}},t_1t_2)=`$ (20)
$`{\displaystyle d(x_1x_2)\mathrm{e}^{ik(x_1x_2)}\mathrm{\Sigma }(t_1,x_1,t_2,x_2)}.`$ (21)
As a rule, we use the lower case to denote operators in the full Wigner representation (momentum, space, time, frequency) and the hat to denote the Wigner representation in space and double-time representation. Note that time arguments of $`\widehat{\sigma }`$ now denote the center-of-mass time, $`\frac{t_1+t_2}{2}`$, and the difference time, $`t_1t_2`$. These two Wigner representations are identical for $`\rho `$ and $`\sigma ^{\mathrm{hf}}`$ which do not depend on the difference time.
In the representation (21) the identity (19) reads
$`\left({\displaystyle \frac{}{t}}+{\displaystyle \frac{k}{m}}{\displaystyle \frac{}{r}}\right)\rho (k,r,t)`$ (24)
$`+2\mathrm{sin}\left({\displaystyle \frac{1}{2}}(_r^\mathit{1}_k^\mathit{2}_k^\mathit{1}_r^\mathit{2})\right)\sigma ^{\mathrm{hf}}(k,r,t)\rho (k,r,t)`$
$`=\mathrm{exp}\left({\displaystyle \frac{i}{2}}\left(_r^\mathit{1}_k^\mathit{2}_k^\mathit{1}_r^\mathit{2}\right)\right){\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑\tau `$
$`\times (`$ $`\widehat{g}^>(k,r,t{\displaystyle \frac{\tau }{2}},\tau )\widehat{\sigma }^<(k,r,t{\displaystyle \frac{\tau }{2}},\tau )`$ (28)
$`+\widehat{\sigma }^<(k,r,t{\displaystyle \frac{\tau }{2}},\tau )\widehat{g}^>(k,r,t{\displaystyle \frac{\tau }{2}},\tau )`$
$`\widehat{g}^<(k,r,t{\displaystyle \frac{\tau }{2}},\tau )\widehat{\sigma }^>(k,r,t{\displaystyle \frac{\tau }{2}},\tau )`$
$`\widehat{\sigma }^>(k,r,t{\displaystyle \frac{\tau }{2}},\tau )\widehat{g}^<(k,r,t{\displaystyle \frac{\tau }{2}},\tau )).`$
Here, the superscripts $`\mathit{1},\mathit{2}`$ denote that the partial derivatives apply only to the first and second function in the product, e.g., $`(_r^\mathit{1}_k^\mathit{2})^3ab=\frac{^3a}{^3r}\frac{^3b}{^3k}`$.
#### 2 Gradient approximation
The power expansion of the goniometric functions in (28) defines the expansion in space gradients. This expansion goes to infinite order, but in the following we will restrict our attention to the linear expansion. It is customary to abbreviate the linear gradient terms with the help of the Poisson brackets,
$$2\mathrm{sin}\left(\frac{1}{2}(_r^\mathit{1}_k^\mathit{2}_k^\mathit{1}_r^\mathit{2})\right)ab\{a,b\}=\frac{a}{k}\frac{b}{r}\frac{a}{r}\frac{b}{k}.$$
(29)
The linear approximation in space gradients is straightforward, however, one has to be careful about time arguments of functions in the collision integral of (28). For example, the first and second product seem to form an anticommutator in which the linear gradients in space cancel. This is not true because of time arguments.
The gradient approximation in time requires a special treatment because of the lower integration limit. Due to this limit the integral does not define a standard matrix product with respect to the time variables. We will give this expansion an explicit notation.
In equilibrium, the functions $`\widehat{g}^{}`$ and $`\widehat{\sigma }^{}`$ do not depend on the center-of-mass time. For slowly evolving systems the center-of-mass dependence is smooth and weak. Accordingly, in the collision integral of (28), we can expand the center-of-mass time dependence around the time $`t`$ in powers of $`\tau `$. Since all functions in (28) have the same center-of-mass time, $`t\frac{\tau }{2}`$, it is possible to write the linear expansion of (28) with respect to time gradients in a compact form
$`{\displaystyle \frac{}{t}}\rho +\{ϵ^{\mathrm{hf}},\rho \}=I+{\displaystyle \frac{}{t}}R,`$ (30)
with $`ϵ^{\mathrm{hf}}=\frac{k^2}{2m}+\sigma ^{\mathrm{hf}}`$. The zero order gradient term, $`I=I^>I^<`$, reads
$`I^<`$ $`=`$ $`\left(1+{\displaystyle \frac{i}{2}}\left(_r^\mathit{1}_k^\mathit{2}_k^\mathit{1}_r^\mathit{2}\right)\right){\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑\tau `$ (31)
$`\times `$ $`\left(\widehat{\sigma }^>(t,\tau )\widehat{g}^<(t,\tau )+\widehat{g}^<(t,\tau )\widehat{\sigma }^>(t,\tau )\right).`$ (32)
$`I^>`$ is obtained from $`I^<`$ via the interchange of particles and holes, $`><`$. In the first order gradient term, $`R=R^>R^<`$, is given by
$`R^<`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(1+{\displaystyle \frac{i}{2}}\left(_r^\mathit{1}_k^\mathit{2}_k^\mathit{1}_r^\mathit{2}\right)\right){\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑\tau \tau `$ (33)
$`\times `$ $`\left(\widehat{\sigma }^>(t,\tau )\widehat{g}^<(t,\tau )+\widehat{g}^<(t,\tau )\widehat{\sigma }^>(t,\tau )\right).`$ (34)
Now we are ready to turn all functions into the Wigner representation for time and energy,
$$\sigma _\omega (k,r,t)=𝑑\tau \mathrm{e}^{i\omega \tau }\widehat{\sigma }(k,r,t,\tau ).$$
(35)
We will write the energy argument as the subscript and keep other arguments implicit in most of the formulas.
After a substitution of $`\sigma `$’s and $`g`$’s of the Wigner representation (35) into (32), one can integrate out the difference time $`\tau `$. The integration over time results in the $`\delta `$ function which describes processes on the energy shell, and the principle value terms which represent off-shell processes,
$$\underset{0}{\overset{\mathrm{}}{}}𝑑\tau \mathrm{e}^{i(\omega \omega ^{})\tau }=\pi \delta (\omega \omega ^{})+i\frac{\mathrm{}}{\omega \omega ^{}}.$$
(36)
The $`\delta `$-function can be readily integrated out leaving both functions with identical energy arguments. The principle value part,$`\mathrm{}`$, is an odd function of energies $`\omega `$ and $`\omega ^{}`$, therefore its non-gradient contributions cancel but the linear gradients survive. Formula (32) thus turns into
$$I^<=\frac{d\omega }{2\pi }\sigma _\omega ^>g_\omega ^<+\frac{d\omega d\omega ^{}}{(2\pi )^2}\frac{\mathrm{}}{\omega ^{}\omega }\{\sigma _\omega ^{}^>,g_\omega ^<\}.$$
(37)
In the correlated part we express the time integral as
$`{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑\tau \tau \mathrm{e}^{i(\omega \omega ^{})\tau }`$ $`=`$ $`i{\displaystyle \frac{}{\omega }}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑\tau \mathrm{e}^{i(\omega \omega ^{})\tau }`$ (38)
$`=`$ $`i\pi \delta ^{}(\omega \omega ^{})+{\displaystyle \frac{\mathrm{}^{}}{\omega \omega ^{}}}.`$ (39)
The meaning of $`\delta ^{}`$ and $`\mathrm{}^{}`$ is seen from comparison with (36). The $`\delta ^{}`$ and $`\mathrm{}^{}`$ are odd and even functions in $`\omega \omega ^{}`$, respectively. The gradient expansion of (34) thus reads
$$R^<=\frac{d\omega d\omega ^{}}{(2\pi )^2}\frac{\mathrm{}^{}}{\omega ^{}\omega }\sigma _\omega ^{}^>g_\omega ^<+\frac{d\omega }{2\pi }\{\sigma _\omega ^>,_\omega g_\omega ^<\}.$$
(40)
The term $`R`$ contributes to the kinetic equation (30) as a gradient correction linear in time. Since the second term of $`R^<`$ in (40) is proportional to space gradients, we neglect this term leading to the second order contribution in gradients.
### B Connection between the Wigner and quasiparticle distributions
At very low temperatures, dissipative processes are known to vanish due to the Pauli exclusion principle. It is desirable to reorganize the kinetic equation so that the scattering integral will vanish in this limit too. As can be seen from (40), the term $`_tR`$ remains finite even for very low temperatures. It is possible to formally remove $`_tR`$ if we shift this term on the drift side of (30) and introduce a new distribution function,
$$f=\rho R.$$
(41)
We have denoted the new function as $`f`$ to anticipate that it is indeed the quasiparticle distribution as will be shown now.
Since we want to arrive at the kinetic equation for $`f`$, it is advantageous to write relation (41) in the opposite way so that we obtain $`\rho `$ as a functional of $`f`$, $`\rho =f+R`$. Using (40) and the particle-hole conjugated term, $`R^>`$, one finds
$$\rho =f\frac{d\omega ^{}d\omega }{(2\pi )^2}\frac{\mathrm{}^{}}{\omega ^{}\omega }\left(\sigma _\omega ^{}^>g_\omega ^<\sigma _\omega ^{}^<g_\omega ^>\right).$$
(42)
This relation provides the so called extended quasiparticle approximation of the Wigner distribution.
#### 1 Wave-function renormalization
To separate the on-shell and off-shell contributions of the correlation term, we use the imaginary part of the selfenergy, $`\gamma =\sigma ^>+\sigma ^<`$ and the spectral function, $`a=g^>+g^<`$ to rewrite (42) as
$$\rho =f\frac{d\omega ^{}d\omega }{(2\pi )^2}\frac{\mathrm{}^{}}{\omega ^{}\omega }\left(\gamma _\omega ^{}g_\omega ^<\sigma _\omega ^{}^<a_\omega \right).$$
(43)
With the help of the Kramers-Kronig relation
$$\frac{d\omega ^{}}{2\pi }\frac{\mathrm{}^{}}{\omega \omega ^{}}\gamma _\omega ^{}=\frac{\sigma _\omega }{\omega },$$
(44)
where $`\sigma `$ is the real part of the selfenergy, we obtain
$$\rho =f+\frac{d\omega }{2\pi }\frac{\sigma _\omega }{\omega }g_\omega ^<+\frac{d\omega ^{}d\omega }{(2\pi )^2}\frac{\mathrm{}^{}}{\omega ^{}\omega }\sigma _\omega ^{}^<a_\omega .$$
(45)
The simplest quasiparticle approximation, $`a_\omega =2\pi \delta (\omega \epsilon )`$ and $`g_\omega ^<=f\mathrm{\hspace{0.17em}2}\pi \delta (\omega \epsilon )`$, in the correlation terms yields
$$\rho =\left(1+\frac{\sigma }{\omega }|_{\omega =\epsilon }\right)f+\frac{d\omega ^{}}{2\pi }\frac{\mathrm{}^{}}{\omega ^{}\epsilon }\sigma _\omega ^{}^<.$$
(46)
In (46) the on-shell (quasiparticle) contribution is reduced by the wave-function renormalization,
$$z=1+\frac{\sigma }{\omega }|_{\omega =\epsilon },$$
(47)
and the term with $`\sigma ^<`$ provides the off-shell contributions.
The exact renormalization is $`z^1=1_\omega \sigma `$, therefore (47) is only the linear approximation in $`_\omega \sigma `$. From the Kramers-Kronig relation (44) one can see that this corresponds to the linear approximation in the scattering rate $`\gamma `$. We will keep all terms only up to this order in $`\gamma `$.
#### 2 Extended quasiparticle approximation
From the relation between the Wigner and the quasiparticle distributions (46) one can deduce an approximative construction of the correlation function $`g^<`$ as a functional of the quasiparticle distribution $`f`$. Since $`\rho =\frac{d\omega }{2\pi }g^<`$, we can write (46) as
$$\frac{d\omega }{2\pi }g_\omega ^<=\frac{d\omega }{2\pi }\left(fz\mathrm{\hspace{0.17em}2}\pi \delta (\omega \epsilon )+\frac{\mathrm{}^{}}{\omega \epsilon }\sigma _\omega ^<\right).$$
(48)
We have multiplied $`f`$ with the $`\delta `$ function and included this term into the integrand.
One can see that the correlation function in the extended quasiparticle approximation ,
$$g_\omega ^<=fz\mathrm{\hspace{0.17em}2}\pi \delta (\omega \epsilon )+\frac{\mathrm{}^{}}{\omega \epsilon }\sigma _\omega ^<,$$
(49)
satisfies (48). Perhaps it is not necessary to recall that the function $`f`$ has been introduced here merely to suppress the non-dissipative part of the collision integral, $`_tR`$. In , the function $`f`$ has been defined as the factor of the singularity of $`g^<`$ which justifies its interpretation in terms of the quasiparticle distribution. As relation (48) shows, both approaches lead to identical results.
In parallel with (49), for the hole correlation function one finds
$$g_\omega ^>=(1f)z\mathrm{\hspace{0.17em}2}\pi \delta (\omega \epsilon )+\frac{\mathrm{}^{}}{\omega \epsilon }\sigma _\omega ^>.$$
(50)
From the spectral identity, $`a=g^>+g^<`$, it follows that (49) and (50) are consistent with the extended quasiparticle approximation of the spectral function
$$a_\omega =z\mathrm{\hspace{0.17em}2}\pi \delta (\omega \epsilon )+\frac{\mathrm{}^{}}{\omega \epsilon }\gamma _\omega .$$
(51)
It is worthy to remark that (51) fulfills the spectral sum rule
$`{\displaystyle \frac{d\omega }{2\pi }a(k,\omega ,r,t)}`$ $`=`$ $`1`$ (52)
as well as the energy weighted sum rule
$`{\displaystyle \frac{d\omega }{2\pi }\omega a(k,\omega ,r,t)}`$ $`=`$ $`{\displaystyle \frac{k^2}{2m}}+\sigma ^{\mathrm{hf}}(k,r,t).`$ (53)
The latter one is violated in the simple quasiparticle picture.
### C Landau-Silin equation
So far we have treated only the time gradients finding that the identity (30) can be expressed as
$$\frac{f}{t}+\{ϵ^{\mathrm{hf}},\rho \}=I.$$
(54)
Now we rearrange this identity into the Landau-Silin equation for $`f`$.
The time-local remainder $`I`$ of the collision integral still includes the space gradients, $`I=I_BI_{}`$. The non-gradient part \[the first term of (37)\] is the Boltzmann-type scattering integral of the Landau-Silin equation
$`I_B`$ $`=`$ $`{\displaystyle \frac{d\omega }{2\pi }(g_\omega ^>\sigma _\omega ^<g_\omega ^<\sigma _\omega ^>)}=z\left((1f)\sigma _\epsilon ^<f\sigma _\epsilon ^>\right)`$ (55)
In the second step we have used (49). The off-shell contributions to (LABEL:o) have not been neglected but cancel exactly.
The gradient part which is the second term of (37) reads
$`I_{}`$ $`=`$ $`{\displaystyle \frac{d\omega d\omega ^{}}{(2\pi )^2}\frac{\mathrm{}}{\omega ^{}\omega }\left(\{\sigma _\omega ^{}^>,g_\omega ^<\}\{\sigma _\omega ^{}^<,g_\omega ^>\}\right)}`$ (57)
$`=`$ $`{\displaystyle \frac{d\omega }{(2\pi )}\left(\{\sigma _\omega ,g_\omega ^<\}\{\sigma _\omega ^<,g_\omega \}\right)}.`$ (58)
Here (44) has been used to remove the $`\omega ^{}`$-integration for $`\sigma ^>+\sigma ^<`$ and a similar Kramers-Kronig relation for the $`\omega `$-integration over the spectral function $`g^>+g^<`$ has been applied. The function $`g=\mathrm{Re}g^{R,A}=\frac{1}{2}(g^R+g^A)`$ is the off-shell part of the propagators.
The mean-field drift $`\{ϵ^{\mathrm{hf}},\rho \}`$ in (54) is not compatible with the quasiparticle distribution in the time derivative and the scattering integral. Moreover, the space gradients from the collision integral, $`I_{}`$, have to be accounted for. Our aim is now to show that all gradients can be collected into the quasiparticle drift,
$$\{ϵ^{\mathrm{hf}},\rho \}+I_{}=\{\epsilon ,f\}.$$
(59)
The quasiparticle energy differs from the mean-field energy by the pole value of the selfenergy, $`\epsilon =ϵ^{\mathrm{hf}}+\sigma _\epsilon `$.
We will proceed in two steps. First we observe that the right hand side of (59) includes only the on-shell contribution, therefore we show that the on-shell part of the left hand side equals the right hand side. Second, we show that the off-shell part of the left hand side of (59) vanishes. The on-shell parts of (49) used in (58) result into
$`I_{}^{\mathrm{on}}`$ $`=`$ $`{\displaystyle 𝑑\omega \{\sigma _\omega ,zf\delta (\omega \epsilon )\}}`$ (60)
$`=`$ $`\{\sigma _\epsilon ,zf\}\sigma _\epsilon ^{}\{\epsilon ,zf\}+zf\{\sigma _\epsilon ^{},\epsilon \}`$ (61)
$`=`$ $`\{\sigma _\epsilon ,zf\}+(1z)\{\epsilon ,zf\}+zf\{z1,\epsilon \}.`$ (62)
Now we add the on-shell part of the commutator $`\{ϵ^{\mathrm{hf}},\rho \}^{\mathrm{on}}=\{ϵ^{\mathrm{hf}},zf\}`$. Using a rearrangement
$`\{ϵ^{\mathrm{hf}},zf\}=\{\epsilon \sigma ,zf\}`$ (63)
$`=\{\epsilon ,f\}+(z1)\{\epsilon ,f\}+f\{\epsilon ,z\}\{\sigma ,zf\}`$ (64)
we obtain
$`\{ϵ^{\mathrm{hf}},\rho \}^{\mathrm{on}}+I_{}^{\mathrm{on}}=\{\epsilon ,f\}\{\epsilon ,(1z)^2f\}.`$ (65)
The last term is of higher order than linear in the damping and can be neglected which confirms the relation (59) already from the on-shell parts.
It remains to prove that all off-shell contributions to the left hand side of (59), $`𝒪=\{ϵ^{\mathrm{hf}},\rho \}^{\mathrm{off}}+I_{}^{\mathrm{off}}`$ cancel, i.e., $`𝒪=0`$. From the off-shell term of (49) we find
$$𝒪=\frac{d\omega }{2\pi }\left(\{ϵ^{\mathrm{hf}}+\sigma _\omega ,\frac{\mathrm{}^{}}{\omega \epsilon }\sigma _\omega ^<\}\{g_\omega ,\sigma _\omega ^<\}\right).$$
(66)
The term with $`ϵ^{\mathrm{hf}}`$ results from $`\{ϵ^{\mathrm{hf}},\rho \}`$, while the others from $`I_{}`$. In the last term we use the extended quasiparticle approximation of the real part of the propagator
$`g_\omega =z{\displaystyle \frac{\mathrm{}}{\omega \epsilon }}+{\displaystyle \frac{\mathrm{}^{}}{\omega \epsilon }}(\sigma _\omega \sigma _\epsilon )`$ (67)
which directly results from (51) and the Kramers-Kronig relation. Using $`\{\frac{\mathrm{}}{\omega \epsilon },\sigma _\omega ^<\}=\{\epsilon ,\frac{\mathrm{}^{}}{\omega \epsilon }\sigma _\omega ^<\}`$ we can rewrite $`𝒪`$ as
$`𝒪=`$ (68)
$`{\displaystyle \frac{d\omega }{2\pi }\left(\{\sigma _\omega \sigma _\epsilon ,\frac{\mathrm{}^{}}{\omega \epsilon }\sigma _\omega ^<\}\{\frac{\mathrm{}^{}}{\omega \epsilon }(\sigma _\omega \sigma _\epsilon ),\sigma _\omega ^<\}\right)}.`$ (69)
(70)
The linear expansion in the vicinity of the pole, $`\sigma _\omega \sigma _\epsilon =(\omega \epsilon )(z1)+o(\gamma ^2)`$, yields
$`𝒪={\displaystyle \frac{d\omega }{2\pi }\{\frac{\mathrm{}^{\prime \prime }}{\omega \epsilon },(z1)\sigma _\omega ^<\}}.`$ (71)
The product $`(z1)\sigma ^<`$ is of second order in $`\gamma `$, i.e., the off-shell contribution $`𝒪`$ is negligible within the assumed accuracy. This completes the proof of relation (59).
In summary, the requirement that the scattering integral vanishes at very low temperatures directly leads to the concept of quasiparticles represented by relation (42) between the Wigner and the quasiparticle distributions
$$\rho =f+\frac{d\omega }{2\pi }\frac{\mathrm{}}{\omega \epsilon }\frac{}{\omega }\left((1f)\sigma _\omega ^<f\sigma _\omega ^>\right).$$
(72)
The space gradients of the collision integral renormalize the mean-field drift into the familiar quasiparticle drift (59). The resulting kinetic equation for the quasiparticle distribution has the structure of the Landau-Silin equation (2)
$$\frac{f}{t}+\frac{\epsilon }{k}\frac{f}{r}\frac{\epsilon }{r}\frac{f}{k}=z((1f)\sigma _\epsilon ^<f\sigma _\epsilon ^>).$$
(73)
We conclude this section that the Levinson type of equation (4) is equivalent to the Landau-Silin equation (2) up to second order in the damping $`\gamma `$ or in the extended quasiparticle picture.
## III Additional retardation of the Levinson equation
The link between the Levinson and Landau-Silin equations shows that the collision integral of the Levinson equation cannot be fully interpreted in terms of the scattering processes but it includes various renormalization features. In the above treatment we have constructed a closed equation for the quasiparticle distribution function while the reduced density is given by an additional functional of this quasiparticle distribution. There arises an important question whether it is not possible to construct the correlation functions, $`g^{>,<}`$, directly from the Wigner distribution so that identity (28) turns into a closed equation for the Wigner distribution. To show the problems and drawbacks with this opposite way we discuss this construction for a homogeneous system.
An approximation giving the Levinson equation is the GKB ansatz which for the homogeneous system reads \[$`\tau >0`$\]
$`\widehat{g}^<(k,t{\displaystyle \frac{\tau }{2}},\tau )`$ $`=`$ $`i\widehat{g}^R(k,t{\displaystyle \frac{\tau }{2}},\tau )\rho (k,t\tau ),`$ (74)
$`\widehat{g}^<(k,t{\displaystyle \frac{\tau }{2}},\tau )`$ $`=`$ $`i\widehat{g}^A(k,t{\displaystyle \frac{\tau }{2}},\tau )\rho (k,t\tau ).`$ (75)
For $`\widehat{g}^>`$ one substitutes $`1\rho `$ for $`\rho `$. Note that this approximation includes the additional retardation between the center-of-mass time $`t\frac{\tau }{2}`$ and the time argument of the Wigner distribution which is $`t\tau `$. This additional retardation causes many problems ranging from instability of numerical treatments over double counts of renormalizations or virial corrections to incorrect interpretations of various correlation phenomena. Here we show how to handle this retardation within the linear approximation.
If one assumes that the propagator $`G^R`$ depends only on the difference time $`\tau `$, i.e., it does not depend on any time dependent external fields or mean fields, it is possible to follow the approach developed above with only minor modifications. To make our discussion specific we limit our attention to homogeneous systems and employ the Galitskii-Feynman T-matrix approximation of the selfenergy,
$`\widehat{\sigma }^<(k,t{\displaystyle \frac{\tau }{2}},\tau )`$ (76)
$`=`$ $`{\displaystyle \frac{dpdq}{(2\pi )^6}\widehat{g}^>(p,t\frac{\tau }{2},\tau )}`$ (77)
$`\times `$ $`{\displaystyle 𝑑\tau _R𝑑\tau _A𝒯^R(k,p,q,\tau _R)𝒯^A(k,p,q,\tau _A)}`$ (78)
$`\times `$ $`\widehat{g}^<(kq,t{\displaystyle \frac{1}{2}}(\tau _R+\tau _A+\tau ),\tau \tau _R+\tau _A)`$ (79)
$`\times `$ $`\widehat{g}^<(p+q,t{\displaystyle \frac{1}{2}}(\tau _R+\tau _A+\tau ),\tau \tau _R+\tau _A).`$ (80)
This selfenergy describes the process in which two particles of initial momenta $`kq`$ and $`p+q`$ are scattered into final states $`k`$ and $`p`$. The function $`𝒯^{R,A}`$ is the retarded/advanced scattering T-matrix, see appendix A. Due to a finite duration of the scattering process, the T-matrices bring yet additional retardation times $`\tau _{R,A}`$. We will show that retardation times $`\tau _{R,A}`$ are responsible for the virial corrections to the density.
### A Decay of propagators and double counts
First we focus on double counts following from the additional retardation resulting from the GKB ansatz. Since these problems appear already for weakly coupled systems, it is sufficient to restrict our attention to the Born approximation, $`𝒯^{R,A}=𝒱\delta (\tau _{R,A})`$. The Levinson equation then reads
$$\frac{}{t}\rho (k,t)=2\mathrm{Re}\frac{dpdq}{(2\pi )^6}𝒱_q^2\underset{0}{\overset{\mathrm{}}{}}𝑑\tau S(t\frac{\tau }{2},\tau )D(t\tau ).$$
(81)
All distributions are collected in
$`D(t)`$ $`=`$ $`(1\rho (k,t))(1\rho (p,t))\rho (kq,t)\rho (p+q,t)`$ (82)
$``$ $`\rho (k,t)\rho (p,t)(1\rho (kq,t))(1\rho (p+q,t)),`$ (83)
all propagators are collected in
$`S(t,\tau )`$ $`=`$ $`\widehat{g}^R(k,t,\tau )\widehat{g}^R(p,t,\tau )`$ (84)
$`\times `$ $`\widehat{g}^A(kq,t,\tau )\widehat{g}^A(p+q,t,\tau ).`$ (85)
We have used that the Wigner distribution is a real function and the complex conjugacy of propagators, $`\widehat{g}^R(k,t,\tau )=\left(\widehat{g}^A(k,t,\tau )\right)^{}`$, to reduce the expressions.
Let us assume for the moment that the propagators do not depend on the center-of-mass time $`t`$. In this case the only $`t`$ dependence in the collision integral of the Levinson equation (81) is due to distributions. Since for the Born approximation all distributions are retarded in an equal way, we can follow a modification of the approach used above.
For slow perturbations one can expand (81) to the lowest order in the memory, $`D(t\tau )=D(t)\tau _tD(t)`$, so that the collision integral splits into the Boltzmann-like scattering integral and the gradient correction, $`_t\rho =+_t`$, with
$``$ $`=`$ $`{\displaystyle \frac{dpdq}{(2\pi )^6}𝒱_q^2D(t)2\mathrm{Re}\underset{0}{\overset{\mathrm{}}{}}𝑑\tau S(\tau )},`$ (86)
$``$ $`=`$ $`{\displaystyle \frac{dpdq}{(2\pi )^6}𝒱_q^2D(t)2\mathrm{Re}\underset{0}{\overset{\mathrm{}}{}}𝑑\tau \tau S(\tau )}.`$ (87)
Please note that $``$ and $``$ are different from (30). Again, from (81) we can define a new function to remove the retardation,
$$=\rho ,$$
(88)
and try to redefine the transport in terms of $``$. This will us lead now to unavoidable double counts.
#### 1 Double counts
The integrals in (86) are controlled by two mechanisms, the de-phasing given by momentum integrals and the decay of propagators. Let us assume for a while that the de-phasing is the dominant part leaving the decay aside. Within the quasiparticle approximation, $`\widehat{g}^R(k,\tau )i\mathrm{e}^{i\epsilon _k\tau }`$, which is free of the decay, the time integrals over propagators yield
$`2\mathrm{Re}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑\tau S_{\mathrm{qp}}(\tau )`$ $`=`$ $`2\pi \delta \left(\epsilon _k+\epsilon _p\epsilon _{kq}\epsilon _{p+q}\right),`$ (89)
$`2\mathrm{Re}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑\tau \tau S_{\mathrm{qp}}(\tau )`$ $`=`$ $`2{\displaystyle \frac{\mathrm{}^{}}{\epsilon _k+\epsilon _p\epsilon _{kq}\epsilon _{p+q}}}.`$ (90)
After approximation (89), $``$ from (86) turns into the Boltzmann-type scattering integral $`I_B`$ of the Landau-Silin equation. Naturally, the on-shell contributions are identical since they do not depend on the retardation.
Using (90) in (87) one finds
$$=2\frac{dpdq}{(2\pi )^6}𝒱_q^2D(t)\frac{\mathrm{}^{}}{\epsilon _k+\epsilon _p\epsilon _{kq}\epsilon _{p+q}}.$$
(91)
The function $``$ is twice the off-shell contribution $`R`$ to the Wigner distribution. To show this, we express (91) in terms of the Born selfenergy (80) and the approximative correlation functions, $`g^<=2\pi \delta (\omega \epsilon )\rho `$, as
$$=2\frac{d\omega ^{}d\omega }{(2\pi )^2}\frac{\mathrm{}^{}}{\omega \omega ^{}}\left(\sigma _\omega ^{}^>g_\omega ^<\sigma _\omega ^{}^<g_\omega ^>\right).$$
(92)
Comparing with the second term of (42) one can see that the additional retardation results in the double count of the correlated part, $`=2R`$. Due to this double count, the function $`=\rho =\rho 2R=fR`$ cannot be interpreted as the distribution of excitations because it has large regions of negative values. Please note a common met pitfall in literature related to this double count given in appendix C.
We note that in numerical treatments of the Levinson equation neglecting or underestimating the decay, leads to a numerical instability of the equation, as it has been reported by Haug . Other numerical solutions show a continuous increase of kinetic energy, e.g. figure 4 of or pulsation modes . The instability of the solution is a problem which shows that the Levinson equation is a difficult tool to handle. The double count of the off-shell contribution resulting in the negative effective occupation factors is likely one of the reasons of this problem. A detailed discussion and numerical comparison of the Levinson equation with the Kadanoff and Baym equation can be found in .
#### 2 Decay of propagators
The double count of correlations is a serious mistake for theories which aim to go beyond the Boltzmann equation. Here we show that the decay of propagators during the collision removes this double count.
We use the simplest approximation of the decay,
$$\widehat{g}^R(k,t,\tau )=i\mathrm{e}^{i\epsilon _k\tau }\mathrm{e}^{\gamma _k\frac{\tau }{2}},$$
(93)
where $`\epsilon _k=\epsilon (k,t\frac{\tau }{2})`$ and $`\gamma _k=\gamma _{\epsilon (k,t\frac{\tau }{2})}(k,t\frac{\tau }{2})`$ is the pole value of $`\gamma =\sigma ^>+\sigma ^<=2\mathrm{I}\mathrm{m}\sigma ^R`$. The time argument shows that external or mean fields can modify the propagation during the collision. The collected propagator thus reads
$$S(t\frac{\tau }{2},\tau )=S_{\mathrm{qp}}(t\frac{\tau }{2},\tau )\mathrm{e}^{(\gamma _k+\gamma _p+\gamma _{kq}+\gamma _{p+q})\frac{\tau }{2}}.$$
(94)
The kinetic equation holds only if the collision duration is short on the scale of the lifetime, $`(2\gamma )^1`$. In this limit we can expand the exponential decay,
$$\mathrm{e}^{(\gamma _k+\gamma _p+\gamma _{kq}+\gamma _{p+q})\frac{\tau }{2}}=1(\gamma _k+\gamma _p+\gamma _{kq}+\gamma _{p+q})\frac{\tau }{2}.$$
(95)
Since we are dealing with the correction terms, we can use a simple approximation of the kinetic equation, $`_t\rho _k=\sigma _k^<\gamma _k\rho `$, to rearrange the contribution of $`\gamma `$’s to its integrand,
$$\frac{\tau }{2}(\gamma _k+\gamma _p+\gamma _{kq}+\gamma _{p+q})=\frac{\tau }{2}\frac{}{t}D\frac{\tau }{2}D_{3\mathrm{p}}.$$
(96)
The last term,
$`D_{3\mathrm{p}}=\sigma _k^>(1\rho _p)\rho _{kq}\rho _{p+q}+(1\rho _k)\sigma _p^>\rho _{kq}\rho _{p+q}`$ (97)
$`+(1\rho _k)(1\rho _p)\sigma _{kq}^<\rho _{p+q}+(1\rho _k)(1\rho _p)\rho _{kq}\sigma _{p+q}^<`$ (98)
$`\sigma _k^<\rho _p(1\rho _{kq})(1\rho _{p+q})\rho _k\sigma _p^<(1\rho _{kq})(1\rho _{p+q})`$ (99)
$`\rho _k\rho _p\sigma _{kq}^>(1\rho _{p+q})\rho _k\rho _p(1\rho _{kq})\sigma _{p+q}^>,`$ (100)
can be neglected leading to the three-particle contribution to the scattering integral. Please note that such neglects have not been necessary in our suggested construction $`\rho [f]`$ outlined in Sec. II.
Within the linear approximation in time gradients one finds that the decay reduces the retardation,
$$\mathrm{e}^{(\gamma _k+\gamma _p+\gamma _{kq}+\gamma _{p+q})\frac{\tau }{2}}D(t\tau )=D(t\frac{\tau }{2},\tau ).$$
(101)
Now the collective propagator and distribution have the same time retardation of the center-of-mass time being equal to $`\frac{\tau }{2}`$. This time shift removes the double count and reproduces the result obtained already in Sec. II.
We note that the original Levinson equation does not include the decay of propagators in the scattering integral. The presented discussion shows that the decay plays the important role for the correct treatment of renormalizations and other contributions of correlations.
### B Collision delay
We have seen that within the Born approximation the retarded collision integral of the Levinson equation can be transformed into the instantaneous scattering integral of the Boltzmann type and quasiparticle renormalizations of the single-particle distribution and energy. Beyond the Born approximation, there is the additional retardation expressed by times $`\tau _{R,A}`$ in the selfenergy (80). This retardation is responsible for the virial correction, in particular for the correlated density $`n_c`$. Since the retardation times $`\tau _{R,A}`$ do not enter the final states of the collision, they cannot be removed in the above manner. On the other hand, the time integrals over $`\tau _{R,A}`$ can be eliminated so that the collision duration is expressed by an effective time-non-locality of the scattering integral. For simplicity of presentation in this section we assume that T-matrices and selfenergies are independent of the center-of-mass time. This will lead us to the basic idea. The complete result with all time and space dependences are given in the next Sec. IV.
Within the linear approximation in time gradients we can treat the contributions of times $`\tau _{R,A}`$ separately from other retardations. We first implement the quasiparticle ansatz and expand in retardations,
$`\widehat{g}^<(kq,t{\displaystyle \frac{1}{2}}(\tau _R+\tau _A+\tau ),\tau \tau _R+\tau _A)`$ (102)
$`\times `$ $`\widehat{g}^<(p+q,t{\displaystyle \frac{1}{2}}(\tau _R+\tau _A+\tau ),\tau \tau _R+\tau _A)`$ (103)
$`=`$ $`\left(1{\displaystyle \frac{1}{2}}\tau {\displaystyle \frac{}{t}}{\displaystyle \frac{1}{2}}(\tau _R+\tau _A){\displaystyle \frac{}{t}}\right)`$ (104)
$`\times `$ $`f_{kq}(t)f_{p+q}(t)\mathrm{e}^{i(\epsilon _{kq}+\epsilon _{p+q})(\tau \tau _R+\tau _A)}.`$ (105)
The contribution of $`\frac{1}{2}\tau \frac{}{t}`$ has been discussed above, now we focus on the term due to $`\frac{1}{2}(\tau _R+\tau _A)\frac{}{t}`$.
Introducing the energy representation of the T-matrix,
$$𝒯^R(\tau _R)=\frac{d\omega _R}{2\pi }t^R(\omega _R)\mathrm{e}^{i\omega \tau _R},$$
(106)
and a complex conjugate for $`𝒯^A`$, the contribution of $`\frac{1}{2}(\tau _R+\tau _A)\frac{}{t}`$ to the selfenergy (80) reads
$`\mathrm{\Delta }\widehat{\sigma }^<(k,t,\tau )`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{dpdq}{(2\pi )^6}\mathrm{e}^{i\epsilon _p\tau }(1f_p)}`$ (107)
$`\times `$ $`{\displaystyle \frac{d\omega _Rd\omega _A}{(2\pi )^2}t^R(k,p,q,\omega _R)t^A(k,p,q,\omega _A)}`$ (108)
$`\times `$ $`{\displaystyle 𝑑\tau _R𝑑\tau _A\mathrm{e}^{i\omega _R\tau _R}\mathrm{e}^{i\omega _A\tau _A}(\tau _R+\tau _A)}`$ (109)
$`\times `$ $`{\displaystyle \frac{}{t}}f_{kq}f_{p+q}\mathrm{e}^{i(\epsilon _{kq}+\epsilon _{p+q})(\tau \tau _R+\tau _A)}.`$ (110)
The integral over $`\tau _{R,A}`$ results in the derived $`\delta `$-functions which can be readily integrated out,
$`\mathrm{\Delta }\widehat{\sigma }^<(k,t,\tau )`$ $`=`$ $`{\displaystyle \frac{dpdq}{(2\pi )^6}\mathrm{e}^{i(\epsilon _p\epsilon _{kq}\epsilon _{p+q})\tau }(1f_p)}`$ (111)
$`\times `$ $`{\displaystyle \frac{1}{2i}}\left({\displaystyle \frac{t^R}{\omega }}t^At^R{\displaystyle \frac{t^A}{\omega }}\right)_{\omega =\epsilon _{kq}+\epsilon _{p+q}}`$ (112)
$`\times `$ $`{\displaystyle \frac{}{t}}f_{kq}f_{p+q}.`$ (113)
The correction (113) is proportional to the collision delay defined from the phase shift,
$$\mathrm{\Delta }_t=\frac{\varphi }{\omega }.$$
(114)
Indeed, writing the T-matrices in the form $`t^{R,A}=|t|\mathrm{e}^{i\varphi }`$, where $`|t|^2`$ is the amplitude of the scattering rate and $`\varphi `$ is the scattering phase shift, we obtain
$$\frac{1}{2i}\left(\frac{t^R}{\omega }t^At^R\frac{t^A}{\omega }\right)=|t|^2\mathrm{\Delta }_t.$$
(115)
This allows us to join the correction with the time-local part of the selfenergy,
$`\widehat{\sigma }^<(k,t,\tau )`$ $`=`$ $`{\displaystyle \frac{dpdq}{(2\pi )^6}|t|^2\mathrm{e}^{i(\epsilon _p\epsilon _{kq}\epsilon _{p+q})\tau }\left(1f_p(t)\right)}`$ (116)
$`\times `$ $`f_{kq}(t\mathrm{\Delta }_t)f_{p+q}(t\mathrm{\Delta }_t).`$ (117)
The selfenergy (117) can be compared with (80). First, the retardation by $`\tau /2`$ has been suppressed because it is eliminated by the method described in Sec. II. Second, the retardations by the internal times of the T-matrices, $`\tau _{R,A}`$, have been reduced from the integral form of (80) to the effective time non-locality of (117). The latter form includes only a single characteristic time $`\mathrm{\Delta }_t`$ which is similar (but not fully identical) to the collision delay defined by Wigner from the scattering phase shift. Note that in form (117) no time integration appears.
An interpretation of the collision delay $`\mathrm{\Delta }_t`$ is troubled by the fact that this effective time does not have to be positive. Apparently, one cannot interpret the collision delay as a mean time scale of the $`\tau _{R,A}`$ integrals in the selfenergy (80). The negative value of $`\mathrm{\Delta }_t`$ appears when the phase shift decreases with increasing energy, as it happens above a bound state or a resonant scattering state. For instance, the collision delay of nucleon-nucleon collisions is negative at low relative energies due to the deuteron bound state, see . The negative value of the collision delay merely tells that particles anti-correlate at a given relative energy.
## IV Nonlocal scattering integral
The selfenergy (117) is easily converted into the scattering integral. The $`|t|^2`$ represents the scattering amplitude in the T-matrix approximation. The exponential function turns into the energy conserving $`\delta `$-function. Compared to the common scattering integral of the Landau-Silin equation, the only difference consists in the collision delay $`\mathrm{\Delta }_t`$ entering the initial states of the collision. Of course, in the previous section we have assumed a homogeneous system with time-independent selfenergies and T-matrices. For a general case more nonlocal corrections appear. These corrections we discuss in this section.
In this section we thus focus on the scattering integral. Our aim is to derive the scattering integral within the same approximations as the skeleton kinetic equation discussed in Sec. II, i.e., keeping all linear gradients. This procedure leads to the non-instantaneous and non-local scattering integral. We use the Galitskii-Feynman T-matrix approximation of the selfenergy which is appropriate for dense systems of particles interacting via short range potentials. The T-matrix approximation for non-equilibrium Green functions is briefly presented in Appendix A.
### A Quasi-classical limit
The consistent treatment of gradient contributions is provided by the quasi-classical limit, i.e., the linear expansion in gradients. The collision delay derived above shows that beside gradient contributions forming the drift in the skeleton equation, there are non-trivial gradient contributions making the scattering integral non-local in time. Transforming the system into a running coordinate framework, one can see that the collision delay has be accompanied by the space non-locality of the collision integral. Accordingly, the gradient expansion has to be applied also to all internal space and time integrations of the scattering integral.
The quasi-classical limit and the limit of small scattering rates explicitly determine how to evaluate the scattering integral from the selfenergy $`\sigma ^<`$. For non-degenerate systems, a very similar scheme was carried out by Bärwinkel . One can see in Bärwinkel’s papers, that the scattering integral is troubled by a large set of gradient corrections. This formal complexity seems to be the main reason why most authors either neglect gradient corrections at all or provide them buried in multi-dimensional integrals . For a degenerate system, the set of gradient corrections to the scattering integral is even larger than for rare gases studied by Bärwinkel, see . To avoid manipulations with long and obscure formulas, the gradient corrections have to be sorted and expressed in a comprehensive form. To this end, we will express all gradient corrections in the form of shifted arguments, as we have done above for the collision delay.
The quasi-classical limit of the scattering integral with all linear gradients kept is a tedious but straightforward algebraic exercise, see Appendix B, giving
$`\sigma ^<(\omega ,k,r,t)`$ (118)
$`={\displaystyle \frac{dp}{(2\pi )^3}\frac{dE}{2\pi }\frac{dq}{(2\pi )^3}\frac{d\mathrm{\Omega }}{2\pi }\left(1\frac{1}{2}\frac{\mathrm{\Delta }_2}{r}\right)}`$ (119)
$`\times \left|t(\omega +E\mathrm{\Delta }_E,k{\displaystyle \frac{1}{2}}\mathrm{\Delta }_K,p{\displaystyle \frac{1}{2}}\mathrm{\Delta }_K,q,r\mathrm{\Delta }_r,t{\displaystyle \frac{1}{2}}\mathrm{\Delta }_t)\right|^2`$ (120)
$`\times g^>(E,p,r\mathrm{\Delta }_2,t)`$ (121)
$`\times g^<(\omega \mathrm{\Omega }\mathrm{\Delta }_E,kq\mathrm{\Delta }_K,r\mathrm{\Delta }_3,t\mathrm{\Delta }_t)`$ (122)
$`\times g^<(E+\mathrm{\Omega }\mathrm{\Delta }_E,p+q\mathrm{\Delta }_K,r\mathrm{\Delta }_4,t\mathrm{\Delta }_t).`$ (123)
With $`g^{>,<}`$ substituted from the quasiparticle approximation, this formula turns into the functional of the quasiparticle distribution. As one can see, the initial states, $`kq`$ and $`p+q`$, include the collision delay $`\mathrm{\Delta }_t`$, as in the homogeneous system describe by (117).
Beside the collision delay, non-local selfenergy (123) includes a number of other gradient corrections which enter nearly all arguments. These are expressed via $`\mathrm{\Delta }`$’s which are derivatives of the scattering phase shift $`\varphi `$, where $`t^{R,A}=|t|\mathrm{e}^{i\varphi }`$, according to the following list
$`\mathrm{\Delta }_2`$ $`=`$ $`{\displaystyle \frac{\varphi }{p}}{\displaystyle \frac{\varphi }{q}}{\displaystyle \frac{\varphi }{k}},`$ (124)
$`\mathrm{\Delta }_3`$ $`=`$ $`{\displaystyle \frac{\varphi }{k}},`$ (125)
$`\mathrm{\Delta }_4`$ $`=`$ $`{\displaystyle \frac{\varphi }{q}}{\displaystyle \frac{\varphi }{k}},`$ (126)
$`\mathrm{\Delta }_t`$ $`=`$ $`{\displaystyle \frac{\varphi }{\omega }},`$ (127)
$`\mathrm{\Delta }_K`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{\varphi }{r}},`$ (128)
$`\mathrm{\Delta }_E`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{\varphi }{t}},`$ (129)
$`\mathrm{\Delta }_r`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left(\mathrm{\Delta }_2+\mathrm{\Delta }_3+\mathrm{\Delta }_4\right).`$ (130)
Finally, there is the energy gain which is discussed in . Numerical values of the delay $`\mathrm{\Delta }_t`$ and the displacements $`\mathrm{\Delta }_{24}`$ for the scattering of isolated particles are presented in .
Sending all $`\mathrm{\Delta }`$’s to zero one recovers the instantaneous and local approximation of the selfenergy, Eq. (123). The non-instantaneous and non-local corrections given by $`\mathrm{\Delta }`$’s appear in three ways: in arguments of the correlation functions, in arguments of the T-matrix, and as a fore-factor. In correlation functions, the $`\mathrm{\Delta }`$’s describe displacements of the initial and final positions of colliding particles, the final state of one of the particles is fixed to the balanced phase-space point $`(k,r,t)`$. The arguments of the T-matrix correspond to the center of arguments of all initial and final states.
The pre-factor has no special physical meaning as it depends on the actual choice of energy and momentum variables. For instance, using in (123) an integral over $`E^{}=E\mathrm{\Delta }_E`$ instead of the integral over $`E`$, a corresponding factor, $`dE/dE^{}=1+_\omega \mathrm{\Delta }_E=1\frac{1}{2}_t\mathrm{\Delta }_t`$, appears.
### B Extended quasiparticle approximation
Now we can complete the derivation of the kinetic equation. Since the scattering integral is already proportional to the ‘small’ scattering rate, we can use from the extended quasiparticle approximation (49) the pole part, $`g^<(\omega ,k,r,t)=f(k,r,t)2\pi \delta (\omega \epsilon (k,r,t))`$, to convert selfenergies $`\sigma ^{>,<}`$ into functionals of the quasiparticle distribution. The selfenergy $`\sigma ^<`$ is given by (123) and $`\sigma ^>`$ is obtained by the interchange $`><`$. The resulting kinetic equation reads
$`{\displaystyle \frac{f_1}{t}}+{\displaystyle \frac{\epsilon _1}{k}}{\displaystyle \frac{f_1}{r}}{\displaystyle \frac{\epsilon _1}{r}}{\displaystyle \frac{f_1}{k}}`$ (131)
$`={\displaystyle \frac{dpdq}{(2\pi )^6}2\pi \delta \left(\epsilon _1+\overline{\epsilon }_2\overline{\epsilon }_3\overline{\epsilon }_42\mathrm{\Delta }_E\right)}`$ (132)
$`\times \left(1{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mathrm{\Delta }_2}{r}}{\displaystyle \frac{\overline{\epsilon }_2}{r}}{\displaystyle \frac{\mathrm{\Delta }_2}{\omega }}\right)`$ (133)
$`\times \left|t(\epsilon _1+\overline{\epsilon }_2\mathrm{\Delta }_E,k{\displaystyle \frac{1}{2}}\mathrm{\Delta }_K,p{\displaystyle \frac{1}{2}}\mathrm{\Delta }_K,q,r\mathrm{\Delta }_r,t{\displaystyle \frac{1}{2}}\mathrm{\Delta }_t)\right|^2`$ (134)
$`\times \left((1f_1)(1\overline{f}_2)\overline{f}_3\overline{f}_4f_1\overline{f}_2(1\overline{f}_3)(1\overline{f}_4)\right).`$ (135)
(136)
The abbreviated notation of $`f`$’s and $`\epsilon `$’s means
$$\begin{array}{ccc}\hfill \epsilon _1& & \epsilon (k,r,t),\hfill \\ \hfill \overline{\epsilon }_2& & \epsilon (p,r\mathrm{\Delta }_2,t),\hfill \\ \hfill \overline{\epsilon }_3& & \epsilon (kq\mathrm{\Delta }_K,r\mathrm{\Delta }_3,t\mathrm{\Delta }_t),\hfill \\ \hfill \overline{\epsilon }_4& & \epsilon (p+q\mathrm{\Delta }_K,r\mathrm{\Delta }_4,t\mathrm{\Delta }_t).\hfill \end{array}$$
(137)
The bar reminds that all $`\mathrm{\Delta }`$’s appear in arguments with negative signs. In the $`\mathrm{\Delta }`$’s and their derivatives the energy $`\omega +E=\epsilon _1+\overline{\epsilon }_2`$ is substituted after all derivatives are taken. Note that the pre-factor has changed compared to (123). In (123), $`\overline{\epsilon }_2`$ depends on $`\mathrm{\Delta }_2`$ which depends on the energy $`E`$. This $`E`$-dependence has resulted in a norm of the singularity.
The non–local kinetic equation (136) represents the main result of this paper. In the following we will discuss some symmetry properties of this equation.
### C Markovian approximation
In the scattering integrals of (137) the collision delay $`\mathrm{\Delta }_t`$ appears together with displacements $`\mathrm{\Delta }_{24}`$. Actual values of these types of corrections are closely linked. In the original Wigner’s approach, the collision delay was identified from the position of the center of wave packets in the final state. Apparently, this position can be described either as the space displacement or via the delay. A corresponding rearrangement of the scattering integral is achieved using the free-space time evolution of the distribution, e.g.,
$`\overline{f}_3`$ $`=`$ $`f(kq\mathrm{\Delta }_K,r\mathrm{\Delta }_3,t\mathrm{\Delta }_t)`$ (138)
$``$ $`f(kq\mathrm{\Delta }_K+F_3\mathrm{\Delta }_t,r\mathrm{\Delta }_3+v_3\mathrm{\Delta }_t,t),`$ (139)
where $`v_3=_k\epsilon _3`$ is a velocity and $`F_3=_r\epsilon _3`$ is a force acting on the particle in the final state. Substituting the right hand side (and similar for $`\overline{f}_4`$ and $`\overline{\epsilon }_{3,4}`$) into (137), one obtains instantaneous scattering integrals. This liberty of rearrangement is only approximative, some quantities are sensitive to such approximations. For instance, the correlated density discussed below is directly proportional to $`\mathrm{\Delta }_t`$ as it measures an effective number of particles in the collision process. On the other hand, the rearrangement into the instantaneous form makes the kinetic equation Markovian what significantly simplifies its numerical treatment, in particular if negative collision delay takes place.
### D Space-time symmetry of non-local corrections
The scattering-out integral in (136) has the same sign of $`\mathrm{\Delta }`$’s as the scattering-in. Accordingly, due to the non-local and non-instantaneous corrections, the scattering-out cannot be obtained from the scattering-in by the space-time symmetry, as it is common in the classical approach to the kinetic equation. This problem has a practical impact. The simulations of pressure with non-instant corrections has been discussed by . In Monte-Carlo simulations based on semi-classical trajectories the space-time symmetry is an unavoidable part of the scattering integral. To understand internal times of collisions and non-local corrections, it is thus necessary to clarify the space-time symmetry.
It is possible to show that the scattering-in and -out are connected by the space-time symmetry, one has to take into account, however, the non-local and non-instantaneous corrections to the T-matrix itself. To this end we employ the optical theorem which can be expressed in two forms,
$$\mathrm{Im}𝒯=𝒯^R𝒜𝒯^A=𝒯^A𝒜𝒯^R.$$
(140)
Here, symbol Im denotes the anti-Hermitian part. The multiplication includes integrals over two-particle states and times. The function $`𝒜_{34}=G_3^>G_4^>G_3^<G_4^<=G^>(3,3^{})G^>(4,4^{})G^<(3,3^{})G^<(4,4^{})`$ is the two-particle spectral function which includes the Pauli blocking of internal states in the Galitskii-Feynman approximation. Both forms of $`\mathrm{Im}𝒯`$ can be derived by algebraic manipulations with the ladder equation for $`𝒯^{R,A}`$ and hold out of equilibrium.
Implementation of the optical theorem (140) requires to rearrange the scattering integrals. This is most conveniently done on the level of the starting non-equilibrium Green functions. The integrand in the right hand side of (19) is composed of two terms, in a reduced notation $`I_l=G_1^>\mathrm{\Sigma }_1^<G_1^<\mathrm{\Sigma }_1^>`$ and $`I_r=\mathrm{\Sigma }_1^<G_1^>\mathrm{\Sigma }_1^>G_1^<`$. For the Galitskii-Feynman selfenergy,
$$I_l=G_1^>G_2^>𝒯^RG_3^<G_4^<𝒯^AG_1^<G_2^<𝒯^RG_3^>G_4^>𝒯^A,$$
(141)
we have to add and subtract $`G_1^<G_2^<𝒯^RG_3^<G_4^<𝒯^A`$ to the second term of (141) to achieve the Galitskii-Feynman two-particle spectral function,
$$I_l=𝒜_{12}𝒯^RG_3^<G_4^<𝒯^AG_1^<G_2^<𝒯^R𝒜_{34}𝒯^A.$$
(142)
In the last term we can apply the optical theorem (140) for the interchange $`𝒯^R𝒜_{34}𝒯^A𝒯^A𝒜_{34}𝒯^R`$, therefore
$$I_l=𝒜_{12}𝒯^RG_3^<G_4^<𝒯^AG_1^<G_2^<𝒯^A𝒜_{34}𝒯^R.$$
(143)
In the same way one can rearrange $`I_r`$.
Technically, the interchange of the retarded and the advanced T-matrices is equivalent to the change of sign of the phase shift $`\varphi `$, $`𝒯^Rt^R=|t|\mathrm{e}^{i\varphi }`$ and $`𝒯^At^A=|t|\mathrm{e}^{i\varphi }`$. The implementation of the optical theorem thus has no effect on the local parts of the scattering integral which depend exclusively on the amplitude $`|t|`$ of the T-matrix. In contrast, all non-local corrections given by $`\mathrm{\Delta }`$’s (124-130) are linear functions of the phase shift $`\varphi `$ and thus reverse their signs. In the kinetic equation (136), the combination corresponding to $`𝒜_{34}`$ is achieved by regrouping distributions as
$`(1f_1)(1\overline{f}_2)\overline{f}_3\overline{f}_4f_1\overline{f}_2(1\overline{f}_3)(1\overline{f}_4)`$ (144)
$`=(1f_1\overline{f}_2)\overline{f}_3\overline{f}_4f_1\overline{f}_2(1\overline{f}_3\overline{f}_4).`$ (145)
The intermediate-state Pauli-blocking factor $`1\overline{f}_3\overline{f}_4`$ is the occupation number corresponding to $`𝒜_{34}`$. Inverting signs of the phase shift, i.e., of $`\mathrm{\Delta }`$’s, related to the scattering integral with the Pauli blocking $`1\overline{f}_3\overline{f}_4`$ one arrives at the kinetic equation
$`{\displaystyle \frac{f_1}{t}}+{\displaystyle \frac{\epsilon _1}{k}}{\displaystyle \frac{f_1}{r}}{\displaystyle \frac{\epsilon _1}{r}}{\displaystyle \frac{f_1}{k}}`$ (146)
$`={\displaystyle \frac{dp}{(2\pi )^3}\frac{dq}{(2\pi )^3}\left(1\frac{1}{2}\frac{\mathrm{\Delta }_2}{r}\frac{\epsilon _2}{r}\frac{\mathrm{\Delta }_2}{\omega }\right)}`$ (147)
$`\times \left|t(\epsilon _1+\overline{\epsilon }_2\mathrm{\Delta }_E,k{\displaystyle \frac{1}{2}}\mathrm{\Delta }_K,p{\displaystyle \frac{1}{2}}\mathrm{\Delta }_K,q,r\mathrm{\Delta }_r,t{\displaystyle \frac{1}{2}}\mathrm{\Delta }_t)\right|^2`$ (148)
$`\times 2\pi \delta \left(\epsilon _1+\overline{\epsilon }_2\overline{\epsilon }_3\overline{\epsilon }_42\mathrm{\Delta }_E\right)(1f_1\overline{f}_2)\overline{f}_3\overline{f}_4`$ (149)
$`{\displaystyle \frac{dp}{(2\pi )^3}\frac{dq}{(2\pi )^3}\left(1+\frac{1}{2}\frac{\mathrm{\Delta }_2}{r}+\frac{\epsilon _2}{r}\frac{\mathrm{\Delta }_2}{\omega }\right)}`$ (150)
$`\times \left|t(\epsilon _1+\epsilon _2+\mathrm{\Delta }_E,k+{\displaystyle \frac{1}{2}}\mathrm{\Delta }_K,p+{\displaystyle \frac{1}{2}}\mathrm{\Delta }_K,q,r+\mathrm{\Delta }_r,t+{\displaystyle \frac{1}{2}}\mathrm{\Delta }_t)\right|^2`$ (151)
$`\times 2\pi \delta \left(\epsilon _1+\epsilon _2\epsilon _3\epsilon _4+2\mathrm{\Delta }_E\right)f_1f_2(1f_3f_4).`$ (152)
(153)
In functions with bars arguments are given by (137), in functions without bars arguments are identical except for the reversed (i.e., positive) signs of all $`\mathrm{\Delta }`$’s.
The kinetic equation with the space-time symmetry of scattering integrals reminds to the Enskog equation for classical hard spheres with respect to two features. First, if one suppresses the collision delay, see (139), the scattering-out becomes the space inversion of the scattering-in. Second, the scattering rate (given by $`|t|^2`$) is centered between initial and final states of collisions.
The kinetic equation (153) has a difficult interpretation because its scattering integrals have positive or negative values following the sign of the Pauli blocking factors $`1f_1f_2=(1f_1)(1f_2)f_1f_2`$ and $`1f_3f_4=(1f_3)(1f_4)f_3f_4`$. This strange behavior appears due to the stimulated transition processes given by the terms $`f_1f_2`$ and $`f_3f_4`$ on top of the standard Pauli blocking, $`(1f_1)(1f_2)`$ and $`(1f_3)(1f_4)`$. While in the local approximation of the scattering integrals the -in and -out stimulated transitions exactly compensate giving no net contribution, the non-local corrections translate this feature known from the equilibrium Green functions on the level of kinetic equation. Our difficulties with the classical interpretation of (153) thus may follow from our low understanding of stimulated processes in general. At very low temperatures, the stimulated transitions become so strong that the above quasiparticle picture breaks down and the system develops super-conductivity having an additional degree of freedom and a different spectrum of single-particle excitations. Although the Galitskii-Feynman approximation is not suited to describe the super-conducting state, it reveals a singularity at $`T_c`$ which signals its onset.
Let us put aside problems with stimulated processes and return to our discussion of characteristic times. For the moment we will assume that the collision delay $`\mathrm{\Delta }_t`$ is positive. In the scattering-in integral which is the first term on the right hand side of (153), the initial condition is at the time $`t\mathrm{\Delta }_t`$, i.e., before the time $`t`$ at which we balance the changes at the phase-space point $`(k,r)`$. This has a clear interpretation, at $`t\mathrm{\Delta }_t`$ the particle of momentum $`kq`$ enters into the collision process from which it is released at $`t`$ into $`(k,r)`$. In the scattering-out process, the particle leaves the cell $`(k,r)`$ at $`t`$ when it enters the collision process. This process will end at $`t+\mathrm{\Delta }_t`$ when we have to check the accessible phase space. While the scattering-out has a natural interpretation, its description violates the causality. Indeed, the occupation factors in future time $`t+\mathrm{\Delta }_t`$ are supposed to decide whether the process comes through or whether it is blocked. We note that the violation of the causality on the microscopic time scale is common in kinetic equations. Already the Fermi Golden Rule includes a time integral into future which eliminates off-shell processes. In fact, it is the Fermi-Golden-Rule-type approximation, in this paper represented by the pole approximation, which makes negative values of the collision delay possible.
For a detailed discussion of the space-time versus particle-hole symmetry for various forms of the two–particle spectral function see .
### E Correlated density
As mentioned in the introduction, in the system described on the level of the Galitskii-Feynman approximation the density of quasiparticles, $`n_f=\frac{dk}{(2\pi )^3}f`$, differs from the density of composing particles $`n=\frac{dk}{(2\pi )^3}\rho `$ by the amount called the correlated density $`n_c=nn_f`$. After some algebra, from relation (42) one finds
$`n_c`$ $`=`$ $`{\displaystyle \frac{dkdpdq}{(2\pi )^8}|t|^2\mathrm{\Delta }_t}`$ (155)
$`\times \delta (\epsilon _1+\epsilon _2\epsilon _3\epsilon _4)f_1f_2(1f_3f_4).`$
The functions $`|t|`$, $`\epsilon `$’s and $`f`$’s are in the local approximation having no $`\mathrm{\Delta }`$’s in arguments. This expression is the generalized Beth-Uhlenbeck formula . The correlated density is proportional to the energy-derivative of the phase shift expressed here via the collision delay.
This correlated density is consistent with the equation of continuity found from the kinetic equation (153). For simplicity we will assume a homogeneous system for the equation of continuity reading, $`_tn=0`$. Taking the momentum integral over (153) one finds
$`{\displaystyle \frac{}{t}}`$ $`{\displaystyle \frac{dk}{(2\pi )^3}f_1}={\displaystyle \frac{}{t}}{\displaystyle \frac{dkdpdq}{(2\pi )^8}|t|^2\mathrm{\Delta }_t}`$ (157)
$`\times \delta (\epsilon _1+\epsilon _2\epsilon _3\epsilon _4)f_1f_2(1f_3f_4).`$
Using (155) in the right hand side, one finds $`_t(n_f+n_c)=0`$. It should be noted that the number of quasiparticles is not conserved, $`_tn_f0`$, because during the time interval $`\mathrm{\Delta }_t`$ the colliding particles are excluded from the single-particle statistics being in two-particle scattering states.
One can see that instantaneous approximations of the kinetic equation cannot capture the correlated density. In the same time, the correct description of the quasiparticle density makes the choice (127) of the collision delay preferable to other characteristic times one can define in the spirit of Wigner from the asymptotic behavior of final states of binary collisions. In Appendix C we comment on some mistakes done in previous studies of the correlated density.
The complete proof of conservation laws with explicit expressions for the correlated pressure and energies can be found in .
## V Summary
In absence of the time operator, there are many definitions of characteristic times one can associate with collisions. Regardless of a definition one uses, if the finite duration of collisions is included in a kinetic equation, this equation is non-Markovian since the initial and final states of the collision are at distinct times. The family of non-Markovian kinetic equation also includes retarded equations of Levinson type in which the collision is expressed in terms of the time integral describing the whole process of the two-particle interaction.
For slowly varying systems, we have shown that one half of the Levinson-type retardation describes the off-shell motion and corresponding renormalizations of single-particle functions. The second half compensates the decay of propagators during the integration and can be eliminated leaving the pole contribution to the scattering integral identical to the one obtained from the quasiparticle approximation.
Finally, we have discussed the collision delay resulting from the energy dependence of the scattering phase shift. We have shown that this non-Markovian correction is connected with other non-local corrections to the scattering integral which gives us a freedom to define the collision duration in different ways. The physical quantity sensitive to the actual choice of the collision delay is the density of quasiparticles, or its complementary quantity, the density of correlated particles. The method presented in this paper, and at the same time the choice of the collision delay, reproduces the correlated density obtained within the generalized Beth-Uhlenbeck approach.
The non-instant and non-local corrections given by the $`\mathrm{\Delta }`$’s do not change the structure and the overall interpretation of the scattering integral but only slightly renormalize its ingredients. The exclusive dependence of the non-local and non-instant corrections on the scattering phase shift confirms results from the theory of gases obtained by very different technical tools. The non-instant and non-local scattering integral in the form (153) parallels the classical Enskog’s equation, therefore it can be treated with numerical tools developed for the theory of classical gases, see e.g. at least as long as $`1f_1f_2`$ remains positive.
The virial corrections to the balance equations appear from intrinsic gradients instead of correlation parts in the equation of the reduced density matrix. In it is shown that the nonlocal kinetic equation leads to complete balance equations for the density, energy and stress tensor which establish conservation laws including correlated parts. The nonlocal kinetic equation derived here requires no more numerical efforts than solving the Boltzmann equation. The numerical solution and application has been demonstrated in .
###### Acknowledgements.
K.M. likes to thank the LPC for a friendly and hospitable atmosphere and Rainer Klages for discussions and useful hints.
## A T-matrix approximation
In order to describe short-ranged two-particle interactions, it is necessary to introduce the standard approximation of the many-body theory, the T-matrix approximation . In the following we give a brief compilation of important formulas.
The sum of ladder diagrams is defined as the causal T-matrix
$$12|𝒯|1^{}2^{}=𝒱(121^{}2^{})+𝒱(1233^{})G(3\overline{1})G(3^{}\overline{2})\overline{1}\overline{2}|𝒯|1^{}2^{}.$$
(A1)
Arguments not present on the left hand side are integrated over. For simplicity, we do not assume the sum over spin and isospin and the anti-symmetrization for identical particles. Since the potential is local in time, we can simplify the T-matrix as $`12|𝒯|1^{}2^{}=x_1x_2t_1|𝒯|x_1^{}x_2^{}t_1^{}\delta (t_1t_2)\delta (t_1^{}t_2^{})`$, and ladder summation (A1) reads
$`x_1x_2t|𝒯|x_1^{}x_2^{}t^{}=𝒱(x_1x_2x_1^{}x_2^{})\delta (tt^{})`$ (A2)
$`+𝒱(x_1x_2x_3x_3^{})x_3x_3^{}t|𝒢|\overline{x}_1\overline{x}_2\overline{t}\overline{x}_1\overline{x}_2\overline{t}|𝒯|x_1^{}x_2^{}t^{},`$ (A3)
where we have introduced the two-particle Green function,
$$x_1x_2t|𝒢|\overline{x}_1\overline{x}_2\overline{t}=G(x_1t\overline{x}_1\overline{t})G(x_2t\overline{x}_2\overline{t}).$$
(A4)
The causal selfenergy then reads
$$\mathrm{\Sigma }(11^{})=i𝑑x_2𝑑x_2^{}G(x_2^{}t_1^{}x_2t_1^+)x_1x_2t_1|𝒯|x_1^{}x_2^{}t_1^{}.$$
(A5)
Using the Langreth-Wilkins rules (which are equivalent to the algebra on the Keldysh contour) we obtain the real-time Green’s functions. From (A5) follows
$`\mathrm{\Sigma }^{}`$ $`=`$ $`𝒯^{}G^{}`$ (A6)
$`\mathrm{\Sigma }^{R/A}`$ $`=`$ $`i\mathrm{\Theta }(t_1t_2)[𝒯^>G^<+𝒯^<G^>]`$ (A7)
$`=`$ $`𝒯^RG^<𝒯^<G^A.`$ (A8)
We have abbreviated the notation, all integrations and variables are the same as in (A5). From (A3) follows
$`𝒯^{}`$ $`=`$ $`𝒯^R𝒢^{}𝒯^A`$ (A9)
$`𝒯^{R/A}`$ $`=`$ $`𝒱+𝒱𝒢^{R/A}𝒯^{R/A}.`$ (A10)
Explicitly (A9) reads
$`<x_1x_2t|𝒯_{ab}^{}|x_1^{}x_2^{}t^{}>`$ (A11)
$`={\displaystyle 𝑑\overline{x}_1𝑑\overline{x}_1^{}𝑑\overline{x}_2𝑑\overline{x}_2^{}𝑑\overline{t}𝑑\overline{t}^{}<x_1x_2t|𝒯_{ab}^R|\overline{x}_1\overline{x}_2\overline{t}>}`$ (A12)
$`\times <\overline{x}_1\overline{x}_2\overline{t}|𝒢_{ab}^{}|\overline{x}_1^{}\overline{x}_2^{}\overline{t}^{}><\overline{x}_1^{}\overline{x}_2^{}\overline{t}^{}|𝒯_{ab}^A|x_1^{}x_2^{}t^{}>.`$ (A13)
In the scattering integral we use the following representation of the two-particle T-matrix,
$`x_1x_2\overline{t}|𝒯|x_1^{}x_2^{}t^{}={\displaystyle \frac{dkdpdqd\mathrm{\Omega }}{(2\pi )^{10}}𝒯(\mathrm{\Omega },k,p,q,r,t)}`$ (A14)
$`\times \mathrm{exp}\left[i\left(kx_1+px_2(kq)x_1^{}(p+q)x_2^{}\mathrm{\Omega }(\overline{t}t^{})\right)\right],`$ (A15)
(A16)
where $`r=\frac{1}{4}(x_1+x_2+x_1^{}+x_2^{})`$ and $`t=\frac{1}{2}(\overline{t}+t^{})`$ are center-of-the-mass coordinate and time.
### 1 Quasiparticle approximation
We introduce the notation for the known quasiparticle approximation without gradient corrections. In the quasiparticle approximation the free two-particle propagator $`𝒢`$ (denoted in the mixed Wigner representation as $`\zeta `$) takes the form
$$\zeta ^R(\mathrm{\Omega },p,p^{},r,t)=\frac{1f(p,r,t)f(p^{},r,t)}{\mathrm{\Omega }\epsilon (p,r,t)\epsilon (p^{},r,t)+i\eta }.$$
(A17)
The selfenergy then reads
$$\sigma ^>(\omega ,k,r,t)=\frac{dp}{(2\pi )^3}t^>(\omega +\epsilon (p,r,t),k,p,0)f(p,r,t),$$
(A18)
and analogously $`\sigma ^<`$. In the rest of this chapter we suppress variables $`r,t`$. From (A10) one finds
$`t^<(\mathrm{\Omega },k,p,0)`$ $`=`$ $`{\displaystyle \frac{dq}{(2\pi )^3}|t^R|^2(\epsilon _{kq}+\epsilon _{p+q},k,p,q)}`$ (A19)
$`\times `$ $`f_{kq}f_{p+q}2\pi \delta (\omega \epsilon _{p+q}\epsilon _{kq}),`$ (A20)
and $`t^>`$ is given by interchange $`f1f`$.
For the determination of the quasiparticle energy one needs the retarded selfenergy
$`\sigma ^R(\omega ,k)`$ $`=`$ $`{\displaystyle \frac{dp}{(2\pi )^3}\frac{d\omega ^{}}{2\pi }\frac{t^>(\omega ^{}+\epsilon _p,k,p,0)f_p}{\omega ^{}\omega i\eta }}`$ (A22)
$`{\displaystyle \frac{dp}{(2\pi )^3}\frac{d\omega ^{}}{2\pi }\frac{t^<(\omega ^{}+\epsilon _p,k,p,0)(1f_p)}{\omega ^{}\omega i\eta }}`$
$`=`$ $`{\displaystyle \frac{dp}{(2\pi )^3}t^R(\omega +\epsilon _p,k,p,0)f_p}`$ (A23)
$``$ $`{\displaystyle \frac{dpdq}{(2\pi )^6}\frac{|t^R(\epsilon _{kq}+\epsilon _{p+q},k,p,q)|^2f_{kq}f_{p+q}}{\epsilon _{kq}+\epsilon _{p+q}\epsilon _p\omega +i\eta }}.`$ (A24)
In the real part of the selfenergy the real part of the T-matrix appears,
$`\mathrm{Re}t^R(\mathrm{\Omega },k,p,0)`$ $`=`$ $`{\displaystyle \frac{dq}{(2\pi )^3}\frac{1f_{kq}f_{p+q}}{\mathrm{\Omega }ϵ_{kq}ϵ_{p+q}}}`$ (A26)
$`\times `$ $`|t^R|^2(ϵ_{kq}+ϵ_{p+q},k,p,q).`$ (A27)
This form follows from the Kramers-Kronig transformation of $`t^>t^<`$.
The ladder equation without gradients reads \[notation is specified by (B5)\]
$`p_1|t^R(q^{})|p_2=p_1|𝒱(q^{})|p_2+{\displaystyle \frac{dp^{}dp^{\prime \prime }}{(2\pi )^6}}`$ (A28)
$`\times p_1|𝒱(q^{})|p^{}p^{}|\zeta ^R(q^{})|p^{\prime \prime }p^{\prime \prime }|t_{sc}^R(q^{})|p_2.`$ (A29)
In the coordinates, $`k=\frac{q^{}}{2}+p_1,p=\frac{q^{}}{2}p_1,q=p_1p_2`$, corresponding to (A16) this equation takes the form
$`t^R(\mathrm{\Omega },k,p,q)=𝒱(k,p,q)`$ (A30)
$`+{\displaystyle \frac{dp^{}dp^{\prime \prime }}{(2\pi )^6}𝒱(k,p,p^{})\zeta ^R(\mathrm{\Omega },kp^{},p+p^{},p^{}p^{\prime \prime })}`$ (A31)
$`\times t^R(\mathrm{\Omega },kp^{\prime \prime },p+p^{\prime \prime },q+p^{\prime \prime }).`$ (A32)
Since $`\zeta ^{}(k,p,q)=(2\pi )^3\delta (q)g^{}(k)g^{}(p)`$, in the quasiparticle approximation one obtains
$`t^R(\mathrm{\Omega },k,p,q)=𝒱(k,p,q)+{\displaystyle \frac{dq^{}}{(2\pi )^3}𝒱(k,p,q^{})}`$ (A34)
$`\times `$ $`{\displaystyle \frac{1f_{kq^{}}f_{p+q^{}}}{\mathrm{\Omega }\epsilon _{p+q^{}}\epsilon _{kq^{}}+i\eta }}t^R(\mathrm{\Omega },kq^{},p+q^{},q+q^{}).`$ (A35)
## B Gradient expansion
Let us consider a product of two-particle functions
$$𝒞(1234)=d3^{}d4^{}𝒜(123^{}4^{})(3^{}4^{}34).$$
(B1)
We transform this product into mixed representation (A16) keeping gradients till the linear order.
Variables in (A16) follow variables of the scattering integral. For the scattering-in, the momenta correspond to the process $`kq,p+qk,p`$. These variables, however, do not form a convenient algebra for the gradient expansion. To this end we use the coordinates
$`\alpha `$ $`=`$ $`12`$ (B2)
$`\beta `$ $`=`$ $`34`$ (B3)
$`\tau `$ $`=`$ $`{\displaystyle \frac{1}{2}}(1+234)`$ (B4)
$`x`$ $`=`$ $`{\displaystyle \frac{1}{4}}(1+2+3+4)(r,t).`$ (B5)
Since the relative coordinates, $`\alpha `$ and $`\beta `$, obey the standard matrix algebra while only the coordinate $`\tau `$ and $`x`$ undergo the gradient expansion, we write these arguments in form $`\alpha |C(\tau ,x)|\beta `$. In this notation the product (B1) reads
$`\alpha |C(\tau ,x)|\beta `$ $`=`$ $`{\displaystyle 𝑑\gamma 𝑑\overline{\tau }\alpha |A(\tau \overline{\tau },x+\frac{1}{2}\overline{\tau })|\gamma }`$ (B6)
$`\times `$ $`\gamma |B(\overline{\tau },x+{\displaystyle \frac{1}{2}}(\tau \overline{\tau }))|\beta .`$ (B7)
Assuming a slow variation of the center-of-mass coordinate $`x`$ we can write this product keeping gradients in $`x`$ up to the linear order. The relative coordinates can be suppressed being independent of this expansion. Via the Fourier transformation of coordinate $`\tau `$,
$$C(\tau ,x)=\frac{d\kappa }{(2\pi )^4}c(\kappa ,x)\mathrm{e}^{i_{n=1}^3\kappa _n\tau _ni\kappa _4\tau _4},$$
(B8)
we express the operator $`𝒞`$ as a function of the sum momentum, $`(\kappa _1,\kappa _2,\kappa _3)=k+p`$, and energy $`\kappa _4=\mathrm{\Omega }`$ of both particles. One can see that the product has the same form as products of single-particle functions,
$$c=ab\frac{i}{2}[a,b]=ab\left(1\frac{i}{2}[\mathrm{ln}a,\mathrm{ln}b]\right),$$
(B9)
where the (generalized) Poisson bracket applies to the sum coordinate and time. The logarithmic form of (B9) is merely a convenient notation because $`a`$ and $`b`$ are matrices in $`\alpha `$ and $`\beta `$. The logarithm is taken of each matrix element of $`a`$ and $`b`$ and summed together with the product $`ab`$.
In parallel with the Dyson equation, there are no linear gradients in the ladder equation (A10). Like the retarded Green function, the T-matrix is given by the matrix inversion, $`𝒯_R^1=𝒱^1𝒢^R`$. Since the T-matrix is symmetric with respect to matrix arguments, $`s|𝒯|s^{}=s^{}|𝒯|s`$ ($`s`$ and $`s^{}`$ are momenta associated with $`\alpha `$ and $`\beta `$, respectively), the matrix inversion does not bring any gradients.
The selfenergy, $`\mathrm{\Sigma }^<(13)=𝒯^<(1234)G^>(42)`$, includes a number of gradient contributions due to internal gradients in $`𝒯^<`$ and a simple non-local correction due the convolution with $`G^>`$.
### 1 Two-particle matrix products
For the calculation of the selfenergy in (B44) we need $`t^<`$ for the zero transferred momentum, $`q=0`$, and its infinitesimal vicinity. Making the gradient expansion of (A9) in terms of the Poisson brackets one finds
$`t^<`$ $`=`$ $`t^R\stackrel{~}{g}^<t^A(1{\displaystyle \frac{i}{2}}([\mathrm{ln}t^R,\mathrm{ln}\stackrel{~}{g}^<][\mathrm{ln}t^A,\mathrm{ln}\stackrel{~}{g}^<]`$ (B11)
$`+[\mathrm{ln}t^R,\mathrm{ln}t^A]))`$
$`=`$ $`|t|^2\stackrel{~}{g}^<\left(1[\varphi ,\mathrm{ln}\stackrel{~}{g}^<][\varphi ,\mathrm{ln}|t|]\right).`$ (B12)
In the second line we have decomposed the T-matrices into amplitudes and phase shifts, $`t^{R,A}=|t|\mathrm{e}^{i\varphi }`$. Note that all gradient corrections depend exclusively on the derivatives of phase shift $`\varphi `$. Using definitions (124-130) and the linear approximation, $`a(x)(1\mathrm{\Delta }_xa)=a(x\mathrm{\Delta })`$, expression (B12) can be given a form
$`t^<(\mathrm{\Omega },k,p,0,r{\displaystyle \frac{1}{2}}\mathrm{\Delta }_2,t)={\displaystyle \frac{dq}{(2\pi )^3}\frac{dQ}{(2\pi )^3}}`$ (B13)
$`\times t^R(\mathrm{\Omega }\mathrm{\Delta }_\omega ,k{\displaystyle \frac{\mathrm{\Delta }_k}{2}},p{\displaystyle \frac{\mathrm{\Delta }_k}{2}},q,r\mathrm{\Delta }_r,t{\displaystyle \frac{\mathrm{\Delta }_t}{2}})`$ (B14)
$`\times \zeta ^<(\mathrm{\Omega }2\mathrm{\Delta }_\omega ,kq\mathrm{\Delta }_k,p+q\mathrm{\Delta }_k,Q,r{\displaystyle \frac{\mathrm{\Delta }_3}{2}}{\displaystyle \frac{\mathrm{\Delta }_4}{2}},t\mathrm{\Delta }_t)`$ (B15)
$`\times t^A(\mathrm{\Omega }\mathrm{\Delta }_\omega ,k{\displaystyle \frac{\mathrm{\Delta }_k}{2}},p{\displaystyle \frac{\mathrm{\Delta }_k}{2}},q+Q,r\mathrm{\Delta }_r,t{\displaystyle \frac{\mathrm{\Delta }_t}{2}}).`$ (B16)
(B17)
Factors of half in the non-local corrections to the amplitude $`|t^{R,A}|`$ result from the fact that this amplitude enters the collision integral in the square.
In the absence of gradients, the two-particle function $`\zeta ^<`$ is singular in the ‘transferred’ momentum $`Q`$ giving the only contribution for $`Q=0`$. When the gradient corrections are already in the explicit form, one can employ this symmetry as it has been used above. Indeed, the complex conjugacy of the advanced and retarded T-matrices, $`t^A=\overline{t}^R`$, requires equal arguments of both functions. Now we show that the infinitesimal vicinity of $`Q=0`$ brings additional gradient contributions.
### 2 Convolution of initial states
The two-particle correlation function $`\zeta ^<`$ representing initial states of the collision in (B17) is not a suitable input for the kinetic equation. We have to express $`\zeta ^<`$ in terms of the convolution of two single-particle functions. This convolution brings gradient corrections by which effective positions of particles entering the collision process become distinct.
¿From the inverse transformation to (A16) we obtain
$`\zeta ^<(\mathrm{\Omega },kp,p+q,Q,r,t)`$ (B18)
$`=2^3{\displaystyle 𝑑x_1𝑑x_2𝑑x_1^{}𝑑x_2^{}𝑑\overline{t}𝑑t^{}\mathrm{exp}[iQ(x_1^{}x_2^{})+i\mathrm{\Omega }(\overline{t}t^{})]}`$ (B19)
$`\times \mathrm{exp}\left[i(kq)(x_1x_1^{})i(p+q)(x_2x_2^{})\right]`$ (B20)
$`\times \delta (x_1+x_2+x_1^{}+x_2^{}4r)\delta (\overline{t}+t^{}2t)`$ (B21)
$`\times G^<(x_1\overline{t},x_1^{}t^{})G^<(x_2\overline{t},x_2^{}t^{}).`$ (B22)
The substitution, $`x_1=r+\alpha +\beta _1/2`$, $`x_1^{}=r+\alpha \beta _1/2`$, $`x_2=r\stackrel{~}{\alpha }+\beta _2/2`$, $`x_2^{}=r\stackrel{~}{\alpha }\beta _2/2`$, $`\overline{t}=t+\tau /2`$, $`t^{}=t\stackrel{~}{\tau }/2`$, shows that the $`\delta `$-functions yield $`\stackrel{~}{\tau }=\tau `$ and $`\stackrel{~}{\alpha }=\alpha `$ and we obtain
$`\zeta ^<(\mathrm{\Omega },kp,p+q,Q,r,t)=2^3{\displaystyle 𝑑\beta _1𝑑\beta _2𝑑\alpha 𝑑\tau \mathrm{exp}[2iQ\alpha ]}`$ (B23)
$`\times \mathrm{exp}\left[i(kq+Q/2)\beta _1i(p+qQ/2)\beta _2+i\mathrm{\Omega }\tau \right]`$ (B24)
$`\times \widehat{g}^<(\beta _1,r+\alpha ,t,\tau )\widehat{g}^<(\beta _2,r\alpha ,t,\tau ),`$ (B25)
where the representation (21) has been used.
Now we linearize the center-of-mass dependence of $`g^<`$’s in $`\alpha `$. The factor $`\alpha `$ can be represented by a derivative in front of the integral, $`\alpha =\frac{i}{2}_Q\frac{i}{4}_q`$. The remaining integration over $`\alpha `$ results in $`\pi ^3\delta (Q)`$. When substituted into (B17), the differential operator in front of the integral is treated with the help of the integration by parts giving
$`\alpha {\displaystyle \frac{1}{2}}(\mathrm{\Delta }_3\mathrm{\Delta }_4)`$ (B26)
and $`Q`$ is integrated out. In other words the $`\alpha `$ in the arguments of $`\widehat{g}`$ in (B25) can be replaced by nonlocal shifts valid up to linear orders in gradients. This leads with (B17) to
$`t^<(\mathrm{\Omega },k,p,0,r{\displaystyle \frac{1}{2}}\mathrm{\Delta }_2,t)={\displaystyle \frac{dq}{(2\pi )^3}}`$ (B27)
$`\times \left|t^R(\mathrm{\Omega }\mathrm{\Delta }_\omega ,k{\displaystyle \frac{\mathrm{\Delta }_k}{2}},p{\displaystyle \frac{\mathrm{\Delta }_k}{2}},q,r\mathrm{\Delta }_r,t{\displaystyle \frac{\mathrm{\Delta }_t}{2}})\right|^2`$ (B28)
$`\times g^<(kq,r\mathrm{\Delta }_3,t\mathrm{\Delta }_t,\mathrm{\Omega }2\mathrm{\Delta }_\omega )`$ (B29)
$`\times g^<(p+q,r\mathrm{\Delta }_4,t\mathrm{\Delta }_t,\mathrm{\Omega }2\mathrm{\Delta }_\omega ).`$ (B30)
(B31)
### 3 Convolution of T-matrix and hole Green function
Next we evaluate the convolution of the $`𝒯`$-matrix with the hole Green function which is required for the selfenergy. Writing the selfenergy in the above representation, \[$`\tau =13`$ and $`x=\frac{1}{2}(1+3)`$\]
$`\mathrm{\Sigma }^<(\tau ,x)`$ $`=`$ $`{\displaystyle 𝑑\alpha 𝑑\beta G^>(\tau \beta +\alpha ,x\frac{1}{2}(\alpha +\beta ))}`$ (B32)
$`\times `$ $`\alpha |𝒯^<(\tau {\displaystyle \frac{1}{2}}(\alpha \beta ),x{\displaystyle \frac{1}{4}}(\alpha +\beta ))|\beta ,`$ (B33)
we see that the convolution couples matrix arguments $`\alpha `$ and $`\beta `$ with the center-of-mass variables. By substitution, $`\lambda =\frac{1}{2}(\alpha +\beta )`$ and $`\mu =\alpha \beta `$, and expansion in gradients we obtain
$`\mathrm{\Sigma }^<(\tau ,x)`$ $`={\displaystyle 𝑑\mu G^>(\mu \tau ,x)}`$ (B40)
$`\times {\displaystyle }d\lambda \lambda +{\displaystyle \frac{\mu }{2}}\left|𝒯^<(\tau {\displaystyle \frac{\mu }{2}},x)\right|\lambda {\displaystyle \frac{\mu }{2}}`$
$`{\displaystyle \frac{1}{2}}{\displaystyle 𝑑\mu G^>(\mu \tau ,x)}`$
$`\times {\displaystyle \frac{}{x}}{\displaystyle 𝑑\lambda \lambda \lambda +\frac{\mu }{2}\left|𝒯^<(\tau \frac{\mu }{2},x)\right|\lambda \frac{\mu }{2}}`$
$`{\displaystyle 𝑑\mu \left(\frac{}{x}G^>(\mu \tau ,x)\right)}`$
$`\times {\displaystyle }d\lambda \lambda \lambda +{\displaystyle \frac{\mu }{2}}\left|𝒯^<(\tau {\displaystyle \frac{\mu }{2}},x)\right|\lambda {\displaystyle \frac{\mu }{2}}.`$
The second and third terms are gradient corrections due to the convolution.
Now we can transform the selfenergy (LABEL:selalp) into the mixed representation. The momentum representation of the matrix algebra is introduced via unity operators, e.g., $`1=\frac{ds}{(2\pi )^3}|ss|`$, with $`\lambda \mu /2|s=\mathrm{e}^{is(\lambda \mu /2)}`$ and $`s^{}|\lambda +\mu /2=\mathrm{e}^{is^{}(\lambda +\mu /2)}`$. For the non-gradient part the integration over $`\lambda `$ results in the known fact that only the diagonal element, $`s=s^{}`$, contributes. For the gradient contribution, $`\lambda (i/2)(_s_s^{})`$, and the diagonal element is taken after the derivatives are performed. Since we have considered the gradient corrections to the $`𝒯`$-matrix already in the last chapter and since the shifts are additive, it is sufficient here to use the zero-order approximation of $`t^<`$,
$$s|t^<|s^{}=\frac{d\overline{s}d\stackrel{~}{s}}{(2\pi )^2}s|t^R|\overline{s}\overline{s}|\zeta ^<|\stackrel{~}{s}\stackrel{~}{s}|t^A|s^{},$$
(B42)
where all functions have the center-of-mass argument $`(\kappa ,x)`$. From $`t^{R,A}=|t|\mathrm{e}^{i\varphi }`$ and condition $`s=s^{}`$ we find that $`\lambda _s\varphi `$. By the substitution into variables of the kinetic equation, $`(q,0)=ss^{}`$, $`(k,\omega )=\kappa /2+s`$, $`(p,0)=\kappa /2s`$ and $`(r,t)=x`$, one confirms that $`_s\varphi =\frac{1}{2}\mathrm{\Delta }_2`$, see (124), and the selfenergy reads
$`\sigma _\omega ^<(k,r,t)`$ $`=`$ $`{\displaystyle \frac{dpd\mathrm{\Omega }}{(2\pi )^4}g_{\mathrm{\Omega }\omega }^>(p,r\mathrm{\Delta }_2,t)\left(1\frac{1}{2}\frac{\mathrm{\Delta }_2}{r}\right)}`$ (B43)
$`\times `$ $`t^<(\mathrm{\Omega },k,p,0,r{\displaystyle \frac{1}{2}}\mathrm{\Delta }_2,t).`$ (B44)
In this expression, $`t^<`$ abbreviates the right hand side of (B42), because the non-local correction $`\mathrm{\Delta }_2`$ is defined only with respect to the integral over internal states of the collision. The norm term, $`1\frac{1}{2}_r\mathrm{\Delta }_2`$ results from the interchange, $`_r(\mathrm{\Delta }_2A)=(_r\mathrm{\Delta }_2)A+\mathrm{\Delta }_2_rA`$, one has to make before the derivative is expressed in terms of the displacement, e.g., $`g^<(r)\mathrm{\Delta }_2_rg^<(r)=g^<(r\mathrm{\Delta }_2)`$.
Together with (B31) we obtain the result (123).
## C Common met pitfall with the correlated density
We note that one of us made a mistake deriving the correlated density directly from the retardation in the Levinson-type equation . The expression found in is identical to the correlated density except for two points, there is a wrong sign and the amplitude is as twice as large. Since this is a pitfall often met in literature we want to give some details here that it can be avoided in future.
The simple expansion of the Levinson equation (4) up to the first order in memory yields
$$\frac{}{t}\rho =+\frac{}{t}.$$
(C1)
To obtain balance equations for the density including virial corrections, one has to integrate (C1) over momentum $`k`$, see . Following the arguments developed for the Boltzmann-type equations, in , the wrong assumption was used that the correlated part of the collision integral, $`\stackrel{~}{n}_c=𝑑kR`$, combines with the left hand side of (C1), $`\stackrel{~}{n}_f=𝑑k\rho `$, to establish the density conservation, $`\frac{}{t}(\stackrel{~}{n}_f+\stackrel{~}{n}_c)=0`$. The correlated density derived in this way reminds the result known from equilibrium . Note that in this picture, the Wigner distribution in left hand side of the Levinson equation is treated as the quasiparticle distribution what is a mistake done usually in density operator studies . The wrong sign follows from this last misinterpretation. By definition the momentum integral over the Wigner distribution yields the full density, $`𝑑k\rho =n`$. The conservation law thus tells that either $`𝑑kR=0`$ or the Levinson equation does not conserve the number of particles.
As we have shown in this paper the Levinson equation contains additional gradient terms which exactly compensate the explicit time gradients (used to derive $`\stackrel{~}{n}_f`$). The above double count follows from the double count of the off-shell contributions discussed in this paper.
Finally, the $`\tau _{R,A}`$ retardation was neglected. Since the correlated density (155) found from the off-shell contribution to the Wigner distribution is consistent with the conservation of the number of particles found from the $`\tau _{R,A}`$ retardation, one can see that these two contributions cancel and the collision integral of the Levinson equation conserves the number of particles, as it should.
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# MEASUREMENT OF HIGGS PROPERTIES AT THE LHC*footnote **footnote *Talk given at the Rencontres de Moriond, QCD and High Energy Hadronic Interactions, Les Arcs, France, March 18-25, 2000.
## I Introduction
One of the prime tasks of the LHC will be to probe the mechanism of electroweak gauge symmetry breaking. Beyond observation of the various CP even and CP odd scalars which nature may have in store for us, this means the determination of the couplings of the Higgs boson to the known fermions and gauge bosons, i.e. the measurement of $`Htt`$, $`Hbb`$, $`H\tau \tau `$ and $`HWW`$, $`HZZ`$, $`H\gamma \gamma `$ couplings, to the extent possible.
Clearly this task very much depends on the expected Higgs boson mass. For $`m_H>200`$ GeV and within the SM, only the $`HZZ`$ and $`HWW`$ channels are expected to be observable, and the two gauge boson modes are related by SU(2). A much richer spectrum of decay modes is predicted for the intermediate mass range, i.e. if a SM-like Higgs boson has a mass between the reach of LEP2 ($`\stackrel{<}{}110`$ GeV) and the $`Z`$-pair threshold. The main reasons for focusing on this range are present indications from electroweak precision data, which favor $`m_H<250`$ GeV, as well as expectations within the MSSM, which predicts the lightest Higgs boson to have a mass $`m_h\stackrel{<}{}130`$ GeV. Recently, an analysis of Higgs coupling measurements at the LHC was completed for this intermediate mass range and in this talk I sumarize the results.
## II Survey of intermediate mass Higgs channels
The total production cross section for a SM Higgs boson at the LHC is dominated by the gluon fusion process, $`ggH`$, which largely proceeds via a top-quark loop. Thus, inclusive Higgs searches will collectively be called “gluon fusion” channels in the following. Three inclusive channels are highly promising for the SM Higgs boson search,
$`ggH\gamma \gamma ,\mathrm{for}m_H\stackrel{<}{}150\mathrm{GeV},`$ (1)
$`ggHZZ^{}4\mathrm{},\mathrm{for}m_H\stackrel{>}{}130\mathrm{GeV},`$ (2)
and
$$ggHWW^{}\mathrm{}\overline{\nu }\overline{\mathrm{}}\nu ,\mathrm{for}m_H\stackrel{>}{}130\mathrm{GeV}.$$
(3)
The $`H\gamma \gamma `$ signal can be observed as a narrow and high statistics $`\gamma \gamma `$ invariant mass peak, albeit on a very large diphoton background. A few tens of $`HZZ^{}4\mathrm{}`$ events are expected to be visible in 100 fb<sup>-1</sup> of data, with excellent signal to background ratios (S/B), ranging between 1:1 and 6:1, in a narrow four-lepton invariant mass peak. Finally, the $`HWW^{}\mathrm{}\overline{\nu }\overline{\mathrm{}}\nu `$ mode is visible as a broad enhancement of event rate in a 4-lepton transverse mass distribution, with S/B between 1:4 and 1:1 (for favorable values of the Higgs mass, around 170 GeV).
Additional, and, as we shall see, crucial information on the Higgs boson can be obtained by isolating Higgs production in weak boson fusion (WBF), i.e. by separately observing $`qqqqH`$ and crossing related processes, in which the Higgs is radiated off a $`t`$-channel $`W`$ or $`Z`$. Specifically, it was recently shown in parton level analyses that the weak boson fusion channels, with subsequent Higgs decay into photon pairs,
$$qqqqH,H\gamma \gamma ,\mathrm{for}m_H\stackrel{<}{}150\mathrm{GeV},$$
(4)
into $`\tau ^+\tau ^{}`$ pairs,
$$qqqqH,H\tau \tau ,\mathrm{for}m_H\stackrel{<}{}140\mathrm{GeV},$$
(5)
or into $`W`$ pairs
$$qqqqH,HWW^{()}e^\pm \mu ^{}/p_T,\mathrm{for}m_H\stackrel{>}{}120\mathrm{GeV},$$
(6)
can be isolated at the LHC. The weak boson fusion channels utilize the significant background reductions which are expected from double forward jet tagging and central jet vetoing techniques, and promise low background environments in which Higgs decays can be studied in detail.
An example of expected events rates (after cuts and including efficiency factors) are summarized in Table I for the $`qqqqH,HWW^{()}e^\pm \mu ^{}/p_T`$ signal. The rates and ensuing statistical errors of the signal cross section are given for 100 fb<sup>-1</sup> of data collcted in both the ATLAS and the CMS detector. The statistical accuracy with which the signal cross sections of the processes in Eqs. (1-6) can be determined is shown in Fig. 1a).
## III Measurement of Higgs properties
In order to translate the cross section measurements of the various Higgs production and decay channels into measurements of Higgs boson properties, in particular into measurements of the various Higgs boson couplings to gauge fields and fermions, it is convenient to rewrite them in terms of partial widths of various Higgs boson decay channels. The Higgs-fermion couplings $`g_{Hff}`$, for example, which in the SM are given by the fermion masses, $`g_{Hff}=m_f(m_H)/v`$, can be traded for $`\mathrm{\Gamma }_f=\mathrm{\Gamma }(H\overline{f}f)`$. Similarly the square of the $`HWW`$ coupling ($`g_{HWW}=gm_W`$ in the SM) or the $`HZZ`$ coupling is proportional to the partial widths $`\mathrm{\Gamma }_W=\mathrm{\Gamma }(HWW^{})`$ or $`\mathrm{\Gamma }_Z=\mathrm{\Gamma }(HZZ^{})`$. $`\mathrm{\Gamma }_\gamma =\mathrm{\Gamma }(H\gamma \gamma )`$ and $`\mathrm{\Gamma }_g=\mathrm{\Gamma }(Hgg)`$ determine the squares of the effective $`H\gamma \gamma `$ and $`Hgg`$ couplings. The Higgs production cross sections are governed by the same squares of couplings, hence, $`\sigma (VVH)\mathrm{\Gamma }_V`$ (for $`V=g,W,Z`$). Combined with the branching fractions $`B(Hii)=\mathrm{\Gamma }_i/\mathrm{\Gamma }`$ the various signal cross sections measure different combinations of Higgs boson partial and total widths, $`\mathrm{\Gamma }_i\mathrm{\Gamma }_j/\mathrm{\Gamma }`$.
The production rate for WBF is a mixture of $`ZZH`$ and $`WWH`$ processes, and we cannot distinguish between the two experimentally. In a large class of models the ratio of $`HWW`$ and $`HZZ`$ couplings is identical to the one in the SM, however, and this includes the MSSM. We therefore assume that 1) the $`HZZ^{}`$ and $`HWW^{}`$ partial widths are related by SU(2) as in the SM, i.e. their ratio, $`z`$, is given by the SM value, $`z=\mathrm{\Gamma }_Z/\mathrm{\Gamma }_W=z_{SM}`$. Note that this assumption can be tested, at the 15-20% level for $`m_H>130`$ GeV, by forming the ratio $`B\sigma (ggHZZ^{})/B\sigma (ggHWW^{})`$.
With $`W,Z`$-universality, the three weak boson fusion cross sections give us direct measurements of three combinations of (partial) widths,
$`X_\gamma ={\displaystyle \frac{\mathrm{\Gamma }_W\mathrm{\Gamma }_\gamma }{\mathrm{\Gamma }}}`$ $`\mathrm{from}`$ $`qqqqH,H\gamma \gamma ,`$ (7)
$`X_\tau ={\displaystyle \frac{\mathrm{\Gamma }_W\mathrm{\Gamma }_\tau }{\mathrm{\Gamma }}}`$ $`\mathrm{from}`$ $`qqqqH,H\tau \tau ,`$ (8)
$`X_W={\displaystyle \frac{\mathrm{\Gamma }_W^2}{\mathrm{\Gamma }}}`$ $`\mathrm{from}`$ $`qqqqH,HWW^{()},`$ (9)
In addition the three gluon fusion channels provide measurements of
$`Y_\gamma ={\displaystyle \frac{\mathrm{\Gamma }_g\mathrm{\Gamma }_\gamma }{\mathrm{\Gamma }}}`$ $`\mathrm{from}`$ $`ggH\gamma \gamma ,`$ (10)
$`Y_Z={\displaystyle \frac{\mathrm{\Gamma }_g\mathrm{\Gamma }_Z}{\mathrm{\Gamma }}}`$ $`\mathrm{from}`$ $`ggHZZ^{()},`$ (11)
$`Y_W={\displaystyle \frac{\mathrm{\Gamma }_g\mathrm{\Gamma }_W}{\mathrm{\Gamma }}}`$ $`\mathrm{from}`$ $`ggHWW^{()}.`$ (12)
When extracting Higgs couplings, the QCD uncertainties of production cross sections enter. These can be estimated via the residual scale dependence of the NLO predictions and are small (of order 5%) for the WBF case, while larger uncertainties of about 20% are found for gluon fusion.
A first test of the Higgs sector is provided by taking ratios of the $`X_i`$’s and ratios of the $`Y_i`$’s. QCD uncertainties, and all other uncertainties related to the initial state, like luminosity and pdf errors, cancel in these ratios. They test $`W,Z`$-universality, and compare the $`\tau \tau H`$ Yukawa coupling with the $`HWW`$ coupling. Typical errors on these cross section ratios are expected to be in the 15 to 20% range. Accepting an additional systematic error of about 20%, a measurement of the ratio $`\mathrm{\Gamma }_g/\mathrm{\Gamma }_W`$, which determines the $`Htt`$ to $`HWW`$ coupling ratio, can be performed, by measuring the cross section ratios $`B\sigma (ggH\gamma \gamma )/\sigma (qqqqH)B(H\gamma \gamma )`$ and $`B\sigma (ggHWW^{})/\sigma (qqqqH)B(HWW^{})`$.
Beyond the measurement of coupling ratios, minimal additional assumptions allow an indirect measurement of the total Higgs width. First of all, the $`\tau `$ partial width, properly normalized, is measurable with an accuracy of order 10%. The $`\tau `$ is a third generation fermion with isospin $`\frac{1}{2}`$, just like the $`b`$-quark. In many models, the ratio of their coupling to the Higgs is given by the $`\tau `$ to $`b`$ mass ratio. In addition to $`W,Z`$-universality we thus assume that (2) $`y=\mathrm{\Gamma }_b/\mathrm{\Gamma }_\tau =y_{SM}`$ and, finally, (3) the branching ratio for unexpected channels is small, i.e. $`ϵ=1(B(Hb\overline{b})+B(H\tau \tau )+B(HWW^{()})+B(HZZ^{()})+B(Hgg)+B(H\gamma \gamma )1`$.
With these three assumptions consider the observable
$`\stackrel{~}{\mathrm{\Gamma }}_W`$ $`=`$ $`X_\tau (1+y)+X_W(1+z)+X_\gamma +Y_W`$ (13)
$`=`$ $`\left(\mathrm{\Gamma }_\tau +\mathrm{\Gamma }_b+\mathrm{\Gamma }_W+\mathrm{\Gamma }_Z+\mathrm{\Gamma }_\gamma +\mathrm{\Gamma }_g\right){\displaystyle \frac{\mathrm{\Gamma }_W}{\mathrm{\Gamma }}}=(1ϵ)\mathrm{\Gamma }_W.`$ (14)
$`\stackrel{~}{\mathrm{\Gamma }}_W`$ provides a lower bound on $`\mathrm{\Gamma }(HWW^{()})=\mathrm{\Gamma }_W`$. Provided $`ϵ`$ is small (within the SM and for $`m_H>110`$ GeV, $`ϵ<0.04`$ and it is dominated by $`B(Hc\overline{c})`$), the determination of $`\stackrel{~}{\mathrm{\Gamma }}_W`$ provides a direct measurement of the $`HWW^{()}`$ partial width. Once $`\mathrm{\Gamma }_W`$ has been determined, the total width of the Higgs boson is given by
$$\mathrm{\Gamma }=\frac{\mathrm{\Gamma }_W^2}{X_W}=\frac{1}{X_W}\left(X_\tau (1+y)+X_W(1+z)+X_\gamma +\stackrel{~}{X}_g\right)^2\frac{1}{(1ϵ)^2}.$$
(15)
The extraction of the total Higgs width, via Eq. (15), requires a measurement of the $`qqqqH,HWW^{()}`$ cross section, which is expected to be available for $`m_H\stackrel{>}{}115`$ GeV. Consequently, errors are large for Higgs masses close to this lower limit but decrease to about 10% for Higgs boson masses around the $`WW`$ threshold. Results are shown in Fig. 1b) and look highly promising.
## IV Summary
With an integrated luminosity of 100 fb<sup>-1</sup> per experiment, the LHC can measure various ratios of Higgs partial widths, with accuracies of order 10 to 20%. This translates into 5 to 10% measurements of various ratios of coupling constants. The ratio $`\mathrm{\Gamma }_\tau /\mathrm{\Gamma }_W`$ measures the coupling of down-type fermions relative to the Higgs couplings to gauge bosons. To the extent that the $`H\gamma \gamma `$ triangle diagrams are dominated by the $`W`$ loop, the width ratio $`\mathrm{\Gamma }_\tau /\mathrm{\Gamma }_\gamma `$ probes the same relationship. The fermion triangles leading to an effective $`Hgg`$ coupling are expected to be dominated by the top-quark, thus, $`\mathrm{\Gamma }_g/\mathrm{\Gamma }_W`$ probes the coupling of up-type fermions relative to the $`HWW`$ coupling. Finally, for Higgs boson masses above $`120`$ GeV, the absolute normalization of the $`HWW`$ coupling is accessible via the extraction of the $`HWW^{()}`$ partial width in weak boson fusion.
These measurements test the crucial aspects of the Higgs sector. The $`HWW`$ coupling, being linear in the Higgs field, identifies the observed Higgs boson as the scalar responsible for the spontaneous breaking of $`SU(2)\times U(1)`$: a scalar without a vacuum expectation value does not exhibit such a trilinear coupling at tree level. The measurement of the ratios of $`g_{Htt}/g_{HWW}`$ and $`g_{H\tau \tau }/g_{HWW}`$ then probes the mass generation of both up and down type fermions. Thus the LHC can do much more than mereley discover the Higgs: it can give us detailed and reasonably precise information on the dynamics of electroweak symmetry breaking.
## Acknowledgments
I would like to thank R. Kinnunen, A. Nikitenko and E. Richter-Was for a most enjoyable collaboration leading to the results summarized here, and the CERN theory group for its hospitality. This work was supported in part by the University of Wisconsin Research Committee with funds granted by the Wisconsin Alumni Research Foundation and in part by the U. S. Department of Energy under Contract No. DE-FG02-95ER40896.
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# Equation of state and critical behavior of polymer models: A quantitative comparison between Wertheim’s thermodynamic perturbation theory and computer simulations.
## I introduction
For a long time, the prediction of the equation of state of polymers from first principles was such a difficult task that it could not be solved without the need of very idealized models with more or less meaningful empirical parameters. Perhaps the most successful of the early attempts was that of Prigogine et al. who considered a lattice and introduced an ad-hoc empirical parameter known as the “number of degrees of freedom per monomer”. Later on, Flory et al. extended this theory to the continuum at the cost of introducing some additional parameters, leading to what is now known as the FOVE theory. Another approach to the problem is based on the polymer+solvent and polymer+vacuum analogy. In this way, the well known Flory-Huggins (FH) and Sanchez-Lacombe equations of state have also been employed to describe the behavior of pure fluids.
On the other hand, the approach from liquid state theory has taken much longer to yield useful results but has now reached a point where a rather satisfactory description of the equation of state of idealized polymer models in the continuum is affordable without the need of any empirical parameters whatsoever. The most popular approaches are the Polymer Reference Interaction Site Model (PRISM) of Curro and Schweizer, the Generalized Dimer Flory theory (GDF) of Honnell and Hall and the Thermodynamic Perturbation Theory of Wertheim (TPT1). The latter has been widely used not only because it yields results that are of similar or superior quality than other alternatives, but because it is the simplest and more tractable of all them and demands a minimum of information.
Originally, this theory was developed to consider fluids of associating hard spheres. Later on, it was realized that in the limit of infinite associating strength, a polydisperse mixture of polymers was recovered. Chapman et al. extended the theory to monodisperse polymer fluids and rewrote the equations in a very convenient notation. Later, by adding a mean field perturbative contribution, Jackson et al. showed that the theory was able to describe the behavior of very different sorts of real fluids. Meanwhile, it was realized that the theory could be used just as well to consider polymers made of attractive monomers. In this way, it is possible to describe real fluids with improved accuracy and no need for a mean field–like perturbation. Since then, the theory has achieved enormous popularity and has been applied to describe polymers of tangent Lennard-Jones beads and square wells, as well as to describe real fluids in chemical engineering applications. Furthermore, modifications of the original theory have been proposed that allow to describe realistic polymer models without the need of empirical parameters.
A very interesting issue both from the practical and theoretical point of view is the behavior of the critical point of polymer fluids as the number of monomers increases. By invoking the polymer+solvent and polymer+vacuum analogy, one would expect from the FH theory that the pure polymer fluid should reach an asymptotic critical temperature, whereas the critical mass density should become vanishingly small. However, in a recent paper, Chatterjee and Schweizer have pointed out that this analogy cannot be taken for granted because the FH scaling predictions are determined by imposing equal chemical potentials, whereas the critical point of pure polymer fluids is related to a phase equilibria that results from the condition of equal pressure at a given temperature.
On the other hand, based on rather limited amount of data, several empirical correlations have been employed to predict the critical properties of substances such as polymethylene. These correlations have predicted widely different behavior, ranging from infinite critical temperature to finite asymptotic critical mass density. More soundly based equations such as the FOVE have been recently employed to support the idea that the critical mass density could reach an asymptotic constant value. However, such an approach relies on extending the applicability of the FOVE theory to densities below which it was not meant to be used.
Renormalization group calculations, however, show that the FH mean field theory yields the correct asymptotic dependence in the long chain length limit, albeit this asymptotic behavior is only reached for extremely long chains and the corrections to scaling are enormous even for typical chain length considered experimentally.
More recently, new experimental techniques have allowed to measure the critical points of longer n-alkanes whose critical temperature lays above the point where thermal decomposition starts. These experiments show that the critical mass density reaches a maximum and then starts to decrease. Simulation results of both realistic alkanes and idealized polymer models give support to this finding.
Surprisingly, very little attention has been devoted to the study of this problem from the point of view of the modern theories such as PRISM, GFD or TPT1. In a recent paper, the PRISM was employed in an attempt to solve the question in the framework of an analytical tractable theory, leading to different predictions depending on the closure employed to solve the PRISM equation. Other previous studies using TPT1 plus a mean field attractive contribution suggested that the critical mass density should vanish in the infinite chain length limit. However, such a conclusion relied on the assumption that the mean field contribution increases linearly with the chain length, a point which at present cannot be taken for granted. Nikitin et al. have recently presented a theory that may be considered as the simplest possible approximation to Wertheim’s theory and arrive to the same conclusion as reference .
In this paper we will try to reach further understanding of this issue. In the next section, we will review the fundamentals of Wertheim’s theory by using somewhat different but more physically appealing arguments suggested by Zhou and Stell. We will then show by means of general arguments how the scaling laws for the critical properties are closely related to the virial coefficients (section III). We will then apply TPT1 to a polymer model and briefly describe how to implement the theory in section IV and in the Appendixes. Section V describes the details of the simulations performed to test the theory. We then devote section VI to present our results and close with a brief conclusion.
## II Equation of state of polymers
### A Preliminary definitions
We consider a monodisperse fluid of polymers made up of $`n`$ monomers each. Non bonded monomers of the same polymer and monomers belonging to different molecules interact through some pair potential, $`u_0`$, while adjacent monomers of the same molecule have an additional bonding potential, $`\mathrm{\Phi }`$, responsible for the connectivity of the polymer. This potential is such that its action vanishes beyond some well defined inter-atomic distance. If eventually, two adjacent monomers were found outside this range, they would no longer be bonded and the n-mer would be considered to have broken. In practice, this can be avoided by making the well depth of the associating potential infinitely large.
Alternatively, we may consider an associating multicomponent mixture of $`n`$ different monomeric species, $`A`$, $`B`$, $`C`$, etc. Each of these species interacts with members of its own species and with members of the remaining species by means of $`u_0`$. Furthermore, the bonding potential $`\mathrm{\Phi }`$ is responsible for associating reactions between monomers of type $`A`$ with monomers of type $`B`$, monomers of type $`B`$ with monomers of type $`C`$ and so on. More concisely, the association reaction taking place is of the form:
$$A+B+C+\mathrm{}ABC\mathrm{}$$
(1)
As the bonding potential is pairwise, the only requirement that is needed to define the ’$`ABC\mathrm{}`$’ complex as an n-mer is that each of the adjacent pairs be found within the range of the bonding potential. Clearly, for an equimolar composition of such a mixture in the limit of complete association, we are lead to a system identical to that described in the preceding paragraph. From a physical point of view, this limit is reached for infinite well depths. However, in what follows it will prove useful to consider that the well depth has some arbitrary finite value. The system is then made of a mixture of free monomers and n-mers whose composition depends on the nature of the bonding potential.
Considering the similarity between the two systems described, we may get an approximation for the equation of state for the chain molecule by studying the behavior of the associating system. In order to do so, we will obtain an expression for the free energy in terms of the degree of association of the mixture. We will then relate the degree of association to the structure of the system and introduce some simplifying assumptions for the n-body correlations. Finally, we will take the limit of complete association and get an equation of state for the chain fluid.
### B The association reaction
Let us consider an associating system as that described before, which is initially prepared by mixing in equal proportions the pure monomers, so that there are $`N`$ monomers of each species inside a volume $`V`$ and the resulting number density of each of the species is $`\rho `$. The system will eventually reach a state of equilibrium, whereby a fraction of the monomers of each species has associated to form n-mers. Let this fraction be $`\alpha `$. Then, the remaining concentration of non bonded monomers of each species is given by $`\rho (1\alpha )`$, while, according to the stoichiometry of the reaction, the concentration of n-mers will be $`\rho \alpha `$ (note that in the limit of complete association $`\rho `$ will actually designate the polymer number density). Schematically, the process can be described as follows:
| | | $`A`$ | $`+`$ | $`B`$ | + | $`C`$ | $`\mathrm{}`$ | $``$ | $`ABC\mathrm{}`$ |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| initially | | $`\rho `$ | | $`\rho `$ | | $`\rho `$ | | | $`0`$ |
| at equilibrium | | $`\rho (1\alpha )`$ | | $`\rho (1\alpha )`$ | | $`\rho (1\alpha )`$ | | | $`\rho \alpha `$ |
(2)
Obviously, the number of n-mers formed is $`V\rho \alpha =N\alpha `$, while the number of remaining free monomers of each species is $`N(1\alpha )`$. The total Helmholtz free energy of the system, $`A=GpV`$ is therefore given by:
$$A(\alpha )=\underset{i=1}{\overset{n}{}}N(1\alpha )\mu _i(\alpha )+N\alpha \mu _{\mathrm{n}\mathrm{mer}}(\alpha )p(\alpha )V$$
(3)
In what follows, we associate a number from $`1`$ to $`n`$ to each of the n species. Therefore, $`\mu _i`$ stands for the chemical potential of species $`i`$ in the mixture of composition given by $`\alpha `$.
It is important to realize that the composition of the mixture as given by the association degree is appropriate to a given model potential. In the above expression for the free energy, we have considered a system of monomers that interact through some reference potential, $`u_0`$, while a bonding potential is responsible for the connectivity of the n-mer. Let us now consider a reference fluid made of monomers that interact through the reference potential but that have no association potential whatsoever. In this case, the free energy is best expressed in terms of the chemical potential of the monomers alone. For the sake of simplicity, we will denote this free energy as $`A(\alpha =0)`$, albeit from a strictly geometrical point of view, n-mers could form eventually even in the absence of a bonding potential:
$$A(\alpha =0)=\underset{i=1}{\overset{n}{}}N\mu _i(\alpha =0)p(\alpha =0)V$$
(4)
The free energy of the associating system measured relative to that of the reference fluid is then given by:
$$\frac{\mathrm{\Delta }A}{N}=\underset{i=1}{\overset{n}{}}[\mu _i(\alpha )\mu _i(\alpha =0)]+\alpha [\mu _{\mathrm{n}\mathrm{mer}}(\alpha )\underset{i=1}{\overset{n}{}}\mu _i(\alpha )][p(\alpha )p(\alpha =0)]\frac{V}{N}$$
(5)
This equation can be further simplified by invoking the condition of chemical equilibrium of the associating mixture, which reads:
$$\underset{i=1}{\overset{n}{}}\mu _i(\alpha )\mu _{\mathrm{n}\mathrm{mer}}(\alpha )=0$$
(6)
Substitution of the above equation into Eq. 5 yields:
$$\frac{\mathrm{\Delta }A}{N}=\underset{i=1}{\overset{n}{}}[\mu _i(\alpha )\mu _i(\alpha =0)][p(\alpha )p(\alpha =0)]\frac{V}{N}$$
(7)
Let us now express the chemical potentials of each of the components in terms of an ideal and an excess contribution:
$$\mu _i(\alpha )=\mu _i^{\mathrm{id}}(\alpha )+\mu _i^{\mathrm{ex}}(\alpha )=\mu _i^0+kT\mathrm{ln}\rho _i(\alpha )+\mu _i^{\mathrm{ex}}(\alpha )$$
(8)
Introducing the above expression for the chemical potential into Eq. 7, we get:
$$\frac{\mathrm{\Delta }A}{N}=kT\mathrm{ln}\underset{i=1}{\overset{n}{}}\frac{\rho _i(\alpha )}{\rho _i(\alpha =0)}+\underset{i=1}{\overset{n}{}}[\mu _i^{\mathrm{ex}}(\alpha )\mu _i^{\mathrm{ex}}(\alpha =0)][p(\alpha )p(\alpha =0)]\frac{V}{N}$$
(9)
If we now recall that $`\rho _i(\alpha )=\rho (1\alpha )`$ and $`\rho _i(\alpha =0)=\rho `$, we finally obtain:
$$\frac{\mathrm{\Delta }A}{N}=nkT\mathrm{ln}(1\alpha )+\underset{i=1}{\overset{n}{}}[\mu _i^{\mathrm{ex}}(\alpha )\mu _i^{\mathrm{ex}}(\alpha =0)][p(\alpha )p(\alpha =0)]\frac{V}{N}$$
(10)
This is an exact equation for the difference in free energy between the associating system (with $`\alpha N`$ n-mers) and the reference non-associating fluid of free monomers. As it stands, it may seem rather useless, because it is a function of the unknown quantities: $`\alpha `$, $`\mu _i^{\mathrm{ex}}(\alpha )`$ and $`p(\alpha )`$. However, we shall see that in the low density limit, the above equation becomes a function of the degree of association only and that this, in turn, may be obtained from knowledge of n-body correlations in the fluid. In a further approximation, the n-body correlations of the fluid will be expressed in terms of n-body correlations of a reference fluid of non–bonded monomers. Finally, by invoking a superposition approximation, we will express the n-body correlations in terms of two body correlations and obtain an expression for the pressure that depends solely on known quantities of a reference fluid of spherical particles.
### C The association reaction in the limit of low density
In the limit of low density, the equation of state of the associating system will depend only on the total number of particles in the fluid, $`N(\alpha )`$:
$$\frac{pV}{kT}=N(\alpha )$$
(11)
$`N(\alpha )`$ is simply obtained by summing the number of particles of each species:
$$N(\alpha )=\underset{i=1}{\overset{n}{}}N(1\alpha )+N\alpha =nN(1\alpha )+N\alpha $$
(12)
On the other hand, the number of particles of the completely un-associated system is simply $`N(\alpha =0)=nN`$. Therefore, the difference in pressure between the non–bonded system and the system with degree of association $`\alpha `$ is:
$$\mathrm{\Delta }p(\alpha )\frac{V}{N}=kT\alpha (n1)$$
(13)
By using this expression and considering that, by definition, the excess chemical potentials vanish in the limit of low density, Eq. 10 becomes:
$$\frac{\mathrm{\Delta }A}{kTN}=n\mathrm{ln}(1\alpha )+\alpha (n1)$$
(14)
This is an exact equation for the difference in free energy in the limit of low density. Applying the standard thermodynamic relationship connecting the pressure with the free energy yields the exact expression for the difference in pressure in the limit of low density:
$$\frac{\mathrm{\Delta }p}{kT\rho }=\rho \frac{\alpha (1n)1}{1\alpha }\frac{\alpha }{\rho }$$
(15)
In what follows, we will consider this expression to be valid in all the density range. The search of an equation of state for the associating system will be thus accomplished if we find an expression for the association degree in terms of known properties. Once this relation has been found, we will obtain an equation of state for the system of n-mers by taking the limit of complete association, i.e., $`\alpha =1`$.
### D Relation between the degree of association and the structure of the fluid
#### 1 Expression for the degree of association in terms of the excess chemical potential of the components
First, consider the equilibrium constant of the reaction, defined as the ratio of the concentration of the products to that of the reactants:
$$K=\frac{\alpha \rho }{(1\alpha )^n\rho ^n}$$
(16)
The connection of the equilibrium constant to the thermodynamics of the process may be obtained by expressing the chemical potential of each of the components as in Eq. 8 and substituting into the condition of chemical equilibrium, Eq. 6. After some simple algebraic manipulations, we are lead to the following expression for the equilibrium constant:
$$kT\mathrm{ln}K+\mu _{\mathrm{n}\mathrm{mer}}^0\underset{i=1}{\overset{n}{}}\mu _i^0+\mu _{\mathrm{n}\mathrm{mer}}^{\mathrm{ex}}(\alpha )\underset{i=1}{\overset{n}{}}\mu _i^{\mathrm{ex}}(\alpha )=0$$
(17)
This expression may be further simplified by considering that, in the limit of low densities, the excess chemical potentials vanish. As a consequence of this, the low density equilibrium constant, $`K_0`$, is given as follows:
$$kT\mathrm{ln}K_0=\underset{i=1}{\overset{n}{}}\mu _i^0\mu _{\mathrm{n}\mathrm{mer}}^0$$
(18)
Substitution of this expression into Eq. 17, leads finally to a simple equation for the equilibrium constant in terms of the excess chemical potential of the components of the mixture:
$$kT\mathrm{ln}\frac{K}{K_0}=\underset{i=1}{\overset{n}{}}\mu _i^{\mathrm{ex}}(\alpha )\mu _{\mathrm{n}\mathrm{mer}}^{\mathrm{ex}}(\alpha )$$
(19)
#### 2 Expression for the structure of the fluid in terms of the excess chemical potential of the components
In order to relate the structure of the fluid to the excess chemical potential of the components of the mixture, (i.e., $`A`$, $`B`$, $`C`$, etc. monomers and n-mers), let us consider the thermodynamic cycle of figure 1.
In the first step of the cycle, an isolated n-mer is dissolved into a fluid mixture with association degree $`\alpha `$. Initially, the total Gibbs free energy of the system is the sum of the free energy of the isolated n-mer, $`G_{\mathrm{n}\mathrm{mer}}`$ and the free energy of the mixture, $`G_{\mathrm{mix}}(\alpha )`$. After dissolving the n-mer, the resulting free energy is that of the original mixture with an extra n-mer, $`G_{\mathrm{mix}+\mathrm{n}\mathrm{mer}}`$. The change in $`G`$ is therefore:
$$\mathrm{\Delta }G_1=G_{\mathrm{mix}+\mathrm{n}\mathrm{mer}}G_{\mathrm{mix}}G_{\mathrm{n}\mathrm{mer}}$$
(20)
In the thermodynamic limit, the difference $`G_{\mathrm{mix}+\mathrm{n}\mathrm{mer}}G_{\mathrm{mix}}`$ becomes equal to the chemical potential of the n-mer in the mixture, while $`G_{\mathrm{n}\mathrm{mer}}`$ may be considered to be the chemical potential of the isolated n-mer (i.e., the free energy difference between a system with a single n-mer and an empty system). It is thus seen that:
$$\mathrm{\Delta }G_1=\mu _{\mathrm{n}\mathrm{mer}}^{\mathrm{ex}}(\alpha )$$
(21)
In a second step, $`n`$ uncorrelated monomers of type $`A`$, $`B`$, $`C`$, etc. dissolved in the mixture are brought together in such a way that the resulting $`ABC\mathrm{}`$ complex forms one of the many possible conformers of the n-mer, say, one such that the vector joining $`B`$ to $`A`$ is $`𝐫_{12}`$, that joining $`C`$ to $`B`$ is $`𝐫_{23}`$, etc. This event will occur according to a probability density given by the n-body correlation function of the mixture, $`\rho ^ng(𝐫_{12},𝐫_{23},\mathrm{},𝐫_{n1,n})`$. The n-mer density is given as an integral of this function over all the conformations compatible with the monomer:
$$\rho _{\mathrm{n}\mathrm{mer}}=\rho ^n_v\mathrm{}_vg(𝐫_{12},𝐫_{23},\mathrm{},𝐫_{n1,n})\mathrm{d}^3𝐫_{12}\mathrm{d}^3𝐫_{23}\mathrm{}\mathrm{d}^3𝐫_{n1,n}$$
(22)
where $`v`$ is the volume within the range of the bonding potential. i.e., any two adjacent monomers whose distance vector is not within this volume are not considered to be bonded.
The process of forming the n-mer from a set of n uncorrelated monomers may be considered as a chemical reaction of the form
| n uncorrelated monomers | $``$ | n correlated monomers |
| --- | --- | --- |
Friedman has shown that such an equation is characterized by an equilibrium constant of the form $`K_{\mathrm{eq}}=\rho _{\mathrm{n}\mathrm{mer}}/\rho ^n`$. Substitution of Eq. 22 into the expression for the equilibrium constant, yields $`K_{\mathrm{eq}}=\rho ^n\mathrm{\Delta }/\rho ^n`$, where,
$$\mathrm{\Delta }=_v\mathrm{}_vg(𝐫_{12},𝐫_{23},\mathrm{},𝐫_{n1,n})\mathrm{d}^3𝐫_{12}\mathrm{d}^3𝐫_{23}\mathrm{}\mathrm{d}^3𝐫_{n1,n}$$
(23)
Now, the free energy of the whole process may be written down as:
$$\mathrm{\Delta }G=\mathrm{\Delta }G_2+kT\mathrm{ln}K_{\mathrm{eq}}$$
(24)
However, when the system reaches equilibrium, $`\mathrm{\Delta }G=0`$, so that we finally obtain:
$$\mathrm{\Delta }G_2=kT\mathrm{ln}\mathrm{\Delta }$$
(25)
Similar arguments as those put through for the first and second steps of the cycle, lead to the conclusion that
$$\mathrm{\Delta }G_3=\underset{i=1}{\overset{n}{}}\mu _i^{\mathrm{ex}}(\alpha )$$
(26)
$$\mathrm{\Delta }G_4=kT\mathrm{ln}\mathrm{\Delta }_0$$
(27)
where $`\mathrm{\Delta }_0`$ is the integral of Eq. 23 evaluated at zero density.
Substitution of Eq. 21, 25, 26, 27 into the net energy balance of the cycle, $`\mathrm{\Delta }G_1+\mathrm{\Delta }G_4=\mathrm{\Delta }G_3+\mathrm{\Delta }G_2`$, leads to the desired equation relating the structure of the fluid with the excess chemical potential of the components of the mixture:
$$kT\mathrm{ln}\frac{\mathrm{\Delta }}{\mathrm{\Delta }_0}=\underset{i=1}{\overset{n}{}}\mu _i^{\mathrm{ex}}(\alpha )\mu _{\mathrm{n}\mathrm{mer}}^{\mathrm{ex}}(\alpha )$$
(28)
A similar equation has been derived in a rather more formal way recently for the particular case of infinitely short ranged association potentials. The derivation we have employed is based on the thermodynamic cycle presented by Zhou and Stell , which allows to extend their previous result to finite range potentials.
#### 3 Expression for the degree of association in terms of the structure of the fluid
Substitution of Eq. 19 into Eq. 28 shows that the equilibrium constant is related to the n-body correlation function of the associating mixture through the following relation:
$$\frac{K}{K_0}=\frac{\mathrm{\Delta }}{\mathrm{\Delta }_0}$$
(29)
Finally, using the expression for $`K`$ in terms of the degree of association, Eq. 16, we obtain the desired equation relating the degree of association with the structure of the system:
$$\frac{\alpha }{(1\alpha )^n\rho ^{n1}}=\frac{K_0}{\mathrm{\Delta }_0}\mathrm{\Delta }$$
(30)
### E The equation of state
Previously, we obtained an approximate equation for the pressure in terms of the association degree of the mixture, Eq. 15. The density derivative of the association degree, required in such an expression is obtained from Eq. 30:
$$\frac{\alpha }{\rho }=\frac{1\alpha }{\rho }\frac{(n1)\alpha +\rho \alpha \frac{\mathrm{ln}\mathrm{\Delta }}{\rho }}{1+\alpha (n1)}$$
(31)
Substitution of this result into Eq. 15 yields an expression for the change in pressure due to the formation of the n-mers from a fluid of non-associated monomers. This expression depends solely on the degree of association and the n-body correlation function of the associating system:
$$\frac{\mathrm{\Delta }p}{kT\rho }=\alpha (n1+\rho \frac{\mathrm{ln}\mathrm{\Delta }}{\rho })$$
(32)
By taking the limit of infinite association, which physically corresponds to infinitely increasing the well depth of the bonding potential, we would arrive at an equation for the pressure of the n-mer fluid relative to that of the monomer reference fluid. In order to do so, however, we would require the n-body correlation function of the associating system, which enters through $`\mathrm{\Delta }`$. Unfortunately, quantitative understanding of such high order correlation functions is far beyond our present knowledge. We will therefore need to make some further approximations in order to get a tractable expression for the pressure.
#### 1 Decoupling of the n-body correlations
In order to simplify the problem of the n-body correlations, we invoke a so called ’linear’ decoupling approximation, which attempts to describe the n-body correlation function in terms of $`n1`$ two body correlation functions:
$$g^{(n)}(𝐫_{12},𝐫_{23},\mathrm{},𝐫_{n1,n})=g^{(2)}(𝐫_{12})g^{(2)}(𝐫_{23})\mathrm{}g^{(2)}(𝐫_{n1,n})$$
(33)
Here it should be understood that $`g^{(2)}(𝐫_{12})`$ stands for the pair correlation function of monomers of type $`A`$ with monomers of type $`B`$ in the multicomponent mixture of associating monomers, $`g^{(2)}(𝐫_{23})`$ stands for the pair correlation function of monomers of type $`B`$ with monomers of type $`C`$ and so on. Still, these two body correlation functions may be quite difficult to obtain. In order to simplify the problem, consider one of these correlation functions, say, $`g^{(2)}(𝐫_{12})`$, in the limit of zero density:
$$g^{(2)}(𝐫_{12};\rho =0)=\mathrm{exp}([u_0(𝐫_{12})+\mathrm{\Phi }(𝐫_{12})]/kT)$$
(34)
It can be seen that, in this limit, the pair correlation function may be exactly expressed in terms of the pair correlation function of a reference system with no bonding potential, $`g_0^{(2)}(𝐫_{12})`$, times the Boltzmann factor of the bonding potential:
$$g^{(2)}(𝐫_{12};\rho =0)=g_0^{(2)}(𝐫_{12};\rho =0)\mathrm{exp}(\mathrm{\Phi }(𝐫_{12})/kT)$$
(35)
Using the linear approximation to the n-body correlations and considering the above equation to hold true at any density, the $`\mathrm{\Delta }`$ integral is simplified considerably, giving:
$$\mathrm{\Delta }=\delta ^{n1}$$
(36)
where $`\delta `$ is defined as:
$$\delta =_vg_0^{(2)}(𝐫_{12})\mathrm{exp}(\mathrm{\Phi }(𝐫_{12})/kT)\mathrm{d}^3𝐫_{12}$$
(37)
#### 2 An equation of state for the n-mer in terms of the thermodynamics and structure of the monomers
As a consequence of the two approximations given for the n-body correlations, we are now able to write down an expression for the pressure of the n-mer in terms of the properties of the reference fluid. Indeed, after setting $`\alpha =1`$, substitution of Eq. 36 into Eq. 32 leads finally to the following result:
$$\frac{\mathrm{\Delta }p}{kT\rho }=[n1][1+\rho \frac{\mathrm{ln}\delta }{\rho }]$$
(38)
Note that as we are considering the limit of complete association, $`\rho `$ is equal to the polymer number density.
Adding the contribution of the reference system to the previous equation, we are now able to write down an equation for the compressibility factor, $`Z=p/kT\rho `$ of the fluid of n-mers:
$$Z_{\mathrm{n}\mathrm{mer}}=nZ_0(n1)[1+\rho \frac{\mathrm{ln}\delta }{\rho }]$$
(39)
where $`Z_0`$ is the compressibility factor of the reference fluid, measured at the same monomer density as the n-mer fluid. This is a rather remarkable equation, as it gives the equation of state of a chain fluid from the properties of a fluid of monomers alone. Different versions of this equation will arise from the different theories available to describe the structure and thermodynamics of the fluid of monomers. In section IV we will consider two such theories in order to describe our model polymer. Let us recall at this point, however, that a simple, qualitative version of Eq.39 may be obtained by simply considering that $`\delta `$ does not depend on the density. In this way, the resulting equation does no longer depend on the structure of the reference fluid. Nikitin et al. have explored this equation using the van der Waals equation of state for $`Z_0`$ and find the same qualitative behavior as is found in this work.
### F Comparison with Wertheim’s theory of association
The arguments we have put through in order to arrive at Eq. 39 are rather physical and intuitive. On the other hand, Wertheim has developed a very general theory of association based on a re-summed cluster expansion, where the significance of each of the approximations is mathematically well understood. It is interesting to compare the results of this rather formal theory with the physically appealing description that we have used, largely based on the work of Zhou and Stell .
In the extension of Wertheim’s theory of association, the compressibility factor of the chain molecule is given as:
$$Z_{\mathrm{n}\mathrm{mer}}=nZ_0(n1)[1+\rho \frac{\mathrm{ln}\kappa }{\rho }]$$
(40)
where $`\kappa `$ is defined as:
$$\kappa =_vg_0^{(2)}(𝐫_{12})[\mathrm{exp}(\mathrm{\Phi }(𝐫_{12})/kT)1]\mathrm{d}^3𝐫_{12}$$
(41)
The limit of complete association requires that $`\mathrm{\Phi }`$ have an infinite well depth, so, within most of the range of the bonding potential, the Boltzmann factor is exceedingly bigger than $`1`$. Therefore, $`\kappa `$ and $`\delta `$ become identical and Eq. 40 is essentially equal to Eq. 39.
Before proceeding to the next section, let us first summarize the approximations invoked to obtain Eq. 39:
1. Assume that the free energy difference between the reference system and the completely associated system takes the form of the low density limit in all the density range.
2. Decouple the n-body correlation function of the associating system into $`n1`$ two body correlation functions through a linear approximation.
3. Assume that the two body correlation function of the associating system is given in terms of the two body correlation function of a reference system with no bonding potential, as suggested by the exact low density limit.
## III Predictions for the scaling laws of the critical properties
We start by assuming that the critical density does become small for large chain lengths, so that one can describe the equation of state in terms of a truncated virial expansion.
$$\frac{p}{kT}=\rho +B_2(T)\rho ^2+B_3(T)\rho ^3$$
(42)
where $`\rho `$ is the polymer number density. By applying the conditions for the critical point of pure fluids, i.e.,
$$\begin{array}{c}\left(\frac{p}{V}\right)_{T_c}=0\\ \left(\frac{^2p}{V^2}\right)_{T_c}=0\end{array}$$
(43)
we obtain a set of equations for the critical temperature and density:
$`B_2(T_c)+\sqrt{3B_3(T_c)}`$ $`=`$ $`0`$ (44)
$`\sqrt{3B_3(T_c)}\rho _c`$ $`=`$ $`1`$ (45)
By making a Taylor expansion on powers of the density, the first and second virial coefficients predicted by Wertheim’s equation are found to be:
$$\begin{array}{c}B_2=n^2\left(b_2\frac{n1}{n}a_2\right)\\ B_3=n^3\left(b_3\frac{n1}{n}a_3\right)\end{array}$$
(46)
where $`b_2`$ and $`b_3`$ are the second and third virial coefficients of the reference fluid of non–bonded monomers, while $`a_2`$ and $`a_3`$ are the zeroth and first order coefficients in a monomer density expansion of $`\mathrm{ln}\delta /\rho `$. Of course, all these quantities are chain length independent.
Now, in order to solve Eq. 44 for the critical temperature we will need to linearize the virial coefficients with respect to the temperature. To do so, let us assume for the time being that there is a finite asymptotic critical temperature in the limit of infinite chain length, which we call $`\mathrm{\Theta }`$, in analogy with the polymer + solvent case. We now make a series expansion of $`B_2`$ and $`B_3`$ in powers of $`\mathrm{\Delta }T=\mathrm{\Theta }T`$ up to first order, and consider the limit of this expression for large $`n`$, leading to
$$\begin{array}{c}B_2(T)=n^2\left(C_2C_2^{}\mathrm{\Delta }T\right)\hfill \\ B_3(T)=n^3\left(C_3C_3^{}\mathrm{\Delta }T\right)\hfill \end{array}$$
(47)
where $`C_2=b_2(\mathrm{\Theta })a_2(\mathrm{\Theta })`$ and $`C_3=b_3(\mathrm{\Theta })a_3(\mathrm{\Theta })`$ while $`C_2^{}`$ and $`C_3^{}`$ are the corresponding derivatives with respect to temperature. Substitution of the linearized virial coefficients into the condition for the critical temperature leads to a quadratic equation for $`\mathrm{\Delta }T`$. Solving for this equation yields $`\mathrm{\Delta }T_c(n)`$, defined as $`\mathrm{\Theta }T_c(n)`$:
$$\mathrm{\Delta }T_c(n)\frac{C_2}{C_2^{}}=\pm \frac{1}{2C_2^{}_{}{}^{}2}\left(12C_2^{}_{}{}^{}2C_312C_2C_2^{}C_3^{}+9C_3^{}_{}{}^{}2\frac{1}{n}\right)^{1/2}\frac{1}{n^{1/2}}\frac{3C_3^{}}{2C_2^{}_{}{}^{}2}\frac{1}{n}$$
(48)
This equation shows that $`\mathrm{\Delta }T_c(n)`$ must reach an asymptotic finite value, since the right hand side term should ultimately vanish for large $`n`$. The requirement for $`T_c(n)`$ to attain a finite asymptotic critical temperature equal to $`\mathrm{\Theta }`$ is then obeyed provided that $`C_2`$ vanishes. If we now notice that
$$\underset{n\mathrm{}}{lim}B_2(\mathrm{\Theta })=n^2C_2$$
(49)
we arrive at the conclusion that indeed $`C_2`$ must vanish at the Boyle temperature of the infinitely long polymer, $`T_B^{\mathrm{}}`$, thus identifying the $`\mathrm{\Theta }`$ temperature with the Boyle temperature of the infinitely long polymer. From the definition of $`C_2`$ we see that this temperature is attained when the following condition is obeyed:
$$b_2(T_B^{\mathrm{}})a_2(T_B^{\mathrm{}})=0$$
(50)
Note also that the leading terms of the expansion (Eq. 48) are of order $`n^{1/2}`$ and $`n^1`$, just as predicted by the FH theory.
The case of the critical polymer density is much simpler. Substitution of the expression for $`B_3`$ in the condition for the critical density shows that:
$$\rho _c(n)n^{3/2}$$
(51)
so that the critical mass density decreases with a power law proportional to $`n^{1/2}`$, as predicted by the FH theory.
It is important to note that the above arguments apply regardless of the specific form in which the reference fluid (thermodynamics and structure) is described. In particular, the simplest implementation of TPT1, proposed by Nikitin et al. considers $`\delta `$ to be a constant. In such a case, $`a_2`$ is zero at all temperatures. However, Eq. 50 shows that this simple version still predicts an asymptotic critical temperature which must obey the condition $`b_2(T_B^{\mathrm{}})=0`$. Obviously, this condition is obeyed for the Boyle temperature of the reference fluid.
Another interesting issue is the apparent universality of the compressibility factor as predicted by the truncated virial expansion of Eq. 42. Indeed, substitution of this equation into the condition for the critical point shows that, apart from Eq. 45 it must also hold that $`\rho _c=B_2/(3B_3)`$. Using this expression for the critical density in the linear term of Eq. 42 and Eq. 45 in the quadratic term, it is seen that both terms cancel each other exactly. Dividing the resulting expression for the pressure by the critical density (Eq. 45) then shows that:
$$Z_c(n)=\frac{p_c}{\rho _ckT_c}=\frac{1}{3}+\frac{B_4}{(3B_3)^{3/2}}+\mathrm{}$$
(52)
Obviously, this result is independent on whatever assumption is made concerning the actual $`n`$ dependence of the virial coefficients and shows that a finite asymptotic critical compressibility factor of about $`1/3`$ is expected in the limit of infinite chain length, irrespective of the nature of the polymer. In the context of TPT1, a constant compressibility factor implies that the critical pressure must decrease as $`n^{3/2}`$.
## IV Application to a polymer model
Let us consider a polymer model as that described in the previous section, with the reference fluid considered to be a truncated and shifted potential of the form:
$$u_0(r)=\{\begin{array}{cc}V_{LJ}(r)V_{LJ}(r_c)\hfill & rr_c\hfill \\ 0\hfill & r>r_c\hfill \end{array}$$
(53)
where $`r_c=22^{1/6}`$ and $`V_{LJ}`$ is the usual Lennard-Jones potential,
$$V_{LJ}(r)=4ϵ\left\{\left(\frac{\sigma }{r}\right)^{12}\left(\frac{\sigma }{r}\right)^6\right\}$$
(54)
As to the bonding potential responsible for the connectivity between adjacent monomers, we will consider the FENE potential, defined in terms of $`R_0`$, the maximum displacement between monomers and $`k_0`$, a sort of elastic constant:
$$\mathrm{\Phi }(r)=\{\begin{array}{cc}k_0R_0^2\mathrm{ln}(1\frac{r^2}{R_0^2})E_b\hfill & 0<r<R_0\hfill \\ 0\hfill & rR_0\hfill \end{array}$$
(55)
In order to ensure permanent connectivity of the n-mer, a constant, $`E_b`$, which is (conceptually) made infinitely large, is added to the actual FENE potential.
In what follows, we will set $`k_0=15ϵ/\sigma ^2`$ and $`R_0=1.5\sigma `$ and use the Lennard-Jones energy and range parameters as energy and length units, respectively. At liquid–like densities the most probable distance between non–bonded monomers is about $`1.12\sigma `$, which is bigger than the most probable distance $`0.96\sigma `$ between bonded monomers. Note at this point that most of the previous applications of Wertheim’s theory have been restricted to bonding potentials of infinitely short range, allowing for a single possible bond length. To our knowledge, only once has the effect of a soft bonding potential been considered previously.
In section II we related the equation of state of such a polymer fluid to the properties of the reference system of LJ monomers. What is now required is a theory for both the thermodynamics and the structure of the monomer fluid. We have obtained the required input from integral equation theory and thermodynamic perturbation theory. Let us consider each of them in turn.
### A The RHNC integral equation theory
In this approach, one attempts to calculate the exact pair correlation function, which is then used to evaluate the mechanical properties of the fluid. The Ornstein-Zernike equation relates the total pair correlation function, $`h(r)=g(r)1`$ to a short range direct correlation function $`c(r)`$:
$$h(r_{12})=c(r_{12})+\rho h(r_{13})c(r_{23})\mathrm{d}^3𝐫_3$$
(56)
Additionally, this integral equation must be provided with a closure that relates $`c(r)`$ to $`h(r)`$. We use the Reference Hyper-netted chain equation of Lado and Ashcroft. Although this set of equations can only be solved numerically and convergence is not a trivial matter, an efficient algorithm due to Labik and Malijevsky makes the calculations affordable with a modest amount of CPU time. Once $`g(r)`$ is known, the pressure of the fluid may be calculated using the standard relation :
$$\frac{p}{k_BT}=\rho \frac{\rho ^2}{6k_BT}r\frac{\mathrm{d}u_0}{\mathrm{d}r}g(r)\mathrm{\hspace{0.33em}\hspace{0.33em}4}\pi r^2dr$$
(57)
Similarly, $`\delta `$ may be calculated using Eq. 37. Actually, solving the integral equation for each of the desired thermodynamic states may result rather cumbersome. In practice we solve the OZ+RHNC equation for several hundreds of state points and fit the pressure and $`\delta `$. Details of the procedure may be found in Appendix A.
### B The perturbation theory of Tang and Lu
Perturbation theory was the first approach to give quantitative results for the thermodynamics of simple fluids at high density. Compared to integral equation theory, it gives similar results at high densities at a smaller computational cost, with the advantage that the free energy is obtained directly, without the need for thermodynamic integration. On the other hand, the traditional perturbation theories of Barker and Henderson and Weeks-Chandler-Andersen are known to be rather poor at low densities because the underlying assumptions of these theories no longer hold true. Fortunately, Tang and Lu have presented rather recently a second order perturbation theory for Lennard-Jones fluids which is very accurate both at low and high densities. The success of this theory relies on a rather good description of the structure of the fluid, which is obtained from the OZ equation, supplemented by a simple closure known as the Mean Spherical Approximation. The use of this closure is very convenient because it has allowed to obtain a very good approximation to the actual free energy with a purely analytical equation. Furthermore, Tang and Lu have been able to obtain also analytic expressions for the pair correlation function using the MSA closure. With minor modifications we were able to extend this theory to our truncated and shifted Lennard-Jones potential and obtain a lengthy but analytic expression for the thermodynamics and structure of our reference fluid. Details of the implementation are explained in Appendix B.
In what follows, we shall present the results of Wertheim’s theory for our polymer model using both the RHNC and the MSA thermodynamic theories for the reference fluid. We will call each of the versions TPT1-RHNC and TPT1-MSA, following the original name for Eq. 39 due to Wertheim.
## V Simulation details
In order to test the TPT1 theory, we have performed extensive computer simulations. Chain length $`n=10`$ was chosen for a detailed comparison to the theory. We have calculated the pressure and chemical potential for 5 temperatures $`T=1.68,2.5,3.0,4.0`$, and 5. The lowest value corresponds to a subcritical isotherm, while the highest value is above the $`\mathrm{\Theta }`$ temperature. The pressure isotherms were evaluated from the virial in NVT Monte Carlo simulations The length $`L`$ of the (cubic) simulation box was fixed to 18 $`\sigma `$ units. The density dependence of the chemical potential was calculated using grand-canonical Monte Carlo simulations. Chain conformations were sampled using local monomer displacements and slithering snake like motions. Particle insertions and deletions were performed following configurational bias grand-canonical acceptance rules. In order to obtain the equation of state we employ cycles of 25 local moves, 25 reptations and 10 CBGC moves. About 40000 such cycles were performed so that at least a few thousand particle insertion-deletion attempts were accepted. The volume of the simulation box was chosen so that an average number of about 50 chain molecules was obtained.
In order to compare the theory to the simulations, we must make sure that the chemical potentials are expressed with respect to the same reference state. In order to do so, we define the excess chemical potential to be the difference between $`\mu `$ and the chemical potential of an ideal gas of chains $`\mu _{\mathrm{id}}`$ with the full intramolecular interactions but no intermolecular interactions.
$$\mu ^{\mathrm{ex}}\mu \mu _{\mathrm{id}}\text{with}\mu _{\mathrm{id}}=\mathrm{ln}\rho (n1)\mathrm{ln}C\mathrm{ln}W_0$$
(58)
The first term denotes the translational entropy. Contributions due to the integration over the momenta are ignored, because they contribute equally to the reference system and the interacting polymer liquid. To determine the ideal gas contribution we construct chains according to the Rosenbluth procedure. The distance $`l`$ between bonded neighbors is chosen according to its Boltzmann weight $`p(|l|)=4\pi l^2\mathrm{exp}((u_0+\mathrm{\Phi })/kT)/C`$ where $`C`$ is the normalization constant. $`W_0`$ denotes the Rosenbluth weight of the chains due to non–bonded interactions, measured at zero density. Once the excess chemical potential has been obtained, it is compared to the excess chemical potential as predicted by the theory, which is evaluated using the standard thermodynamic relation:
$$\frac{\mu ^{\mathrm{ex}}}{k_BT}=\frac{A^{ex}}{Nk_BT}+Z1$$
(59)
The grand-canonical ensemble allows also for an accurate measurement of the phase diagram, because the order parameter (i.e., the density) is not conserved and density fluctuations are efficiently equilibrated. We monitor the probability distribution $`P(\rho )`$ of the density. Close to two phase coexistence, the probability distribution is bimodal: one peak corresponds to the vapor, the other corresponds to the liquid. The coexistence chemical potential $`\mu _{\mathrm{coex}}`$ is fixed by the condition of equal weight in both peaks:
$$_0^\rho ^{}d\rho P(\rho )\stackrel{!}{=}_\rho ^{}^{\mathrm{}}d\rho P(\rho )\text{with}\rho ^{}=_0^{\mathrm{}}d\rho \rho P(\rho )$$
(60)
Far below the critical points the probability between the two peaks is very low, and we use a re-weighting scheme as to encourage the system to “tunnel” between the two phases. To this end we add a term $`k_BT\mathrm{ln}W(\rho )`$ to the original Hamiltonian. Choosing $`W(\rho )P(\rho )`$ the system visits all densities with roughly equal probability. The probability distribution of the grand-canonical ensemble is obtained via re-weighting the distribution in the simulations $`P_{\mathrm{MC}}`$ according to $`P(\rho )=P_{\mathrm{MC}}(\rho )W(\rho )`$. At very low temperatures the density of the liquid in coexistence with its vapor becomes very high and the configurational bias scheme becomes quite inefficient. Since the density of the vapor is very low, however, its pressure is vanishingly small. Hence, we employed NpT simulations at zero pressure to obtain the liquid density at coexistence.
At the critical point the correlation length of density fluctuations diverges and universal behavior is expected. For finite chain length the unmixing transition exhibits 3D Ising universal behavior. We have located the critical point for chain length $`n=1,10,20,40,`$ and $`60`$ by mapping the symmetrized order parameter distribution $`P_{\mathrm{sym}}(\rho )=[P(\rho )+P(\rho _c\rho )]/2`$ onto the universal scaling function of the 3D Ising model. This symmetrization reduces field mixing corrections which are antisymmetric in $`\rho \rho _c`$ to leading order. Normalizing $`P_{\mathrm{sym}}`$ to unit variance and norm we eliminates all non–universal factors. The results of this mapping are presented in Fig.2, where we have used system sizes $`L=11.3,13.8,18,22.5,`$ and $`27`$ for chain length $`n=1,10,20,40`$, and $`60`$, respectively. This method gives an accurate location of the critical temperature and density (finite size corrections to the critical density of the order $`L^{(1\alpha )/\nu }`$ are neglected). The locations of the critical points are collected in Tab. I.
## VI Results and Discussion
Let us first examine the thermodynamic data for the chains of 10 monomers. Fig.3 shows the predictions of TPT1 for several pressure isotherms ($`kT/ϵ`$=5, 4, 3, 2.5 and 1.68) compared with simulation results. Both the RHNC and MSA versions of the theory are seen to give rather good estimates; at the highest temperatures, far above the estimated $`\mathrm{\Theta }`$ point of our model (see below) as well as at the lowest, a subcritical isotherm. Overall, the RHNC version seems to describe the isotherms slightly better.
Results for the excess chemical potential of the chains are shown in Fig.4. The agreement is also quite satisfactory, though the results are slightly worse than for the pressure isotherms, specially at the lowest temperatures and densities. Indeed, the main assumption of the theory, that the local environment of a monomer in the polymer fluid is similar to that of the monomer fluid breaks down in the low density limit. The fluid is then made of isolated clusters of $`n`$ monomers, rather than of single monomers uniformly distributed in space. Likewise, the theory is unable to describe the density dependence of the single chain internal energy and entropy.
The liquid–vapor coexistence curve of the 10-mer as obtained from simulation and theory is shown in Fig.5. Both the RHNC and MSA versions overestimate the critical temperature as obtained from simulation by about $`15\%`$. Of course, this is expected for any classical theory. On the other hand, far away from the critical point, results from both versions of the theory are seen to yield fair agreement with simulation. The MSA version is somewhat more convenient, however, because it allows to calculate the coexistence at low temperatures with no additional cost, while it becomes rather problematic to calculate the coexistence for the RHNC version below the reference fluid critical temperature. The reason for this is that the RHNC integral equation presents a region of no solutions below this point, so that the resulting equation of state is no longer defined inside the liquid-vapor envelope.
We have also investigated the critical points of chains of 20, 40 and 60 n-mers, in an attempt to study the behavior of TPT1 for longer chains. Table I gives a summary of the simulation results, obtained by finite size scaling, together with predictions from TPT1-RHNC and TPT1-MSA. Both versions overestimate the critical temperatures by about 15% for all chain lengths studied. However, the MSA and RHNC predictions seem to converge as the chain length increases. On the other hand, the critical monomer densities are always underestimated, though the MSA version seems to give much better agreement than the RHNC version. In the latter theory the density decreases much too fast compared to the MC results. The overall behavior of the critical parameters is illustrated in Fig. 6, where both $`T_c`$ and $`\rho _c`$ are plotted against $`n^{1/2}`$, the predicted asymptotic scaling law for both of these properties. It is seen that for chain lengths up to 60 monomers, the critical properties are far from reaching their asymptotic behavior, so that the simulations do not allow as to asses unambiguously the predicted scaling laws.
Although the calculation of the critical point of fluids larger than about 100 monomers by computer simulation becomes prohibitively expensive, we can estimate the $`\mathrm{\Theta }`$ point of our polymer model by an analysis of the temperature dependence of the polymer extension. Fig. 7 shows a plot of the mean squared end to end distance divided by $`n1`$ as a function of temperature for various chain lengths. In the infinite chain length limit, the intercept of two such plots occurs at the $`\mathrm{\Theta }`$ point of the polymer model. Extrapolation of the results gives as an estimate $`\mathrm{\Theta }3.3`$. As to the theory, fitting the critical temperature predicted by TPT1-RHNC to a power law of the form $`T_c=T_c^{\mathrm{}}+bn^{1/2}+cn^1`$ in the range $`10^2`$ to $`10^7`$ gives $`T_c^{\mathrm{}}=3.44`$. On the other hand, by searching for the root in Eq. 50, we find that TPT1-MSA predicts $`T_c^{\mathrm{}}=3.14`$. Assuming that the $`\mathrm{\Theta }`$ point is indeed the critical point of the infinitely long chain, as suggested by the considerations of Section III, it would seem that TPT1 is capable of giving an excellent prediction for the $`\mathrm{\Theta }`$ point of the polymer, even though the actual prediction may vary somewhat depending on the theory used to describe the reference fluid. Remarkably, considering that $`\delta `$ is density independent and using the simple van der Waals equation of state, Nikitin et al. have shown that TPT1 predicts an asymptotic critical temperature $`T_c^{\mathrm{}}=\frac{27}{8}T_c^0`$. As $`T_c^0`$, the critical temperature of the reference fluid, is approximately $`1.0`$ LJ reduced units, we find that the simplest TPT1 approach already gives an excellent prediction for the $`\mathrm{\Theta }`$ point of $`T_c^{\mathrm{}}=3.375`$. This is, however, somewhat fortuitous as it was shown in Section III that this TPT1-van-der-Waals approach of Nikitin actually predicts that $`\mathrm{\Theta }`$ is equal to the Boyle temperature of the monomer fluid, $`T_B^0`$, which is about $`T_B^0=2.58`$. Thus, the relatively good estimate of $`T_c^{\mathrm{}}`$ turns out to be a consequence of the over prediction of $`T_B^0`$ implicit in the van der Waals equation of state.
Recently, Chatterjee and Schweizer have analyzed the behavior of the critical point of infinite chain lengths using the PRISM theory. For two of the closures employed, the same behavior as that predicted by TPT1 is observed, at least concerning i) vanishing critical density and ii) finite critical temperature. The power laws are, however, different. The critical monomer density is predicted to vanish with a weaker dependence which may be either $`1/3`$ or $`1/4`$ depending on whether the RMPY/HTA or the MSA closures are used. The critical temperature is predicted to reach a finite critical value with the same exponents as the critical density. However, there was no a priori reason for choosing one closure over the others and several different trends could be obtained depending on the closure that was used. It is pleasing to see that TPT1 is able to give a unique conclusion, independent of the molecular theory used to describe the monomer fluid.
## VII conclusion
In this paper we have used the formalism of Zhou and Stell to extend the TPT1 theory to polymers with variable bond-length. By using rigorous molecular theories for the reference fluid of non–bonded monomers we have been able to explore two versions of the theory that allow to give a good description of the fluid without the need of any empirical data for the polymer. Comparison with numerical simulations show that both the RHNC version and the MSA version give good agreement with the simulation data. At low temperatures, the RHNC version seems to be more reliable, while at high temperatures there is apparently little difference. The MSA version is seen to give fair predictions for the critical points of the longer chains, with the advantage that it is almost analytic. The $`\mathrm{\Theta }`$ point of the model is predicted in very good agreement by the RHNC version, as well as by the MSA version.
Concerning the critical behavior of long chains, it has been shown that TPT1 predicts an approach of the critical temperature to the $`\mathrm{\Theta }`$ point with a power law of $`n^{1/2}`$. This is the correct behavior for the infinitely long polymer chain, as mean field behavior must be recovered in the infinitely long chain limit. The $`\mathrm{\Theta }`$ point has been shown to be the Boyle temperature of the infinitely long polymer while the critical mass density is predicted to vanish with a power law of $`n^{1/2}`$. All of these predictions concerning the scaling behavior of the critical points of the pure monomer fluid are seen to agree exactly with the mean field (Flory–Huggins) predictions for the critical behavior of polymer+solvent mixtures. This gives further support to the idea that pure polymer equations of state may be used as effective equations of state for the polymer+continuum-solvent system and vice versa, polymer+solvent equations of state may be used for the polymer+vacuum case; a formal prove of this intuitively appealing idea is, however, difficult.
In a subsequent paper we use the implementation of TPT1 proposed here to describe the equation of state, together with a self consistent field theory to study the surface and interfacial properties of a polymer+solvent system.
### Acknowledgments
Helpful correspondence with Y.Tang is acknowledged. The authors also benefited from stimulating discussions with A. Milchev and V. Ivanov. Financial support is due to grant Bi314/17 of the DFG and to project No.PB97-0329 of the Spanish DGICYT. L.G.M. wishes to thank the Universidad Complutense for the award of a predoctoral grant and for founding a stay in Mainz.
## VIII Appendix A: Fitting the RHNC data for the monomer fluid
Rather than solving the integral equation of each of the state points desired, which would be rather expensive, we followed a similar approach as that used by Johnson et al. to describe the Lennard-Jones fluid. We solved the RHNC integral equation for a set of 756 states in the range $`1.1T6.0`$ and $`0<\rho <0.85`$, we calculated the resulting compressibility factor (see Eq. 57) and then fit the data to a Modified-Benedict-Webb-Rubin equation of state, given by:
$$\rho kT(Z1)=\underset{i=1}{\overset{8}{}}a_i\rho ^{i+1}+F\underset{i=1}{\overset{6}{}}b_i\rho ^{2i+1}$$
(61)
where $`F=\mathrm{exp}(\gamma \rho ^2)`$, while the exact form of the $`a_i`$ and $`b_i`$ coefficients is given in Table II. Once the fit to $`Z`$ is performed, the free energy may be determined from:
$$A/N=\underset{i=1}{\overset{8}{}}\frac{a_i\rho ^i}{i}+\underset{i=1}{\overset{6}{}}b_iG_i$$
(62)
The $`G_i`$ coefficients obey the following recursive relation:
$$G_i=\frac{F\rho ^{2(i1)}2(i1)G_{i1}}{2\gamma }$$
(63)
with the first term given by $`G_1=(1F)/(2\gamma )`$. The parameters obtained for the fit are collected in Table III.
Contrary to the approach of Johnson et at., we calculate the thermodynamics using the RHNC theory, rather than computer simulations. In this way we are able to save several orders of magnitude of CPU time. The effective diameter required in the RHNC equation is determined as suggested by Lado et al., while the effective hard sphere bridge function is calculated from the parameterization of Labik and Malijevsky.
Also required is the associating strength $`\delta `$ of the fluid. This is obtained by solving equation 37 for each of the 756 state points. Rather than fitting $`\delta `$, which is a difficult task, as it varies several orders of magnitude in the range $`1.1T6.0`$, we fit the dimensionless ratio $`\delta (\rho )/\delta (\rho =0)`$. The actual functional form used is:
$$\delta /\delta _0=1+\underset{i=1}{\overset{5}{}}\underset{j=1}{\overset{5}{}}a_{ij}\rho ^iT^{(1j)}$$
(64)
The parameters for this fit are gathered in Table IV.
## IX Appendix B: The perturbation theory of Tang and Lu
In order to describe the thermodynamics of the reference fluid we perform a Barker-Henderson decomposition of the monomer fluid pair potential such that $`u_0`$ is described in terms of a repulsive reference potential $`w_{\mathrm{ref}}`$, which is made of the positive region of $`u_0`$ and a perturbation, $`w_{\mathrm{per}}`$, which is made of the negative part of the potential:
$$w_{\mathrm{ref}}(r)=\{\begin{array}{cc}u_0(r)\hfill & rt\sigma \hfill \\ 0\hfill & r>t\sigma \hfill \end{array}\text{and}w_{\mathrm{per}}(r)=u_0(r)w_{\mathrm{ref}}(r)=\{\begin{array}{cc}0\hfill & rt\sigma \hfill \\ u_0(r)\hfill & r>t\sigma \hfill \end{array}$$
(65)
where $`t=1.0013`$ defines the value of $`r`$ where $`u_0`$ becomes negative (recall that we are considering a cut and shifted potential). We now couple the perturbation potential to the reference potential with a coupling parameter, $`\lambda `$, so that the actual potential is recovered for $`\lambda =1`$:
$$w(r;\lambda )=w_{\mathrm{ref}}(r)+\lambda w_{\mathrm{per}}(r)$$
(66)
At this stage we recall the fundamental functional expression that relates the Helmholtz free energy with the radial distribution function,
$$\frac{\delta A}{\delta w(r;\lambda )}=\frac{1}{2}N\rho g(r;\lambda )$$
(67)
Integration of this equation following the rules of functional calculus, leads to an expression relating the free energy of the monomer fluid with that of the reference fluid:
$$\frac{A}{N}\frac{A_{\mathrm{ref}}}{N}=\frac{1}{2}\rho _{\lambda =0}^{\lambda =1}_{t\sigma }^{\mathrm{}}g(r;\lambda )w_{\mathrm{per}}4\pi r^2drd\lambda $$
(68)
It is then assumed that the radial distribution function may be expanded as a series in powers of $`\lambda `$ of the form $`g(r;\lambda )=g_0(r)+g_1(r)\lambda +\mathrm{}`$. Truncation of the series to first order then yields:
$$\frac{A}{N}\frac{A_{\mathrm{ref}}}{N}=\frac{1}{2}\rho _{t\sigma }^{\mathrm{}}g_0(r)w_{\mathrm{per}}(r)4\pi r^2dr+\frac{1}{4}\rho _{t\sigma }^{\mathrm{}}g_1(r)w_{\mathrm{per}}(r)4\pi r^2dr$$
(69)
Following Barker and Henderson we now choose to describe the reference potential by means of a hard sphere fluid of appropriate hard sphere diameter, $`d`$. This choice is justified because the reference potential is made essentially of the repulsive part of the potential. In this way, $`\beta A_{\mathrm{ref}}/N`$ and $`g_0`$ may be considered to be the free energy and radial distribution function of an effective hard sphere fluid, while $`g_1`$ may be solved using the MSA closure. However, the integrals appearing in the previous equation are still quite tedious to calculate and a further approximation allows to get fully analytical results. This is done by fitting the actual monomer potential, $`u_0`$ to a Two Yukawa potential, following the procedure of ref. :
$$u(r)^{TY}=k_0ϵ\frac{e^{z_1(rt\sigma )}}{r}+k_0ϵ\frac{e^{z_2(rt\sigma )}}{r}$$
(70)
Actually, the fit needs to be performed only for values of $`r`$ greater than $`d`$ and the resulting function (with $`k_0=2.4405`$, $`z_1=3.492456`$ and $`z_2=13.109857`$) is virtually identical to the true potential. However, substitution of Eq. 70 in the first and second order contributions of Eq. 69 and solving $`g_1`$ for $`u^{TY}`$ rather than for $`u_0`$, a very accurate expression for Eq. 69 may be obtained which is fully analytic.
This is done by rearranging Eq. 69 into the form:
$$\frac{\beta A}{N}=a_0+a_1+a_2$$
(71)
where $`a_0`$ is given by the Carnahan-Starling equation of state:
$$a_0=\frac{4\eta 3\eta ^2}{(1\eta )^2}$$
(72)
while
$$a_1=\frac{1}{2}\rho _d^{\mathrm{}}(g_01)w_{per}^{TY}4\pi r^2dr+\frac{1}{2}\rho _d^{\mathrm{}}w_{per}4\pi r^2dr+g_0(d)_{t\sigma }^dw_{per}4\pi r^2dr$$
(73)
and
$$a_2=\frac{1}{4}\rho _d^{\mathrm{}}g_1(r)w_{per}^{TY}4\pi r^2dr+\frac{1}{4}\rho g_1(d)_{t\sigma }^dw_{per}4\pi r^2dr$$
(74)
In the last two equations it is understood that $`w_{\mathrm{per}}^{TY}`$ is the perturbation potential when expressed as in Eq. 70; the actual $`w_{\mathrm{per}}`$ potential is used where ever possible; $`g_0=1`$ and $`g_1=0`$ beyond the cutoff distance of the potential while $`g_0`$ and $`g_1`$ are considered to remain constant in the range $`[d,t\sigma ]`$. The only difference between the above expressions and those obtained by Tang et al. for the true Lennard-Jones potential are found in the trivial integrals of the form $`w_14\pi r^2dr`$ because both the integration limits and the perturbation potential differ. Solving $`a_1`$ and $`a_2`$ yields:
$`a_1`$ $`=`$ $`{\displaystyle \frac{12\eta \beta ϵ}{d^3}}[k_1({\displaystyle \frac{L(z_1d)}{z_1^2(1\eta )^2Q(z_1d)}}{\displaystyle \frac{1+z_1d}{z_1^2}})`$ (79)
$`k_2({\displaystyle \frac{L(z_2d)}{z_2^2(1\eta )^2Q(z_2d)}}{\displaystyle \frac{1+z_2d}{z_2^2}})]`$
$`+48\eta \beta ϵW_{\mathrm{cs}}(r_c,d)48\eta \beta ϵg_0(d)W_{\mathrm{cs}}(t\sigma ,d)`$
$`a_2`$ $`=`$ $`{\displaystyle \frac{6\eta \beta ^2ϵ^2}{d^3}}[{\displaystyle \frac{k_1^2}{2z_1Q^4(z_1d)}}+{\displaystyle \frac{k_2^2}{2z_2Q^4(z_2d)}}`$ (84)
$`{\displaystyle \frac{2k_1k_2}{(z_1+z_2)Q^2(z_1d)Q^2(z_2d)}}]`$
$`24\eta \beta ^2ϵ^2\left[{\displaystyle \frac{k_1/d}{Q^2(z_1d)}}{\displaystyle \frac{k_2/d}{Q^2(z_2d)}}\right]W_{\mathrm{cs}}(t\sigma ,d)`$
where $`\eta =\pi /6d^3\rho `$ and $`k_i=k_0e^{z_i(t\sigma d)}`$. The Barker-Henderson diameter, defined as
$$d=_0^{\mathrm{}}(1e^{\beta w_{\mathrm{ref}}(r)})dr$$
(85)
is parameterized using the formula proposed in Ref. :
$$d/\sigma =2^{1/6}\left[1+\left(1+\frac{T+c_2T^2+c_3T^4}{c_1}\right)^{1/2}\right]^{1/6}$$
(86)
where $`c_1=1.150167`$, $`c_2=0.046498`$, and $`c_3=0.0004477054`$.
On the other hand, $`Q`$ is defined as:
$$Q(t)=\frac{S(t)+12\eta L(t)e^t}{(1\eta )^2t^3}$$
(87)
where
$$S(t)=(1\eta )^2t^3+6\eta (1\eta )t^2+18\eta ^2t12\eta (1+2\eta )$$
(88)
and
$$L(t)=(1+\eta /2)t+1+2\eta $$
(89)
The reference radial distribution at contact is given by:
$$g_0(d)=\frac{1+\eta /2}{(1\eta )^2}$$
(90)
while $`W_{cs}`$ is to constant factors, the definite integral of the perturbation potential, $`w_{\mathrm{per}}`$:
$`W_{cs}(x,y)`$ $`=`$ $`\{{\displaystyle \frac{1}{3}}[\left({\displaystyle \frac{\sigma }{x}}\right)^3\left({\displaystyle \frac{\sigma }{y}}\right)^3]{\displaystyle \frac{1}{9}}[\left({\displaystyle \frac{\sigma }{x}}\right)^9\left({\displaystyle \frac{\sigma }{y}}\right)^9]`$ (93)
$`{\displaystyle \frac{1}{3}}[\left({\displaystyle \frac{\sigma }{r_c}}\right)^{12}\left({\displaystyle \frac{\sigma }{r_c}}\right)^6][\left({\displaystyle \frac{x}{\sigma }}\right)^3\left({\displaystyle \frac{y}{\sigma }}\right)^3]\}{\displaystyle \frac{\sigma ^3}{d^3}}`$
The compressibility factor may be determined by density differentiation of the free energy, which leads to:
$$Z=\frac{pV}{NkT}=Z_0+Z_1+Z_2$$
(94)
where
$$Z_0=\frac{1+\eta +\eta ^2\eta ^3}{(1\eta )^3}$$
(95)
$`Z_1`$ $`=`$ $`a_1{\displaystyle \frac{12\eta ^2\beta ϵ}{d^3}}[k_1({\displaystyle \frac{(5/2+\eta /2)z_1d+4+2\eta }{z_1^2(1\eta )^3Q(z_1d)}}{\displaystyle \frac{L(z_1d)Q_\eta ^{}(z_1d)}{z_1^2(1\eta )^2Q^2(z_1d)}})`$ (100)
$`k_2({\displaystyle \frac{(5/2+\eta /2)z_2d+4+2\eta }{z_2^2(1\eta )^3Q(z_2d)}}{\displaystyle \frac{L(z_2d)Q_\eta ^{}(z_2d)}{z_2^2(1\eta )^2Q^2(z_2d)}})]`$
$`16\eta \beta ϵd{\displaystyle \frac{g_0(d)}{d}}W_{\mathrm{cs}}(t\sigma ,d)`$
$`Z_2`$ $`=`$ $`a_2+{\displaystyle \frac{12\eta ^2\beta ^2ϵ^2}{d^3}}[{\displaystyle \frac{k_1^2Q_\eta ^{}(z_1d)}{z_1Q^5(z_1d)}}+{\displaystyle \frac{k_2^2Q_\eta ^{}(z_2d)}{z_2Q^5(z_2d)}}`$ (105)
$`{\displaystyle \frac{2k_1k_2[Q_\eta ^{}(z_1d)Q(z_2d)+Q(z_1d)Q_\eta ^{}(z_2d)}{(z_1+z_2)Q^3(z_1d)Q^3(z_2d)}}]`$
$`+48\eta ^2\beta ^2ϵ^2\left[{\displaystyle \frac{k_1/d}{Q^3(z_1d)}}Q_\eta ^{}(z_1d){\displaystyle \frac{k_2/d}{Q^3(z_2d)}}Q_\eta ^{}(z_2d)\right]W_{\mathrm{cs}}(t\sigma ,d)`$
while
$$d\frac{g_0(d)}{d}=\frac{3\eta (5+\eta )}{2(1\eta )^3}$$
(106)
and
$$Q_\eta ^{}(t)=\frac{6(1\eta )t^2+36\eta t12(1+5\eta )+12[(1+2\eta )t+1+5\eta ]e^t}{(1\eta )^3t^3}$$
(107)
In order to calculate the associating strength, $`\delta `$, the radial distribution function could have been assumed to be $`g=g_0+g_1`$, which is already a rather good approximation. However, we use the SEXP (simplified Exponential) approximation which considerably improves the estimate of $`g`$ around $`\sigma `$ with no additional information. According to this approximation, $`g=g_0e^{g_1}`$. Once $`g`$ is known, $`\delta `$ is calculated directly by invoking Eq. 37. The integral must be performed numerically but would have been analytic if the bond length was held fixed.
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# The neutrino star in the bulk Universe
## Introduction
Recently there has been considerable interest in the field theories with large extra spacetime dimensions. In the comparison to the standard Kaluza-Klein theory these extra dimensions may be restricted only to the gravity sector of theory while the Standard Model (SM) fields are assumed to be localized on the 4-dimensional spacetime 1 ,4 ,ADD3 .This is a promising scenario from the phenomenological point of view because it shift the energy scale of unification from $`10^{19}GeV`$ to $`1100TeV`$. It has been recently shown L1 that this framework can be embedded into string models, where the fundamental Planck scale can be identified with the string scale which could be as low as the weak scale. The extra dimensions have the potential to lower the unification scale as well DDG1 .
The aim of this paper is to examine the degenerate neutrino star originated from the extra dimensional theory. The neutrino star ( neutrino ball ) was a subject of interest in the theory of the Standard Model hd . The fermion star in the extra dimensional theory was also the a subject of interest nk .The neutrino star model is capable to explain the nature of the object in Sgr A in the center of the Galaxy torr .
## The bulk neutrino extension of the electroweak theory
Much of the interesting phenomenology of brane-world models is associated with the Kaluza-Klein theories tap that originate from large, gravity-only additional dimensions. The higher-dimensional bulk fermions lead to Kaluza-Klein towers of standard model singlets that may be interpreted as sterile neutrinos dinest ; ADDM ; dvali ; BCS . In this paper we shall consider the six-dimensional Kaluza - Klein theory. Let us now consider the action in the six-dimensional spacetime:
$$𝒮=d^6x\sqrt{g_6}L=𝒮_g+𝒮_F=d^6x\sqrt{g_6}(L_g+L_F),$$
(1)
where $`g_6=det(g_{MN})`$ and $`M=\{\mu ,i\}`$, $`N=\{\nu ,j\}`$ with $`x^M=\{x^\mu ,y^i\},i=1,2`$. The metrical tensor in the six-dimensional spacetime can be written:
$$g_{MN}=\left(\begin{array}{cc}e^{2\xi (x)/f_0}\overline{g}_{\mu \nu }& 0\\ 0& \delta _{ij}e^{+2\xi (x)/f_0}\end{array}\right).$$
(2)
According to the above definition we can write:
$$\sqrt{g_6}=\sqrt{\overline{g}}e^{2\xi (x)/f_0}.$$
(3)
We consider here the Lagrangian of the field as follows:
$`L=L_g+L_F,`$ (4)
$`L_g={\displaystyle \frac{1}{2\kappa _6}}R,`$ (5)
$`L_F=L_{brane}\delta (y)+L_{bulk},`$ (6)
where $`\kappa _6`$ is the six-dimensional gravitational coupling. Its natural interpretation originates from the distance scaling in the four-dimensional spacetime. Let us compactify the six-dimensional spacetime $`(_6_4\times 𝒮^1\times 𝒮^1)`$ to the four-dimensional Minkowski one. In this paper we assume that the extra dimensions are compactified on the $`2`$ dimensional torus with a single radius $`r_2`$. The six-dimensional action may be rewritten as:
$$𝒮=d^4xd^2y\sqrt{g_6}L=d^4x\sqrt{\overline{g}},$$
(7)
where $`d^2y=(2\pi r_2)^2`$. The six-dimensional gravitional coupling $`\kappa _6=8\pi G_6`$ is convenient to define as
$$G_6^1=\frac{4\pi }{(2\pi )^2}M^4,$$
where $`M`$ \- is the energy scale of the compactification ($`10100TeV`$). Cosmological consideration hall gives the bound $`M>100TeV`$ what corresponds $`r_2<5.1\times 10^5mm`$. If we define the four-dimensional coupling constant $`\kappa =8\pi G_N=8\pi M_{Pl}^2`$ we get
$$M_{Pl}^2=4\pi M^4r_2^2.$$
(8)
The parameter $`f_0`$ is defined as:
$$f_0^2=\frac{1}{\pi }M_{Pl}^2$$
(9)
to produce the $`1/2`$ term in for the dilaton field kinetic term. The model can be motivated by the ten-dimensional string theory with the string scale $`M_s`$. For simplicity we assume that the first $`2`$ dimensions have the same size characterized by the radius $`r_2`$. This is likewise assumed for the remaining $`4`$ dimensions where the corresponding radius is called $`r_4<r_2`$. The four-dimensional Planck constant
$$(\frac{M_{Pl}}{M_s})^2=16\pi (2\pi r_2M_s)^2(2\pi r_4M_s)^4$$
(10)
is connected to the string scale $`M_s`$. In the result of the comparison to (8) we have
$$M=\sqrt{4\pi }(2\pi r_4M_s)M_s.$$
(11)
The paremeters of the model karmen are presented on the Table 1.
In this section we shall extend the Standard Model minimally with the bulk neutrino $`\psi (x,y)`$. The lagrangian of the neutrino sector of the model is then:
$$L_{brane}=i\overline{L}_f\widehat{\mathrm{\Gamma }}^\mu _\mu L_f\frac{1}{M_s}(h_{f,a}\overline{\psi }_aL_fH_0+h.c.),$$
and the fermion field are
$$L_f=\left(\begin{array}{c}\nu _f^0\\ 0\end{array}\right)_Lf=e,\nu ,\tau .$$
(12)
The Higgs field will have only the residual form
$$H_0=\frac{1}{\sqrt{2}}\left(\begin{array}{c}0\\ \upsilon \end{array}\right)$$
(13)
after the spontaneous symmetry breaking. The additional bulk neutrino is described by the Lagrange function
$`_{bulk}`$ $`=`$ $`i\overline{\psi }\widehat{\mathrm{\Gamma }}^N_N\psi m_D(\overline{\psi }_L\psi _R+h.c),`$ (14)
where $`\psi _{R,L}=(I\pm \mathrm{\Gamma }^7)\psi `$. The system has global $`U(1)`$ lepton symmetry which generates the lepton charge $`Q_L`$. The bulk neutrino field can be decomposed into four-dimensional Dirac spinors $`\psi _a`$, where $`a=\{+,\}`$. The gamma matrices $`\mathrm{\Gamma }`$ are defined as
$$\{\widehat{\mathrm{\Gamma }}^M,\widehat{\mathrm{\Gamma }}^N\}=2g^{M,N}I.$$
(15)
In the flat four-dimensional spacetime, when $`\overline{g}_{\mu \nu }=\eta _{\mu \nu }`$ they may be defined as
$$\widehat{\mathrm{\Gamma }}^\mu =e^{\xi /f_0}\mathrm{\Gamma }^\mu ,\mathrm{\Gamma }^\mu =\gamma ^\mu I_{(2)},$$
(16)
$$\widehat{\mathrm{\Gamma }}^i=e^{2\xi /f_0}\mathrm{\Gamma }^i,\mathrm{\Gamma }^i=i\gamma ^5\sigma ^i$$
and
$`\{\gamma ^\mu ,\gamma ^\nu \}=2\eta ^{\mu \nu }I,`$
$`\{\gamma ^i,\gamma ^j\}=2\delta _{i,j}I,`$
where
$$\mathrm{\Gamma }^7=\overline{\mathrm{\Gamma }}^0\mathrm{}\overline{\mathrm{\Gamma }}^5=\gamma ^5\sigma ^3,\gamma ^5=i\gamma ^0\gamma ^1\gamma ^2\gamma ^3.$$
Using the metric tensor six-dimensional form (2) and the (16) we can calculate the Lagrange function in the four-dimensional Minkowski spacetime
$`L=i\overline{L}\gamma ^\mu D_\mu L+e^{2\xi /f_0}{\displaystyle \underset{a,\underset{¯}{n}}{}}i\overline{\psi }_{a,\underset{¯}{n}}(x)\gamma ^\mu D_\mu \psi _{a,\underset{¯}{n}}(x)`$ (17)
$`m_D{\displaystyle \underset{a,\underset{¯}{n}}{}}\overline{\psi }_{a,\underset{¯}{n},L}\psi _{a,\underset{¯}{n},R}e^{\xi /f_0}{\displaystyle \underset{a,\underset{¯}{n}}{}}\overline{\psi }_{a,\underset{¯}{n}}(x)(\gamma ^iu_{\underset{¯}{n}}^i)\psi _{a,\underset{¯}{n}}(x),`$
where
$$u_{\underset{¯}{n}}^i=\frac{1}{r_2}n^i,m_{\underset{¯}{n}}^2=\underset{¯}{u}_{\underset{¯}{n}}^2=m\frac{1}{r_2^2}(n_1^2+n_2^2).$$
(18)
The bulk neutrino is decomposed as
$$\psi _a(x,y)=\frac{e^{\xi /f_0}}{(2\pi r_2)}\underset{\underset{¯}{n}}{}\psi _{a,\underset{¯}{n}}(x)exp(\frac{1}{r_2}i\underset{¯}{n}\underset{¯}{y}).$$
(19)
Each of these four-dimensional Dirac spinors can be decomposed into Weyl spinors
$$\psi _{a,\underset{¯}{n}}(x)=\left(\begin{array}{c}\xi _{a,\underset{¯}{n}}^c\\ \eta _{a,\underset{¯}{n}}\end{array}\right).$$
(20)
After electroweak symmetry breaking we introduce the mass
$$m_{f,a}=\frac{h_{f,a}v}{2\pi (M_sr_2)}10^1MeV.$$
(21)
The exact diagonalization of the mass matrix gives three exactly massles Weyl fermions
$$\nu _f=U_{f,f^{}}\nu _f^{}^0+V_{f,a}(\underset{¯}{0})\eta _a(\underset{¯}{0})$$
(22)
and two Dirac spinors for each mode number $`\underset{¯}{n}`$ with mass
$$M_{\underset{¯}{n}}=\sqrt{m_D^2+m_{\underset{¯}{n}}^2}.$$
(23)
In general karmen , there is a superposition of the electroweak neutrinos $`\nu _f^0`$ on the brane and the Kaluza-Klein bulk neutrinos $`\eta _a(\underset{¯}{n})`$.
## The neutrino star
An alternative model for the supermassive compact object in the center of our Galaxy has been recently proposed by Tsiklauri and Viollier rm tsi . The main ingredient of the proposal is that the dark matter at the center of the galaxy is non-baryonic composed with massive neutrinos or gravitinos. Such neutrino balls could have formed in early epochs, during a first-order phase transition in the Standard Model.
In case of the spherically symmertic gravitational field we have:
$$\overline{g}_{\mu \nu }=\left(\begin{array}{cccc}e^{\nu (r)}& & & \\ & e^{\lambda (x)}& & \\ & & r^2& \\ & & & r^2\mathrm{sin}^2\theta \end{array}\right).$$
(24)
In the similar way the dilaton filed $`\xi (r)`$ will be dependent on the radius $`r`$. As $`\xi (r)f_0=M_{Pl}/\sqrt{\pi }10^{19}GeV`$ we shall neglect the dilaton filed $`\xi (r)`$ in the first approximation. In this paper we present numerical results describing the structure of neutrino star. It is possible to describe a static spherical star solving the Oppenheimer-Tolman-Volkoff (OTV) equations (more general case with the dilaton filed $`\xi (r)`$ is presented in Appendix I).
$$\frac{dP(r)}{dr}=\frac{G}{r^2}(\rho (r)+\frac{P(r)}{c^2})\frac{(m(r)+\frac{4\pi }{c^2}P(r)r^3)}{(1\frac{2Gm(r)}{c^2r})},$$
(25)
$$\frac{dm(r)}{dr}=4\pi r^2\rho (r).$$
(26)
Having solved the OTV equation the pressure $`p(r)`$, mass $`m(r)`$ and density $`\rho (r)`$ were obtained. To obtain the total radius $`R`$ of the star the fulfillment of the condition $`p(R)=0`$ is necessary. Our aim is to achieve the equation of state of neutrino star matter at finite temperature. In such a case the physical system can be defined by the thermodynamic potential $`\mathrm{\Omega }`$ fet
$$\mathrm{\Omega }=k_BTlnTr(e^{\beta (H\mu Q_L)}),$$
(27)
where $`\beta =1/k_BT`$, $`k_B`$ is the Boltzmann constant, $`Q_L`$ lepton charge. The chemical potential $`\mu `$ reflects the lepton number conservation. $`H`$ stands for the Hamiltonian of the physical system. All needed averages are calculated with the Hamiltonian $`H`$. We define the density of energy and pressure by the energy - momentum tensor
$$T_{\mu \nu }=(P+ϵ)u_\mu u_\nu Pg_{\mu \nu },$$
(28)
where $`u_\mu `$ is a unite vector ($`u_\mu u^\mu =1`$). So, the calculations give
$$ϵ(x_F,T)=\rho c^2=ϵ_{F,bulk}+ϵ_{F,brane},$$
(29)
$$P(x_F,T)=P_{F,bulk}+P_{F,brane},$$
(30)
where
$$ϵ_F=ϵ_0\chi (x_F,T),$$
(31)
$$P_F=P_0\varphi (x_F,T).$$
(32)
The fact that the bulk Dirac neutrinos are massive means that they play the same role like ions in a white dwarf or neutrons in a neutron star. We have
$`\chi _{bulk}(x_F,T)={\displaystyle \frac{1}{^{\pi ^2}}}{\displaystyle _0^{\mathrm{}}}dzz^2\sqrt{z^2+1}\{{\displaystyle \frac{1}{\mathrm{exp}((\sqrt{1+z^2}\mu ^{})/\tau )+1}}`$ (33)
$`+{\displaystyle \frac{1}{\mathrm{exp}((\sqrt{1+z^2}+\mu ^{})/\tau )+1}},`$
$`\varphi _{bulk}(x_F,T)={\displaystyle \frac{1}{3\pi ^2}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{z^4dz}{\sqrt{z^2+1}}}\{{\displaystyle \frac{1}{\mathrm{exp}((\sqrt{1+z^2}\mu ^{})/\tau )+1}}`$ (34)
$`+{\displaystyle \frac{1}{\mathrm{exp}((\sqrt{1+z^2}+\mu ^{})/\tau )+1}}\},`$
where $`\tau =(k_BT)/M_2`$,
$$\mu ^{}=\mu /M_2=\sqrt{1+x_F^2}$$
(35)
and
$$x_F=k_F/M_s.$$
(36)
Similarly to the paper toki we have introduced (35,36) the dimensionless “Fermi” momentum even at finite temperature which exactly corresponds to the Fermi momentum at zero temperature. For the massless brane neutrinos we have
$`\chi _{brane}(x_F,T)={\displaystyle \frac{\gamma }{2\pi ^2}}{\displaystyle 𝑑zz^3\{\frac{1}{(e^{(zx_F)/\tau }+1)}+\frac{1}{(e^{(z+x_F)/\tau }+1)}\}}`$ (37)
$`={\displaystyle \frac{3\gamma }{\pi ^2}}t^4{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dzz^3}{e^{(zx_F)/\tau }+1}}\left\{Li_4(e^{x_F/\tau })+Li_4(e^{x_F/\tau })\right\},`$ (38)
where the presure is
$$\varphi _{brane}(x_F,T)=\frac{1}{3}ϵ_{brane}(x_F,T).$$
(39)
Both $`ϵ_F`$ and $`P_F`$ depend on the neutrino chemical potential $`\mu `$ or Fermi momentum $`x_F`$. This parametric dependence on $`\mu `$ (or $`x_F`$) defines the equation of state.
Similarly to the paper toki we have introduced the dimensionless ’Fermi’ momentum even at finite temperature which exactly corresponds to the Fermi momentum at zero temperature. Both $`ϵ_F`$ and $`P_F`$ depend on the neutrino chemical potential $`\mu `$ or Fermi momentum $`x_F`$. This parametric dependence on $`\mu `$ (or $`x_F`$) defines the equation of state.
When Fermi momentum reaches the second Kaluza-Klein level $`m_D\sqrt{2}`$ the number of avaliable neutrino modes will change. For the high bulk neutrino density in the ultrarelativistic limit all Kaluza-Klein modes should be included. In that limit energy density of the bulk neutrinos is equal to
$`\chi _{bulk}(x_F,T)={\displaystyle \frac{\gamma }{3\pi ^2}}(r_2M_2){\displaystyle 𝑑zz^5\{\frac{1}{(e^{(zx_F)/\tau }+1)}+\frac{1}{(e^{(z+x_F)/\tau }+1)}\}}`$ (40)
$`={\displaystyle \frac{40\gamma }{\pi }}(r_2M_2)t^6{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dzz^5}{e^{(zx_F)/\tau }+1}}\left\{Li_6(e^{x_F/\tau })+Li_6(e^{x_F/\tau })\right\}.`$ (41)
In the ultrarelativistic limit the equation of state is
$$\varphi _{bulk}(x_F,T)=\frac{1}{5}ϵ_{bulk}(x_F,T).$$
(42)
The equations (25,26) are easy integrated numerically, For example, for the neutron star with the central density $`\rho _c=\mathrm{\hspace{0.17em}10}^4g/cm^3`$ and temperature $`T=50keV`$ the star density and pressure profile are presented on the Fig.1 and Fig. 2. Similarly to a structure of white dwarf massive sterile neutrinos like ions contribute to the density of the star , while the massles neutrinos like electrons contribute to the pressure. This feature will be more visible with the increasing temperature.
The parameters of the maximum mass configuration are:
$$M_{max}=2.3\times 10^4M_{},R=1.2\times 10^6km.$$
(43)
This fact is easy to notice on the mass-radius diagram (Fig. 3). In general, we have all family of neutrino stars, depending on growing neutrino Fermi momentum. This result concerns the Ferni momentum below the second Kaluza-Klein level $`m_D\sqrt{2}`$ . The ulrarelativistic limit gives higher masses $`M_{max}1.6\times 10^6M_{}`$. The density profile in this limit is presented on the Fig 1 (magenta). In limit of bulk neutrino ball (Appendix I) the neutrino star properties are presented in Table 2.
and the neutrino star mass dependence from the central density $`\rho _c`$ (Fig. 4).
## Appendix
The Einstein equations in in the six-dimensional spacetime (the metric tensor (2,24) ) can be written as
$$e^{\lambda (r)}(\frac{\nu ^{}(r)}{r}+\frac{1}{r^2}2\frac{\xi ^{}(r)^2}{f_0^2})\frac{1}{r^2}=\frac{8\pi G_6}{c^4}e^{2\xi (r)/f_0}P(r),$$
(44)
$$e^{\lambda (r)}(\frac{\lambda ^{}(r)}{r}+\frac{1}{r^2}+2\frac{\xi ^{}(r)^2}{f_0^2})\frac{1}{r^2}=\frac{8\pi G_6}{c^2}e^{2\xi (r)/f_0}\rho (r).$$
(45)
The Einstein equations give also the equation for the dilaton field:
$`\xi ^{\prime \prime }(r)+{\displaystyle \frac{2}{r}}\xi ^{}(r){\displaystyle \frac{1}{2}}\xi ^{}(r)(\lambda ^{}(r)\nu ^{}(r)){\displaystyle \frac{1}{f_0}}\xi ^{}(r)^2+`$ (46)
$`f_0({\displaystyle \frac{1}{2r^2}}(e^{\lambda (r)}1)+{\displaystyle \frac{1}{2r}}(\lambda ^{}(r)\nu ^{}(r))+{\displaystyle \frac{1}{8}}(\lambda ^{}(r)\nu ^{}(r))\nu ^{}(r)`$
$`{\displaystyle \frac{1}{4}}\nu ^{\prime \prime }(r)+{\displaystyle \frac{4\pi G_6}{c^4}}e^{\lambda (r)2\xi (r)/f_0}P_s(r)).`$
The continuity equation $`T_{;M}^{MN}=0`$ gives
$$\nu ^{}(r)=\frac{2P^{}(r)}{(P+c^2\rho )}2\frac{\xi ^{}(r)}{f_0}\frac{(P2P_sc^2\rho )}{(P+c^2\rho )}.$$
(47)
We assume that the energy-momentum tensor has a diagonal form $`diag(T_{\mu \nu })=\{\epsilon ,P,P,P,P_s,P_s\}`$. Integrating of equation (45) yields we can write:
$$e^{\lambda (r)}=\frac{\frac{dr}{r}((1\frac{8\pi G}{c^2}r^2\rho (r)e^{\frac{2\xi (r)}{f_0}})e^{{\scriptscriptstyle 𝑑r({\scriptscriptstyle \frac{1}{r}}+{\scriptscriptstyle \frac{2r}{f_0^2}}\xi ^{}(r)^2)}})}{\mathrm{exp}(𝑑r(\frac{1}{r}+\frac{2r}{f_0^2}\xi ^{}(r)^2))}.$$
(48)
If we define
$$e^{\lambda (r)}=1\frac{2G_6}{c^2r}m(r),$$
(49)
then we can obtain the generalized Oppenheimer-Tolman-Volkoff equation:
$`{\displaystyle \frac{dP(r)}{dr}}={\displaystyle \frac{G_6}{r^2}}{\displaystyle \frac{(\rho (r)+\frac{P(r)}{c^2})(m(r)+\frac{4\pi }{c^2}P(r)r^3e^{2\xi (r)/f_0})}{(1\frac{2G_6}{c^2r}m(r))}}+`$ (50)
$`({\displaystyle \frac{\xi ^{}(r)^2}{f_0^2}}{\displaystyle \frac{\xi ^{}(r)}{^{f_0}}}{\displaystyle \frac{(P2P_sc^2\rho )}{(P+c^2\rho )}})(c^2\rho +P).`$ (51)
In vacuum $`P=0,P_s=0`$, of course we have well known Schwarzschild solution:
$`e^{\lambda (r)}=(1{\displaystyle \frac{r_g}{r}}),`$
$`e^{\nu (r)}=(1{\displaystyle \frac{r_g}{r}}),`$
$`\xi (r)=0.`$
Neglecting the dilaton field $`\xi (r)`$ the Oppenheimer-Tolman-Volkoff (25).
The theory of neutrinos, bound by gravity, can be easily sketched considering a Thomas-Fermi model for fermions bilic . We can set the Fermi energy equal to the gravitational potential which binds the system, and see that the number density is a function of the gravitational potential. Such a gravitational potential will obey a Poisson equation, where neutrinos (and anti-neutrinos) are the source term. Including gravity the local equilibrium condition demands
$$\frac{1}{r^2}\frac{d}{dr}(\frac{r^2}{\rho (r)}\frac{dp}{dr})=4\pi G_N\rho (r).$$
(52)
Inside the ball the pressure of the massive bulk neutrinos is $`P=\rho /5`$ while $`P=\rho /3`$ for the brane neutrinos in the relativistic limit in high temperature. Defining
$$\rho (r)=\rho _ce^\phi $$
(53)
with
$$\kappa =24\pi G_N\rho _c$$
(54)
we have the Liuoville equation
$$\mathrm{}\phi =\kappa e^\phi .$$
(55)
The Laplace operator is
$$\mathrm{}=\frac{1}{r^2}\frac{d}{dr}(r^2\frac{d}{dr})+\mathrm{}$$
(56)
Using the thin-wall approximation one can obtain the following expression
$$F_0(r)=e^\varphi =\frac{1}{\mathrm{cosh}^2(\frac{\sqrt{2\kappa }}{2}r)}.$$
(57)
The equation (57) allows to estimate the mass of the neutrino ball which is given by the dependence
$$M(R)=4\pi \rho _c_0^{\mathrm{}}\frac{drr^3}{\mathrm{cosh}^2(\frac{\sqrt{2\kappa }}{2}r)}$$
(58)
with the radius scale
$$r_0=\sqrt{\frac{\kappa }{2}}.$$
(59)
The bulk neutrino density profile (57) is presented on the Fig. 1 (the dotted line).
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# Confidence Intervals and Upper Bounds for Small Signals in the Presence of Background Noise
## 1 Introduction
Finding confidence intervals or upper limits for small signals has recently attracted a great deal of attention. The paper by Feldman and Cousins Cousins-Feldman rekindled the interest in this area by providing a unified approach, that is, it based the choice of quoting an upper limit or a two-sided confidence interval on the data alone, without the experimenter having to make this decision. The unified approach uses the Neyman construction together with an ordering principle based on likelihood ratios. Subsequent papers such as Giunti Giunti and Roe and Woodroofe Roe-Woodroofe gave variations of this method, basing the ordering on other quantities. Common to all these methods is the need to have a fairly precise knowledge of the background rate, for example, from Monte Carlo simulations. Unfortunately, as we will see in section 4, these methods can fail when used in the presence of background uncertainty. A possible remedy is discussed in Cousins and Highland Cousins-Highland where it is suggested to treat the background uncertainty as a systematic error. We will instead treat the background uncertainty as a statistical error and develop a method that is suitable in this situation. Our method is based on the likelihood ratio test statistic, together with an adjustment for those situations where there is very little (or no) signal observed. We will show that this method yields the correct coverage rate and that it has good power.
The background rate has to be estimated either from the data or through Monte Carlo. Both of those methods have their strengths and their weaknesses. Using data requires choosing sidebands, and implicitly makes the assumption that the density generating the background events is the same in the signal region as it is in the sidebands. This leads to a quandary: If we choose a small sideband, this assumption seems more reasonable, but then we will also see fewer background events and therefore have a higher statistical uncertainty in the estimate of the background rate. Choosing a large region might yield higher statistics but makes the assumption of a linear background more tenuous. An alternative way to estimate the background rate is by Monte Carlo. One problem with this approach is that a good Monte Carlo is often difficult to do because we can never be completely sure that we have modeled all the relevant effects correctly. This is particularly true when searching for small signals since the efforts to reduce the background often mean one is probing the tails of distributions which may be difficult to model. Also, in High Energy Physics today running a complete Monte Carlo simulation of an experiment can be a formidable task from a computational point of view, and it might not be possible to run enough Monte Carlo to effectively eliminate the uncertainty in the background estimate. For these reasons, it is in many cases not possible to ignore the uncertainty in the background rate. Our method is well equipped to deal with the uncertainty that comes from limited statistics in both situations, those where the background is estimated from the data as well as those where the background is estimated using Monte Carlo.
In some cases there are two different sources of background. It turns out that our method can be extended fairly easily to this situation also, regardless of the method of estimation used.
## 2 A Description of the Method
In this section we will outline the basic ideas of this new method. We will need the following notation. Assume that we observe $`x`$ events in a suitably chosen signal region and a total of $`y`$ events in the background region. Here the background region can be chosen fairly freely and need not be contiguous. Furthermore, we assume that the ratio of the size of the background region to the size of the signal region is $`\tau `$. For example, if we use two background regions of the same size as the signal region we get $`\tau =2`$. Then a probability model for the data is given by
$$XPois(\mu +b),YPois(\tau b)$$
where $`\mu `$ is the signal rate, $`b`$ is the background rate and $`Pois`$ is the usual Poisson distribution. We can assume $`X`$ and $`Y`$ to be independent and so
$$P_{\mu ,b}(X=x,Y=y)=\frac{(\mu +b)^x}{x!}e^{(\mu +b)}\frac{(\tau b)^y}{y!}e^{\tau b}$$
### 2.1 A Confidence Region for $`\mu `$ and $`b`$
A common technique for constructing confidence regions (or intervals for one-dimensional parameters) is to find a corresponding hypothesis test and then to invert the test. We will start with a simultaneous hypothesis test for $`\mu `$ and $`b`$ with the null hypothesis $`H_0:\mu =\mu _0,b=b_0`$. The steps are as follows:
1. List all observations $`(u_i,v_i),i=1,..,K`$, together with their probabilities, which are given by
$$P_{\mu _0,b_0}(X=u_i,Y=v_i)=\frac{(\mu _0+b_0)^{u_i}}{u_i!}e^{(\mu _0+b_0)}\frac{(\tau b_0)^{v_i}}{v_i!}e^{\tau b_0}$$
Here we use the values $`\mu _0`$ and $`b_0`$ specified in the null hypothesis. We list only those observations which have a probability above a certain small threshold. In our algorithm we require the probabilities to be larger than $`10^6`$.
2. Sort all observations from the most likely to the most unlikely.
3. Find the partial sums from the largest to the $`k^{th}`$ observation until you reach $`1\alpha ,`$ if the desired level of the test is $`\alpha .`$
4. If the observed $`(x,y)`$ appears in the list of possible observations before $`1\alpha `$ is reached, accept the null hypothesis, otherwise reject it.
This is in effect a Neyman construction like the one used in Feldman and Cousins Cousins-Feldman , only we use the probabilities as the ordering quantities. The test simply checks whether or not the observation $`(x,y)`$ is compatible with the rates $`\mu _0`$ and $`b_0`$ specified in the null hypothesis.
Using the likelihood ratios as the ordering quantity in a way similar to Feldman and Cousins Cousins-Feldman would in principle be superior, but unfortunately does not yield a viable method here because the list of ”likely” observations is infinite.
The inversion of the hypothesis test then involves a search through all pairs $`(\mu ,b)`$. If a certain pair leads to the acceptance of the null hypothesis, we add it to the confidence region, otherwise we do not. As an example we have figure 1 where we show the confidence regions obtained for three observations. The boundaries of the confidence regions are somewhat ragged due to the discrete nature of the Poisson random variable.
### 2.2 A Confidence Interval for $`\mu `$
Again we will start with a hypothesis test, but this time we will only fix the signal rate $`\mu `$. The null hypothesis then becomes $`H_0:\mu =\mu _0`$.
A popular test in Statistics for any kind of hypothesis test is the likelihood ratio test, which is based on the likelihood ratio test statistic $`\mathrm{\Lambda }`$ given in our problem by:
$$\mathrm{\Lambda }(\mu _0;x,y)=\frac{\mathrm{max}\left\{l(\mu _0,b;x,y):b0\right\}}{\mathrm{max}\left\{l(\mu ,b;x,y):\mu 0,b0\right\}}$$
Here $`l(\mu ,b;x,y)=P_{\mu ,b}(X=x,Y=y)`$ is the likelihood function of $`\mu `$ and $`b`$ given the observation$`(x,y)`$. This test statistic can be thought of as the ratio of the best explanation for the data if $`H_0`$ is true and the best explanation for the data if no assumption is made on $`\mu `$. The denominator is simply the likelihood function evaluated at the usual maximum likelihood estimator. To find the numerator we have to find the maximum likelihood estimator of the background rate $`b`$ assuming that the signal rate is known to be $`\mu _0.`$
$$\frac{}{b}\mathrm{log}l(\mu _0,b;x,y)=\frac{x}{\mu _0+b}1+\frac{y}{b}\tau 0$$
$$\widehat{b}=\frac{x+y(1+\tau )\mu _0+\sqrt{\left(x+y(1+\tau )\mu _0\right)^2+4(1+\tau )y\mu _0}}{2(1+\tau )}$$
$`l(\mu ,\widehat{b};x,y)`$ is called the profile likelihood function of $`\mu `$. The usefulness of the likelihood ratio test statistic lies in the fact that approximately we have
$$2\mathrm{log}\mathrm{\Lambda }(\mu _0;x,y)\chi ^2(d)$$
that is $`2\mathrm{log}\mathrm{\Lambda }`$ has an approximate Chi-Square distribution with $`d`$ degrees of freedom, where $`d`$ is the number of parameters in the model minus the number of parameters specified in the null hypothesis. Here we have $`d=1`$. Because the profile likelihood $`l`$ differs from the likelihood ratio $`\mathrm{\Lambda }`$ only by a constant independent of $`\mu _0`$, we can concentrate on the profile likelihood. For more details on the likelihood ratio test statistic see Casella and Berger Casella-Berger . For information on the profile likelihood see Bartlett Barlett , Lawley Lawley and Murphy and Van Der Vaart Murphy and Van Der Vaart . In figure 2 we have the profile likelihood function for the case $`x=6`$, $`y=2`$, $`\tau =2`$. To find a $`(1\alpha )100\%`$ confidence interval we start at the minimum, which of course is at the usual maximum likelihood estimator, and then move to the left and to the right until the function rises by the $`\alpha `$ percentile of a $`\chi ^2`$ distribution with $`1`$ degree of freedom.
The method here uses an approximation to the profile likelihood function by a quadratic function. Unfortunately, in cases where the number of observations in the signal region is small compared to the number of background events, the profile likelihood function becomes almost linear and this approximation does not work. For those cases we will use the following method:
We overlay the confidence region described previously with the profile likelihood curve $`(\mu ,\widehat{b})`$. Then we find the smallest value of $`\mu `$ that is on this curve but not in the confidence region. Figure 3 illustrates this method. Clearly, in the case of fewer observations in the signal region than in the background region only an upper bound will be quoted. We will use this second method whenever the profile likelihood function has a positive derivative at $`\mu =0`$. It turns out that the limits obtained by these two methods are compatible. Figure 4 shows the upper bound for a number of cases with the limits obtained by the confidence region method drawn as squares and the limits using the likelihood ratio method as diamonds. The transition from one method to the other is quite smooth.
## 3 Extensions of the Method
### 3.1 Estimating background from Monte Carlo
Our method can also be applied when a Monte Carlo with limited statistics is used to estimate the background rate. Assume we run the Monte Carlo $`n`$ times and observe a total of $`y`$ events. In the data we see $`x`$ events in the signal region. Then a probability model for this situation is given by:
$$XPois(\mu +b),YPois(nb)$$
We notice, of course, that this is actually the same model as the one previously used, only with an $`n`$ instead of a $`\tau `$. Therefore, we can use our method for this situation without any changes.
### 3.2 Include a second background source
Sometimes there is a second source of background present in the data, this one characterized by the fact that it only appears in the signal region. An example is an invariant mass histogram where some of the background comes from the misidentification of pions with muons. Say we run a Monte Carlo $`n`$ times and observe a total of $`z`$ events of this type. In the data we have $`x`$ observations in the signal region and $`y`$ observations in a suitably chosen background region. The probability model for this case is given by :
$$\begin{array}{c}XPois(\mu +b+\eta ),YPois(\tau b)\\ ZPois(n\eta )\end{array}$$
where $`\mu `$ is the signal rate, $`b`$ is the rate of the first background source and $`\eta `$ is the rate of the second background source. Then:
$$P_{\mu ,b,\eta }(X=x,Y=y,Z=z)=\frac{(\mu +b+\eta )^x}{x!}e^{(\mu +b+\eta )}\frac{(\tau b)^y}{y!}e^{\tau b}\frac{(n\eta )^z}{z!}e^{n\eta }$$
We can extend our method to this situation in a fairly straightforward manner. First, we have to find the profile likelihood function, which leads to the following nonlinear system of equations:
$$\frac{\mathrm{log}l}{b}=\frac{x}{\mu +b+\eta }1+\frac{y}{b}\tau =0$$
$$\frac{\mathrm{log}l}{\eta }=\frac{x}{\mu +b+\eta }1+\frac{z}{\eta }n=0$$
It turns out that $`\widehat{\eta }`$ can be found as the second largest root of the cubic equation $`ax^3+bx^2+cx+d=0`$ with
| $`a=(1+n)(n\tau )`$ |
| --- |
| $`b=\left(x+z(1+n)\mu \right)(\tau n)(1+n)(z+y)`$ |
| $`c=xz+(\tau n)z\mu +z^2+yz(1+n)z\mu `$ |
| $`d=\mu z^2`$ |
and then $`\widehat{b}`$ is given by
$$\widehat{b}=\frac{x+y(1+\tau )(\mu +\widehat{\eta })+\sqrt{\left(x+y(1+\tau )(\mu +\widehat{\eta })\right)^2+4(1+\tau )y(\mu +\widehat{\eta })}}{2(1+\tau )}$$
The same problem as before arises again in this situation: in the case of few events in the signal region the profile likelihood function is nearly linear. We can use the same remedy as before by instead finding the intersection of the boundary of the three-dimensional confidence region in $`(\mu ,b,\eta )`$ space with the profile likelihood curve $`(\mu ,\widehat{b},\widehat{\eta })`$.
## 4 Performance of this Method
In this section we will study the true coverage and the power of this method. For comparison we will use the unified approach of Feldman and Cousins Cousins-Feldman . Although that method was not designed to deal with uncertainty in the background, it is the standard method used in calculating confidence intervals in High Energy Physics at this time according to the Review of Particle Physics Particle Data Group . The coverage rates shown here were obtained through exact computation, without approximation. In figure 5 we have the case of one background region of equal size to the signal region, and finding $`90\%`$ confidence intervals. The true background rates are $`b=1,3`$ and $`5`$ and the signal rates go from $`0`$ to $`5.`$ Clearly our method has much better coverage than Feldman and Cousins Cousins-Feldman ; although, due to the approximation used in the method, the coverage is sometimes slightly worse than the nominal one. The ragged appearance of the graph is due to the discrete nature of the Poisson distribution and in general is unavoidable. Figure 6 has the case of two background regions and a $`99\%`$ confidence interval. Again the new method performs very well, whereas Feldman and Cousins Cousins-Feldman does not achieve the nominal coverage.
Next, we illustrate the performance of this method when the background rate is estimated by Monte Carlo. In figure 7 we compute the coverage rates for the cases $`\mu =0`$ and $`b=0.5,1.5`$ and $`2.5`$. The Monte Carlo is run $`n`$ times where we let $`n`$ range from $`1`$ to $`10`$. Our method performs satisfactorily for all cases. This plot is also interesting because it gives some insight into the number of runs of the Monte Carlo needed before one can assume that the background rate is known.
The performance of the extension of our method to the case of two different background sources discussed in section 3.2 has to be studied using mini Monte Carlo because the number of possible observations $`(x,y,z)`$ becomes quite large. We have run a variety of these mini Monte Carlo studies. In Figure 8 we show the results for the case $`\tau =2`$, $`n=10`$ and a $`90\%`$ confidence interval. The true coverage rates of our method appear to be in line with the nominal ones, again with a few cases where the coverage is slightly worse due to the approximation used in the method. Due to the discreteness of the Poisson distribution the method is also quite conservative for many situations.
As Feldman and Cousins Cousins-Feldman noted, there has been some criticism of their limits in cases where an experiment finds fewer events than are expected from the background rate. For example, an experiment with an expected background rate of $`b=0.5`$ and no observed events would quote a Feldman-Cousins upper limit of $`1.94`$; whereas an experiment with an expected background rate of $`b=1.5`$ and no observed events would quote a Feldman-Cousins upper limit of $`1.33`$. Our tables show a similar effect. Feldman and Cousins Cousins-Feldman explain this apparent inconsistency in some detail, and we are in full agreement with their reasoning.
It may also come as somewhat of a surprise that our method yields limits that are sometimes smaller than the limits in Feldman and Cousins Cousins-Feldman .Contrary to the instance discussed in the previous paragraph, this happens in some cases where more events are observed in the signal region than in the sidebands. For example, consider the following situation: say we use one sideband of equal size to the signal region, that is we have $`\tau =1`$, and we observe $`x=1`$ events in the signal region and $`y=0`$ events in the sideband. Then the $`90\%`$ upper limit using our method is found to be $`3.65`$, whereas assuming $`b=0`$ and using Feldman and Cousins method gives a $`90\%`$ upper limit of $`4.36`$, or about $`19\%`$ larger. So here we get a smaller limit despite the fact that we have an additional uncertainty. This apparent paradox can be explained as follows: Let us consider the question whether the single event observed in the signal region is a signal or a background event. Using Feldman and Cousins Cousins-Feldman the answer is clear: having assumed $`b=0`$ we know that there is no background, therefore this event has to be from the signal. In fact, it would be quite reasonable in this situation to actually quote a lower limit larger than $`0`$, because we already have proof that $`\mu >0`$. On the other hand, in our method $`y=0`$ does not imply $`b=0`$, it only means that $`b`$ is not too large. It is therefore still possible that the event in the signal region is in fact a background event, and that $`\mu =0`$. It should come as no surprise, then, that in those cases we are quoting a smaller upper limit for $`\mu `$ than Feldman and Cousins Cousins-Feldman . Notice that in theory it would still be preferable to have absolute knowledge of the background rate because one could then make an actual observation of the signal rate rather than having to settle for setting a limit. In practice it is usually impossible to have such precise knowledge and, of course, nobody would claim a discovery based on just one event.
## 5 Conclusion
The methods for quoting limits for rare decays, mainly Feldman and Cousins Cousins-Feldman and its variants, all suffer from the requirement that the background source be known with a high degree of precision. We have described a new method based on the likelihood ratio test which treats the background uncertainty as a statistical error. The performance of the method is shown to be quite good, with the true coverage rates close to the nominal ones and good power. The method can be used in situations where the background rate has been estimated from data sidebands as well as where it has been obtained from Monte Carlo. The method can also be extended to cases where a second background source solely appearing in the signal region is present.
In the appendix we provide tables for the confidence intervals and the experimental sensitivity as discussed in Feldman and Cousins Cousins-Feldman and in Review of Particle Physics Particle Data Group for the cases $`\tau =1,2`$ and $`\alpha =0.9,0.99`$. The currently used algorithm for computing the limits becomes unreliable in some extreme cases, and in those cases we use NA in the tables. A FORTRAN program for the computation of the limits for any other case as well as for the extension discussed in section 4.2. can be obtained by writing to w\_rolke@rumac.upr.clu.edu.
## 6 Appendix
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# Theorem 1
SIMULTANEOUS APPROXIMATION AND
INTERPOLATION OF FUNCTIONS
ON CONTINUA IN THE COMPLEX PLANE
VLADIMIR V. ANDRIEVSKII
IGOR E. PRITSKER
RICHARD S. VARGA
footnotetext: This research was supported in part by the National Science Foundation grant DMS-9707359.
## Abstract
We construct polynomial approximations of Dzjadyk type (in terms of the $`k`$-th modulus of continuity, $`k1`$) for analytic functions defined on a continuum $`E`$ in the complex plane, which simultaneously interpolate at given points of $`E`$. Furthermore, the error in this approximation is decaying as $`e^{cn^\alpha }`$ strictly inside $`E`$, where $`c`$ and $`\alpha `$ are positive constants independent of the degree $`n`$ of the approximating polynomial.
Key words: polynomial approximation, interpolation, analytic functions, quasiconformal curve.
AMS subject classification: 30E10, 41A10
1. Introduction and main results
Let $`E\text{C}`$ be a compact set with connected complement $`\mathrm{\Omega }:=\overline{\text{C}}E`$, where $`\overline{\text{C}}:=\text{C}\{\mathrm{}\}`$ is the extended complex plane. Denote by $`A(E)`$ the class of all functions continuous on $`E`$ and analytic in $`E^0`$, the interior of $`E`$ (the case $`E^0=\mathrm{}`$ is not excluded). Let $`\text{P}_n,n\text{N}_0:=\{0,1,2,\mathrm{}\},`$ be the class of complex polynomials of degree at most $`n`$. For $`fA(E)`$ and $`n\text{N}_0`$, define
$$E_n(f,E):=\underset{p\text{P}_n}{inf}fp_E,$$
where $`||||_E`$ denotes the uniform norm on $`E`$. By Mergelyan’s theorem (see ), we have that
$$\underset{n\mathrm{}}{lim}E_n(f,E)=0(fA(E)).$$
The following assertion on “simultaneous approximation and interpolation” quantifies a result of Walsh \[38, p. 310\]: Let $`z_1,\mathrm{},z_NE`$ be distinct points, $`fA(E)`$. Then for any $`n\text{N}:=\{1,2,\mathrm{}\},nN1`$, there exists a polynomial $`p_n\text{P}_n`$ such that
(1.1)
$$fp_n_EcE_n(f,E),$$
$$p_n(z_j)=f(z_j)(j=1,\mathrm{},N),$$
where $`c>0`$ is independent of $`n`$ and $`f`$.
A suitable polynomial has the form
$$p_n(z)=p_n^{}(z)+\underset{j=1}{\overset{N}{}}\frac{q(z)}{q^{}(z_j)(zz_j)}(f(z_j)p_n^{}(z_j)),$$
where
$$q(z):=\underset{j=1}{\overset{N}{}}(zz_j),$$
and $`p_n^{}\text{P}_n`$ satisfies
$$fp_n^{}_E=E_n(f,E).$$
It is natural to ask whether it is possible to interpolate the function $`f`$ as before at arbitrary prescribed points and to simultaneously approximate it in an even stronger sense than in (1.1). The theorem of Gopengauz about simultaneous polynomial approximation of real functions continuous on the interval $`[1,1]`$ and their interpolation at $`\pm 1`$ is an example of such result. For recent accounts of improvements and generalizations of this remarkable statement (for real functions) we refer the reader to , and .
We shall make use of the D-approximation (named after Dzjadyk, who found in the late 50’s - early 60’s a constructive description of Hölder classes requiring a nonuniform scale of approximation) as a substitute for (1.1). There is an extensive bibliography devoted to this subject (see, for example, the monographs , , , and ). In the overwhelming majority of the results on D-approximation, $`E`$ is a continuum (one of the rare exceptions is the recent interesting paper ). In it is shown that, for the D-approximation to hold for a continuum $`E`$, it is sufficient and under some mild restrictions also necessary that $`E`$ belongs to the class $`H^{}`$, which can be defined as follows (cf. and ).
From now on we assume that $`E`$ is a continuum with diam$`E>0`$, connected complement $`\mathrm{\Omega }`$ and boundary $`L:=E`$. In the sequel, we denote by $`\alpha ,\beta ,c,c_1,\mathrm{}`$ positive constants (possibly different at different occurences) that either are absolute or depend on parameters not essential for the arguments; otherwise, such a dependence will be indicated.
We say that $`EH`$ if any points $`z,\zeta E`$ can be joined by an arc $`\gamma (z,\zeta )E`$ whose length $`|\gamma (z,\zeta )|`$ satisfies the condition
(1.2)
$$|\gamma (z,\zeta )|c|z\zeta |,c=c(E)1.$$
Let us compactify the domain $`\mathrm{\Omega }`$ by prime ends in the Caratheodory sense (see ). Let $`\stackrel{~}{\mathrm{\Omega }}`$ be this compactification, and let $`\stackrel{~}{L}:=\stackrel{~}{\mathrm{\Omega }}\mathrm{\Omega }`$. Assuming that $`EH`$, then all the prime ends $`Z\stackrel{~}{L}`$ are of the first kind, i.e., they have singleton impressions $`|Z|=zL`$. The circle $`\{\xi :|\xi z|=r\},\mathrm{\hspace{0.17em}0}<r<\frac{1}{2}`$diam$`E`$, contains one arc, or finitely many arcs, dividing $`\mathrm{\Omega }`$ into two subdomains: an unbounded subdomain and a bounded subdomain such that $`Z`$ can be defined by a chain of cross-cuts of the bounded subdomain. Let $`\gamma _Z(r)`$ denote that one of these arcs for which the unbounded subdomain is as large as possible (for given $`Z`$ and $`r`$). Thus, the arc $`\gamma _Z(r)`$ separates the prime end $`Z`$ from $`\mathrm{}`$ (cf. , ).
If $`0<r<R<\frac{1}{2}`$ diam$`E`$, then $`\gamma _Z(r)`$ and $`\gamma _Z(R)`$ are the sides of some quadrilateral $`Q_Z(r,R)\mathrm{\Omega }`$ whose other two sides are parts of the boundary $`L`$. Let $`m_Z(r,R)`$ be the module of this quadrilateral, i.e., the module of the family of arcs that separate the sides $`\gamma _Z(r)`$ and $`\gamma _Z(R)`$ in $`Q_Z(r,R)`$ (see , ).
We say that $`EH^{}`$ if $`EH`$ and if there exist constants $`c=c(E)<\frac{1}{2}`$diam$`E`$ and $`c_1=c_1(E)`$ such that
(1.3)
$$|m_Z(|z\zeta |,c)m_𝒵(|z\zeta |,c)|c_1$$
for any pair of prime ends $`Z,𝒵\stackrel{~}{L},`$ with their impressions $`z=|Z|,\zeta =|𝒵|`$ satisfying $`|z\zeta |<c`$.
In particular, $`H^{}`$ includes domains with quasiconformal boundaries (see , ) and the classes $`B_k^{}`$ of domains introduced by Dzjadyk . For a more detailed investigation of the geometric meaning of conditions (1.2) and (1.3), see .
We will be studying functions defined by their $`k`$-th modulus of continuity $`(k\text{N})`$. There is a number of different definitions of these moduli in the complex plane (see , , , ). The definition by Dyn’kin is the most convenient for our purpose here.
From now on, suppose that $`EH^{}`$. Set
$$D(z,\delta ):=\{\zeta :|\zeta z|\delta \}(z\text{C},\delta >0).$$
The quantity
$$\omega _{f,k,z,E}(\delta ):=E_{k1}(f,ED(z,\delta )),$$
where $`fA(E),k\text{N},zE,\delta >0`$, is called the $`k`$-th local modulus of continuity, and
$$\omega _{f,k,E}(\delta ):=\underset{zE}{sup}\omega _{f,k,z,E}(\delta )$$
is called the $`k`$-th (global) modulus of continuity of $`f`$ on $`E`$. It is known (see ) that the behavior of this modulus is essentially the same as in the classical case of the interval $`E=[1,1]`$. In particular,
(1.4)
$$\omega _{f,k,E}(t\delta )ct^k\omega _{f,k,E}(\delta )(t>1,\delta >0).$$
We denote by $`A^r(E),r\text{N}`$, the class of functions $`fA(E)`$ which are $`r`$-times continuously differentiable on $`E`$, where we set $`A^0(E):=A(E)`$.
By definition, the function $`w=\mathrm{\Phi }(z)`$ maps $`\mathrm{\Omega }`$ conformally and univalently onto $`\mathrm{\Delta }:=\{w:|w|>1\}`$ and is normalized by the conditions
(1.5)
$$\mathrm{\Phi }(\mathrm{})=\mathrm{},\mathrm{\Phi }^{}(\mathrm{})>0.$$
The same symbol $`\mathrm{\Phi }`$ denotes the homeomorphism between the compactification $`\stackrel{~}{\mathrm{\Omega }}`$ of $`\mathrm{\Omega }`$ and $`\overline{\mathrm{\Delta }}`$, which coincides with $`\mathrm{\Phi }(z)`$ in $`\mathrm{\Omega }`$. Let $`\mathrm{\Psi }:=\mathrm{\Phi }^1`$. We define the distance to the level curves of $`\mathrm{\Phi }(z)`$
$$L_\delta :=\{\zeta :|\mathrm{\Phi }(\zeta )|=1+\delta \}(\delta >0)$$
by
$$\rho _\delta (z):=\text{dist}(z,L_\delta )(z\text{C},\delta >0),$$
where
$$\text{dist}(\zeta ,B):=inf\{|\zeta z|:zB\}(\zeta \text{C},B\text{C}).$$
###### Theorem 1
Let $`EH^{},fA(E),k\text{N}`$, and let $`z_1,\mathrm{},z_NE`$ be distinct points. Then for any $`n\text{N},nN+k`$, there exists a polynomial $`p_n\text{P}_n`$ such that
(1.6)
$$|f(z)p_n(z)|c_1\omega _{f,k,E}(\rho _{1/n}(z))(zL),$$
(1.7)
$$p_n(z_j)=f(z_j)(j=1,\mathrm{},N)$$
with $`c_1`$ independent of $`n`$.
Moreover, if $`E^0\mathrm{}`$ and if for any $`0<\delta <1`$, there is a constant $`c_2`$ such that
(1.8)
$$\underset{0}{\overset{\delta }{}}\omega _{f,k,E}(t)\frac{dt}{t}c_2\omega _{f,k,E}(\delta ),$$
then, in addition to (1.6) and (1.7),
(1.9)
$$fp_n_Kc_3\mathrm{exp}(c_4n^\alpha )$$
for every compact set $`KE^0`$, where the constants $`c_3,c_4`$ and $`0<\alpha 1`$ are independent of $`n`$.
A polynomial $`p_n`$ satisfying (1.6) is called a D-approximation of the function $`f`$ (D-property of $`E`$, Dzjadyk type direct theorem). For $`k>1`$, (1.6) generalizes the corresponding direct theorems of Belyi and Tamrazov (when $`E`$ is a quasidisk) and Shevchuk (when $`E`$ belongs to the Dzjadyk class $`B_k^{}`$). More detailed history can be found in these papers.
It was first noticed by Shirokov that the rate of D-approximation may admit significant improvement strictly inside $`E`$. Saff and Totik proved that if $`L`$ is an analytic curve, then an exponential rate is achievable strictly inside $`E`$, while on the boundary the approximation is “near-best”. However, even for domains with piecewise smooth boundary without cusps (and therefore belonging to $`H^{}`$), the error of approximation strictly inside $`E`$ cannot be better than $`e^{cn^\alpha }`$ (cf. (1.9)), where $`\alpha `$ may be arbitrarily small (see , ). In the results from , , containing estimates of the form (1.9), it is usually assumed that $`\mathrm{\Omega }`$ satisfies a wedge condition. For a continuum $`EH^{}`$, this condition can be violated.
Keeping in mind the Gopengauz result , we generalize Theorem 1 to the case of the Hermite interpolation and simultaneous approximation of a function $`fA^r(E)`$ and its derivatives. For simplicity we formulate and prove this assertion only for the case of boundary interpolation points and without the analog of (1.9).
###### Theorem 2
Let $`EH^{},fA^r(E),r\text{N},k\text{N}`$, and let $`z_1,\mathrm{},z_NE`$ be distinct points. Then for any $`n\text{N},nNr+k`$, there exists a polynomial $`p_n\text{P}_n`$ such that for $`l=0,\mathrm{},r`$,
(1.10)
$$|f^{(l)}(z)p_n^{(l)}(z)|c\rho _{1/n}^{rl}(z)\omega _{f^{(r)},k,E}(\rho _{1/n}(z))(zL),$$
and
(1.11)
$$p_n^{(l)}(z_j)=f^{(l)}(z_j)(j=1,\mathrm{},N),$$
with $`c`$ independent of $`n`$.
Our next goal is to allow the number of interpolation nodes $`N`$ to grow infinitely with the degree of approximating polynomial $`n`$. It is well known that we cannot take $`N1`$ equal to $`n`$, preserving uniform convergence (cf. Faber’s theorem claiming that for $`E=[1,1]`$ there is no universal set of nodes such that the Lagrange interpolating polynomials converge to every continuous function in uniform norm). However, it was first observed by Bernstein that for any continuous function on $`E=[1,1]`$ and any small $`\epsilon >0`$, there exists a sequence of polynomials interpolating in the Chebyshev nodes and uniformly convergent on $`[1,1]`$, such that $`n(1+\epsilon )N`$. This result was developed in several directions. In particular, Erdős (see and ) found a necessary and sufficient condition on the system of nodes, for this type of simultaneous approximation and interpolation to be valid.
We generalize the results of Bernstein and Erdős in the following Theorem. In order to accomplish this, we specify the choice of points $`z_1,\mathrm{},z_N`$ in an optimal fashion from the point of view of interpolation theory. Namely, we require that the discrete measure
$$\mu _N=\frac{1}{N}\underset{j=1}{\overset{N}{}}\delta _{z_j},$$
where $`\delta _z`$ denotes the unit mass placed at $`z`$, is close to the equilibrium measure for $`E`$ (for details, see ). Fekete points (see , ) are natural candidates for this purpose.
A Jordan curve is called quasiconformal if it is an image of the unit circle under a quasiconformal homeomorphism of the complex plane onto itself, with infinity as a fixed point (see for details).
###### Theorem 3
Let $`E`$ be a closed Jordan domain bounded by a quasiconformal curve $`L`$. Let $`f,r,k`$ be as in Theorem 1 and let $`z_1,\mathrm{},z_NE`$ be the points of an $`N`$-th Fekete point set of $`E`$. Then for any $`\epsilon >0`$ there exists a polynomial $`p_n\text{P}_n,n(1+\epsilon )N,`$ satisfying conditions (1.6) and (1.7). Moreover, if (1.8) holds then in addition to (1.6) and (1.7) we have (1.9), and the constants $`c_1,c_3,c_4`$ and $`\alpha `$ are independent of $`N`$.
2. Auxiliary results
In this section, we give some results from -, , which are needed for the proofs of the above theorems and which characterize the properties of the mappings $`\mathrm{\Phi }`$ and $`\mathrm{\Psi }`$ in the case $`EH^{}`$. For $`a>0`$ and $`b>0`$, we will use the expression $`ab`$ (order inequality) if $`acb`$. The expression $`ab`$ means that $`ab`$ and $`ba`$ simultaneously. The distance $`\rho _\delta (z)`$ to the level lines of $`\mathrm{\Phi }`$ is, for any $`zL`$, a normal majorant (in the terminology of ), i.e.,
(2.1)
$$\rho _{2\delta }(z)\rho _\delta (z)(\delta >0).$$
Let $`z,\zeta L,\delta >0`$. The condition $`|z\zeta |\rho _\delta (z)`$ yields
(2.2)
$$\rho _\delta (\zeta )\rho _\delta (z).$$
If $`L`$ is a quasiconformal curve, $`zL,\zeta \mathrm{\Omega }`$ and if $`|z\zeta |\rho _\delta (z)`$, then the inequality
(2.3)
$$\frac{\rho _\delta (z)}{|z\zeta |}\left(\frac{\delta }{|\mathrm{\Phi }(z)\mathrm{\Phi }(\zeta )|}\right)^\alpha $$
holds with some $`\alpha =\alpha (E)`$.
One of the fundamental problems that, as a rule, is encountered in the construction of approximations by polynomials, is the problem of approximating the Cauchy kernel $`1/(\zeta z),zE,\zeta \overline{\mathrm{\Omega }}`$, by polynomial kernels of the form
(2.4)
$$K_n(\zeta ,z)=\underset{j=0}{\overset{n}{}}a_j(\zeta )z^j.$$
The most general kernels of such type, the functions $`K_{r,m,k,n}(\zeta ,z)`$, were introduced by Dzjadyk (see \[13, Chapter 9\] or \[7, Chapter 3\]). Taking them as a basis for our discussion, we can establish the following result (cf. \[3, Lemma 9\]).
###### Lemma 1
Let $`EH^{}`$, and let $`m,r\text{N}`$. Then for any $`n\text{N}`$ there exists a polynomial kernel of the form (2.4) such that the following relations hold for $`l=0,\mathrm{},r`$, $`zL`$ and $`\zeta \overline{\mathrm{\Omega }}`$ with $`d(\zeta ,E)3`$:
$$\left|\frac{^l}{z^l}\left(\frac{1}{\zeta z}K_n(\zeta ,z)\right)\right|\frac{c_1}{|\zeta z|^{l+1}}\left(\frac{\rho _{1/n}(z)}{|\zeta z|+\rho _{1/n}(z)}\right)^m,$$
(2.5)
$$\left|\frac{^l}{z^l}K_n(\zeta ,z)\right|\frac{c_2}{(|\zeta z|+\rho _{1/n}(z))^{l+1}},$$
where $`c_j=c_j(m,r,E),j=1,2.`$
In order to improve the approximation properties of the polynomial kernel $`K_n(\zeta ,z)`$ inside of $`E`$, we use an idea from \[31, Theorem 2\], completing it by the following geometrical fact. Let
$$d(\zeta ,B):=\text{dist}(\zeta ,B)=inf\{|\zeta z|:zB\}(\zeta \text{C},B\text{C}).$$
###### Lemma 2
Let $`EH^{},E^0\mathrm{}`$. For any $`\zeta \overline{\mathrm{\Omega }}`$ with $`d(\zeta ,L)3`$, there exists a Jordan domain $`G_\zeta `$ with the following properties:
(i) $`\zeta G_\zeta ,E\overline{G_\zeta }`$;
(ii) $`\text{diam}G_\zeta c;`$
(iii) $`G_\zeta `$ is $`K`$-quasiconformal.
Here, the constants $`c>\text{diam}E`$ and $`K1`$ are independent of $`\zeta `$.
Proof. If $`\zeta \mathrm{\Omega }`$ we set $`𝒵:=\zeta `$; if $`\zeta L`$ we denote by $`𝒵\stackrel{~}{L}`$ the prime end whose impression coincides with $`\zeta `$ (or any of such prime ends). Let
$$\mathrm{\Gamma }_\zeta :=\{\xi \mathrm{\Omega }:\text{arg}\mathrm{\Phi }(\xi )=\text{arg}\mathrm{\Phi }(𝒵)\}.$$
By virtue of \[4, Lemma 1 and Lemma 2\],
(2.6)
$$d(z,L)|z\zeta |(z\mathrm{\Gamma }_\zeta ),$$
and for any $`z_1,z_2\mathrm{\Gamma }_\zeta `$ the length of the part of $`\mathrm{\Gamma }_\zeta `$ between these points satisfies
(2.7)
$$|\mathrm{\Gamma }_\zeta (z_1,z_2)||z_1z_2|.$$
A result of Rickman (see also \[7, p. 144\]) together with (2.7) imply that $`\mathrm{\Gamma }_\zeta `$ is $`K_1`$-quasiconformal with some $`K_11`$ independent of $`\zeta `$, i.e., there exists a $`K_1`$-quasiconformal mapping $`F:\overline{\text{C}}\overline{\text{C}}`$ such that
$$F(\zeta )=0,F(\mathrm{})=\mathrm{},F(\mathrm{\Gamma }_\zeta )=\{w:w>0\}.$$
We can assume that $`|F(z_0)|=1`$ for a fixed $`z_0E^0`$. We recall the following well-known property of quasiconformal automorphisms of the complex plane (see, for example, \[7, p. 98\]): If $`|\xi _1\xi _2||\xi _1\xi _3|`$ then
(2.8)
$$|F(\xi _1)F(\xi _2)||F(\xi _1)F(\xi _3)|$$
and vice versa.
According to (2.6) and (2.8) there are constants $`c_1`$ and $`c_2`$ such that
$$F(E)G_\zeta ^{}:=\{w=re^{i\theta }:\mathrm{\hspace{0.17em}0}r<c_1,c_2<|\theta |\pi \}.$$
By the Ahlfors criterion (see , \[20, p. 100\]), $`G_\zeta ^{}`$ is $`K_2`$-quasiconformal with $`K_2=K_2(c_1,c_2)1`$. Therefore, by (2.8) the domain $`G_\zeta :=F^1(G_\zeta ^{})`$ satisfies the conditions (i)-(iii) with $`K=K_1K_2`$.
$`\mathrm{}`$
Let $`E,\zeta `$ and $`G_\zeta `$ be as in Lemma 2 and let $`z_0E^0`$ be fixed. Consider the conformal mapping $`\mathrm{\Phi }_\zeta :\overline{\text{C}}\overline{G_\zeta }\mathrm{\Delta }`$ normalized as in (1.5), and the conformal mapping $`\varphi _\zeta :G_\zeta \{w:|w\frac{1}{2}|<\frac{1}{2}\}`$ normalized by the conditions
$$\varphi _\zeta (z_0)=\frac{1}{2},\varphi _\zeta (\zeta )=1.$$
Next, we use results from the theory of local distortion, under conformal mappings of an arbitrary simply connected domain onto a canonical one, developed by Belyi (see also ).
Lemma 2 as well as \[8, Theorem 1 and Theorem 6\] imply that the functions $`\mathrm{\Phi }_\zeta ^1`$ and $`\varphi _\zeta `$ satisfy a Hölder condition (with constants independent of $`\zeta `$). Therefore, by \[8, Theorem 4\] for any $`M\text{N}`$ there exists a polynomial $`t_M(\zeta ,z)\text{P}_M`$ (in $`z`$) such that
$$\varphi _\zeta t_M(\zeta ,)_{\overline{G_\zeta }}\frac{c_1}{M^\beta }$$
with some $`c_1`$ and $`\beta `$ independent of $`\zeta `$. We can assume that $`t_M(\zeta ,\zeta )=1`$.
Now for $`n\text{N}`$, we set
$$M:=\left[\frac{n^{1/(1+\beta )}}{2}\right],N:=[n^{\beta /(1+\beta )}]$$
(here $`[x]`$ denotes the Gauss bracket of $`x`$, the largest integer not exceeding $`x`$) and we note that, for the polynomial
$$u_{n/2}(\zeta ,z):=t_M^N(\zeta ,z),$$
the inequality
(2.9)
$$u_{n/2}(\zeta ,)_E\left(1+\frac{c_1}{M^\beta }\right)^N1$$
holds, as well as for any compact set $`KE^0`$ and $`\alpha :=\beta /(1+\beta )`$,
(2.10)
$$u_{n/2}(\zeta ,)_K(1c_2)^Ne^{cn^\alpha },$$
where the constants $`c_2<1`$ and $`c`$ are independent of $`\zeta `$.
Hence, the function defined by
$$T_n(\zeta ,z):=\frac{1u_{n/2}(\zeta ,z)}{\zeta z}+u_{n/2}(\zeta ,z)K_{[n/2]}(\zeta ,z),$$
where $`K_{[n/2]}(\zeta ,z)`$ is the polynomial kernel from Lemma 1, is a polynomial (in $`z`$) of degree at most $`n`$. According to Lemma 1, (2.9) and (2.10), it satisfies for $`\zeta \overline{\mathrm{\Omega }},d(\zeta ,L)3`$, arbitrary but fixed $`m\text{N}`$ and each compact set $`KE^0`$ the following conditions:
$$\left|\frac{1}{\zeta z}T_n(\zeta ,z)\right|=|u_{n/2}(\zeta ,z)|\left|\frac{1}{\zeta z}K_{[n/2]}(\zeta ,z)\right|$$
(2.11)
$$\{\begin{array}{cc}\frac{1}{|\zeta z|}\left(\frac{\rho _{1/n}(z)}{|\zeta z|+\rho _{1/n}(z)}\right)^m,\hfill & \text{ if }zL,\hfill \\ e^{cn^\alpha },\hfill & \text{ if }zK.\hfill \end{array}$$
In addition,
(2.12)
$$|T_n(\zeta ,z)|\frac{1}{|\zeta z|}(zE,\zeta \overline{\mathrm{\Omega }},d(\zeta ,L)3).$$
We will also need the continuous extension of an arbitrary function $`FA(E)`$ into the complex plane which preserves the smoothness properties of $`F`$. The corresponding construction, proposed by Dyn’kin , , is based on the Whitney partition of unity (see ) and local properties of the $`k`$-th modulus of continuity of $`F`$. A slight modification of the reasoning in , and gives the following result (cf. \[7, pp. 13-15\]).
###### Lemma 3
Let $`EH^{}`$. Any $`FA(E)`$ can be continuously extended to the complex plane (we preserve the notation $`F`$ for the extension) such that:
(i) $`F(z)=0`$ for $`z`$ with $`d(z,E)3`$, i.e., $`F`$ has compact support;
(ii) for $`z\text{C}E`$,
$$\left|\frac{F(z)}{\overline{z}}\right|c_1\frac{\omega _{F,k,z^{},E}(23d(z,E))}{d(z,E)},$$
where $`z^{}E`$ is an arbitrary point among those ones which are closest to $`z`$, $`c_1=c_1(k,\text{diam}E)`$;
(iii) if $`\zeta E,z\text{C},|z\zeta |<\delta ,\mathrm{\hspace{0.17em}0}<\delta <\frac{1}{2}`$diam$`E`$, then
$$|F(z)P_{F,k,\zeta ,E,\delta }(z)|c_2\omega _{F,k,\zeta ,E}(25\delta ),$$
where $`P_{F,k,\zeta ,E,\delta }(z)\text{P}_{k1}`$ is the (unique) polynomial such that
$$FP_{F,k,\zeta ,E,\delta }_{ED(\zeta ,\delta )}=\omega _{F,k,\zeta ,E}(\delta ),$$
and $`c_2=c_2(k)`$;
(iv) if $`F`$ satisfies a Lipschitz condition on $`E`$, i.e.,
$$|F(z)F(\zeta )|c|z\zeta |(z,\zeta E),$$
then the extension satisfies the same condition for $`z,\zeta \text{C}`$, with $`c_3=c_3(c,\text{diam}E,k)`$ instead of $`c`$.
3. Proof of Theorem 1
We fix a point $`z_0E`$ and consider a primitive of $`f`$:
(3.1)
$$F(\zeta ):=\underset{\gamma (z_0,\zeta )}{}f(\xi )𝑑\xi (\zeta E),$$
where $`\gamma (z_0,\zeta )E`$ is an arbitrary rectifiable arc joining $`z_0`$ and $`\zeta `$.
On writing for $`zL,\zeta E`$ with $`|\zeta z|\delta `$,
$`F(\zeta )`$ $`=`$ $`F(z)+{\displaystyle \underset{\gamma (z,\zeta )}{}}f(\xi )𝑑\xi `$
$`=`$ $`\nu _\delta (\zeta ,z)+{\displaystyle \underset{\gamma (z,\zeta )}{}}\left(f(\xi )P_{f,k,z,E,c\delta }(\xi )\right)𝑑\xi ,`$
where $`c1`$ is the constant from (1.2), we obtain
$$\omega _{F,k+1,z,E}(\delta )F\nu _\delta (,z)_{ED(z,\delta )}\delta \omega (\delta ),$$
where $`\omega (\delta ):=\omega _{f,k,E}(\delta )`$. Using Lemma 3, we can extend $`F`$ continuously to C, so that $`F`$ has compact support and satisfies
(3.2)
$$\left|\frac{F(\zeta )}{\overline{\zeta }}\right|\omega (d(\zeta ,L)),$$
for $`\zeta \mathrm{\Omega }^{}:=\{\zeta \overline{\mathrm{\Omega }}:d(\zeta ,L)3\}.`$ Moreover, for $`zL,\zeta \text{C}`$ with $`|z\zeta |\delta <\frac{1}{2}`$diam$`E`$, we have
(3.3)
$$|F(\zeta )\nu _\delta (\zeta ,z)|\delta \omega (\delta ).$$
Indeed, since for $`\zeta ED(z,\delta )`$,
$`|\nu _\delta (\zeta ,z)P_{F,k+1,z,E,\delta }(\zeta )|`$
$``$ $`|F(\zeta )\nu _\delta (\zeta ,z)|+|F(\zeta )P_{F,k+1,z,E,\delta }(\zeta )|\delta \omega (\delta ),`$
we have by the Bernstein-Walsh lemma \[38, p. 77\]
$$\nu _\delta (,z)P_{F,k+1,z,E,\delta }_{D(z,\delta )}\delta \omega (\delta ).$$
Hence (3.3) follows from the last inequality and assertion (iii) of Lemma 3.
Next, we consider the most complicated case, that is, $`E^0\mathrm{}`$ and (1.8) holds. We introduce the polynomial kernel $`Q_{n/2}(\zeta ,z):=T_{[n/2]}(\zeta ,z)`$, which by (2.11) and (2.12) satisfies
(3.4)
$$\frac{1}{\zeta }Q_{n/2}(\zeta ,)_Ke^{cn^\alpha }(\zeta \mathrm{\Omega }^{})$$
on each compact set $`KE^0`$, and
(3.5)
$$\left|\frac{1}{\zeta z}Q_{n/2}(\zeta ,z)\right|\frac{1}{|\zeta z|}\left(\frac{\rho _{1/n}(z)}{|\zeta z|+\rho _{1/n}(z)}\right)^k(zL),$$
(3.6)
$$|Q_{n/2}(\zeta ,z)|\frac{1}{|\zeta z|}(zE).$$
Further, we consider the polynomial
$$t_n(z)=\frac{1}{\pi }\underset{\mathrm{\Omega }^{}}{}\frac{F(\zeta )}{\overline{\zeta }}Q_{n/2}^2(\zeta ,z)𝑑m(\zeta )(zE),$$
where $`dm(\zeta )`$ means integration with respect to the two-dimensional Lebesgue measure (area). Let $`zL,D:=D(z,\rho ),\sigma :=D,\rho :=\rho _{1/n}(z)`$. According to assertion (iv) of Lemma 3, $`F`$ is an ACL-function (absolutely continuous on lines parallel to the coordinate axes) in C. Hence Green’s formula can be applied here (see ) to obtain
(3.7) $`f(z)t_n(z)`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle \underset{\mathrm{\Omega }^{}D}{}}{\displaystyle \frac{F(\zeta )}{\overline{\zeta }}}\left(Q_{n/2}^2(\zeta ,z){\displaystyle \frac{1}{(\zeta z)^2}}\right)𝑑m(\zeta )`$
$`+`$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle \underset{D}{}}{\displaystyle \frac{F(\zeta )}{\overline{\zeta }}}Q_{n/2}^2(\zeta ,z)𝑑m(\zeta )`$
$`+`$ $`f(z){\displaystyle \frac{1}{2\pi i}}{\displaystyle \underset{\sigma }{}}{\displaystyle \frac{F(\zeta )}{(\zeta z)^2}}𝑑\zeta `$
$`=`$ $`U_1(z)+U_2(z)+U_3(z).`$
The first two integrals in (3.7) can be estimated in an appropriate way by passing to polar coordinates and using (1.4), (1.8), (3.2), (3.5) as well as (3.6):
(3.8)
$$|U_1(z)|\underset{\rho }{\overset{c}{}}\omega (t)\frac{\rho ^{k+1}}{t^{k+2}}𝑑t\omega (\rho )\rho \underset{\rho }{\overset{c}{}}\frac{dt}{t^2}\omega (\rho ),$$
(3.9)
$$|U_2(z)|\underset{0}{\overset{\rho }{}}\frac{\omega (t)}{t}𝑑t\omega (\rho ).$$
In order to estimate the third term in (3.7), we note that
$$|f(z)(\nu _\rho )_\zeta ^{}(z,z)|=|f(z)P_{f,k,z,E,c\rho }(z)|\omega (c\rho )\omega (\rho ),$$
so that by (3.3):
(3.10)
$$|U_3(z)||f(z)(\nu _\rho )_\zeta ^{}(z,z)|+\frac{1}{2\pi }\left|\underset{\sigma }{}\frac{F(\zeta )\nu _\rho (\zeta ,z)}{(\zeta z)^2}𝑑\zeta \right|\omega (\rho ).$$
Comparing (3.7)-(3.10), we obtain that
(3.11)
$$|f(z)t_n(z)|\omega (\rho _{1/n}(z))(zL).$$
The estimate
(3.12)
$$ft_n_Ke^{cn^\alpha },$$
for any compact set $`KE^0`$, follows immediately from (3.2) and (3.4) by a straight-forward modification of the above reasoning.
To satisfy the interpolation condition (1.7), we argue as follows. Let $`n>2N`$. We consider the polynomials
$$V_{n/2+1}(\zeta ,z):=\{\begin{array}{cc}1(\zeta z)Q_{n/2}(\zeta ,z),\hfill & \text{ if }\zeta L,zE,\hfill \\ 1,\hfill & \text{ if }\zeta E^0,zE,\hfill \end{array}$$
and
$$u_n(z):=\underset{j=1}{\overset{N}{}}\frac{q(z)}{q^{}(z_j)(zz_j)}(f(z_j)t_n(z_j))V_{n/2+1}(z_j,z).$$
By (3.4), (3.5), (3.11) and (3.12),
$$|u_n(z)|\{\begin{array}{cc}\underset{j}{\overset{}{}}\omega (\rho _{1/n}(z_j))\left(\frac{\rho _{1/n}(z)}{|zz_j|+\rho _{1/n}(z)}\right)^k,\hfill & \text{ if }zL,\hfill \\ e^{cn^\alpha },\hfill & \text{ if }zK,\hfill \end{array}$$
where $`_j^{}`$ means the sum in all $`j`$ with $`z_jL`$. To show that
$$p_n(z):=t_n(z)+u_n(z)$$
satisfies (1.6), (1.7) and (1.9), it is sufficient to prove that the inequality
(3.13)
$$\omega (\rho _{1/n}(\zeta ))\left(\frac{\rho _{1/n}(z)}{|z\zeta |+\rho _{1/n}(z)}\right)^k\omega (\rho _{1/n}(z))$$
holds for any $`z,\zeta L`$.
This relation is trivial if $`|\zeta z|\rho _{1/n}(\zeta )`$ (cf. (2.2)). Hence we may assume that $`|\zeta z|>\rho _{1/n}(\zeta )`$. Then by (1.4),
$$\omega (\rho _{1/n}(\zeta ))\left(\frac{\rho _{1/n}(z)}{|z\zeta |+\rho _{1/n}(z)}\right)^k\omega (|\zeta z|)\left(\frac{\rho _{1/n}(z)}{|\zeta z|}\right)^k\omega (\rho _{1/n}(z)),$$
which completes the proof of (3.13).
Note that we used assumption (1.8) only for the estimation of $`U_2(z)`$ in (3.9). If we are interested only in relations (1.6) and (1.7), then we need to choose in the above reasoning $`Q_{n/2}(\zeta ,z)=K_{[n/2]}(\zeta ,z)`$, where $`K_n(\zeta ,z)`$ is the polynomial kernel from Lemma 1. Then, instead of (3.9), we obtain by (2.5) that
$$|U_2(\zeta ,z)|\underset{0}{\overset{\rho }{}}\omega (t)\frac{tdt}{\rho ^2}\omega (\rho ),$$
and (1.8) becomes superfluous.
$`\mathrm{}`$
4. Proof of Theorem 2
Since the scheme of this proof is the same as in the proof of Theorem 1, we describe it only briefly. We begin with the Taylor formula for a primitive $`F`$ defined by (3.1):
$$F(\zeta )=F(z)+\underset{j=1}{\overset{r}{}}\frac{f^{(j1)}(z)}{j!}(\zeta z)^j+\frac{1}{r!}\underset{\gamma (z,\zeta )}{}(\zeta \xi )^rf^{(r)}(\xi )𝑑\xi ,$$
where $`z,\zeta E`$ and an arc $`\gamma (z,\zeta )E`$ joins these points and satisfies (1.2). Therefore, we have for $`zL,\zeta E`$ with $`|z\zeta |\delta `$,
$$F(\zeta )=\kappa _\delta (\zeta ,z)+\frac{1}{r!}\underset{\gamma (z,\zeta )}{}(z\xi )^r\left(f^{(r)}(\xi )P_{f^{(r)},k,z,E,c\delta }(\xi )\right)𝑑\xi ,$$
where $`c1`$ is the constant from (1.2) and $`\kappa _\delta (\zeta ,z)`$ is a polynomial (in $`\zeta `$) of degree $`k+r`$. Using Lemma 3, we extend $`F`$ continuously, so that $`F`$ has compact support and satisfies
$$\left|\frac{F(\zeta )}{\overline{\zeta }}\right|d(\zeta ,L)^r\omega (d(\zeta ,L))(\zeta \mathrm{\Omega }^{}:=\{\zeta \overline{\mathrm{\Omega }}:d(\zeta ,L)3\}),$$
$$|F(\zeta )\kappa _\delta (\zeta ,z)|\delta ^{r+1}\omega (\delta )(zL,\zeta \text{C},|\zeta z|\delta ),$$
where $`\omega (\delta ):=\omega _{f^{(r)},k,z,E}(\delta )`$.
Next, we introduce the polynomial
$$t_n(z)=\frac{1}{\pi }\underset{\mathrm{\Omega }^{}}{}\frac{F(\zeta )}{\overline{\zeta }}\frac{}{z}K_n(\zeta ,z)𝑑m(\zeta )(zE),$$
where $`K_n(\zeta ,z)`$ is the polynomial kernel from Lemma 1 (with $`m=2r`$).
Let $`l=0,\mathrm{},r`$ and let $`z,D`$ as well as $`\sigma `$ be the same as in (3.7).
By Green’s formula, we have that
$`f^{(l)}(z)t_n^{(l)}(z)`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle \underset{\mathrm{\Omega }^{}D}{}}{\displaystyle \frac{F(\zeta )}{\overline{\zeta }}}{\displaystyle \frac{^{l+1}}{z^{l+1}}}\left(K_n(\zeta ,z){\displaystyle \frac{1}{\zeta z}}\right)𝑑m(\zeta )`$
$`+`$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle \underset{D}{}}{\displaystyle \frac{F(\zeta )}{\overline{\zeta }}}{\displaystyle \frac{^{l+1}}{z^{l+1}}}K_n(\zeta ,z)𝑑m(\zeta )`$
$`+`$ $`f^{(l)}(z){\displaystyle \frac{1}{2\pi i}}{\displaystyle \underset{\sigma }{}}F(\zeta ){\displaystyle \frac{^{l+1}}{z^{l+1}}}{\displaystyle \frac{1}{\zeta z}}𝑑\zeta .`$
Reasoning as in the proof of (3.11), we obtain that
(4.1)
$$|f^{(l)}(z)t_n^{(l)}(z)|\rho _{1/n}^{rl}(z)\omega (\rho _{1/n}(z))(zL).$$
Further, we assume that $`n>2N(r+1)`$ and introduce the auxiliary polynomials
$$V_{n/2}(\zeta ,z):=1+\frac{(\zeta z)^{r+1}}{r!}\frac{^r}{z^r}K_{[n/2]}(\zeta ,z)$$
and
$$u_n(z):=\underset{j=1}{\overset{N}{}}\frac{q^{r+1}(z)}{(zz_j)^{r+1}}V_{n/2}(z_j,z)\underset{s=0}{\overset{r}{}}A_{j,s}(zz_j)^s,$$
where
$$A_{j,s}:=\underset{\nu =0}{\overset{s}{}}\frac{1}{\nu !(s\nu )!}\left(f^{(\nu )}(z_j)t_n^{(\nu )}(z_j)\right)\left(\frac{^{s\nu }}{z^{s\nu }}\frac{(zz_j)^{r+1}}{q^{r+1}(z)}\right)|_{z=z_j}.$$
According to the Hermite interpolation formula (see ), we have
$$u_n^{(l)}(z_j)=f^{(l)}(z_j)t_n^{(l)}(z_j)(j=1,\mathrm{},N).$$
Therefore the polynomial
$$p_n:=u_n+t_n$$
satisfies the interpolation condition (1.11).
Since
$$|A_{j,s}|\rho _{1/n}^{rs}(z_j)\omega (\rho _{1/n}(z_j)),$$
we obtain by Lemma 1 for any $`zL`$,
(4.2) $`|u_n(z)|`$ $``$ $`{\displaystyle \underset{j=1}{\overset{N}{}}}\left({\displaystyle \frac{\rho _{1/n}(z)}{|zz_j|+\rho _{1/n}(z)}}\right)^{2r}{\displaystyle \underset{s=0}{\overset{r}{}}}\rho _{1/n}^{rs}(z_j)\omega (\rho _{1/n}(z_j))|zz_j|^s`$
$``$ $`\rho _{1/n}^r(z)\omega (\rho _{1/n}(z)),`$
where we used (2.2) and the following inequality: for $`z,\zeta L`$ with $`|\zeta z|\rho _{1/n}(z)`$,
$$\left|\frac{\rho _{1/n}(z)}{z\zeta }\right|^{2r}|z\zeta |^r\omega (|z\zeta |)\rho _{1/n}^r(z)\omega (\rho _{1/n}(z)).$$
By a theorem of Tamrazov (see also \[7, p. 187\]), (4.2) yields
(4.3)
$$|u_n^{(l)}(z)|\rho _{1/n}^{rl}(z)\omega (\rho _{1/n}(z)).$$
Combining (4.1) and (4.3), we obtain (1.10).
$`\mathrm{}`$
5. Proof of Theorem 3
We use the same scheme as in the proof of Theorem 1. Let (1.8) hold. We construct a polynomial $`t_N\text{P}_N`$ such that
(5.1)
$$|f(z)t_N(z)|\omega (\rho _{1/N}(z))(zL),$$
where $`\omega (\delta ):=\omega _{f,k,E}(\delta )`$, and
(5.2)
$$ft_N_Ke^{cN^\alpha }$$
for any compact set $`KE^0`$.
Let $`m:=[\epsilon N]`$. Consider the polynomial
$$u_{N+m}(z):=\underset{j=1}{\overset{N}{}}\frac{q(z)}{q^{}(z_j)(zz_j)}(f(z_j)t_N(z_j))V_{m+1}(z_j,z),$$
where
$$V_{m+1}(\zeta ,z):=1(\zeta z)Q_m(\zeta ,z)(\zeta L,zE),$$
and $`Q_m(\zeta ,z):=T_m(\zeta ,z)`$ is a polynomial of degree at most $`m`$ (in $`z`$) satisfying the inequalities (cf. (2.11))
(5.3)
$$\left|\frac{1}{\zeta z}Q_m(\zeta ,z)\right|\frac{1}{|\zeta z|}\left(\frac{\rho _{1/m}(z)}{|\zeta z|+\rho _{1/m}(z)}\right)^{k+l}(z,\zeta L)$$
(the choice of $`l=l(E)>0`$ will be specified below) and
(5.4)
$$\frac{1}{\zeta }Q_m(\zeta ,)_Ke^{cm^\alpha }(\zeta L)$$
on each compact set $`KE^0`$.
Let $`zL,\mathrm{\Phi }(z)=e^{i\theta _0},\mathrm{\Phi }(z_j)=e^{i\theta _j},`$
$$0\theta _1<\theta _2<\mathrm{}<\theta _N<\theta _{N+1}:=\theta _1+2\pi .$$
It is proved in that
(5.5)
$$|\theta _{j+1}\theta _j|\frac{1}{N}(j=1,\mathrm{},N).$$
We rename the points $`\{e^{i\theta _j}\}_1^N`$ by $`\{e^{i\theta _j^{}}\}_1^\mu `$, $`\{e^{i\theta _j^{\prime \prime }}\}_1^\nu `$ and $`\{e^{i\theta _j^{\prime \prime \prime }}\}_1^{N\mu \nu }`$ in such a way that
$$|\theta _0\theta _j^{}|\frac{1}{m}(j=1,\mathrm{},\mu ),$$
and $`\theta _j=\theta _j^{\prime \prime },\theta _j^{\prime \prime \prime }`$ satisfy
$$|\theta _0\theta _j|>\frac{1}{m},(\theta _j\{\theta _1^{},\mathrm{},\theta _\mu ^{}\}),$$
$$\theta _0<\theta _1^{\prime \prime }<\theta _2^{\prime \prime }<\mathrm{}<\theta _\nu ^{\prime \prime }\pi +\theta _0,$$
$$\theta _0\pi <\theta _{N\mu \nu }^{\prime \prime \prime }<\mathrm{}<\theta _1^{\prime \prime \prime }<\theta _0.$$
Equation (5.5) implies that
$$\mu \frac{1}{\epsilon },\nu N\mu \nu N.$$
Furthermore, for the function
$$h(\theta ,\theta _0):=(f(\mathrm{\Psi }(e^{i\theta }))t_N(\mathrm{\Psi }(e^{i\theta })))V_{m+1}(\mathrm{\Psi }(e^{i\theta }),\mathrm{\Psi }(e^{i\theta _0}))$$
we have by (1.4), (2.2), (5.1) and (5.3),
(5.6)
$$|h(\theta _j^{},\theta _0)|\omega (\rho ),$$
(5.7)
$$|h(\theta _j^{\prime \prime },\theta _0)|\omega (|zz_j^{\prime \prime }|)\left(\frac{\rho }{|zz_j^{\prime \prime }|}\right)^{k+l}\omega (\rho )\left(\frac{\rho }{|zz_j^{\prime \prime }|}\right)^l,$$
(5.8)
$$|h(\theta _j^{\prime \prime \prime },\theta _0)|\omega (|zz_j^{\prime \prime \prime }|)\left(\frac{\rho }{|zz_j^{\prime \prime \prime }|}\right)^{k+l}\omega (\rho )\left(\frac{\rho }{|zz_j^{\prime \prime \prime }|}\right)^l,$$
where $`\rho :=\rho _{1/m}(z),z_j^{\prime \prime }:=\mathrm{\Psi }(e^{i\theta _j^{\prime \prime }}),z_j^{\prime \prime \prime }:=\mathrm{\Psi }(e^{i\theta _j^{\prime \prime \prime }})`$.
It follows from (5.2) and (5.4) that the polynomial
(5.9)
$$p_{[(1+\epsilon )N]}(z):=t_N(z)+u_{N+m}(z)$$
satisfies (1.7) and (1.9).
We choose $`l`$ so that
$$\left|\frac{\rho }{\zeta z}\right|^l\left(\frac{1}{m|\mathrm{\Phi }(\zeta )\mathrm{\Phi }(z)|}\right)^2,$$
for $`\zeta L`$ with $`|\zeta z|>\rho `$ (cf. (2.3)).
Since
$`|u_{N+m}(z)|`$ $``$ $`{\displaystyle \underset{j=1}{\overset{\mu }{}}}|h(\theta _j^{},\theta _0)|+{\displaystyle \underset{j=1}{\overset{\nu }{}}}|h(\theta _j^{\prime \prime },\theta _0)|+{\displaystyle \underset{j=1}{\overset{N\mu \nu }{}}}|h(\theta _j^{\prime \prime \prime },\theta _0)|`$
$``$ $`\omega (\rho )\left(1+{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{1}{j^2}}\right)\omega (\rho ),zL,`$
by (5.6)-(5.8), we obtain the desired inequality (1.6) by (2.1) and (5.1), for $`p_{[(1+\epsilon )N]}`$ given by (5.9).
Taking in the above argument $`Q_m(\zeta ,z):=K_m(\zeta ,z)`$, we obtain equations (1.6) and (1.7) even without assumption (1.8).
$`\mathrm{}`$
GSF-Forschungszentrum, Institut für Biomathematik und Biometrie, Ingolstädter Landstr. 1, D-85764 Neuherberg, Germany;
e-mail: mgk002@eo-dec-mathsrv.ku-eichstaett.de
Department of Mathematics, 401 Mathematical Sciences, Oklahoma State University, Stillwater, OK 74078-1058, U.S.A.;
e-mail: igor@math.okstate.edu
Institute for Computational Mathematics, Kent State University, Kent, OH 44242, U.S.A.;
e-mail: varga@mcs.kent.edu
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# 1 Introduction
## 1 Introduction
The naturalness of the standard model (SM) <sup>?</sup> may require the existence of low energy supersymmetry (SUSY) <sup>?</sup>. The flavor puzzle, namely the fermion masses, mixings and CP violation, of the SM needs new physics to be understood. It would be nice if SUSY also provides (partial) understanding to the flavor problem <sup>?</sup>.
Let us look at the fermion mass pattern. The fact is that the third generation is much heavier than the second generation which is also much heavier than the first. Does this imply a family symmetry? We assume that the answer is yes. Let us consider the charged leptons. By assuming a $`Z_3`$ cyclic symmetry among the SU(2) doublets $`L_i`$ ($`i=1,2,3`$) of the three generations <sup>?</sup>, the Yukawa interactions result in a democratic mass matrix. This mass matrix is of rank $`1`$. Therefore only the tau lepton gets mass, the muon and electron are still massless.
The essential point is how the family symmetry breaks. Naively the symmetry breaking can be achieved by introducing more Higgs fields. We consider this problem within SUSY. We have observed that SUSY naturally provides Higgs like fields, which are the scalar neutrinos. Furthermore, if the vacuum expectation values (VEVs) of the sneutrinos are non-vanishing, $`v_i0`$, the R-parity violating interactions $`L_iL_jE_k^c`$, with $`E_k^c`$ denoting the anti-particle superfields of the SU(2) singlet leptons, contribute to the fermion masses, in addition to the Yukawa interactions. This made us to propose that the family symmetry is broken by the sneutrino VEVs <sup>?</sup>.
A remark should be made here. The sneutrino VEV is not enough to the electroweak symmetry breaking (EWSB), the Higgs fields are still necessary. As will be seen, $`v_i`$ is typically ($`510`$) GeV. Such a VEV however, is too large to accommodate the neutrino oscillation data. The point we made is that such a sneutrino VEV breaks a family symmetry which can be useful for the understanding of the charged lepton masses.
Let us focus on the lepton sector. The $`Z_3`$ family symmetry results in that Yukawa interactions only give $`\tau `$ lepton mass, $`m_\tau yv_d`$, where $`y`$ is the coupling constant, and the Higgs VEV $`v_dv_u100`$ GeV. Taking $`y10^2`$ gives realistic $`m_\tau `$. The muon mass is due to the family symmetry breaking $`v_i0`$. To make the model phenomenologically acceptable, we found $`v_310`$ GeV, and the trilinear R-parity violating couplings $`\lambda 10^2`$. The muon mass is then $`m_\mu \lambda v_3100`$ MeV.
## 2 Challenging Problems
Immediately the following questions are raised to challenge the above described scenario:
1. Is $`10`$ GeV $`v_3`$ safe? Namely do we have unacceptable Majoron?
2. Is $`m_{\nu _\tau }`$ too large? Roughly it is expected to be $`g_2^2v_3^2/M_Z100`$ MeV $`1`$ GeV, with $`g_2`$ being the coupling of the SU(2) interaction.
3. Is $`\lambda 10^2`$ contradict with the experimental data on the rare decays $`\mu 3e`$ and $`\mu e\gamma `$?
To get rid of the above difficulties, the following things are assumed respectively. First, a large (weak scale) mixing mass term of the Higgs scalar and slepton $`\stackrel{~}{L}_3`$, $`B\mu _3\stackrel{~}{H}_u\stackrel{~}{L}_3`$ is introduced. This term breaks SUSY softly. Because it breaks the lepton number explicitly, no massless Majoron which would be the Goldstone particle corresponding to the spontaneous lepton number violation, appears. Note that the trilinear R-parity violating interactions themselves are not enough to keep the would be Majoron from being light. Second, in the superpotential, the bilinear mixing mass terms, like $`\mu _3H_uL_3`$ and $`\mu H_uH_d`$ should be small. We take $`\mu _3=0`$ and $`\mu =0`$. In this way, $`m_{\nu _\tau }=0`$ at tree level. This may avoid the difficulty of Question 2. Note that the EWSB is achieved by an alternative superpotential. Third, the family symmetry $`Z_3`$ in the ordinary Yukawa and trilinear R-parity violating interactions is adopted. Because of this, $`\mu 3e`$ and $`\mu e\gamma `$ do not occur. This symmetry can be regarded as accidental at low energy.
The relevant superpotential is
$$\begin{array}{ccc}𝒲\hfill & =\hfill & y_j\left(\underset{i}{\overset{3}{}}L_i\right)H_dE_j^c+\lambda _j\left(L_1L_2+L_2L_3+L_3L_1\right)E_j^c\hfill \\ & & +\lambda ^{}X\left(H_uH_d\mu ^2\right),\hfill \end{array}$$
(1)
where the last term is for the EWSB with $`\lambda ^{}`$ being the coupling constant, $`X`$ a singlet superfield and $`\mu `$ the weak scale. The baryon number conservation is assumed <sup>?</sup>. Eq. (1) reproduces the lepton masses as we have planned. It would be easier to work in the mass eigenstates of Yukawa interactions,
$$𝒲=y^\tau L_\tau H_dE_\tau ^c+L_eL_\mu \left(\lambda _\mu E_\mu ^c+\lambda _\tau E_\tau ^c\right)+\lambda ^{}X\left(H_uH_d\mu ^2\right),$$
(2)
where the left-handed leptons are
$$\begin{array}{ccc}\tau _L\hfill & =\hfill & \frac{1}{\sqrt{3}}\left(e_1+e_2+e_3^{}\right),\hfill \\ \mu _L\hfill & =\hfill & \frac{1}{\sqrt{2}}\left(e_1e_2\right),\hfill \\ e_L\hfill & =\hfill & \frac{1}{\sqrt{6}}\left(e_1+e_22e_3^{}\right),\hfill \end{array}$$
(3)
with ($`\nu _i`$, $`e_i`$) being the fermionic component of $`L_i`$. ($`\nu _i^{}`$, $`e_i^{}`$) means the physical leptons after considering the mixing with the neutralinos and charginos. From Eq. (2), we see that $`\mu 3e`$ does not occur. Instead, the rare decays $`\tau 2e\mu `$ and $`3\mu `$ have branching ratios $`10^7`$ if $`m_{\stackrel{~}{\nu }_i}100`$ GeV.
At the quantum level, a comparatively large neutrino mass is inevitably induced due to the large lepton number violating effect in $`B\mu _3`$ mass. It occurs at the one-loop level with a Zino exchange. Therefore $`m_{\nu _\tau }0`$ at one loop which will be studied further. Is it natural in a theory in which $`B\mu _3`$ is large, $`\mu _3`$ is vanishingly small and $`m_{\nu _\tau }`$ is consistent with experiment?
## 3 A Model of GMSB
With the framework of gauge mediated SUSY breaking (GMSB) <sup>?</sup>, the scenario asked in the last question can be realized naturally <sup>?</sup>. Lepton number violation is introduced originally in the messenger sector. It is then communicated to the SM sector including the related soft SUSY breaking terms. We will make use of the observation of the $`\mu `$-problem in GMSB <sup>?</sup>. It was noted that both $`\mu `$ term and its corresponding soft breaking $`B\mu `$ term can be generated at one loop. Either $`\mu `$ is at the weak scale and $`B\mu `$ is unnaturally large, or $`B\mu `$ is at the weak scale and $`\mu `$ is very small. This is not a problem in our model, because the EWSB given in Eq. (1) does not need the $`\mu `$ term. However, we apply similar observation to the discussion of the mixing of $`H_u`$ and $`L_3`$.
The messengers are introduced as follows with the SU(3)$`\times `$SU(2)$`\times `$U(1) quantum numbers,
$$S,S^{}=(1,2,1),\overline{S},\overline{S^{}}=(1,2,1),$$
(4)
and
$$T,T^{}=(3,1,2/3),\overline{T},\overline{T^{}}=(\overline{3},1,2/3).$$
(5)
Two more gauge singlets are introduced $`Y`$ for the SUSY breaking and $`V`$ for the lepton number violation. The superpotential is then
$$𝒲_{\mathrm{total}}=𝒲+𝒲_1+𝒲_2,$$
(6)
where
$$\begin{array}{ccc}𝒲_1\hfill & =\hfill & m_1\left(\overline{S^{}}S+S^{}\overline{S}\right)+m_2\left(\overline{T^{}}T+T^{}\overline{T}\right)+m_3S\overline{S}+m_4T\overline{T}+m_5V^2\hfill \\ & & +Y\left(\lambda _1S\overline{S}+\lambda _2T\overline{T}+\lambda _3V^2\mu _1^2\right),\hfill \end{array}$$
$$𝒲_2=V\left(\lambda _5H_uS+\lambda _6L_3\overline{S}\right),$$
(7)
with $`\mu _1`$ being the SUSY breaking scale. The lepton number violation lies in $`𝒲_2`$. By integrating out the heavy messengers, the effective Lagrangian related to lepton number violation is
$$_{\mathrm{eff}}^{\mathit{}}=\mu _3L_3H_u|_{\theta \theta }+B\mu _3\stackrel{~}{L_3}\stackrel{~}{H_u}+\mathrm{h}.\mathrm{c}.,$$
(8)
where
$$\mu _3\frac{\lambda _5\lambda _6}{16\pi ^2}\frac{\mu _1^2}{m_3},B\mu _3\frac{\lambda _5\lambda _6}{16\pi ^2}\left(\frac{\mu _1^2}{m_3}\right)^2.$$
(9)
Or
$$B\mu _3=\mu _3\frac{\mu _1^2}{m_3}.$$
(10)
This is what we needed. $`\mu _1^2/m_3`$ is usually around $`100`$ TeV. $`\mu _3`$ is very small if $`B\mu _3`$ is taken to be the weak scale.
From the scalar potential, the sneutrino VEV is obtained as
$$v_1=0,v_2=0,v_3=\frac{B\mu _3v_u}{M_A^2+\frac{1}{2}M_Z^2\mathrm{cos}2\beta }.$$
(11)
Numerically, $`v_310`$ GeV if $`M_A300`$ GeV and $`B\mu _3\left(60\mathrm{GeV}\right)^2`$. Nonzero $`v_3`$ implies the mixing between neutrino $`\nu _3`$ and the neutralinos. In addition, $`B\mu _3`$ causes comparatively large $`\nu _3`$-Higgsino mixing at one loop,
$$m_{3H}\mu _3+\frac{g_2^2}{16\pi ^2}\frac{B\mu _3}{M_{\stackrel{~}{Z}}}0.040.1\mathrm{GeV}.$$
(12)
The $`\tau `$ neutrino mass should be obtained from the full neutralino mass matrix,
$$i\left(\nu _3H_d^0H_u^0\stackrel{~}{Z}X\right)\left(\begin{array}{ccccc}0& 0& m_{3H}& av_3& 0\\ 0& 0& 0& av_d& \lambda ^{}v_u\\ m_{3H}& 0& 0& av_u& \lambda ^{}v_d\\ av_3& av_d& av_u& M_{\stackrel{~}{Z}}& 0\\ 0& \lambda ^{}v_u& \lambda ^{}v_d& 0& 0\end{array}\right)\left(\begin{array}{c}\nu _3\\ H_d^0\\ H_u^0\\ \stackrel{~}{Z}\\ X\end{array}\right)+\mathrm{h}.\mathrm{c}.,$$
(13)
with $`a=\left({\displaystyle \frac{g_1^2+g_2^2}{2}}\right)^{1/2}`$. It gives
$$m_{\nu _\tau }\frac{m_{3H}v_3}{M_Z}\left(110\right)\mathrm{MeV}.$$
(14)
This heavy $`\nu _\tau `$ can decay to $`e^+e^{}\nu _e`$. We later noted that $`\nu _\tau `$ can decay to gravitino $`+`$ photon with a longer lifetime. In writing down the expression of the physical $`\nu _\tau `$ and $`\tau `$ states in Refs. <sup>?</sup> and <sup>?</sup>, the gaugino masses were neglected. In fact, $`\nu _3`$ mixes with photino.
## 4 Discussion
The idea about fermion masses presented in this model essentially depends on SUSY. A model of large sneutrino VEV exists. ($`110`$) MeV $`\nu _\tau `$ is a consequence of this VEV. It should be noted that $`L_3`$ and $`H_d`$ appear in the superpotential in different ways, so that the VEV cannot be rotated away through redefining Higgs superfield.
The atmospheric neutrino anomaly must be explained by introducing a sterile neutrino which is also necessary from the constraint of the Big-Bang Neucleosynthesis. The extension of this idea to the quark sector can be found in Ref. <sup>?</sup>.
The $`B\mu _3`$ term breaks the family symmetry explicitly. It would be more appealing if the family symmetry breaking is spontaneous in some clever model.
## 5 Acknowledgments
I would like to thank Profs. Dongsheng Du and H.S. Song for collaborations.
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# 1 Introduction
## 1 Introduction
Non-commutative field theories have recently received a great deal of attention , stemming in part from the fact that they arise as low-energy descriptions of string backgrounds with anti-symmetric tensor fields . Renormalizability of non-commutative theories remains an open question.<sup>1</sup><sup>1</sup>1As this paper was nearing completion, we received , which argues that a four-dimensional non-commutative Wess-Zumino model is renormalizable. The argument hinges on controlling the infrared singularities which give rise to the interesting phase structure explored in this paper. Thus there is little overlap between and our work. Standard approaches to demonstrating perturbative renormalizability (see for instance ) encounter difficulties because of infrared singularities. These singularities are not associated with any massless propagating fields in the theory, but instead arise through loop effects. They can have dramatic physical consequences: for instance, in a theory whose classical action is that of a massive scalar with cubic interactions, the $`\varphi =0`$ vacuum becomes not just globally unstable (on account of an effective potential which is unbounded below), but locally unstable, as if the scalar had become tachyonic .
The main interest in this paper will be in scalar theories with $`\varphi ^4`$ interactions. For the most part we will restrict attention to Euclidean signature and to even dimensions. Following , we will briefly review in Section 3.2 how non-planar one-loop graphs lead to a singularity in the one particle irreducible (1PI) two-point function: $`\mathrm{\Gamma }^{(2)}(p)\mathrm{}`$ as $`p0`$. What this amounts to physically is long-range frustration: $`\varphi (x)\varphi (0)`$ oscillates in sign for large $`x`$. A natural expectation, given such a correlator, is that the usual Ising-type phase transition, to an ordered phase with $`\varphi 0`$, will be modified to a transition to a phase where $`\varphi `$ varies spatially. This is indeed what we will find in Section 4: more particularly, we will find a fluctuation-driven first order transition to a stripe phase, where only one momentum mode of the scalar field condenses.<sup>2</sup><sup>2</sup>2In it was argued that a translationally invariant ordered phase with massless Goldstone bosons is impossible in continuum renormalized perturbation theory. This can be regarded as a hint of an exotic ordered phase such as the stripe phase that we find. In Section 4.6 we will consider more complicated ordered phases, where more than one momentum mode condenses. In the perturbative regime of Section 4, it turns out that stripe phases are favored; however, it appears that as couplings are increased, the system alternates between preferring the condensation of one or several momentum modes.
The overall picture we will find for the phase diagram is a first order line terminating on one end at a Lifshitz point, where first order behavior merges back into the second order transition of the Ising model; and terminating on the other end at a critical point which arises in a planar version of the commutative theory. We review in Section 2 the relationship between phase structure and renormalizability, and argue via scaling that the existence of the critical point in the planar theory should imply the renormalizability of the non-commutative theory.
The main method we use to establish the existence of phase transitions is a self-consistent Hartree treatment of one-loop graphs. This same method is our primary tool in demonstrating the validity of the scaling arguments that guarantee renormalizability. Thus, these arguments are airtight only for the large $`N`$ limit of a non-commutative $`O(N)`$ vector model, where the self-consistent one-loop Hartree treatment becomes exact.<sup>3</sup><sup>3</sup>3Strictly speaking, this is true only in the disordered phase and at critical points at the boundary of this phase. In the ordered phase the distinction between the self-consistent $`N=1`$ problem and the large $`N`$ limit could be important, see e.g. and especially . We expect to return to this question in the current context in the future. More sophisticated methods are called for to decide the validity of scaling relations in more general non-commutative field theories. As explained in Section 5, it is possible to arrange the quartic couplings in the $`O(N)`$ theory so that there are no divergences at all at leading order in $`N`$, independent of the dimension. Our renormalizability arguments apply, however, for arbitrary quartic couplings in the $`O(N)`$ theory.
Other authors have discussed finite temperature effects in non-commutative field theories, using the Matsubara formalism with periodic Euclidean time. The aim of this paper is rather different: we work with non-commutative field theories on uncompactified flat space, and varying “temperature” is regarded as equivalent to changing the bare mass.
We conclude in Section 7 with a discussion of the relevance of our results to string theory and to quantum hall systems.
## 2 Phase Structure and Renormalizability
We begin with an extremely brief recapitulation of the Wilsonian connection between the phase structure of a (cutoff) field theory and that of taking its continuum limit—which is the problem of renormalizability. As we are interested in this paper in the renormalization of scalar fields, we will recall the lore on commuting scalar fields.
In the Wilsonian approach, we study the Euclidean field theory governed by the action
$$S=d^dx\left[\frac{1}{2}(\varphi )^2+\frac{1}{2}m^2\varphi ^2+\frac{g^2}{4}\varphi ^4\right]$$
(1)
as a statistical mechanics problem—i.e. we introduce a momentum cutoff $`\mathrm{\Lambda }`$ and measure all dimensionful quantities in its units. This leaves us with a problem with one degree of freedom per dimensionless volume ($`\mathrm{\Lambda }^d`$ in physical units) and dimensionless couplings $`m^2/\mathrm{\Lambda }^2`$ and $`g^2/\mathrm{\Lambda }^{4d}`$, commonly labelled $`r`$ and $`u`$ in the statistical mechanics/condensed matter literature (see for instance ). We next search for lines of continuous phase transitions, or critical surfaces, in the $`r,u`$ plane; the critical surface in this problem (which is in the universality class of the Ising model) is sketched in Figure 1. On this surface, the correlation length, in dimensionless units, $`\xi `$ (the “lattice” correlation length, literally so if the cutoff is implemented via discretization) diverges which is the sine qua non of taking a continuum limit. Having located a critical surface, three continuum limits are possible at each point on it—a massless limit obtained by sitting exactly on the critical surface, and two massive limits in which $`r(\mathrm{\Lambda })`$ and $`u(\mathrm{\Lambda })`$ are chosen to approach the critical surface from either phase as the cutoff is taken to infinity while keeping $`\mathrm{\Lambda }/\xi (r,u)=M_R`$ fixed and equal to a renormalized mass. That such limits can be taken, requires scaling in the statistical mechanics. Also, the phenomenon of universality will imply that some domain of a critical surface will exhibit the same long distance correlations and hence give rise to the same continuum limit.
In the above description we have used the language of phase structure. A more powerful account is that of renormalization group (RG) flows which we have sketched in $`d=4`$ (Figure 2a) and $`d<4`$ (Figure 2b). In the RG description we are interested in fixed points of infinite correlation length. In $`d=4`$ we see that all critical theories flow into the gaussian fixed point $`(r,u)=(0,0)`$ whence the (strict) continuum limit is trivial, while in $`d<4`$ the non-trivial Wilson-Fisher fixed point leads to an interacting continuum limit. This information is not available from a phase diagram alone. The fixed point analysis also implies scaling and hence guarantees a continuum limit.
We belabor this point because in the balance of the paper, we will deal in phase diagrams and not RG flows. We will find critical points and will argue that these guarantee the existence of continuum limits (renormalizability) of non-commuting scalar theories. In a specific large $`N`$ theory, we will be able to show this explicitly, and we do not doubt that the claim is correct. Nevertheless, it does not have the full generality of an RG analysis and it would be nice to carry out such an analysis, even perturbatively as done for the commutative $`\varphi ^4`$ theory in $`d=4`$ by Polchinski .
## 3 One Loop Action for Non-Commutative Scalars
### 3.1 Generalities
Our starting point is the non-commutative scalar field theory specified by the action
$$S=d^dx\left[\frac{1}{2}(\varphi )^2+\frac{1}{2}m^2\varphi ^2+\frac{g^2}{4}\varphi ^4\right]$$
(2)
where $`\varphi ^4=\varphi \varphi \varphi \varphi `$ and the star product is defined as usual by
$$(fg)(x)=e^{\frac{i}{2}\mathrm{\Theta }^{\mu \nu }_{x^\mu }_{y^\nu }}f(x)g(y)|_{y=x}.$$
(3)
The anti-symmetric matrix $`\mathrm{\Theta }^{\mu \nu }`$, whose relation to the star product can be most simply expressed through
$$[x^\mu ,x^\nu ]_{}x^\mu x^\nu x^\nu x^\mu =i\mathrm{\Theta }^{\mu \nu },$$
(4)
will for most of this paper be assumed to be of the form
$$\mathrm{\Theta }^{\mu \nu }=\theta \left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\mathrm{𝟏}_{d/2}$$
(5)
for even dimensions $`d`$. We will comment in Section 4.5 on the more complicated cases where $`\mathrm{\Theta }^{\mu \nu }`$ has unequal eigenvalues or $`d`$ is odd.
A useful extension of (3) is
$$(f_1f_2\mathrm{}f_{\mathrm{}})(x)=e^{\frac{i}{2}_{i<j}\mathrm{\Theta }^{\mu \nu }_{x_i^\mu }_{x_j^\nu }}f_1(x_1)f_2(x_2)\mathrm{}f_{\mathrm{}}(x_{\mathrm{}})|_{x_i=x}.$$
(6)
The result (6) is easiest to obtain via Fourier analysis, using the basic result
$$e^{ip_1x}e^{ip_2x}=e^{\frac{i}{2}p_1p_2}e^{i(p_1+p_2)x},$$
(7)
where by definition $`pq=\mathrm{\Theta }^{\mu \nu }p_\mu p_\nu `$. (Condensed matter readers should note that this is the lowest Landau level algebra of density operators ).
Like matrix multiplication, the star product is non-commutative. However, a product of exponentials $`e^{ip_kx}`$, can be reordered cyclically if the $`p_k`$ sum to zero. As a result it is necessary to specify the cyclic order of vertices in writing down Feynman rules, just as in large $`N`$ theories. It is well known that planar amplitudes of the field theory (2) are independent of $`\mathrm{\Theta }^{\mu \nu }`$ (and hence are the same as in the commutative theory where $`\mathrm{\Theta }^{\mu \nu }=0`$), up to an overall phase. The effects of non-planar diagrams were first studied systematically in .
We will also consider two variants of (2), namely complex scalars, where
$$S=d^dx\left[\varphi \varphi ^{}+m^2\varphi \varphi ^{}+g^2\varphi \varphi ^{}\varphi \varphi ^{}+g^2\varphi \varphi \varphi ^{}\varphi ^{}\right],$$
(8)
and the non-commutative $`O(N)`$ vector model,
$$S=d^dx\left[\frac{1}{2}(\varphi _i)^2+\frac{1}{2}m^2\varphi _i^2+\frac{g^2}{4}\varphi _i\varphi _i\varphi _j\varphi _j+\frac{g^2}{4}\delta ^{ik}\delta ^{jl}\varphi _i\varphi _j\varphi _k\varphi _l\right]$$
(9)
The field theory (8) has been studied previously .
### 3.2 One loop diagrams
In this section we review the results of which will be relevant for our calculations. The one loop corrections to the propagator of the scalar field $`\varphi `$ split into planar and non-planar parts:
| $`\mathrm{\Sigma }_{\mathrm{planar}}(p)`$ | $`=2g^2{\displaystyle \frac{d^dk}{(2\pi )^d}\frac{1}{k^2+m^2}}`$ |
| --- | --- |
| $`\mathrm{\Sigma }_{\mathrm{non}\mathrm{planar}}(p)`$ | $`=g^2{\displaystyle \frac{d^dk}{(2\pi )^d}\frac{e^{ipk}}{k^2+m^2}}.`$ |
(10)
Using the Schwinger parametrization
$$\frac{1}{k^2+m^2}=_0^{\mathrm{}}𝑑\alpha e^{\alpha (k^2+m^2)},$$
(11)
one easily extracts
| $`I_2(p)`$ | $`={\displaystyle ^\mathrm{\Lambda }}{\displaystyle \frac{d^dk}{(2\pi )^d}}{\displaystyle \frac{e^{ipk}}{k^2+m^2}}`$ |
| --- | --- |
| | $`={\displaystyle \frac{1}{(4\pi )^{d/2}}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{d\alpha }{\alpha ^{d/2}}}e^{\alpha m^2\frac{pp}{4\alpha }\frac{1}{\mathrm{\Lambda }^2\alpha }}`$ |
| | $`=(2\pi )^{d/2}m^{\frac{d2}{2}}\left(pp+{\displaystyle \frac{4}{\mathrm{\Lambda }^2}}\right)^{\frac{2d}{4}}K_{\frac{d2}{2}}\left(m\sqrt{pp+{\displaystyle \frac{4}{\mathrm{\Lambda }^2}}}\right),`$ |
(12)
where, following , we have introduced an ultraviolet regulator $`\mathrm{\Lambda }`$ and defined a symmetric product
$$pq=p_\mu \mathrm{\Theta }^{\mu \nu }\mathrm{\Theta }_\nu {}_{}{}^{\lambda }q_{\lambda }^{}.$$
(13)
$`K_\nu (x)`$ denotes a modified Bessel function.
The end result is a corrected propagator of the following form:
$$\mathrm{\Gamma }^{(2)}(p)=p^2+m^2+2g^2I_2(0)+g^2I_2(p).$$
(14)
The parallel results for the complex scalar and the $`O(N)`$ vector model are
| complex scalar: | $`\mathrm{\Gamma }^{(2)}(p)=p^2+m^2+(2g^2+g^2)I_2(0)+g^2I_2(p)`$ |
| --- | --- |
| $`O(N)`$ model: | $`\mathrm{\Gamma }^{(2)}(p)=p^2+m^2+g^2NI_2(0)+g^2NI_2(p)+O(1),`$ |
(15)
where for the $`O(N)`$ model we have only evaluated the graphs which contribute at leading order in large $`N`$. The salient property of $`\mathrm{\Gamma }^{(2)}(p)`$ for our subsequent discussion is that it grows for small $`p`$. If we think of absorbing the planar one-loop contribution to $`\mathrm{\Gamma }^{(2)}(p)`$ into the definition of mass (for instance, define a renormalized mass through $`M^2=m^2+2g^2I_2(p^2)`$ for the real scalar theory), and then removing the cutoff while holding $`M^2`$ fixed, then $`\mathrm{\Gamma }^{(2)}(p)`$ is finite for finite $`p`$ but goes to $`+\mathrm{}`$ as $`p0`$. This means that the low-momentum modes are extremely stiff; thus they seem likely never to participate in a phase transition. If there is a phase transition, it should involve the momentum modes where $`\mathrm{\Gamma }^{(2)}(p)`$ is smallest, since these are the ones most likely to be destabilized as $`M^2`$ becomes negative. The details of how this happens turn out to be somewhat intricate, and the next section is devoted to sorting them out.
## 4 Phase Structure for Non-commutative scalars
### 4.1 Generalities
The non-commutative action (2) is characterized by three dimensionless parameters: $`m^2/\mathrm{\Lambda }^2`$, $`g^2/\mathrm{\Lambda }^{4d}`$ and $`\theta \mathrm{\Lambda }^2`$. Accordingly we need to establish a phase diagram in a three dimensional space. In what follows we will typically work at a fixed dimensionless coupling and take two dimensional cuts through parameter space instead. The resulting phase diagram is sketched in Figure 3, and the rest of this section is devoted to the arguments that give rise to it.
### 4.2 The Planar Limit: $`\theta \mathrm{\Lambda }^2=\mathrm{}`$
It has been noted previously that the large $`\theta `$ limit picks out planar diagrams . In the context of a cutoff field theory, this statement can be made precise.
Consider the theory at $`m^2/\mathrm{\Lambda }^2>0`$ so that all diagrams are infrared finite in addition to being ultraviolet finite on account of the cutoff. The perturbative expansion for the free energy $`F=\mathrm{log}Z`$ is the sum of all connected vacuum diagrams (no external legs). Of these, the planar graphs involve no phase factors stemming from the non-commutativity while the non-planar graphs involve at least one. In the limit $`\theta \mathrm{\Lambda }^2\mathrm{}`$ the latter vanish on account of the infinitely rapid oscillations of the integrand. Consequently, the high temperature expansion of $`F`$ in the maximally non-commuting scalar theory is identical to that of the planar theory, defined as the sum over planar diagrams alone.<sup>4</sup><sup>4</sup>4For us, high temperature is synonymous with large positive $`m^2`$, while low temperature means large negative $`m^2`$.
This statement can be extended to selected classes of fields. The perturbative expansion for the free energy in a field is,
| $`F[J]`$ | $`=\mathrm{log}Z[J]`$ |
| --- | --- |
| | $`=\mathrm{log}Z[0]+{\displaystyle \underset{n}{}}{\displaystyle _{p_1}}\mathrm{}{\displaystyle _{p_n}}G_c(p_1,\mathrm{}.,p_n)J(p_1)\mathrm{}.J(p_n)`$ |
(16)
where $`G_c(p_1,\mathrm{}.,p_n)`$, inclusive of the momentum conserving delta function, is the connected $`n`$-point Green function in momentum space. On account of the vanishing of the non-planar graphs alluded to earlier,
$$\underset{\theta \mathrm{\Lambda }^2\mathrm{}}{lim}e^{\frac{i}{2}_{i<j}p_ip_j}G_c(p_1,\mathrm{}.,p_n;\theta \mathrm{\Lambda }^2)=G_c^{\mathrm{planar}}(p_1,\mathrm{}.,p_n)$$
(17)
for each $`n`$. For fields that are modulated only in one direction, the momenta $`p_1`$ through $`p_n`$ are parallel whence the explicit phase factor for planar diagrams vanishes. Hence in the limit $`\theta \mathrm{\Lambda }^2\mathrm{}`$, $`F[J]`$ is still given as a sum of planar diagrams alone. Note that for more complicated field configurations, the equivalence is no longer true.
Thus far we have made statements in the high temperature phase. The planar theory, inheriting the Ising symmetry ($`\varphi \varphi `$) of the full theory, will exhibit a twofold degenerate broken symmetry phase. This will be reached on traversing a continuous phase transition at a critical $`m_c^2/\mathrm{\Lambda }^2<0`$ which will be mean field in $`d>4`$, mean field corrected by logarithms in $`d=4`$, but likely in a different universality class from the standard Ising transition in $`d<4`$ (see below). An immediate deduction from the equivalence of the infinite $`\theta \mathrm{\Lambda }^2`$ and planar theories in a field is that we may analytically continue their common free energy into the low temperature phase (bypassing the critical point), thereby establishing the equivalence of their low temperature thermodynamics as well. Altogether, we may conclude that the infinite $`\theta \mathrm{\Lambda }^2`$ theory has a critical point at the same $`m_c^2/\mathrm{\Lambda }^2`$ as the planar theory, which is exactly the same transition. This critical point will play a central role in achieving a continuum limit for the noncommutative theory.
Finally we should note the well known result that the planar theory is also the $`N\mathrm{}`$ limit of a hermitian matrix model with an appropriately generalized version of (1) as its action. Via this route, there are exact results on the phase transition in $`d=1`$ <sup>5</sup><sup>5</sup>5It is common to refer to the $`c=1`$ matrix model as string theory in two dimensions, on account of the anomaly-induced dynamics for the Liouville field. From the point of view of the matrix model, however, $`d=1`$. that show that the transition is different from that of the Ising model, and an approximate RG treatment that has yielded exponents in $`2<d<4`$ as well . There does not appear to be a treatment of the broken symmetry phase in this formulation of the problem—at issue is how the Ising symmetry of the planar theory is embedded in what appears prima facie to be a much larger symmetry group in the matrix model.
### 4.3 First Order Transition at Large $`\theta \mathrm{\Lambda }^2`$
Having established that there is a critical point at $`\theta \mathrm{\Lambda }^2=\mathrm{}`$, we will now show that this terminates a line of fluctuation driven first order transitions at large $`\theta \mathrm{\Lambda }^2`$. That this should happen, is a fairly general expectation from the work of Brazovskii when combined with the observation that the one loop $`\mathrm{\Gamma }^{(2)}(p)`$ has a minimum at non-zero $`p`$ at large $`\theta \mathrm{\Lambda }^2`$ (the latter condition is implicit in their analysis), and we will review this physics below. Following Brazovskii, we will employ the self-consistent Hartree approximation. This approximation is sensible in $`d4`$, but does not capture the planar theory transition in $`d=2`$. Consequently, we will not treat that case in detail, although our later arguments on renormalizability will apply there as well.
We will present two separate arguments for the first order transition. The first will involve an asymptotic construction of the solution to the self-consistent problem—this will follow the original analysis as closely as possible but with complications that we detail below. This will also require that we take a double limit in which $`\theta \mathrm{\Lambda }^2`$ is taken large but $`g^2`$ is taken to zero—i.e. we will be working in the vicinity of the massless Gaussian theory. The second will be a more indirect argument which will be carried out at fixed $`g^2`$ and will consist of showing that the free energy of two solutions must cross in a certain region of parameter space.
#### 4.3.1 Take I
We turn to the direct construction. To keep the discussion simple, we will continue to assume even Euclidean dimension $`d`$, and $`\mathrm{\Theta }^{\mu \nu }`$ with maximal rank and eigenvalues $`\pm \theta `$. We will also omit inessential factors of order unity.
A self-consistent treatment of the one-loop correction to the propagator leads to
$$\mathrm{\Gamma }^{(2)}(p)=p^2+m^2+2g^2\frac{d^dk}{(2\pi )^d}\frac{1+\frac{1}{2}e^{ikp}}{\mathrm{\Gamma }^{(2)}(k)}$$
(18)
Having $`\mathrm{\Theta }^{\mu \nu }`$ of maximal rank, with eigenvalues $`\pm \theta `$, leads to a considerable simplification: $`\mathrm{\Gamma }^{(2)}(p)`$ is a function only of $`p^2`$. This can be shown inductively, order by order in $`g^2`$. Suppose that $`\mathrm{\Gamma }^{(2)}(p)`$ is $`SO(d1)`$ to $`\mathrm{}`$-th order in $`g^2`$. The $`(\mathrm{}+1)`$-st order expression for $`\mathrm{\Gamma }^{(2)}(p)`$ is then invariant on $`SO(d1)`$ transformations of $`p_\mu \mathrm{\Theta }^{\mu \nu }`$, and so can only be a function of $`pp=\theta ^2p^2`$. Equation (18) can thus be re-expressed as
$$\mathrm{\Gamma }^{(2)}(p)=p^2+M^2+g^2_0^{\mathrm{}}k^{d1}𝑑k\frac{J_{(d2)/2}(\theta kp)}{(\theta kp)^{(d2)/2}}\frac{1}{\mathrm{\Gamma }^{(2)}(k)}.$$
(19)
Here $`m^2`$ is the bare mass, and
$$M^2=m^2+2g^2\frac{d^dk}{(2\pi )^d}\frac{1}{\mathrm{\Gamma }^{(2)}(k)}$$
(20)
is the mass corrected by the planar graph, Figure 4a. The non-commutativity does not affect the large $`p`$ behavior, $`\mathrm{\Gamma }^{(2)}(p)p^2`$. Thus the integral in (20) must be regulated, for instance with an explicit cutoff as in Section 3.2. No other ultraviolet divergences will arise in the following computations, so we may in effect take the cutoff to infinity and treat $`M^2`$ as the finite quantity which we dial to produce a transition. This leads to an argument for renormalizability, as we shall explain in Section 6.
The last term in (18) has a contribution from large momenta $`k`$ which goes as $`1/p^{d2}`$ as $`p0`$. This prevents condensation of very low-momentum modes of $`\varphi `$. Thus, if there is a phase transition at all, the ordered state must break translation invariance.
From now on, let us restrict attention to $`d=4`$. At the end of this section we will remark on the extension to other dimensions. Proceeding naively, we could solve (19) to the first non-trivial order in $`g`$ by replacing $`\mathrm{\Gamma }^{(2)}(k)`$ on the right hand side by $`k^2+M^2`$. Then as $`M^2`$ is decreased below $`c_1\frac{g}{\theta }`$, for some (easily computable) constant $`c_1`$ of order unity, $`\mathrm{\Gamma }^{(2)}(p)`$ becomes negative near $`p=p_c\sqrt{g/\theta }`$, signalling a second order phase transition to an ordered state where some momentum modes of $`\varphi `$ condense. Our aim in the rest of this subsection is to show that this naive result is altered in a self-consistent analysis to a rather more interesting conclusion: there is a first order transition to a stripe pattern at $`M^2=c_1\frac{g}{\theta }c_2\frac{g^{7/3}}{\theta }`$, where $`c_2`$ is another constant of order unity.
First observe that (19) makes the second order transition impossible. If we were to suppose that $`\mathrm{\Gamma }^{(2)}(p_c)=0`$ for some $`M^2`$ and some $`p_c`$, then the integral in (19) diverges.<sup>6</sup><sup>6</sup>6To be more precise, the integral in (19) diverges if $`\mathrm{\Gamma }^{(2)}(p)(p^2p_c^2)^2`$ for $`p`$ near $`p_c`$: that is, $`\mathrm{\Gamma }^{(2)}(p)`$ cannot at the same time remain smooth and have a zero. To cement the conclusion that there are no second order phase transitions, one must check that it is also impossible for $`\mathrm{\Gamma }^{(2)}(p)`$ to have a cusp form such as $`|p^2p_c^2|^{2\nu }`$ with $`\nu <1`$. Such a form can eliminate the divergence in (19). But there is still a contradiction: $`d\mathrm{\Gamma }^{(2)}(p)/dp`$ diverges as $`pp_c`$, but an integral expression for $`d\mathrm{\Gamma }^{(2)}(p)/dp`$ remains finite. The possibility of a first order transition was first realized by Brazovskii , given a suitable effective Lagrangian, with a minimum in the inverse propagator at non-zero $`p_c`$. In , such an effective Lagrangian was simply assumed, and a one-loop fluctuation analysis with this starting point was shown to lead to a first order transition. This same method was further developed in . In our case, the basic Lagrangian does not have an appropriate inverse propagator with a minimum at nonzero $`p_c`$—as we have seen, this only arises in the 1PI effective action after a one-loop analysis. One cannot feed $`\mathrm{\Gamma }^{(2)}(p)`$ directly into a Brazovskiian analysis: all fluctuations are supposed to be incorporated into $`\mathrm{\Gamma }^{(2)}(p)`$ already. However, we shall see that a self-consistent treatment based on (19) reproduces the essential features of , and does support the conclusion that a first order transition takes place.
Although $`\mathrm{\Gamma }^{(2)}(p)`$ can never vanish, it should be possible to make its minimum as small as we please by decreasing $`M^2`$ sufficiently. We then have the approximate forms
$`\mathrm{\Gamma }^{(2)}(p)`$ $`p^2`$ for large $`p`$ $`\mathrm{\Gamma }^{(2)}(p)`$ $`\xi _0^2(p^2p_c^2)^2+r`$ for $`pp_c`$ $`\mathrm{\Gamma }^{(2)}(p)`$ $`{\displaystyle \frac{g^2}{(\theta p)^2}}`$ for small $`p`$. (21)
Two regions of the integral in (19) will then make significant contributions: the large $`p`$ region and the $`pp_c`$ region. The full integral can be approximated as a sum of the contributions from these regions. The parameters $`p_c`$ and $`\xi _0`$ can then be determined self-consistently: for $`pp_c`$,
| $`\mathrm{\Gamma }^{(2)}(p)`$ | $`p^2+M^2+g^2\left({\displaystyle _0^{\mathrm{}}}k^3𝑑k{\displaystyle \frac{J_1(\theta kp)/\theta kp}{k^2}}+p_c^2{\displaystyle \frac{J_1(\theta p_c^2)}{\theta p_c^2}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{kdk}{r+\xi _0^2(k^2p_c^2)^2}}\right)`$ |
| --- | --- |
| | $`p^2+M^2+{\displaystyle \frac{g^2}{(\theta p)^2}}+{\displaystyle \frac{g^2p_c^2}{\xi _0\sqrt{r}}}`$ |
| $`p_c`$ | $`={\displaystyle \frac{1}{\xi _0}}=\sqrt{{\displaystyle \frac{g}{\theta }}}r=2p_c^2+M^2+{\displaystyle \frac{g^3/\theta }{\xi _0\sqrt{r}}}.`$ |
(22)
The last line of (22) follows from matching the previous line to the approximate form $`\mathrm{\Gamma }^{(2)}(p)\xi _0^2(p^2p_c^2)^2`$. We have again suppressed inessential factors of order unity in (22), and we have assumed $`g1`$, so that $`\theta p_c^21`$. This latter inequality will be essential to our further analysis. It turns out to be difficult to get the factors of order unity right, because $`\xi _0`$ is only on the order of $`1/p_c`$: the breadth of the minimum of $`\mathrm{\Gamma }^{(2)}(p)`$, in the approximation scheme we have used, is comparable to its distance from the origin.
So far we have done all calculations as expansions around the disordered phase, where $`\varphi =0`$. The putative ordered phase is a stripe pattern: $`\varphi =\varphi _0=A\mathrm{cos}p_cx`$. (Strictly speaking, we do not expect an exact $`\mathrm{cos}p_cx`$ dependence, only a function of $`x`$ with period $`2\pi /p_c`$. But because the momentum modes near $`p_c`$ make the dominant contribution, neglecting higher harmonics should not change the qualitative picture. One should in principle consider sums of the form $`\mathrm{cos}p_cx+\mathrm{cos}p_cy`$, as well as more complicated superpositions; however, as we shall see in Section 4.6, the simple stripe solution is favored at small coupling). Writing $`\varphi =\varphi _0+\eta `$ and expanding in fluctuations of $`\eta `$, one obtains
| $``$ | $`={\displaystyle \frac{1}{2}}(\varphi _0)^2+{\displaystyle \frac{g^2}{4}}\varphi _0^4+{\displaystyle \frac{1}{2}}(\eta )^2+{\displaystyle \frac{g^2}{4}}\eta ^4`$ |
| --- | --- |
| | $`+\varphi _0\eta +g^2\varphi _0^3\eta +g^2\varphi _0\eta ^3+g^2\varphi _0^2\eta ^2+{\displaystyle \frac{g^2}{2}}\varphi _0\eta \varphi _0\eta .`$ |
(23)
We will treat perturbatively all terms with factors of both $`\eta `$ and $`\varphi _0`$. This is justified if $`A`$ is small, which is not completely obvious given that the transition is first order. But for small $`g`$, $`A`$ is small very close to the phase transition—at least, this is an assumption which can easily be verified at the end of the computation.
Classically, one would demand that the terms in (23) linear in $`\eta `$ should vanish: this corresponds to summing the graphs in Figure 5a and 5b. At the one-loop level, the graphs in Figure 5c and 5d also contribute, and the end result is
$$h\frac{1}{2}\frac{d\mathrm{\Phi }}{dA}=\mathrm{\Gamma }^{(2)}(p_c)A\frac{g^2}{2}A^3,$$
(24)
where $`\mathrm{\Phi }`$ is the free energy per unit volume. Adding the graphs in Figures 5e and 5f to the self-consistent one-loop graphs that contributed to (18), one finds only a slight modification of the propagator in the disordered phase:<sup>7</sup><sup>7</sup>7Actually, the propagator is not diagonal in momentum space once the graphs in Figure 5e and 5f are added. Non-diagonal terms would have to be taken into account if we wanted to determine the exact form of the ordered phase; however they should not change the qualitative features of the phase transition. To (18) one must simply add the terms $`g^2A^2(1+\frac{1}{2}\mathrm{cos}pp_c)`$. The first of these terms comes from the graph in Figure 5e, while the second is from the graph in Figure 5f. For momenta $`p`$ on the order of $`p_c`$, the argument of the cosine is small (again assuming $`\theta p_c^2`$ is small), so for such momenta we may alter (22) to
| $`\mathrm{\Gamma }^{(2)}(p)`$ | $`p^2+M^2+{\displaystyle \frac{g^2}{(\theta p)^2}}+{\displaystyle \frac{g^2p_c^2}{\xi _0\sqrt{r}}}+g^2A^2`$ |
| --- | --- |
| $`p_c`$ | $`={\displaystyle \frac{1}{\xi _0}}=\sqrt{{\displaystyle \frac{g}{\theta }}}r=2p_c^2+M^2+{\displaystyle \frac{g^3/\theta }{\xi _0\sqrt{r}}}+g^2A^2.`$ |
(25)
To summarize the situation so far: a self-consistent analysis of the graphs in Figures 4 and 5 led to an expression for the minimum, $`r=\mathrm{\Gamma }^{(2)}(p_c)`$, of the 1PI propagator $`\mathrm{\Gamma }^{(2)}(p)`$; and an expression for the tadpole for the magnitude $`A`$ of the condensate. These expressions are
$$r=\tau +g^2A^2+\frac{\alpha }{\sqrt{r}}\frac{1}{2}\frac{\mathrm{\Phi }}{A}=A\left(r\frac{g^2}{2}A^2\right),$$
(26)
where we have defined
$$\tau =2p_c^2+M^2\alpha =\frac{g^{7/2}}{\theta ^{3/2}}p_c=\sqrt{\frac{g}{\theta }}.$$
(27)
The various factors of order unity which we have neglected can be absorbed by rescaling $`g`$, $`p_c`$, and $`\alpha `$.<sup>8</sup><sup>8</sup>8In point of fact, it is difficult to compute some of these factors, since they arise from the integral equation (19). We thank D. Priour for consultations on the possibility of treating (19) numerically. None of these rescalings affects the conclusion that there is a first order transition. However, the rest of the analysis is somewhat delicate, and we must from now on be meticulous about factors. The calculation we are about to sketch appeared in with some slight errors, which were corrected in , modulo a typo in the sign of the last term in their (2.13).
A stable phase must have a vanishing tadpole. This happens either when $`A=0`$ (the disordered phase) or when $`r\frac{g^2}{2}A^2=0`$ (the ordered phase). The respective values of $`r`$, as well as the value of $`A`$ in the ordered phase, can be determined from
$`r_o+{\displaystyle \frac{\alpha }{\sqrt{r_o}}}+\tau =0r_d{\displaystyle \frac{\alpha }{\sqrt{r_d}}}\tau =0`$ $`A_o={\displaystyle \frac{1}{g}}\sqrt{r_o{\displaystyle \frac{\alpha }{\sqrt{r_o}}}\tau }.`$ (28)
Note that the ordered phase only exists for $`\tau <3(\alpha /2)^{2/3}1.89\alpha ^{2/3}`$. In this range, there are two solutions for $`r_o`$ (both positive): let us call these $`r_{o1}`$ and $`r_{o2}`$. The free energy difference is conveniently calculated as
| $`\mathrm{\Delta }\mathrm{\Phi }`$ | $`=\mathrm{\Phi }_o\mathrm{\Phi }_d={\displaystyle _{r_d}^{r_o}}𝑑r{\displaystyle \frac{dA}{dr}}{\displaystyle \frac{d\mathrm{\Phi }}{dA}}`$ |
| --- | --- |
| | $`={\displaystyle \frac{g^2}{4}}A_o^4+{\displaystyle \frac{r_o^2r_d^2}{2g^2}}+{\displaystyle \frac{\alpha }{g^2}}(\sqrt{r_o}\sqrt{r_d})`$ |
| | $`={\displaystyle \frac{1}{g^2}}\left[{\displaystyle \frac{r_o^2+r_d^2}{2}}+\alpha (\sqrt{r_o}\sqrt{r_d})\right].`$ |
(29)
This equation applies equally to the two solutions for $`r_o`$. By subtracting (29) with $`r_o=r_{o1}`$ from the same equation with $`r_o=r_{o2}`$, one can verify that the larger solution for $`r_o`$ leads to the lower free energy. One can also show, directly from (29), that the disordered phase is favored for $`2.03\alpha ^{2/3}<\tau <1.89\alpha ^{2/3}`$, and the ordered phase is favored for $`\tau <2.03\alpha ^{2/3}`$.<sup>9</sup><sup>9</sup>9In principle one should be able to prove these assertions using the formula for solutions of a cubic. However, since the problem has only one parameter, $`\tau /\alpha ^{2/3}`$, we have found it more convenient to proceed numerically.
In terms of the original variables, and again in $`d=4`$, the phase transition to an ordered phase occurs when $`M^2=c_1\frac{g}{\theta }c_2\frac{g^{7/3}}{\theta }`$, for some constants $`c_1`$ and $`c_2`$ of order unity. It is straightforward to show that $`c_1=1/\pi `$ and $`\theta p_c^2=g/2\pi `$. But to obtain an accurate value for $`c_2`$ requires numerics. The $`g^{7/3}`$ correction is a measure of the extent to which the system avoids the second order behavior expected from a naive one-loop analysis. Higher loop corrections are expected to shift the critical value of $`M^2\theta `$ by $`O(g^3)`$.
#### 4.3.2 Take II
We now consider a second route to establishing a first order transition at large $`\theta \mathrm{\Lambda }^2`$ which exploits the intimate connection between the phases of the noncommutative theory and their limits in the planar theory. The basic idea is this: As shown earlier, one cannot have a continuous transition once the disordered phase propagator has a minimum away from zero which is the case once $`\theta \mathrm{\Lambda }^2`$ is reduced from infinity. While this allows the high temperature solution to exist below the critical temperature of the planar theory, there must be a second solution that grows out of the ordered solution of the planar theory but is now modulated at a very long wavelength. If we take $`\theta \mathrm{\Lambda }^2`$ sufficiently large at fixed $`m^2<m_c^2`$, the ordered solution will win for it is the correct solution at $`\theta \mathrm{\Lambda }^2=\mathrm{}`$. Hence there must be a first order line emanating from the planar theory critical point.
To construct a proof on these lines it would appear that we could get away by comparing free energies exactly at $`\theta \mathrm{\Lambda }^2=\mathrm{}`$ for continuity would give us a range of $`\theta \mathrm{\Lambda }^2`$ over which the ordering would still hold. This almost works—while there is no problem with the ordered solution, which connects to the ordered solution of the planar theory, the delicate point is that the $`\theta \mathrm{\Lambda }^2\mathrm{}`$ limit of the disordered solution is not itself a solution of the planar theory self-consistency equation. Consequently we need to take the limit carefully, which can be done, modulo a weak and entirely plausible assumption, without actually solving the equations. (For the simpler Brazovskii problem, a similar strategy is entirely successful. This is detailed in the Appendix.)
To begin, let us consider the limiting ordered free energy. This is the free energy of the planar theory,
$$F_{\mathrm{ord}}=(\frac{m^2}{2}\varphi _o^2+\frac{g^2}{4}\varphi _o^4)\frac{(m^2+g^2\varphi _o^2)^2}{8g^2}+\frac{1}{2}\frac{d^dk}{(2\pi )^d}\mathrm{log}[k^2+g^2\varphi _o^2]$$
(30)
where $`\varphi _o=\varphi `$ and is determined by $`m^2`$ via,
$$m^2+4g^2\varphi _o^2+8g^2\frac{d^dk}{(2\pi )^d}\frac{1}{k^2+g^2\varphi _o^2}=0.$$
(31)
In writing this compact form, we have used the relation $`\mathrm{\Sigma }=m^2+4g^2\varphi _o^2`$ valid for the fluctuation propagator in the ordered phase.
We next turn to the disordered phase self-consistency equation, now written in terms of the self energy:
$$\mathrm{\Sigma }(p)=2g^2\frac{d^dk}{(2\pi )^d}\frac{1+\frac{1}{2}e^{ikp}}{k^2+m^2+\mathrm{\Sigma }(k)}.$$
(32)
We will assume that the integrals are regulated in the ultraviolet, but the the precise choice of regulator will not be crucial. Consider the following assertions:
(i) At any $`\mathrm{\Lambda }`$ and $`\theta \mathrm{\Lambda }^2`$, $`\mathrm{\Sigma }(0)=\frac{3}{2}\mathrm{\Sigma }(\mathrm{})`$. This follows from the rapid oscillation of the phase factor as $`p\mathrm{}`$. At increasingly large $`\theta \mathrm{\Lambda }^2`$, this fall in $`\mathrm{\Sigma }(p)`$ will happen increasingly rapidly. Hence if we examine $`\mathrm{\Gamma }^{(2)}(p)=p^2+m^2+\mathrm{\Sigma }(p)`$ it must develop a minimum away from $`p=0`$. As already noted this minimum will prevent a continuous transition.
(ii) For stability, $`\mathrm{\Gamma }^{(2)}(p)=k^2+m^2+\mathrm{\Sigma }(p)>0`$.
(iii) It follows that $`\mathrm{\Sigma }(\mathrm{})`$ and $`\mathrm{\Sigma }(0)`$ are positive.
(iv) The combination $`m^2+\mathrm{\Sigma }(p)`$ must change sign between zero and infinity whenever $`m^2<m_c^2`$ and $`\theta \mathrm{\Lambda }^2`$ is sufficiently large. The proof is by contradiction. At $`p=0`$ it is positive by (ii). Were it to remain bounded away from zero we could bound
$$\frac{1}{k^2+m^2+\mathrm{\Sigma }(k)}<\frac{1}{k^2}$$
(33)
and deduce that $`\mathrm{\Sigma }(\mathrm{})<m_c^2`$ and thence that $`\mathrm{\Sigma }(\mathrm{})+m^2<m^2m_c^2<0`$ which contradicts the hypothesis.
(v) The minimum value of $`m^2+\mathrm{\Sigma }(p)`$ must vanish as $`\theta \mathrm{\Lambda }^2\mathrm{}`$. This follows from (i) and (ii). The minimum value will be achieved (or arbitrarily closely approximated in the case of purely monotone decreasing $`\mathrm{\Sigma }(p)`$) ever closer to $`p=0`$ in this limit. To avoid violating (ii), it must vanish.
Thus far we have not made any assumptions. Now we need to make a mild assumption to proceed further.
(vi) We assume that $`m^2+\mathrm{\Sigma }(\mathrm{})`$ vanishes in the large $`\theta \mathrm{\Lambda }^2`$ limit as well. If $`\mathrm{\Sigma }(p)`$ is monotone decreasing, this follows from (v). Otherwise is is strongly indicated by continuity from higher temperatures. For $`m^2>m_c^2`$ the limit is a positive constant, being the renormalized mass of the planar theory, which vanishes at $`m^2=m_c^2`$. It seems highly unlikely that at any fixed $`\theta \mathrm{\Lambda }^2\mathrm{}`$ the large momentum behavior will be non-monotone as a function of temperature whence the conclusion. In sum, we assume that $`\mathrm{\Sigma }(p)`$ goes pointwise to $`m^2`$. Note that this is not a solution of the limit of (32).
(vii) The remaining task is to evaluate the free energy in this limit. The free energy of the disordered solution is given by,
$$F_{\mathrm{dis}}=\frac{1}{4}\frac{d^dk}{(2\pi )^d}\frac{\mathrm{\Sigma }(k)}{k^2+m^2+\mathrm{\Sigma }(k)}+\frac{1}{2}\frac{d^dk}{(2\pi )^d}\mathrm{log}[k^2+m^2+\mathrm{\Sigma }(k)].$$
(34)
As the limit of the disordered solution is not itself a solution, we have to be a bit careful about evaluating $`F_{\mathrm{dis}}`$ as $`\theta \mathrm{\Lambda }^2\mathrm{}`$. Setting $`\mathrm{\Sigma }(p)+m^2=0`$ in the logarithm is unproblematic. The remaining integral can be evaluated by breaking it up as
$$\frac{d^dk}{(2\pi )^d}\frac{m^2+\mathrm{\Sigma }(k)}{k^2+m^2+\mathrm{\Sigma }(k)}\frac{d^dk}{(2\pi )^d}\frac{m^2}{k^2+m^2+\mathrm{\Sigma }(k)}$$
(35)
wherein the first term vanishes at large $`\theta \mathrm{\Lambda }^2`$ and the second can be evaluated via the self-consistency equation. Finally,
$$\underset{\theta \mathrm{\Lambda }^2\mathrm{}}{lim}F_{\mathrm{dis}}=\frac{(m^2)^2}{8g^2}+\frac{1}{2}\frac{d^dk}{(2\pi )^d}\mathrm{log}[k^2].$$
(36)
It is not hard to see that this is higher than the free energy of the ordered solution for $`m^2<m_c^2`$ by expanding the former to leading order in $`\varphi _o^2`$, which is what we set out to prove.
### 4.4 Small $`\theta \mathrm{\Lambda }^2`$: Ising Transition and Lifshitz Point
We now turn to the opposite limit, namely that in which $`\theta \mathrm{\Lambda }^2`$ is sufficiently small—we will quantify that below. The analysis in this limit is straightforward. Exactly at $`\theta \mathrm{\Lambda }^2=0`$ we recover the commuting theory which has the standard Ising critical point. At small, non-zero values of $`\theta \mathrm{\Lambda }^2`$ we can expand the exponential phase factor in the quartic vertex in powers of $`_{i<j}p_ip_j`$. The leading term is the quartic interaction of the commuting theory, and the additional terms are clearly irrelevant at its critical fixed point in $`d4`$. This statement is true below $`d=4`$ as well in the $`ϵ`$ expansion, for the additional terms possess a momentum structure that is not generated by the interactions already present at the Wilson-Fisher fixed point.<sup>10</sup><sup>10</sup>10This assumes we continue the wedge structure in some fashion between dimensions four and two.
This RG statement can be understood more prosaically from the one-loop computation reviewed in Section 3. On examining the minimum in $`\mathrm{\Gamma }^{(2)}(k)`$ we find that it continues to be at $`k=0`$ at small $`\theta \mathrm{\Lambda }^2`$ whence the phase transition would still be into the uniform Ising symmetry breaking state. As a bonus, we discover that there is a critical value $`(\theta \mathrm{\Lambda }^2)_c1/g`$ at which the minimum starts to move away from $`k=0`$ and that initially it grows as the square root of the deviation in $`\theta \mathrm{\Lambda }^2`$. While these precise claims will be modified in an actual theory of this region (which the straight one loop computation is not, being as it is only valid for $`m^2/\mathrm{\Lambda }^2>0`$), the general feature that the Ising transition to a uniform phase will give way to a first order transition (expected by continuity with what happens at large $`\theta \mathrm{\Lambda }^2`$) to a modulated phase at a Lifshitz point should be correct. Unfortunately, our self consistent treatment is incapable of producing a Lifshitz point at a finite temperature ($`m^2/\mathrm{\Lambda }^2>\mathrm{}`$) in $`d=4`$ (which is the lower critical dimension in this approximation) and so the Ising transition and the first order line meet only at $`m^2/\mathrm{\Lambda }^2=\mathrm{}`$. Further investigation of this region is clearly desirable.
### 4.5 Odd dimensions
The phase transition we have described is dimension-sensitive. In $`d=2`$, there can be no long-range order: the stripe phase is unstable to infrared fluctuations.<sup>11</sup><sup>11</sup>11This statement is essentially an application of the Coleman-Mermin-Wagner theorem; however we will check it explicitly in Section 4.6. In $`d=3`$, with only $`\theta ^{12}0`$, the story is slightly more complicated, since $`\mathrm{\Gamma }^{(2)}(p)`$ is no longer a function only of $`p^2`$. The momentum modes for which $`\mathrm{\Gamma }^{(2)}(p)`$ is a minimum lie on a ring in the $`p_1`$$`p_2`$ plane. One may nevertheless argue as before that a second order phase transition is impossible: the integral equation for $`\mathrm{\Gamma }^{(2)}(p)`$ implies that it is smooth; but then if $`\mathrm{\Gamma }^{(2)}(p)`$ is to have a zero along the ring, it must be a quadratic zero, and such a zero renders the integral infinite for all momenta. A weakly first order transition to a stripe phase is the expected behavior (weak because only logarithmic divergences make the transition first order). The main point that must be checked is that the stripe phase is not destroyed by infrared fluctuations. This we will do in section 4.6.
Suppose we add one more commuting dimension: that is, we work in $`d=4`$ but with only $`\mathrm{\Theta }^{12}`$ and $`\mathrm{\Theta }^{21}`$ nonzero. $`\mathrm{\Gamma }^{(2)}(p)`$ should again have a minimum for $`p`$ on a ring in the $`p_1`$$`p_2`$ plane. Since this ring is now codimension three, it is possible for $`\mathrm{\Gamma }^{(2)}(p)`$ to go smoothly to zero in the self-consistent Hartree treatment. So a second order transition becomes possible. What happens with a generic $`\mathrm{\Theta }^{\mu \nu }`$ is less obvious, but again a second order transition seems possible. One might be able to rule out a second order transition on the basis of an instability in the one-loop corrected quartic interactions. It would be interesting to work out this case in more detail.
### 4.6 Ordering beyond Stripes
Our analysis of the ordering in non-commutative scalar theories has thus far been restricted to stripe phases. Following Brazovskii, we would expect these to be the correct solution in the self-consistent Hartree approximation at weak coupling. However this ignores two orthogonal possibilities. The first is the option of going between the stripe phase and the disordered phase by analogs of the smectic and nematic phases in liquid crystal physics.<sup>12</sup><sup>12</sup>12We should note, however, that in the extensively studied problem of A–B diblock copolymer melts there appears to be no evidence for intermediate phases. This is conceivable on symmetry grounds but as it is beyond the reach of our techniques, we are unable to say anything definite on this score other than to note that the destruction of long range stripe order in $`d=2`$ (see below) would imply that any ordered phase would necessarily be more symmetric.
The second possibility is that, at lower temperatures, the stripe phase might give way to one in which two (or more) different momentum modes of $`\varphi `$ might condense, leading to a pattern of squares or rhombi. An interesting question here is whether the noncommutativity can influence the state selection in an interesting way. The purpose of this section is to argue that stripes are indeed preferred over rhombi for small $`\theta p_c^2`$, but that rhombi (or more complicated patterns) become favored as $`\theta p_c^2`$ is increased. Since $`\theta p_c^2`$ is small when the coupling $`g^2`$ is small, the more interesting ordered phases can only arise at strong coupling. We will also remark on the first order phase transition to a triangular crystal that arises when a $`\varphi ^3`$ interaction is present. The analysis will be very rough in that we will use a model lagrangian that will incorporate the soft modes and nonlinearities in the problem at the tree level instead of attempting a full fluctuation analysis. However, the conclusion that stripes give way to more complicated patterns as $`\theta p_c^2`$ is increased should be robust, since it relies simply on the phases that emerge from the star product.
As a preliminary, and to introduce our model, let us demonstrate that long range order for a stripe phase is impossible in two dimensions. This is a well-established result , so we will be brief. In the case of even dimension $`d`$, and assuming $`\theta ^{\mu \nu }`$ has all eigenvalues equal to $`\pm \theta `$, non-commutativity contributes to the analysis only by generating a minimum in the propagator for non-zero momentum $`p_c`$. Thus we can work with an effective action of the Brazovskiian form :
$$S_{\mathrm{eff}}=d^dx\left(\frac{1}{2}\kappa _1\left[(^2+p_c^2)\varphi \right]^2+\frac{1}{2}\kappa _2\varphi ^2+\frac{1}{4}\kappa _4\varphi ^4\right).$$
(37)
A variant for complex scalars is
$$S_{\mathrm{eff}}=d^dx\left(\kappa _1|(^2+p_c^2)\varphi |^2+\kappa _2|\varphi |^2+\frac{1}{2}\kappa _4|\varphi |^4\right).$$
(38)
It is assumed that $`\kappa _1>0`$ and $`\kappa _4>0`$. When $`\kappa _2<0`$, there is pattern formation. The complex case is simpler, since the classical minima of (38) are plane waves, $`\varphi =Ae^{ipx}`$, with $`|A|=\sqrt{\kappa _2/\kappa _4}`$ and $`|p|=p_c`$. We will focus on the complex case at the outset, and return to real scalars, and to $`\varphi ^3`$ interactions, near the end of this section.
The precise claim about absence of long-range order in the model (38) is that at any finite temperature, in $`d3`$, long-range order is destroyed by infrared fluctuations. To see this, suppose $`\varphi =A(x)e^{ip_cx^1}`$ where $`A(x)`$ is assumed to be slowly varying. Plugging this ansatz into (38), one finds
$$S_{\mathrm{eff}}=d^dx\left(\kappa _1|(2ip_c_{x^1}+_x^{}^2)A|^2+\kappa _2|A|^2+\frac{1}{2}\kappa _4|A|^4\right).$$
(39)
The massless modes are those where $`A=\sqrt{\kappa _2/\kappa _4}e^{i\theta (x)}`$. We will use $`x^{}`$ to indicate the directions perpendicular to the unit vector $`\widehat{x}^1`$. To quadratic order in $`\theta `$, (39) reduces to
$$S_{\mathrm{eff}}=d^dx\left(\left|\frac{\kappa _1\kappa _2}{\kappa _4}\right|\left[4p_c^2(_1\theta )^2+(_{}^2\theta )^2\right]\frac{1}{2}\frac{\kappa _2^2}{\kappa _4}\right).$$
(40)
Because the action is fourth-order in derivatives, Coleman-Mermin-Wagner arguments are stronger than usual: an estimate of the fluctuations gives
$$\theta ^2\frac{d^dq}{4p_c^2q_1^2+q_{}^4},$$
(41)
which is infrared divergent for $`d3`$.
The above argument is not substantially modified in $`d=2`$ by non-commutativity because the infrared divergence in (41) refers to physics at far larger length scales than the scale of non-commutativity or of frustration. However, in $`d=3`$, the premise of the model is wrong: non-commutativity can only exist for two of the three coordinates, say $`x^1`$ and $`x^2`$. If a stripe phase forms in the $`x^2`$ direction, the dispersion relation (up to dimensionful constants) is $`\omega (q)q_1^4+q_2^2+q_3^2`$. The integral $`d^3q/\omega (q)`$ is now finite in the infrared, so the stripe phase is indeed stable, and our remarks at the end of Section 4.3.1 were justified.
Let us now proceed to the comparison of energetics for stripes/rolls and rhombi, in a non-commutative version of the Brazovskii model,<sup>13</sup><sup>13</sup>13We thank E. Witten for a conversation in which the simplest case of this argument was worked out for small $`\theta p_c^2`$.
| $`S_{\mathrm{eff}}`$ | $`={\displaystyle d^dx_{\mathrm{eff}}}`$ |
| --- | --- |
| $`_{\mathrm{eff}}`$ | $`=\left(\kappa _1|(^2+p_c^2)\varphi |^2+\kappa _2|\varphi |^2+{\displaystyle \frac{1}{2}}\kappa _4\varphi \varphi \varphi ^{}\varphi ^{}+{\displaystyle \frac{1}{2}}\kappa _4^{}\varphi \varphi ^{}\varphi \varphi ^{}\right),`$ |
(42)
with $`\kappa _1>0`$, $`\kappa _2<0`$, and $`\kappa _4+\kappa _4^{}>0`$. For simplicity, we will assume that $`\varphi `$ is modulated at most in two directions, $`x^1`$ and $`x^2`$, which are rotated onto one another by the action of $`\theta ^{\mu \nu }`$. More complicated situations can be imagined since we must assume $`d>2`$ for the analysis to proceed at all; however, most of the interesting physics should refer to dimensions paired together by $`\theta ^{\mu \nu }`$.
Clearly, there are still roll solutions, $`\varphi =Ae^{ipx}`$, to the equations of motion. Rhombi would arise from an ansatz $`\varphi =A_1e^{ip_1x}+A_2e^{ip_2x}`$ with $`|p_1|=|p_2|=p_c`$, but with $`p_1`$ and $`p_2`$ linearly independent. This ansatz is not a solution to the equation of motion, but if we can show that it has a lower energy than the roll solution, it is reasonable evidence that rhombic patterns form. Plugging the trial function into (42), we obtain
| $`|\varphi |^2`$ | $`=|A_1|^2+|A_2|^2`$ |
| --- | --- |
| $`\varphi \varphi \varphi ^{}\varphi ^{}`$ | $`=|A_1|^4+|A_2|^4+2(1+\mathrm{cos}p_1p_2)|A_1A_2|^2`$ |
| $`\varphi \varphi ^{}\varphi \varphi ^{}`$ | $`=|A_1|^4+|A_2|^4+4|A_1A_2|^2`$ |
| $`_{\mathrm{eff}}`$ | $`=\kappa _2(|A_1|^2+|A_2|^2)+{\displaystyle \frac{1}{2}}(\kappa _4+\kappa _4^{})(|A_1|^2+|A_2|^2)^2`$ |
| | $`+(\kappa _4\mathrm{cos}p_1p_2+\kappa _4^{})|A_1A_2|^2,`$ |
(43)
where $`U`$ indicates that $`U`$ is to be averaged over space.
Clearly, rolls are favored over rhombi when $`\kappa _4^{}>|\kappa _4|`$. In order for rhombi to be favored over rolls without making $`\kappa _4+\kappa _4^{}<0`$ (which would destabilize the theory altogether), we need $`\kappa _4>\kappa _4^{}`$ and $`p_1p_2`$ in the vicinity of $`\pi `$. Because
$$p_1p_2=\theta p_c^2\mathrm{sin}angle(p_1,p_2),$$
(44)
we need $`\theta p_c^2\theta \lambda `$ to be finite.
Suppose we indeed arrange values of the $`\kappa _i`$ where rhombi are favored over rolls. The simplest case is $`\kappa _4>0`$ but $`\kappa _4^{}=0`$. The curious fact is that a square lattice appears almost never to be the preferred pattern: if $`\theta p_c^2<\pi /2`$, rolls are preferred; whereas if $`\theta p_c^2>\pi /2`$, then among the trial wave-functions studied so far the one with the lowest energy has $`angle(p_1,p_2)=\mathrm{sin}^1(2\theta p_c^2/\pi )`$. We have not found actual solutions to the equations of motion with the point symmetries of a rhombic lattice. Because of the $`\varphi ^4`$ terms, any such solution would necessarily be a combination of infinitely many plane waves with wave-vectors on the reciprocal lattice. It is straightforward though tedious work to optimize a variational ansatz which is the sum of finitely many plane waves. However there is not much point in going through the exercise because the form of (38) is only intended to capture the behavior of momenta close to $`|p|=p_c`$.
Finally, let us turn to the case of a real scalar. We will continue to use the approach of a trial wave-function composed of plane waves with all momenta on the ring $`|p|=p_c`$. Plugging the ansatz
$$\varphi =A_1e^{ip_1x}+A_2e^{ip_2x}+\mathrm{c}.\mathrm{c}.,$$
(45)
into the effective action
| $`S_{\mathrm{eff}}`$ | $`={\displaystyle d^dx_{\mathrm{eff}}}`$ |
| --- | --- |
| $`_{\mathrm{eff}}`$ | $`={\displaystyle \frac{1}{2}}\kappa _1\left[(^2+p_c^2)\varphi \right]^2+{\displaystyle \frac{1}{2}}\kappa _2\varphi ^2+{\displaystyle \frac{1}{4}}\kappa _4\varphi \varphi \varphi \varphi ,`$ |
(46)
with $`\kappa _1>0`$, $`\kappa _2<0`$, and $`\kappa _4>0`$, one obtains the following:
| $`\varphi ^2`$ | $`=2(|A_1|^2+|A_2|^2)`$ |
| --- | --- |
| $`\varphi \varphi \varphi \varphi `$ | $`=6(|A_1|^4+|A_2|^4)+(16+8\mathrm{cos}p_1p_2)|A_1A_2|^2`$ |
| $`_{\mathrm{eff}}`$ | $`=\kappa _2(|A_1|^2+|A_2|^2)+{\displaystyle \frac{3}{2}}\kappa _4(|A_1|^2+|A_2|^2)^2+\kappa _4(1+2\mathrm{cos}p_1p_2)|A_1A_2|^2.`$ |
(47)
The sign on the last term can be negative, and the conclusion is that on the space of trial wave-functions that we have examined, rhombi are preferred over stripes if $`p_1p_2>2\pi /3`$ is finite. This last inequality is possible when $`\theta p_c^2>2\pi /3`$. In a similar way, one can consider the condensation of three momentum modes and show that for $`\theta p_c^2`$ in a neighborhood of $`2\pi /\sqrt{3}`$, a hexagonal pattern is preferred over both rhombi and stripes. Again, although these trial wave-functions are not solutions to the equations of motion, it is reasonable to expect solutions with the same symmetries to exist and to compete with one another in a similar way. We have considered condensing momentum modes which all lie in a single plane on which $`\theta _{\mu \nu }=\theta ϵ_{\mu \nu }`$. In $`d=4`$ and higher, there are more complicated possibilities where momentum modes condense in many directions. We will not attempt to classify all the possible ordered phases.
One can also consider adding a cubic term to (46):
| $`S_{\mathrm{eff}}`$ | $`={\displaystyle d^dx_{\mathrm{eff}}}`$ |
| --- | --- |
| $`_{\mathrm{eff}}`$ | $`={\displaystyle \frac{1}{2}}\kappa _1\left[(^2+p_c^2)\varphi \right]^2+{\displaystyle \frac{1}{2}}\kappa _2\varphi ^2+{\displaystyle \frac{1}{3}}\kappa _3\varphi \varphi \varphi +{\displaystyle \frac{1}{4}}\kappa _4\varphi \varphi \varphi \varphi .`$ |
(48)
As explained in , a $`\varphi ^3`$ interaction contributes a term $`I_2(p)`$ to $`\mathrm{\Gamma }^{(2)}(p)`$; thus if the cubic interaction is strong as compared to the quartic interaction, $`\mathrm{\Gamma }^{(2)}(p)`$ is monotonic, and there is no reason to think that there are spatially non-uniform phases at all. For the purposes of this discussion, let us assume then that the cubic interaction is weak but nonzero.
Given that the leading nonlinearity is cubic, the natural expectation in two dimensions would be that commutative versions of (48) exhibit a first order transition to hexagonal crystals as $`\kappa _2`$ is lowered. This persists in the presence of non-commutativity, as the following computation shows. The ansatz for $`\varphi `$ is
$$\varphi =A_1e^{ip_1x}+A_2e^{ip_2x}+A_3e^{i(p_1+p_2)x}+\mathrm{c}.\mathrm{c}.,$$
(49)
with $`|p_1|=|p_2|=|p_1+p_2|=p_c`$. In two dimensions, the only way that this can be arranged is for $`p_1`$, $`p_2`$, and $`p_1p_2`$ to point to the vertices of an equilateral triangle centered on the origin. The relevant spatial averages are
| $`\varphi ^2`$ | $`=2{\displaystyle \underset{i}{}}|A_i|^2`$ |
| --- | --- |
| $`\varphi \varphi \varphi `$ | $`=3\mathrm{cos}{\displaystyle \frac{p_1p_2}{2}}(A_1A_2A_3+A_1^{}A_2^{}A_3^{})`$ |
| $`\varphi \varphi \varphi \varphi `$ | $`=6\left({\displaystyle \underset{i}{}}|A_i|^2\right)^2+4(1+2\mathrm{cos}p_1p_2){\displaystyle \underset{i<j}{}}|A_iA_j|^2`$ |
| $`_{\mathrm{eff}}`$ | $`=\kappa _2{\displaystyle \underset{i}{}}|A_i|^2+\kappa _3\mathrm{cos}{\displaystyle \frac{p_1p_2}{2}}(A_1A_2A_3+A_1^{}A_2^{}A_3^{})`$ |
| | $`+{\displaystyle \frac{3}{2}}\kappa _4\left({\displaystyle \underset{i}{}}|A_i|^2\right)^2+\kappa _4(1+2\mathrm{cos}p_1p_2){\displaystyle \underset{i<j}{}}|A_iA_j|^2.`$ |
(50)
As long as the cubic term is present, a first order transition to a hexagonal lattice is the expected behavior. To be more precise, the expectation is that with $`\kappa _2`$ sufficiently small but positive, there is a minimum where $`_{\mathrm{eff}}`$ is negative, and where all three $`A_i`$ are nonzero. To see that this must be so, set $`A_i=An_i`$ with $`_in_i^2=1`$. Then
| $`_{\mathrm{eff}}`$ | $`=\stackrel{~}{\kappa }_2A^2+\stackrel{~}{\kappa }_3A^3+\stackrel{~}{\kappa }_4A^4`$ |
| --- | --- |
| $`\stackrel{~}{\kappa }_2`$ | $`=\kappa _2`$ |
| $`\stackrel{~}{\kappa }_3`$ | $`=2\kappa _3\mathrm{cos}{\displaystyle \frac{p_1p_2}{2}}n_1n_2n_3`$ |
| $`\stackrel{~}{\kappa }_4`$ | $`=\kappa _4\left[{\displaystyle \frac{3}{2}}+(1+2\mathrm{cos}p_1p_2){\displaystyle \underset{i<j}{}}n_i^2n_j^2\right].`$ |
(51)
One may easily show that $`\stackrel{~}{\kappa }_4>0`$, so there is no runaway behavior. It is straightforward to show that $`_{\mathrm{eff}}`$ attains negative values precisely if $`\sqrt{\stackrel{~}{\kappa }_2\stackrel{~}{\kappa }_4}<2|\stackrel{~}{\kappa }_3|`$. This is the desired result: even assuming a small cubic interaction, the first order transition at finite positive $`\kappa _2`$ swamps the would-be second order behavior at $`\kappa _2=0`$. The only exception is when $`\mathrm{cos}\frac{p_1p_2}{2}=0`$: at this point, the cubic term vanishes. A solution with $`|A_1|=|A_2|=|A_3|`$ is still preferred, but the phase of $`A_1A_2A_3`$ is a flat direction. On general Landau-Ginzburg theory grounds one might expect a second order point at $`\mathrm{cos}\frac{p_1p_2}{2}=0`$ separating the two first order lines corresponding, respectively, to phase locking of the three plane waves with $`A_1A_2A_3`$ positive or negative; however it seems likely that fluctuations again drive the transition first order. In higher dimensions the story is again more complicated. The usual expectation is a lattice with the maximal number of equilateral triangles; in the presence of non-commutativity these lattices are likely to have deformations and preferred orientations. We leave an investigation of the possibilities for future work.
## 5 The non-commutative $`O(N)`$ vector model
Naively, there are two reasons to think that quantum field theories defined on non-commutative spaces are under better control perturbatively than their commutative counterparts. First, non-planar loop diagrams include an oscillatory factor in the integrand which generically cures ultraviolet divergences . Second, compactifying the non-commutative position space also entails a compactification in momentum space : this is one of several manifestations of the interplay between ultraviolet and infrared effects. But there are also reasons to think that such quantum field theories are not under good perturbative control. First and perhaps most seriously, it has not been shown that the special form of the tree level lagrangian (polynomial in derivatives and fields, but with all ordinary products replaced by star products) is preserved by quantum corrections beyond one loop. Second, planar divergences are just as bad as for commutative field theories.
The latter two difficulties can be resolved for special quantum field theories. Consider the $`O(N)`$ vector model:
$$S=d^dx\left[\frac{1}{2}(\varphi _i)^2+\frac{1}{2}m^2\varphi _i^2+\frac{g^2}{4}\varphi _i\varphi _i\varphi _j\varphi _j+\frac{g^2}{4}\delta ^{ik}\delta ^{jl}\varphi _i\varphi _j\varphi _k\varphi _l\right],$$
(52)
which, in the large $`N`$ limit, is dominated by bubble graphs in its disordered phase. If $`g^2=0`$, the loop integrals in these graphs converge even without an ultraviolet cutoff, due to the oscillatory factor in the loop integrands. Thus there are no counterterms in the large $`N`$ limit: the theory is perfectly finite. At subleading orders in $`N`$, counterterms do arise, and all the usual questions arise regarding whether the special form of the lagrangian is preserved in the quantum theory. However it can perhaps be regarded as progress toward deciding the issue of perturbative renormalizability that these difficulties can be suppressed by $`1/N`$ in an appropriate ’t Hooftian limit.
Returning to the bubble sum for (52) with $`g^2=0`$, this can be reduced to an integral equation—the 1PI two-point function is given implicitly by
$$\mathrm{\Gamma }^{(2)}(p)=p^2+m^2+\frac{d^dk}{(2\pi )^d}\frac{g^2Ne^{ipk}}{\mathrm{\Gamma }^{(2)}(k)}.$$
(53)
In other words, the self-consistent Hartree approximation in the disordered phase is exact at large $`N`$. The existence of an ordered phase, following our earlier caveat, remains an open question.
It may be that (52) with $`g^2=0`$ is a special case of a more general strategy for generating quantum field theories with improved convergence properties in a special limit—the main ingredient being two conflicting notions of planarity, one from the index structure of the fields, and one from the star product. However, we are not guaranteed to obtain a finite quantum field theory (even in special limits) through this trick, as there may be graphs which are planar in both senses. For example, the lagrangian
$$=\frac{1}{2}tr(\mathrm{\Phi })^2+\frac{1}{2}m^2tr\mathrm{\Phi }^2+\frac{g^2}{4}\delta ^{i_2k_1}\delta ^{k_2j_1}\delta ^{j_2l_1}\delta ^{l_2i_1}\mathrm{\Phi }_{i_1i_2}\mathrm{\Phi }_{j_1j_2}\mathrm{\Phi }_{k_1k_2}\mathrm{\Phi }_{l_1l_2},$$
(54)
where $`\mathrm{\Phi }`$ is an $`N\times N`$ hermitian matrix, specifies a quantum field theory with two conflicting notions of planarity; but there are still certain planar diagrams diverge, for instance the one-loop correction to the quartic interaction vertex. This divergence requires a counterterm of the form $`tr(\mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi }\mathrm{\Phi })`$, so we learn that the special form of the lagrangian (54) is not preserved by quantum corrections.
## 6 Continuum limits
We turn now the question of taking a continuum limit, i.e. the question of renormalizability. We will examine this in the disordered phase of the theory within our self-consistent Hartree approach. Formally, we will consider the renormalization of the infinite $`N`$ theory, (52), with $`g^2N`$ and $`g^2N`$ finite. In this approximation, the only non-trivial correlation function that enters the theory is $`\mathrm{\Gamma }^{(2)}(p)`$ but on account of the noncommutativity it has significant momentum dependence that makes the procedure somewhat more complicated than it might seem at first sight. Nevertheless, the claim is that there are non-trivial continuum limits, and therefore interacting renormalizable field theories, in any even dimension. It was conjectured in that non-commutative $`\varphi ^4`$ theory should be renormalizable in $`d=4`$, despite the infrared divergences. The results of this section do not amount to a demonstration of this claim, because we persist in working at infinite $`N`$; however we hope that an extension to finite $`N`$ may be possible. We should note though, that our Wilsonian attempt to renormalize by means of the planar theory critical point is closely connected to the the planar subtraction algorithm suggested in .
As a warmup recall the problem of the commuting, infinite $`N`$, $`O(N)`$ vector model. Here,
$$\mathrm{\Gamma }^{(2)}(p)=p^2+m^2+\mathrm{\Sigma }$$
(55)
where
$$\mathrm{\Sigma }=\frac{g^2}{2}^\mathrm{\Lambda }\frac{d^dk}{(2\pi )^d}\frac{1}{k^2+m^2+\mathrm{\Sigma }},$$
(56)
where we have indicated the regulation of the integral explicitly. In $`2<d<4`$, this theory has a critical point $`m^2=m_c^2(\mathrm{\Lambda })`$ at which the renormalized mass $`m^2+\mathrm{\Sigma }`$ vanishes. Explicitly,
$$m_c^2=\frac{g^2}{2}^\mathrm{\Lambda }\frac{d^dk}{(2\pi )^d}\frac{1}{k^2},$$
(57)
but that is not central.
From the existence of the critical point it follows that we may choose $`m^2(\mathrm{\Lambda })`$ in order that $`m^2(\mathrm{\Lambda })+\mathrm{\Sigma }(g^2,\mathrm{\Lambda })=M^2`$ is held fixed,<sup>14</sup><sup>14</sup>14This is the point in the argument where the existence of a critical point in the bare theory is crucial. As we know that we can tune the bare mass to get $`M^2`$ to vanish at any value of the cutoff, it follows that we can hold it fixed at a specified value as well. whereupon we may take $`\mathrm{\Lambda }`$ to infinity and obtain a renormalized propagator $`\mathrm{\Gamma }_R^{(2)}(p)=p^2+M^2`$. In this example the end product of this mass renormalization is trivial, but the underlying lattice problem is not. A non-trivial correlation length exponent $`\nu `$ is hidden in the relation between the bare mass and the renormalized mass.
Turning now to the noncommuting infinite $`N`$ theory, we can divide $`\mathrm{\Sigma }(p)`$ into the parts coming from the planar and non-planar diagrams and rewrite the self consistency equation in two parts:
$$M^2=m^2+2g^2N^\mathrm{\Lambda }\frac{d^dk}{(2\pi )^d}\frac{1}{k^2+M^2+\mathrm{\Sigma }_{\mathrm{np}}(k)}$$
(58)
and
$$\mathrm{\Sigma }_{\mathrm{np}}(p)=g^2N^\mathrm{\Lambda }\frac{d^dk}{(2\pi )^d}\frac{e^{ikp}}{k^2+M^2+\mathrm{\Sigma }_{\mathrm{np}}(k)}.$$
(59)
Note that $`M^2m^2\mathrm{\Sigma }(\mathrm{})`$ introduced earlier.
Now, the planar theory has a critical point at $`m^2=m_c^2`$ and $`\theta \mathrm{\Lambda }^2=\mathrm{}`$ where $`M^2`$ vanishes. Keeping a fixed dimensionful $`\theta `$ then sends $`\theta \mathrm{\Lambda }^2`$ to infinity automatically, while the critical point enables us to choose $`m^2(\mathrm{\Lambda };\theta )`$ such that $`M^2`$ is held fixed as $`\mathrm{\Lambda }\mathrm{}`$.
In this limit we are able to remove the cutoff from (59) as well which still defines a finite $`\mathrm{\Sigma }_{\mathrm{np}}(p)`$ and thereby a finite renormalized two point function $`\mathrm{\Gamma }_R^{(2)}(p)=p^2+M^2+\mathrm{\Sigma }_{\mathrm{np}}(p)`$. Note that by dimensional analysis $`\mathrm{\Sigma }_{\mathrm{np}}`$ has the scaling form $`\mathrm{\Sigma }_{\mathrm{np}}=M^2f(p/M,\theta M^2)`$. Note also, that $`\mathrm{\Sigma }_{\mathrm{np}}\frac{1}{(\theta p)^{d2}}`$ is divergent in the infrared.<sup>15</sup><sup>15</sup>15Note that this is the same integral equation that arose in (53), albeit with $`M^2`$ replacing $`m^2`$.
In this fashion we see that the existence of the planar theory critical point does allow a continuum limit to be taken at $`N=\mathrm{}`$. There is one more limiting theory in this case, the massless purely planar theory obtained by sitting exactly at the critical point. (If the $`N=\mathrm{}`$ theory has a broken symmetry phase then a third continuum limit would be feasible from within that phase.) We expect that the planar limit does not change much between $`N=\mathrm{}`$ and $`N=1`$ in $`d4`$. That suggests that more progress could be made in establishing renormalizability via an $`1/N`$ expansion—the same logic underlies the comparable demonstration for commutative scalar theories. However, we have not examined this question in any detail.
Note also the interesting feature that the continuum limit is nontrivial in the infrared while the ultraviolet behavior is free. The latter is consistent with the triviality of commuting scalar theories in $`d4`$ although in this case it cannot be distinguished from the vanishing of the anomalous dimension of the scalar field at $`N=\mathrm{}`$. At any rate, the point is that the continuum limit will be nontrivial even above $`d=4`$ (say in $`d=6`$) for finite noncommutativity is, in a sense, a relevant perturbation at the planar theory fixed point. This is in contrast to the situation in the commuting case where only the free massive theory is possible.
Finally, we should note that a different set of continuum limits is possible near the Lifshitz point in the $`N=1`$ theory. These would entail keeping $`\theta \mathrm{\Lambda }^2`$ finite as $`\mathrm{\Lambda }\mathrm{}`$ and hence a vanishing $`\theta `$. We have not studied these, but our large $`N`$ approach could be extended above $`d=4`$ should interest in these theories be warranted.
## 7 Concluding Remarks
The purpose of this paper has been to investigate unusual phase structure in the simplest interacting non-commutative field theory, namely $`\varphi ^4`$ theory. We expect the existence of stripe phases to be quite common in non-commutative theories, the reason being that $`\mathrm{\Gamma }^{(2)}(p)`$ typically has singular behavior as $`p0`$ when the cutoff is removed. If $`\mathrm{\Gamma }^{(2)}(p)\mathrm{}`$ as $`p0`$, the theory is sick in the sense that for no value of bare masses have we found a stable vacuum. If $`\mathrm{\Gamma }^{(2)}(p)+\mathrm{}`$, then some version of the arguments of Section 4 should establish a first order transition to a stripe phase.
It is natural to ask, what manifestation might this transition find in string theory? At present, we have no definite answer; but let us remark on one obvious venue where such a transition might be expected to arise. D-branes in bosonic string theory can, in a rough approximation, be thought of as classical lumps of an open string tachyon field $`T`$ (see for example ). The potential for the tachyon field is cubic, so non-commutativity (in the form of a $`B_{\mu \nu }`$-field) merely destabilizes the $`T=0`$ vacuum further. Unstable D-branes in type II superstrings, however, have a quartic potential for the tachyon field, and it is possible that a stripe phase will arise for sufficiently large $`B_{\mu \nu }`$. If a stripe phase indeed exists, it would be quite a peculiar string background, not readily comprehensible in classical terms. Classically speaking, the stripe phase would amount to an alternating arrangement of stable D-branes and their anti-branes, parallel to one another. Such an arrangement seems obviously unstable toward the branes and anti-branes collapsing into one another. There are however some caveats: first, in working with the non-commutative field theory as an approximation to the string dynamics, we are by assumption working in a limit where the closed string interactions are suppressed. So the collapse of branes into anti-branes could have a much longer time scale than the transition to a stripe phase from a $`T=0`$ phase. Second, working directly with a field theory is suspect because the cutoff ($`1/\sqrt{\alpha ^{}}`$) is on the same order as the tachyon mass. A full string theory computation, with finite $`\alpha ^{}`$, seems to be the only wholly reliable approach.
We should note that our description of stripe phases may not incorporate an important piece of the physics of unstable D-branes: namely that the open string degrees of freedom are believed to be confined when the tachyon field has condensed. The mechanism for this confinement is not yet wholly clear (see for two interesting proposals), but it could play a role in the stability of various phases.
A relatively well-understood aspect of unstable D-branes in a strong background $`B_{\mu \nu }`$ field is condensation into a semiclassical configuration known as the non-commutative soliton . It has been argued that this configuration represents the collapse of a bosonic Dp-brane into a D(p-2)-brane; see also for closely related work. The role of these objects in our quantum considerations is not entirely clear. For one thing, the stable solitons/instantons at infinite $`\theta `$ found in by leaving out the gradient term in the action, have exactly zero action in our case which makes it imperative to keep the gradient term in the analysis. Should such an improved computation be feasible, we would still suspect that in the large $`N`$ limit these objects will be unstable to unwinding via other directions in field space and hence do not play a role.
In the $`N=1`$ case, the solitons/instantons could exist. As instantons we would have to consider whether they invalidate our conclusions on the nature of the ordered phases. This seems unlikely, at least at weak coupling where the wavelength of our stripe phases is much larger, by a factor of $`1/\sqrt{g}`$, than the size of the instantons in the classical analysis. More generally, they are local distortions of the condensate, but do not (unlike, say vortices in an XY model) affect the ordering at large distances. Hence we expect that they will renormalize the properties of the ordered phases but not lead to any singular effects.
In contrast, actual solitons in a quantum theory in odd dimensions would rely on the existence of an ordered phase. In the classical analysis, this is assumed to be a uniform condensate but as we have seen this is no longer the case in the quantum theory. Again, at weak coupling, it is plausible that the solitons are essentially undistorted on account of their size, but their dynamics would clearly be affected by motion in an inhomogeneous background.
It is perhaps worth recalling some features of quantum Hall physics that may find a formulation as noncommutative field theories. The first is the physics of the $`\nu =1/2`$ state which has an elegant microscopic interpretation in terms of dipoles<sup>16</sup><sup>16</sup>16For a recent discussion including an invocation of non-commutative geometry, see . with evident parallels to the discussion of Bigatti and Susskind . We note that stripe phases do occur in quantum Hall systems at half filling (in high Landau levels) but hasten to add that any connection between them, the dipolar $`\nu =1/2`$ theory and our noncommutative results is purely speculative at this point.
An equally speculative connection is to a widely studied model of the integer quantum Hall effect, namely Pruisken’s non-linear sigma model. It has several incarnations, but the simplest has a lagrangian of the form
$$=\frac{1}{4}\sigma _{xx}^{(0)}tr(T)^2+\frac{1}{8}\sigma _{xy}^{(0)}tr\left(T[_xT,_yT]\right),$$
(60)
where $`T`$ takes values in the coset $`U(2N)/(U(N)\times U(N))`$. The second term is topological. The model arises from a replica treatment of disorder in a theory of non-interacting electrons in a magnetic field. In the end, physical quantities must be computed in a $`N0`$ limit. In the original formulation, the RG flow in the $`(\sigma _{xy},\sigma _{xx})`$ plane was argued to have fixed points at $`\sigma _{xx}=0`$ and $`\sigma _{xy}𝐙`$ and at $`\sigma _{xx}=1/2`$ and $`\sigma _{xy}𝐙+1/2`$.<sup>17</sup><sup>17</sup>17However, Zirnbauer has argued persuasively that the fixed point theories have instead a Wess-Zumino-Witten form. The first set of fixed points represent the quantum Hall plateaux, and the second set control the critical behavior of transitions between plateaux. The behavior $`\sigma _{xx}0`$ in the infrared is the analog of confinement in this model, and the special fixed points at non-zero $`\sigma _{xx}`$ are believed to arise from instanton effects, similar to ’t Hooft’s treatment of deconfinement in QCD at $`\theta =\pi `$.
It is suggestive that precisely the $`U(2N)/(U(N)\times U(N))`$ coset arises as the vacuum manifold in the condensation of unstable D-branes in type II string theory. However it is difficult to see how the topological term $`tr\left(T[_xT,_yT]\right)`$ can arise in the string effective action. This term is the leading order term in a derivative expansion of $`tr\left(T^3T^3\right)`$; however, from a type II string theory point of view, neither the mixing of star products with ordinary products inside a single trace, nor the cubic form of this potential, is expected. Nevertheless, strictly from a field theory point of view, it would be interesting to ask whether replacing the topological term with $`tr\left(T^3T^3\right)`$ preserves the universality class. We hope to return to these issues in the future.
## Acknowledgements
We would like to thank D. Huse for many useful discussions and for pointers to the literature on the Brazovskii transition. The research of S.S.G. was supported in part by DOE grant DE-FG02-91ER40671 while S.L.S. is grateful to the NSF (grant DMR 99-7804), the Alfred P. Sloan Foundation and the David and Lucille Packard Foundation for support.
## Appendix: The Brazovskii Transition by Non-Brazovskiian Means
Here we will give a version of the argument used in Section 4.3.2 for the transition studied originally by Brazovskii . The relevant action is (37), which we rewrite as
| $`S_{\mathrm{eff}}`$ | $`={\displaystyle d^dx_{\mathrm{eff}}}`$ |
| --- | --- |
| $`_{\mathrm{eff}}`$ | $`={\displaystyle \frac{\kappa }{2}}((^2+p_c^2)\varphi )^2+{\displaystyle \frac{1}{2}}m^2\varphi ^2+{\displaystyle \frac{g^2}{4!}}\varphi ^4.`$ |
(61)
The traditional route to showing that there is a first order transition is a weak coupling analysis. Here we will work at fixed coupling but instead take $`p_c`$ to be the “small parameter”. For this to work we will need that the theory at $`p_c=0`$ possess a continuous transition to a broken symmetry phase, else nothing is gained. In the self-consistent Hartree approximation of Brazovskii for the high temperature phase,
$$\mathrm{\Sigma }=\frac{g^2}{2}\frac{d^dk}{(2\pi )^d}\frac{1}{\kappa (k^2p_c^2)^2+m^2+\mathrm{\Sigma }}$$
(62)
this requires that we take $`d>4`$.<sup>18</sup><sup>18</sup>18Alternatively, we can modify the Brazovskii action by replacing $`(k^2p_c^2)^2`$ by $`(kp_c)^2`$ which preserves the feature of a codimension one surface of soft modes while retaining the connection to a critical point in $`d>2`$. In these dimensions, the critical point of the uniform ($`p_c=0`$) theory is at $`m_c^2=\frac{g^2}{2}\frac{d^dk}{(2\pi )^d}\frac{1}{\kappa k^4}`$.
Now consider $`p_c>0`$ at fixed $`g^2`$. By the standard argument sketched in Section 4.3.1, there is always a disordered solution at any $`m^2/\mathrm{\Lambda }^2>\mathrm{}`$, with a free energy (written in terms of self-consistent parameters)
$$F=\frac{\mathrm{\Sigma }^2}{2g^2}+\frac{1}{2}\frac{d^dk}{(2\pi )^d}\mathrm{log}[\kappa (k^2p_c^2)^2+m^2+\mathrm{\Sigma }],$$
(63)
whence a continuous transition out of the high temperature phase is impossible.
Next, we follow the disordered solution as we take $`p_c`$ to zero. For $`m^2>m_c^2`$ the solution smoothly connects to the disordered solution of the uniform theory with $`\mathrm{\Sigma }+m^2>0`$. But for $`m^2<m_c^2`$ this is not possible as there is only an ordered solution. So in the vicinity of the $`p_c=0`$ line we must have two solutions, where the second evolves out of the broken symmetry solution of the uniform theory as $`p_c`$ is tuned away from zero.
The remaining task is to show that for $`m^2<m_c^2`$ there is always a value of $`p_c`$ below which the ordered solution wins—this then shows, as in the example of the noncommutative theory in Section 4, that there is a first order line in the $`m^2,p_c`$ plane (see Figure 6). For this again it suffices to make the comparison in the limit $`p_c0`$ on the grounds that free energies are continuous.
The disordered solution is readily shown to require that $`m^2+\mathrm{\Sigma }0`$ as $`p_c0`$ whenever $`m^2<m_c^2`$. Hence its free energy has the limit,
$$F=\frac{(m^2)^2}{2g^2}+\frac{1}{2}\frac{d^dk}{(2\pi )^d}\mathrm{log}[\kappa k^4].$$
(64)
The free energy of the ordered solution at $`p_c=0`$ is
$$F=(\frac{m^2}{2}\varphi _o^2+\frac{g^2}{4!}\varphi _o^4)\frac{(m^2+\frac{g^2}{6}\varphi _o^2)^2}{2g^2}+\frac{1}{2}\frac{d^dk}{(2\pi )^d}\mathrm{log}[\kappa k^4+\frac{g^2}{3}\varphi _o^2]$$
(65)
where the parameters satisfy,
$$m^2+\frac{g^2}{6}\varphi _o^2+\frac{g^2}{2}\frac{d^dk}{(2\pi )^d}\frac{1}{\kappa k^4+\frac{g^2}{3}\varphi _o^2}=0$$
(66)
and
$$\mathrm{\Sigma }=\frac{g^2}{3}\varphi _o^2.$$
(67)
An explicit comparison, easily done perturbatively in $`\varphi _o^2`$, shows that the ordered solution has lower energy, which is what we set out to show.
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# Monodromy Filtration and Positivity
We study Deligne’s conjecture on the monodromy weight filtration on the nearby cycles in the mixed characteristic case, and reduce it to the nondegeneracy of certain pairings in the semistable case. We also prove a related conjecture of Rapoport and Zink which uses only the image of the Cech restriction morphism, if Deligne’s conjecture holds for a general hyperplane section. In general we show that Deligne’s conjecture is true if the standard conjectures hold.
monodromy filtration, polarization, standard conjecture
Introduction
<sup>1</sup><sup>1</sup>1
Date : Sept. 22, 2003, v.11
Let $`f:XS`$ be a projective morphism of complex manifolds with relative dimension $`n`$ where $`S`$ is an open disk and $`f`$ is semistable (i.e. $`X_0:=f^1(0)`$ is a reduced divisor with normal crossings whose irreducible components are smooth). J. Steenbrink constructed a limit mixed Hodge structure by using a resolution of the nearby cycle sheaf on which the monodromy weight filtration can be defined. This limit mixed Hodge structure coincides with the one obtained by W. Schmid using the $`\text{SL}_2`$-orbit theorem, because the weight filtration coincides with the (shifted) monodromy filtration, see and also , , , \[24, 2.3\], etc.
A similar construction was then given by M. Rapoport and T. Zink in the case $`X`$ is projective and semistable over a henselian discrete valuation ring $`R`$ of mixed characteristic. Here we may assume $`f`$ semistable because the non semistable case can be reduced to the semistable case by replacing the discrete valuation field with a finite extension if necessary. Then Deligne’s conjecture on the monodromy filtration \[3, I\] is stated as
0.1. Conjecture. The obtained weight filtration coincides with the monodromy filtration shifted by the degree of cohomology.
This conjecture is proved so far in the case $`n2`$ by (using ) and in some other cases (see e.g. , ). We cannot apply the same argument as in the complex analytic case, because we do not have a good notion of positivity for $`l`$-adic sheaves. In the positive characteristic case, however, the conjecture was proved by Deligne (assuming $`f`$ is the base change of a morphism to a smooth curve over a finite field). In this paper we apply the arguments in , to the mixed characteristic case, and show that Conjecture (0.1) is closely related to the standard conjectures which would give a notion of positivity for the pairings of correspondences (but not for the pairings of cohomology groups in general).
Let $`Y^{(i)}`$ denote the disjoint union of the intersections of $`i`$ irreducible components of the special fiber $`X_0`$. Let $`k`$ be the residue field of $`R`$ which is assumed to be a finite field. Let $`\overline{k}`$ be an algebraic closure of $`k`$, and put $`Y_{\overline{k}}^{(i)}=Y^{(i)}_k\overline{k}`$. Let $`l`$ be a prime number different from the characteristic of $`k`$. Then the $`E_1`$-term of the weight spectral sequence of Rapoport and Zink is given by direct sums of the $`l`$-adic cohomology groups of $`Y_{\overline{k}}^{(i)}`$ which are Tate-twisted appropriately. Its differential $`d_1`$ consists of the Cech restriction morphisms and the co-Cech Gysin morphisms which are denoted respectively by $`\rho `$ and $`\gamma `$ in this paper. The primitive cohomology of $`Y_{\overline{k}}^{(i)}`$ has a canonical pairing induced by Poincaré duality together with the hard Lefschetz theorem (choosing and fixing an ample divisor class of $`f)`$.
0.2. Theorem. Conjecture (0.1) is true if the restrictions of the canonical pairing to the intersections of the primitive part with $`\text{Im}\rho `$ and with $`\text{Im}\gamma `$ are both nondegenerate.
So the problem is reduced to the study of the canonical pairing on the primitive part. There are some examples satisfying the assumption of (0.2), see (2.8). If $`n=3`$, the converse of (0.2) is also true and the hypothesis on $`\text{Im}\gamma `$ is always satisfied in this case (using , ). Note that (0.1) may apparently depend on the choice of $`l`$, and we can prove (0.1) for certain $`l`$ in a simple case where every eigenvalue of the Frobenius action has multiplicity $`1`$, see (2.7) and (5.4). (In this case we can determine the endomorphism ring of a simple motive, see (5.2).) However, the general case seems to be related closely to the standard conjectures.
A conjecture similar to (0.2) but using only the restriction morphisms was noted by Rapoport and Zink in the introduction of . We can prove this conjecture under an inductive hypothesis as follows:
0.3. Theorem. Assume that Conjecture (0.1) holds for a general hyperplane section of the generic fiber, and that the restriction of the modified pairing in to $`\text{Im}\rho `$ or the restriction of the canonical pairing to the intersection of the primitive part with $`\text{Im}\rho `$ is nondegenerate. Then (0.1) is true.
In the complex analytic case, the hypothesis of (0.2) is trivially satisfied because of the positivity of polarizations of Hodge structures. We can argue similarly if we have a kind of “positivity” in characteristic $`p>0`$. The positivity for a zero-dimensional variety (or a motive) is clear, because the pairing is defined over the subfield $``$ of $`_l)`$. In the one-dimensional case, this notion is provided by the theory of Riemann forms for abelian varieties (see also , ) combined with work of Deligne on the compatibility of the Weil pairing and Poincaré duality (see also (3.4) below). These were used in an essential way for the proof of Conjecture (0.1) for $`n=2`$ in . However, for higher dimensional varieties, we do not know any notion of positivity except the standard conjecture of Hodge index type . Using the theory on the standard conjectures (loc. cit.) we show
0.4. Theorem. Assume that the standard conjectures hold for $`Y^{(i)}`$, $`Y^{(i)}\times _kY^{(i)}`$, and the numerical equivalence and the homological equivalence coincide for $`Y^{(i)}\times _kY^{(i+1)}(i>0)`$. Then Conjecture (0.1) is true.
I am quite recently informed that if the generic fiber can be uniformized by the Drinfeld upper half space, Conjecture (0.1) is proved by T. Ito using Theorem (0.2).
Part of this work was done during my stay at the University of Leiden. I thank Professor J. Murre for useful discussions, and the staff of the institute for the hospitality. I also thank Takeshi Saito for a useful discussion about the consequence of de Jong’s theory of alternations.
In Sect. 1 and Sect. 2, we review the theory of graded or bigraded modules of Lefschetz-type and prove (0.2) and (0.3). In Sect. 3, we review the work of Deligne on the compatibility of the Weil pairing and Poincaré duality. In Sect. 4, we review the standard conjectures and prove (0.4). In Sect. 5, we study the Frobenius action, and prove (0.1) in some simple cases.
1. Graded Modules of Lefschetz-Type
We first review a theory of morphisms of degree $`1`$ between graded modules of Lefschetz-type . A typical example is given by the restriction or Gysin morphism associated to a morphism of smooth projective varieties whose relative dimension is $`1`$.
1.1. In this and the next sections we denote by $`\mathrm{\Lambda }`$ a field. By a graded $`\mathrm{\Lambda }[L]`$-module we will mean a finite dimensional graded vector space $`M^{}`$ over $`\mathrm{\Lambda }`$ having an action of $`L`$ with degree $`2`$. We call $`M^{}`$ $`n`$-symmetric if
$$L^j:M^{nj}\stackrel{}{}M^{n+j}\text{for}j>0.$$
Here the Tate twists are omitted to simplify the notation. We say that $`M^{}`$ is a graded module of Lefschetz-type if it is a $`0`$-symmetric graded $`\mathrm{\Lambda }[L]`$-module. Then we have the Lefschetz decomposition
$$M^j=_{i0}L^i{}_{0}{}^{}M_{}^{j2i}\text{with}{}_{0}{}^{}M_{}^{j}=\text{Ker}L^{j+1}M^j.$$
We say that $`f:M^{}N^{}`$ is a morphism of degree $`m`$ between graded modules of Lefschetz-type, if $`f(M^j)N^{j+m}`$ and $`f`$ is $`\mathrm{\Lambda }[L]`$-linear. Shifting the degree of $`N^{}`$, it is identified with a graded morphism of a $`0`$-symmetric graded $`\mathrm{\Lambda }[L]`$-module to a $`(m)`$-symmetric graded $`\mathrm{\Lambda }[L]`$-module. Considering the image of the primitive part, we see that a morphism of degree $`0`$ preserves the Lefschetz decomposition and there is no nontrivial morphism of negative degree (i.e. a morphism of a $`0`$-symmetric module to an $`n`$-symmetric module for $`n>0`$ is trivial). For a morphism of degree $`1`$, we see that
$$f({}_{0}{}^{}M_{}^{j}){}_{0}{}^{}N_{}^{j+1}+L{}_{0}{}^{}N_{}^{j1}.$$
$`(\mathrm{1.1.1})`$
Let $`f:M^{}N^{}`$ be a morphism of degree $`1`$ between graded modules of Lefschetz-type. We define
$$\text{Im}^0f=_j(_{i0}L^i(\text{Im}f{}_{0}{}^{}N_{}^{j})),\text{Im}^1f=\text{Im}f/\text{Im}^0f.$$
Then $`\text{Im}^0f`$ is $`0`$-symmetric, and its primitive part is given by $`\text{Im}f{}_{0}{}^{}N_{}^{j}`$. It has been remarked by T. Ito that the proof of the next lemma easily follows from the definition itself.
1.2. Lemma. The quotient module $`\text{Im}^1f`$ is $`1`$-symmetric, i.e. we have the bijectivity of
$$L^{j+1}:(\text{Im}^1f)^j(\text{Im}^1f)^{j+2}\text{for}j0,$$
$`(\mathrm{1.2.1})`$
where the upper index $`j`$ means the degree $`j`$ part, etc.
Proof. The surjectivity follows from the $`0`$-symmetry of $`M`$. To show the injectivity, let $`n=_{i0}L^in_i(\text{Im}f)^j`$ with $`n_i{}_{0}{}^{}N_{}^{j2i}`$, and assume $`L^{j+1}n(\text{Im}^0f)^{j+2}`$. Since $`L^{j+1}n=_{i>0}L^{j+i+1}n_i`$ and $`\text{Im}^0f`$ is $`0`$-symmetric, we may assume $`n_i=0`$ for $`i>0`$ by modifying $`n`$ modulo $`\text{Im}^0f`$ if necessary. Then $`n\text{Im}^0f`$ by definition, and the assertion follows.
1.3. Proposition. Let $`M^{},N^{}`$ be graded $`\mathrm{\Lambda }[L]`$-modules of Lefschetz-type, and
$$f:M^{}N^{},g:N^{}M^{}$$
be morphisms of degree one. Assume there are nondegenerate pairings of $`\mathrm{\Lambda }`$-modules
$$\mathrm{\Phi }_M:M^j_\mathrm{\Lambda }M^j\mathrm{\Lambda },\mathrm{\Phi }_N:N^j_\mathrm{\Lambda }N^j\mathrm{\Lambda },$$
such that $`\mathrm{\Phi }_N(fid)`$ coincides with $`\mathrm{\Phi }_M(idg)`$ up to a nonzero multiple constant and that $`\mathrm{\Phi }_M(idL)=\mathrm{\Phi }_M(Lid)`$ (and the same for $`\mathrm{\Phi }_M)`$. Then
$$dim(\text{Im}^0f)^j=dim(\text{Im}^1g)^{j+1},dim(\text{Im}^0g)^j=dim(\text{Im}^1f)^{j+1}.$$
Proof. Let $`a_j,b_j,c_j,d_j`$ denote respectively the above dimensions so that $`a_j=a_j`$, $`b_j=b_j`$, etc. By the duality of $`f`$ and $`g`$, we have
$$a_{j+1}+d_j=dim(\text{Im}f)^{j+1}=dim(\text{Im}g)^j=b_{j1}+c_j=b_{j+1}+c_j.$$
If we put $`p_j=a_jb_j`$, $`q_j=c_jd_j`$, we get $`p_{j+1}=q_j`$ for any $`j`$. So $`p_j=q_j=0`$ by the symmetry of $`p_j`$, $`q_j`$, and the assertion follows.
1.4. Lemma. With the notation and assumption of (1.3), assume that $`n\text{Im}f{}_{0}{}^{}N_{}^{j}`$ vanishes if $`\mathrm{\Phi }_N(f(m),L^jn)=0`$ for any $`mM^{j1}`$. Then
$$\text{Ker}g\text{Im}^0f=0.$$
$`(\mathrm{1.4.1})`$
Proof. Let $`n=_{0ir}L^in_i`$ with $`n_i{}_{0}{}^{}N_{}^{j2i}`$, and assume $`g(n)=0`$. Since
$$L^{j+2r}g(n_r)=L^{j+r}g(n)=0,$$
the assertion is reduced by induction on $`r`$ to the injectivity of
$$gL^j:\text{Im}f{}_{0}{}^{}N_{}^{j}M^{j+1}.$$
But this injectivity is clear, because $`\mathrm{\Phi }_N(f(m),L^jn)`$ is equal to $`\mathrm{\Phi }_M(m,g(L^jn))`$ up to a nonzero multiple constant. So the assertion follows.
1.5. Proposition. With the notation and assumption of (1.3), assume the vanishing of $`fgf`$ and the injectivity of
$$f:\text{Im}^0gN^{},g:\text{Im}^0fM^{}.$$
$`(\mathrm{1.5.1})`$
Then the compositions
$$\text{Im}^0g\stackrel{𝑓}{}\text{Im}f\text{Im}^1f,\text{Im}^0f\stackrel{𝑔}{}\text{Im}g\text{Im}^1g$$
$`(\mathrm{1.5.2})`$
are isomorphisms, and we have canonical decompositions
$$\text{Im}f=\text{Im}^0f\text{Im}^1f,\text{Im}g=\text{Im}^0g\text{Im}^1g,$$
$`(\mathrm{1.5.3})`$
such that the restriction of $`g`$ to $`\text{Im}^0f`$ is injective, and that to $`\text{Im}^1f`$ is zero, and similarly for the restriction of $`f`$. Furthermore, we have canonical isomorphisms
$$\text{Im}fg=\text{Im}^1f,\text{Im}gf=\text{Im}^1g.$$
$`(\mathrm{1.5.4})`$
Proof. The hypothesis implies $`f(\text{Im}^0g)\text{Im}^0f=g(\text{Im}^0f)\text{Im}^0g=0`$, because the vanishing of $`fgf`$ is equivalent to that of $`gfg`$ by duality. So the assertion follows from (1.3).
1.6. Proposition. With the notation and the assumptions of (1.3), assume further that (1.4.1) holds, $`\text{Im}fg`$ is $`1`$-symmetric, $`gfg=0`$, and
$$(\text{Ker}g\text{Im}f)^j=(\text{Im}fg)^j\text{for}j0.$$
$`(\mathrm{1.6.1})`$
Then (1.6.1) holds also for $`j=0`$.
Proof. Since $`\text{Ker}g\text{Im}f\text{Im}fg`$, it is enough to show the coincidence of the images of both sides of (1.6.1) in $`\text{Im}^1f`$ using (1.4.1). We can identify
$$(\text{Ker}g\text{Im}f)/\text{Im}fg$$
$`(\mathrm{1.6.2})`$
with a graded $`\mathrm{\Lambda }[L]`$-submodule of $`\text{Im}^1f/\text{Im}fg`$, and the last module is $`1`$-symmetric by hypothesis and (1.2). Furthermore (1.6.2) vanishes except for the degree $`0`$ by hypothesis. So it vanishes at any degree, and the assertion follows.
2. Bigraded Modules of Lefschetz-Type
We prove (0.2–3) after reviewing a theory of bigraded modules of Lefschetz-type , . These modules appear in the $`E_1`$-term of the Steenbrink-type spectral sequence associated to a semi-stable degeneration.
2.1. Let $`M^,`$ be a bigraded $`\mathrm{\Lambda }[N,L]`$-module of Lefschetz-type, i.e. it is a finite dimensional bigraded vector space over $`\mathrm{\Lambda }`$ having commuting actions of $`N,L`$ with bidegrees $`(2,0)`$ and $`(0.2)`$ respectively such that
$$N^i:M^{i,j}\stackrel{}{}M^{i,j}(i>0),L^j:M^{i,j}\stackrel{}{}M^{i,j}(j>0).$$
Put $`M_i^j=M^{i,j}`$, and $`{}_{(0)}{}^{}M_{i}^{j}=\text{Ker}N^{i+1}M_i^j`$ for $`i0`$, and $`0`$ otherwise. Then we have the Lefschetz decomposition for the first index:
$$M_i^j=_{a0}N^a{}_{(0)}{}^{}M_{i+2a}^{j}.$$
$`(\mathrm{2.1.1})`$
We define $`C_i^j={}_{(0)}{}^{}M_{i}^{j}`$ and
$$C_{i,a}^j=N^aC_{i+a}^jM_{ia}^j,$$
so that $`C_{i,a}^j=0`$ for $`i<0`$ or $`a<0`$, and $`C_{i,0}^j=C_i^j`$. Then
$$N^{ia}:C_{i,a}^j\stackrel{}{}C_{a,i}^j(i>a),L^j:C_{i,a}^j\stackrel{}{}C_{i,a}^j(j>0).$$
$`(\mathrm{2.1.2})`$
Let $`d`$ be a differential of bidegree $`(1,1)`$ on $`M^,`$ which commutes with $`N,L`$ and satisfies $`d^2=0`$. By (1.1.1) applied to the action of $`N`$, we have a decomposition $`d=d^{}+d^{\prime \prime }`$ such that
$$d^{}:C_{i,a}^jC_{i1,a}^{j+1},d^{\prime \prime }:C_{i,a}^jC_{i,a+1}^{j+1}$$
are differentials which anti-commute with each other. Let $`C_i^{}=_jC_i^j`$. We have morphisms of degree $`1`$ between graded $`\mathrm{\Lambda }[L]`$-modules
$$\gamma _i:C_i^{}C_{i1}^{}(i>0),\rho _i:C_i^{}C_{i+1}^{}(i0)$$
such that $`d^{},d^{\prime \prime }`$ are identified with $`\gamma _{ia},\rho _{ia}`$ respectively.
We define $`H_i^j=Z_i^j/B_i^j`$ with
$$Z_i^j=\text{Ker}(d:M_i^jM_{i1}^{j+1}),B_i^j=\text{Im}(d:M_{i+1}^{j1}M_i^j).$$
Then we have the induced morphism $`N:H_i^jH_{i2}^j`$, and similarly for $`Z_i^j,B_i^j`$. For a positive integer $`i`$ and an integer $`j`$, we will consider the condition for the bijectivity of the morphism
$$N^i:H_i^jH_i^j\text{for}i>0.$$
$`(\mathrm{2.1.3})`$
For this we will assume that the action of $`L`$ induces bijections
$$L^j:H_i^j\stackrel{}{}H_i^j\text{for}j>0,$$
$`(\mathrm{2.1.4})`$
and that both terms of (2.1.3) have the same dimension (so that bijectivity is equivalent to injectivity and to surjectivity). Using (2.1.4), the last assumption is satisfied if we have a self-duality of $`M^,`$.
Since $`M_{i1}^{}=C_{i1}^{}NM_{i+1}^{}`$ and $`N^i:M_i^j\stackrel{}{}M_i^j(i>0)`$, we have a morphism
$$\stackrel{~}{d}:=dN^i:Z_i^jC_{i1}^{j+1}\text{for}i>0.$$
2.2. Proposition. (i) The surjectivity of (2.1.3) is equivalent to
$$\stackrel{~}{d}(Z_i^j)=(\text{Im}\gamma _i\rho _{i1})^{j+1}.$$
$`(\mathrm{2.2.1})`$
(ii) If (2.1.3) for $`(i+2,j)`$ is surjective, then the surjectivity of (2.1.3) is further equivalent to
$$\gamma _i(\text{Ker}\rho _i)^j=(\text{Im}\gamma _i\rho _{i1})^{j+1}.$$
$`(\mathrm{2.2.2})`$
(iii) If furthermore (2.1.3) for $`(i+2,j)`$ is surjective and (2.1.3) for $`(i+1,j+1)`$ is injective, then the surjectivity of (2.1.3) is further equivalent to
$$(\text{Ker}\rho _{i1}\text{Im}\gamma _i)^{j+1}=(\text{Im}\gamma _i\rho _{i1})^{j+1}.$$
$`(\mathrm{2.2.3})`$
(iv) If (2.1.3) for $`(i+1,j1)`$ is surjective, then the injectivity of (2.1.3) is equivalent to
$$(\text{Ker}\gamma _i\text{Im}\rho _{i1})^j=(\text{Im}\rho _{i1}\gamma _i)^j(=(\text{Im}\gamma _{i+1}\rho _i)^j).$$
$`(\mathrm{2.2.4})`$
Proof. We have by definition $`\stackrel{~}{d}(B_i^j)=(\text{Im}\gamma _i\rho _{i1})^{j+1}`$. Then the first assertion is clear because the surjectivity is equivalent to $`Z_i^j=B_i^j+N^iZ_i^j`$.
For the second, take $`m=_{a0}N^am_aN^iZ_i^j`$ with $`m_aC_{i+2a}^j`$ to calculate $`\stackrel{~}{d}(Z_i^j)`$. We may assume that $`m_0\text{Ker}\rho _i`$ applying (2.2.1) for $`i+2`$ to $`_{a1}N^{a1}m_a`$ and modifying $`m`$ by an element of $`d(C_{i+1}^{j1})`$ because this does not change $`dm`$. So the second assertion is clear.
Before showing the third assertion, we see that
$$(\text{Ker}\gamma _i\text{Im}\rho _{i1})^j/\stackrel{~}{d}(Z_{i1}^{j1})\stackrel{}{}(N^iB_i^jZ_i^j)/B_i^j,$$
because $`N^iB_i^j=(\text{Im}\rho _{i1})^j+B_i^j`$. Then the last assertion follows from (2.2.1).
For the third assertion, take $`\gamma _im(\text{Ker}\rho _{i1})^{j+1}`$, where $`mC_i^j`$ with $`\rho _{i1}\gamma _im=0`$. Then $`\rho _im\text{Ker}\gamma _{i+1}`$, and we may assume $`\rho _im=0`$ using (2.2.4) for $`(i+1,j+1)`$ and modifying $`m`$ by $`\gamma _{i1}m^{}`$ because it does not change $`\gamma _im`$. So the assertion follows from (2.2.2).
2.3. Application. With the notation of the introduction we may assume
$$C_i^j=H^{ni+j}(Y_{\overline{k}}^{(i+1)},_l(i)),H_i^j=\text{Gr}_{n+j+i}^WH^{n+j}(X_{\overline{K}},_l),$$
$`(\mathrm{2.3.1})`$
and $`\rho _i,\gamma _{i+1}`$ are respectively the Cech restriction and co-Cech Gysin morphisms, which are dual of each other up to a sign. In particular, $`\rho _j\gamma _{j+1}\rho _j=\gamma _{j+2}\rho _{j+1}\rho _j=0`$, i.e. the first hypothesis of Proposition (1.5) is satisfied. Note that $`N`$ is the logarithm of the monodromy, and $`L`$ is given by the ample divisor class of $`f`$, see , . So the bijectivity of (2.1.4) follows from the hard Lefschetz theorem for the generic fiber together with the strict compatibility of the weight filtration .
2.4. Conjecture of Rapoport and Zink. For a smooth projective $`\overline{k}`$-variety $`Y`$ with an ample line bundle $`L`$, we have a canonical pairing on $`H^j(Y,_l)`$ by Poincaré duality and the Hard Lefschetz theorem . Using further the Lefschetz decomposition and modifying the sign as in (4.1), we get a modified pairing on $`H^j(Y,_l)`$ as in . Rapoport and Zink noted there that Conjecture (0.1) would be verified if the restriction of this modified pairing to $`\text{Im}\rho `$ is nondegenerate. Note that the hypothesis of (1.4) is satisfied under the above hypothesis (where $`f=\rho _{i1}`$ and $`g=\gamma _i`$) because the modified pairing coincides with the canonical pairing up to a sign if one factor belongs to the primitive part.
2.5. Proof of (0.2) and (0.3). By , the $`E_1`$-term of the weight spectral sequence has a structure of bigraded $`_l`$-modules of Lefschetz-type, see (2.3). Then (2.2.3–4) follows from (1.4–5), and Theorem (0.2) follows.
For Theorem (0.3) we will show by decreasing induction on $`i`$ that (2.1.3) is bijective (or equivalently, injective) and $`\text{Im}\gamma _i\rho _{i1}`$ is $`1`$-symmetric. Assume the two assertions are true for $`i+r`$ with $`r>0`$. Considering the restriction morphism to a general hyperplane section of the generic fiber and using the weak Lefschetz theorem, we see that (2.1.3) is injective for $`j<0`$, because the restriction morphism is strictly compatible with the weight filtration. Then (2.1.3) is injective also for $`j>0`$ by (2.1.4), and it is injective for any $`j`$ by (2.2) (iv) and (1.6) where $`f=\rho _{i1}`$ and $`g=\gamma _i`$. So it remains to show that $`\text{Im}\gamma _i\rho _{i1}`$ is $`1`$-symmetric.
Since the injectivity of the first morphism of (1.5.1) follows from (1.4), the assertion is reduced by (1.5) to the injectivity of $`\rho _{i1}:\text{Im}^0\gamma _iC_i^{}[1]`$, i.e. to the vanishing of
$$\text{Ker}(\rho _{i1}:\text{Im}^0\gamma _i(\text{Im}\rho _{i1}\gamma _i)[1]).$$
$`(\mathrm{2.5.1})`$
Here $`[m]`$ means a shift of degree for an integer $`m`$, i.e. $`(M[m])^i=M^{i+m}`$. We see that (2.5.1) is $`0`$-symmetric, because $`\text{Im}^0\gamma _i`$ and $`(\text{Im}\rho _{i1}\gamma _i)[1]`$ are (using inductive hypothesis). Furthermore (2.5.1) is a submodule of $`\text{Ker}\rho _{i1}\text{Im}\gamma _i`$, and the latter is identified by (2.2.3–4) with
$$\text{Im}(\gamma _i:(\text{Im}\rho _{i1})[1]\text{Im}\gamma _i)=(\text{Im}\rho _{i1}/\text{Im}\rho _{i1}\gamma _i)[1].$$
The last module is an extension of a $`2`$-symmetric module by a $`1`$-symmetric module, because $`\text{Im}\rho _{i1}\gamma _i\text{Im}^0\rho _{i1}=0`$ by (1.4) and $`\text{Im}\rho _{i1}\gamma _i(=\text{Im}\gamma _{i+1}\rho _i)`$ is $`1`$-symmetric by inductive hypothesis. But there is no nontrivial morphism of a $`0`$-symmetric module to an $`n`$-symmetric module for $`n>0`$, see (1.1). So the assertion follows.
2.6. First cohomology case. The restriction of the pairing to the image of $`\rho _{i1}`$ in $`H^1(Y_{\overline{k}}^{(i+1)},_l)`$ is always nondegenerate by the abelian-positivity in Sect. 3, because $`\rho _{i1}`$ is associated to the morphism of Picard varieties using the Tate module.
2.7. Three-dimensional case. Assume $`n=3`$. Then we can show that Conjecture (0.1) is equivalent to the nondegeneracy of the restrictions of the canonical pairing to
$$H^3(Y_{\overline{k}}^{(1)},_l)^{\text{prim}}\text{Im}\gamma ,H^2(Y_{\overline{k}}^{(2)},_l)^{\text{prim}}\text{Im}\rho .$$
$`(\mathrm{2.7.1})`$
Indeed, using the abelian-positivity in Sect. 3, the arguments in (2.5) show that these conditions in this case are equivalent respectively to the bijectivity of
$$N:H_1^1H_1^1,N:H_1^0H_1^0.$$
Since the first morphism is bijective by using a general hyperplane section, the nondegeneracy of the pairing for the intersection with $`\text{Im}\gamma `$ in (0.2) is always true for $`n=3`$. Note that the bijectivity of the first morphism cannot be proved by using simply the abelian positivity in Sect. 3 unless we know that $`\text{Ker}\gamma H^1(Y_{\overline{k}}^{(2)},_l)`$ corresponds to an abelian subvariety of the Picard variety.
If $`\text{Im}\rho H^2(Y_{\overline{k}}^{(2)},_l)`$ is contained in the subspace generated by algebraic cycle classes, then Conjecture (0.1) can be reduced to $`D(Y^{(1)})`$ or $`A(Y^{(1)},L)`$ in the notation of (4.1) (or to the Tate conjecture for divisors on $`Y^{(1)}`$) because $`D(Y^{(2)})`$, $`D(Y^{(1)})`$ for divisors and $`I(Y^{(2)},L)`$ are known.
If, for every eigenvalue of the Frobenius action on $`\text{Im}\rho H^2(Y_{\overline{k}}^{(2)},_l)^{\text{prim}}`$, its multiplicity as an eigenvalue of the Frobenius action on $`H^2(Y_{\overline{k}}^{(2)},_l)^{\text{prim}}`$ is $`1`$, then Conjecture (0.1) can be proved for certain prime numbers $`l`$, see (5.4).
2.8. Example. Let $`VS:=\text{Spec}R`$ be a smooth projective morphism of relative dimension $`n+1`$. For $`i=0,\mathrm{},r`$, let $`Z_i`$ be a smooth hypersurface of $`V`$ defined by a section $`P_i`$ of a relative ample line bundle $`L_i`$. Assume $`L_0=_{1ir}L_i`$, and $`_{0ir}Z_i`$ is a divisor with normal crossings relative to $`R`$ (i.e. the fiber over $`\text{Spec}k`$ is a divisor with normal crossings). Let $`\pi `$ be a generator of the maximal ideal of $`R`$, and define
$$P=P_1\mathrm{}P_r+P_0\pi .$$
Let $`X`$ be the hypersurface of $`V`$ defined by $`P`$. (This is an analogue of for $`r=2`$.) Then Conjecture (0.1) is true for a semi-stable model of the generic fiber of $`X`$ although $`X`$ is not semistable in general.
Indeed, consider an iteration of blow-ups $`\sigma _j:V^{(j)}V^{(j1)}`$ along the proper transform of $`Z_0Z_j`$ for $`j=1,\mathrm{},r1`$, where $`V^{(0)}=V`$. We have local coordinates $`x_0,\mathrm{},x_n`$ over $`R`$ such that $`P`$ is locally given by
$$ux_1\mathrm{}x_m+x_0\pi ,$$
where $`x_i`$ is the restriction of $`P_i`$ for $`im`$ and $`u`$ is locally invertible. Then the proper transform of $`P`$ by the blow-up along $`Z_0Z_1`$ is locally given by
$$ux_2\mathrm{}x_m+x_0^{}\pi \text{or}ux_1^{\prime \prime }x_2\mathrm{}x_m+\pi ,$$
where $`(x_0^{},x_1,x_2,\mathrm{},x_n)`$ and $`(x_0,x_1^{\prime \prime },x_2,\mathrm{},x_n)`$ are local coordinate systems on open subvarieties $`U^{}`$, $`U^{\prime \prime }`$ of $`V^{(1)}`$ such that $`x_0=x_0^{}x_1`$ on $`U^{}`$ and $`x_1=x_1^{\prime \prime }x_0`$ on $`U^{\prime \prime }`$. Here the pull-back of $`x_i`$ is also denoted by $`x_i`$ to simplify the notation. Since $`U^{\prime \prime }`$ does not intersect the proper transform of $`Z_0Z_j`$ for $`j>1`$, we can proceed by induction, and get a semi-stable model. Its generic fiber is same as that of $`X`$ because the intersection of the center of the blow-up with the generic fiber of (the proper transform of) $`X`$ is a locally principal divisor on the generic fiber. Furthermore the $`Y^{(i)}`$ are lifted to smooth projective schemes over $`R`$ so that the assumption of (0.2) is verified by using Hodge theory.
3. Abelian-Positivity
3.1. Canonical pairing. Let $`A`$ be an abelian variety over a field $`k`$. We denote by $`A^{}`$ its dual variety, and by $`T_lA_{\overline{k}}`$ the Tate module of $`A_{\overline{k}}:=A_k\overline{k}`$ where $`l`$ is a prime number different form $`\text{char}k`$. Using the Kummer sequence, we have a canonical isomorphism
$$H^1(A_{\overline{k}},\mu _n)=A^{}(\overline{k})_n(:=\text{Ker}(n:A^{}(\overline{k})A^{}(\overline{k}))),$$
$`(\mathrm{3.1.1})`$
where $`n`$ is an integer prime to $`\text{char}k`$. Then, passing to the limit, we get
$$H^1(A_{\overline{k}},_l(1))=T_lA_{\overline{k}}^{}.$$
$`(\mathrm{3.1.2})`$
Since the left-hand side of (3.1.1) is identified with
$$\text{Hom}(A(\overline{k})_n,\mu _n),$$
using torsors , we get the canonical pairing of Weil (see also , ):
$$A(\overline{k})_nA^{}(\overline{k})_n\mu _n,T_lA_{\overline{k}}__lT_lA_{\overline{k}}^{}_l(1).$$
$`(\mathrm{3.1.3})`$
To get a pairing of $`T_lA_{\overline{k}}`$, we take a divisor $`D`$ on $`A`$ which induces a morphism
$$\phi _D:AA^{},$$
$`(\mathrm{3.1.4})`$
such that $`\phi _D(a)A^{}(\overline{k})`$ for $`aA(\overline{k})`$ is given by $`T_a^{}D_{\overline{k}}D_{\overline{k}}`$, where $`T_a`$ is the translation by $`a`$ (see loc. cit.) Note that $`\phi _D`$ depends only on the algebraic equivalence class of $`D`$, and $`\phi _D`$ is an isogeny if $`D`$ is ample.
If the pairing is induced by an ample divisor on $`A`$, its restriction to $`T_lB_{\overline{k}}`$ for any abelian subvariety $`B`$ of $`A`$ is nondegenerate, because we have the commutative diagram:
$$\begin{array}{ccc}B& & A\\ \phi _{D|_B}& & \phi _D& & \\ B^{}& & A^{}\end{array}$$
$`(\mathrm{3.1.5})`$
Note that this holds only for subgroups of $`T_lA_{\overline{k}}`$ corresponding to abelian subvarieties.
We say that a pairing of a $`_l`$-module $`V`$ with a continuous action of $`G:=\text{Gal}(\overline{k}/k)`$ is abelian-positive if there exists an abelian variety with an ample divisor $`D`$ such that $`V`$ is isomorphic to $`T_lA_{\overline{k}}__l_l`$ up to a Tate twist as a $`_l[G]`$-module and the pairing corresponds to the one on $`T_lA_{\overline{k}}`$ defined by the canonical pairing and $`\phi _D`$. Note that abelian-positive pairings (having the same weight) are stable by direct sums.
3.2. Compatibility of the cycle classes. Let $`A`$ be an abelian variety over $`k`$, $`X`$ a smooth projective variety over $`k`$, and $`D`$ a divisor on $`A\times _kX`$ such that its restriction to $`\{0\}\times X`$ is rationally equivalent to zero. Let $`P`$ be the Picard variety of $`X`$. Then $`D`$ induces a morphism of abelian varieties
$$\mathrm{\Psi }_D:AP,$$
such $`\mathrm{\Psi }_D(a)P(\overline{k})`$ for $`aA(\overline{k})`$ is defined by the restriction of $`D_{\overline{k}}`$ to $`\{a\}\times X_{\overline{k}}`$, see (and also , ).
Let $`cl(D)^{1,1}H^1(A_{\overline{k}},R^1(pr_1)_{}\mu _n)`$ denote the $`(1,1)`$-component of the cycle class of $`D`$, where $`n`$ is an integer prime to $`\text{char}k`$, and $`pr_1`$ is the first projection. Assume $`\text{NS}(X)`$ is torsion-free. Then $`R^1(pr_1)_{}\mu _n`$ is a constant sheaf on $`A_{\overline{k}}`$ with fiber $`H^1(X_{\overline{k}},\mu _n)=P(\overline{k})_n`$, and we get
$$cl(D)^{1,1}H^1(A_{\overline{k}},R^1(pr_1)_{}\mu _n)=\text{Hom}(A(\overline{k})_n,P(\overline{k})_n).$$
3.3. Theorem (Deligne). The induced morphism $`\mathrm{\Psi }_D:A(\overline{k})_nP(\overline{k})_n`$ coincides with $`cl(D)^{1,1}`$.
(The proof is essentially the same as in .)
3.4. Corollary (Deligne). Let $`C`$ be a smooth projective curve over a field $`k`$, and $`J`$ its Jacobian. Then we have a canonical isomorphism $`H^1(C_{\overline{k}},_l(1))=T_lJ(\overline{k})`$ such that Poincaré duality on $`H^1(C_{\overline{k}},_l(1))`$ is identified with the pairing of $`T_lJ(\overline{k})`$ given by the canonical pairing (3.1.3) together with (3.1.4) for the theta divisor on $`J`$.
Proof. We may assume that $`C`$ has a $`k`$-rational point replacing $`k`$ with a finite extension $`k^{}`$ and $`C`$ with $`C_kk^{}`$ if necessary. Choosing a $`k`$-rational point of $`C`$, we have a morphism $`f:CJ`$. It is well-known that this induces isomorphisms
$`f^{}:H^1(J_{\overline{k}},_l)\stackrel{}{}H^1(C_{\overline{k}},_l),`$
$`f_{}:H^1(C_{\overline{k}},_l)\stackrel{}{}H^{2g1}(J_{\overline{k}},_l(g1)).`$
These are independent of the choice of the $`k`$-rational point of $`C`$, because a translation on $`J`$ acts trivially on the cohomology of $`J`$. Since $`f_{}`$ and $`f^{}`$ are dual of each other, the canonical pairing on $`H^1(C_{\overline{k}},_l)`$ is identified by $`f^{}`$ with the pairing on $`H^1(J_{\overline{k}},_l)`$ given by Poincaré duality and $`f_{}f^{}`$.
Let $`\mathrm{\Gamma }_1=m^{}\mathrm{\Theta }pr_1^{}\mathrm{\Theta }pr_2^{}\mathrm{\Theta }`$, where $`\mathrm{\Theta }`$ is the theta divisor, and $`m:J\times _kJJ`$ is the multiplication. Let $`\mathrm{\Gamma }_2\text{CH}^g(J\times _kJ)`$ be the diagonal of $`f(C)J`$ so that $`f_{}f^{}`$ is identified with $`\mathrm{\Gamma }_2`$. Since the canonical pairing (3.1.3) can be identified with Poincaré duality (see (3.6)), the assertion is reduced by (3.3) (where $`A=X=J`$ and $`D=\mathrm{\Gamma }_1`$) to that the actions of the correspondences
$`(\mathrm{\Gamma }_1)_{}:H^{2g1}(J_{\overline{k}},_l)H^1(J_{\overline{k}},_l(1g)),`$
$`(\mathrm{\Gamma }_2)_{}:H^1(J_{\overline{k}},_l(1g))H^{2g1}(J_{\overline{k}},_l)`$
are inverse of each other up to sign, or equivalently, that the action of the composition of $`\mathrm{\Gamma }_2\mathrm{\Gamma }_1`$ on $`H^{2g1}(J_{\overline{k}},_l)`$ is the multiplication by $`1`$. Here it is enough to show the assertion for the action on the Albanese variety of $`J`$, see (3.7). For $`a,bJ(\overline{k})`$, we see that the image of $`[a][b]`$ by the action of $`\mathrm{\Gamma }_2\mathrm{\Gamma }_1`$ is given by
$$f_{}f^{}(T_a^{}\mathrm{\Theta }T_b^{}\mathrm{\Theta }).$$
$`(\mathrm{3.4.1})`$
Let $`C^{(j)}`$ denote the $`j`$-th symmetric power of $`C`$. Then $`f`$ induces $`f^{(j)}:C^{(j)}J`$, and $`f^{(g)}`$ is birational . So there is a nonempty Zariski-open subset $`U`$ of $`J`$ such that for $`aU(\overline{k})`$, there exists uniquely $`\{c_1,\mathrm{},c_g\}C^{(g)}(\overline{k})`$ satisfying
$$a\underset{ij}{}f(c_i)=f(c_j).$$
$`(\mathrm{3.4.2})`$
Since $`\mathrm{\Theta }=f^{(g1)}(C^{(g1)})`$, it implies that $`T_a^{}\mathrm{\Theta }^{}f(C)=\{f(c_j)\}`$, where $`\mathrm{\Theta }^{}=(1)^{}\mathrm{\Theta }`$. But $`\phi _\mathrm{\Theta }=\phi _\mathrm{\Theta }^{}`$, because the action of $`1:JJ`$ on $`\text{NS}(J)`$ is the identity. So the assertion follows from (3.4.1–2).
3.5. First cohomology case. Let $`Y`$ be a smooth projective variety with an ample divisor class $`L`$. Then the canonical pairing on $`H^1(Y_{\overline{k}},_l)`$, defined by Poincaré duality and $`L`$, is abelian-positive. Indeed, we may assume $`L`$ is very ample, and take $`C`$ a smooth closed subvariety of dimension $`1`$ which is an intersection of general hyperplane sections (replacing $`k`$ with a finite extension if necessary). Then the composition of the restriction and Gysin morphisms
$$H^1(Y_{\overline{k}},_l)H^1(C_{\overline{k}},_l)H^{2n1}(Y_{\overline{k}},_l(n1))$$
coincides with $`L^{n1}`$, where $`n=dimX`$. So the pairing on $`H^1(Y_{\overline{k}},_l)`$ is identified with the restriction of the natural pairing on $`H^1(C_{\overline{k}},_l)`$ to $`H^1(Y_{\overline{k}},_l)`$.
3.6. Complement to the proof of (3.4), I. Let $`X`$ be an irreducible smooth projective variety over a field $`k`$ having a $`k`$-rational point $`0`$. Let $`P_X`$ be the Picard variety of $`X`$. Then we have a canonical morphism $`\text{Alb}:XP_X^{}`$ sending $`0`$ to $`0`$ so that $`P_X^{}`$ is identified with the Albanese variety $`\text{Alb}_X`$ of $`X`$. (Indeed, for an abelian variety $`A`$ and a morphism $`f:(X,0)(A,0)`$, we have $`f^{}:A^{}P_X`$ and $`f^{}:P_X^{}A`$ so that $`f^{}\text{Alb}=f`$, using the theory of divisorial correspondences.)
Let $`n=dimX`$, and $`l`$ a prime number different from the characteristic of $`k`$. Let $`V_lM=T_lM__l_l`$ for an abelian group $`M`$. Then, using the Kummer sequence together with the arguments in (3.1), we have canonical isomorphisms
$`V_lP_X(\overline{k})`$ $`=H^1(X_{\overline{k}},_l(1)),`$
$`V_l\text{Alb}_X(\overline{k})`$ $`=\text{Hom}(V_lP_X(\overline{k}),_l(1))=H^{2n1}(X_{\overline{k}},_l)(n),`$
where the last isomorphism follows from the first together with Poincaré duality, and the last term is identified with $`H_1(X_{\overline{k}},_l):=\text{Hom}(H^1(X_{\overline{k}},_l),_l)`$.
3.7. Complement to the proof of (3.4), II. Let $`X,Y`$ be irreducible smooth projective varieties over a field $`k`$ having a $`k`$-rational point $`0`$. Let $`\mathrm{\Gamma }\text{CH}^m(X\times _kY)`$ with $`m=dimY`$. Then $`\mathrm{\Gamma }`$ induces morphisms
$$(X,0)(\text{Alb}_Y,0),\mathrm{\Gamma }_{}:\text{Alb}_X\text{Alb}_Y,$$
and the induced morphism $`\mathrm{\Gamma }_{}:V_l\text{Alb}_X(\overline{k})V_l\text{Alb}_Y(\overline{k})`$ is identified by (3.6) with
$$\mathrm{\Gamma }_{}:H^{2n1}(X_{\overline{k}},_l)(n)H^{2m1}(Y_{\overline{k}},_l)(m),$$
which is given by $`cl(\mathrm{\Gamma })^{1,2m1}H^1(X_{\overline{k}},_l)H^{2m1}(Y_{\overline{k}},_l)(m)`$, the $`(1,2m1)`$-component of the cycle class of $`\mathrm{\Gamma }_{}`$, see also . Although this is well-known to specialists, its proof does not seem to be completely trivial, and we give here a short sketch for the convenience of the reader.
First we may replace $`Y`$ with $`\text{Alb}_Y`$ (by composing $`\mathrm{\Gamma }`$ with the graph of $`Y\text{Alb}_Y)`$ so that the assertion is reduced to the case $`Y`$ is an abelian variety $`A`$. Consider $`\text{Alb}(\mathrm{\Gamma })\text{CH}^m(X\times _kA)`$ which is a graph of a morphism of $`X`$ to $`A`$, and is defined by using the additive structure of $`A`$ (applied to the restriction of $`\mathrm{\Gamma }`$ to the generic fiber of $`X\times _kAX)`$. Then the assertion is easily verified if $`\mathrm{\Gamma }`$ is replaced by $`\text{Alb}(\mathrm{\Gamma })`$, because $`\text{Alb}(\mathrm{\Gamma })`$ is then extended to an element of $`\text{CH}^m(\text{Alb}_X\times _kA)`$ which is a graph of a morphism of abelian varieties.
So the assertion is reduced to that $`cl(\mathrm{\Gamma })^{1,2m1}=0`$ if $`\text{Alb}(\mathrm{\Gamma })=0`$. Here we may assume $`k`$ is algebraically closed. Replacing $`X`$ with a variety which is étale over $`X`$, we may assume that $`\mathrm{\Gamma }`$ is a linear combination of the graphs of morphisms of $`X`$ to $`A`$ (where $`X`$ is smooth and irreducible, but may be nonproper). By induction on the number of the components of $`\mathrm{\Gamma }`$, we may assume that
$$\mathrm{\Gamma }=\mathrm{\Gamma }_{g_1+g_2}\mathrm{\Gamma }_{g_1}\mathrm{\Gamma }_{g_2}\mathrm{\Gamma }_0$$
for morphisms $`g_i:XA(i=1,2)`$, where $`\mathrm{\Gamma }_{g_i}`$ denotes the graph of $`g_i`$.
Let $`\stackrel{~}{\mathrm{\Gamma }}`$ denote the pullback of the cycle $`\mathrm{\Gamma }_{g_1}\mathrm{\Gamma }_0`$ by the projection $`A\times _kX\times _kAX\times _kA`$ sending $`(a,x,b)`$ to $`(x,a+b)`$. Then $`\mathrm{\Gamma }_{g_1}\mathrm{\Gamma }_0`$ and $`\mathrm{\Gamma }_{g_1+g_2}\mathrm{\Gamma }_{g_2}`$ coincide with the pull-backs of $`\stackrel{~}{\mathrm{\Gamma }}`$ by the inclusions $`X\times _kAA\times _kX\times _kA`$ sending $`(x,a)`$ to $`(0,x,a)`$ and $`(g_2(x),x,a)`$ respectively. Since these inclusions are sections of the projection to the second and third factors, it is enough to show that the Künneth component of the cycle class of $`\stackrel{~}{\mathrm{\Gamma }}`$ in
$$H^1(A\times _kX,_l)H^{2m1}(A,_l)(m)$$
comes from $`H^1(X,_l)H^{2m1}(A,_l)(m)`$ by $`pr_2\times id`$, where $`pr_2:A\times _kXX`$ is the second projection. The last assertion is easily verified by using the Künneth decomposition.
4. Standard Conjectures
4.1. Let $`Y`$ be an equidimensional smooth projective variety over a field $`k`$. We fix an ample divisor class $`L`$ of $`Y`$. Then $`L`$ acts on étale cohomology. Let $`n=dimY`$. By the hard Lefschetz theorem , we have
$$L^j:H^{nj}(Y_{\overline{k}},_l)\stackrel{}{}H^{n+j}(Y_{\overline{k}},_l(j))(j>0),$$
which implies the Lefschetz decomposition
$$H^j(Y_{\overline{k}},_l)=_{i0}L^iH^{j2i}(Y_{\overline{k}},_l(i))^{\text{prim}}.$$
This induces a morphism
$$\mathrm{\Lambda }:H^j(Y_{\overline{k}},_l)H^{j2}(Y_{\overline{k}},_l(1)),$$
such that for $`mH^j(Y_{\overline{k}},_l)^{\text{prim}}`$, we have $`\mathrm{\Lambda }(L^im)=L^{i1}m`$ if $`i>0`$, and $`0`$ otherwise. The standard conjecture $`B(Y)`$ asserts that $`\mathrm{\Lambda }`$ is algebraically defined as an action of a correspondence.
We define
$$A^j(Y)=\text{Im}(cl:\text{CH}^j(Y)_{}H^{2j}(Y_{\overline{k}},_l(j))),$$
$`(\mathrm{4.1.1})`$
so that we have the injective morphism $`A^j(Y)H^{2j}(Y_{\overline{k}},_l(j))`$. Note that $`A^j(Y)`$ may apparently depend on the choice of $`l`$. The standard conjecture $`A(Y,L)`$ asserts the isomorphism
$$L^{n2i}:A^i(Y)\stackrel{}{}A^{ni}(Y)(0i<n/2).$$
This is independent of $`L`$ if the $`A^i(Y)`$ are finite dimensional, because it is equivalent to the equality of the dimensions by the hard Lefschetz theorem for the étale cohomology groups. The conjecture $`A(Y,L)`$ follows from $`B(Y)`$, and implies that the Lefschetz decomposition is compatible with the subspace $`A^j(Y)`$ so that
$$A^j(Y)=_{i0}L^iA^{ji}(Y)^{\text{prim}}.$$
The standard conjecture of Hodge index type $`I(Y,L)`$ asserts that the pairing
$$(1)^jL^{n2j}a,b\text{for}a,bA^j(Y)^{\text{prim}}$$
is positive definite for $`0jn/2`$.
We will denote by $`D(Y)`$ the conjecture which asserts the coincidence of the homological equivalence and the numerical equivalences for the cycles on $`Y`$ (i.e. the canonical pairing between $`A^j(Y)`$ and $`A^{nj}(Y)`$ is nondegenerate for $`A^j(Y)`$). This is equivalent to $`A(Y,L)`$ under the assumption $`I(Y,L)`$, and implies the injectivity of
$$A^j(Y)_{}_lH^{2j}(Y_{\overline{k}},_l(j)),$$
$`(\mathrm{4.1.2})`$
and also the independence of $`A^j(Y)`$ of the choice of $`l`$. It is known that $`D(Y)`$ is true for divisors (by Matsusaka), and $`I(Y,L)`$ is true for surfaces (by Segre ), see , (and also , ).
By the Lefschetz decomposition, we have an isomorphism
$${}_{}{}^{}:H^{n+j}(Y_{\overline{k}},_l)\stackrel{}{}H^{nj}(Y_{\overline{k}},_l(j))$$
such that for $`mH^i(Y_{\overline{k}},_l(a))^{\text{prim}}`$, we have
$${}_{}{}^{}(L^am)=(1)^{i(i+1)/2}L^{nia}m.$$
Combined with Poincaré duality, this induces a pairing on $`H^{}(Y_{\overline{k}},_l)`$ defined by $`m,{}_{}{}^{}n`$ for $`m,nH^j(Y_{\overline{k}},_l)`$.
For a nonzero correspondence $`\mathrm{\Gamma }A^n(Y\times _kY)\text{End}(H^{}(Y_{\overline{k}},_l))`$, we define $`\mathrm{\Gamma }^{}`$ to be the composition of , $`{}_{}{}^{t}\mathrm{\Gamma }`$ and , where $`{}_{}{}^{t}\mathrm{\Gamma }`$ is the transpose of $`\mathrm{\Gamma }`$. Then $`\mathrm{\Gamma }^{}`$ is algebraic and
$$\text{Tr}(\mathrm{\Gamma }^{}\mathrm{\Gamma })>0,$$
$`(\mathrm{4.1.3})`$
if $`B(Y)`$ and $`I(Y\times _kY,L1+1L)`$ are satisfied, see . (This is an analogue of the positivity of the Losati involution.) Note that $`H^{}(Y_{\overline{k}},_l)H^{}(Y_{\overline{k}},_l)`$ is the direct sum of
$$S_{i,j}:=\underset{a=0}{\overset{i}{}}\underset{b=0}{\overset{j}{}}L^aH^{ni}(Y_{\overline{k}},_l)^{\text{prim}}L^bH^{nj}(Y_{\overline{k}},_l)^{\text{prim}}.$$
For $`i=j`$, the primitive part of $`S_{i,i}`$ of degree $`2n`$ for the action of $`L1+1L`$ is isomorphic to $`H^{ni}(Y_{\overline{k}},_l)^{\text{prim}}H^{ni}(Y_{\overline{k}},_l)^{\text{prim}}`$, and this isomorphism is compatible with the canonical pairing up to a nonzero multiple constant.
If $`B(Y)`$ and $`I(Y\times _kY,L1+1L)`$ are true, then $`A^n(Y\times _kY)`$ is a semisimple algebra and there are projectors to primitive parts (and $`D(Y\times _kY)`$ holds), see , , . We have the projector to each cohomology group by the algebraicity of the Künneth components using the Frobenius morphism (because $`k`$ is assumed to be a finite field). By Jannsen it is actually enough to assume $`D(Y\times _kY)`$ for the semisimplicity of $`A^n(Y\times _kY)`$.
We will denote by
$$\iota :A^n(Y\times _kY)\text{End}(H^{}(Y_{\overline{k}},_l))$$
the canonical injection induced by (4.1.1).
4.2. Remarks. (i) Assume $`D(Y\times _kY)`$ holds. Then $`A:=A^n(Y\times _kY)`$ is a direct product of full matrix algebras over skew fields by (together with Wedderburn’s theorem). For an element $`f`$ of $`A`$, there exists an idempotent $`\pi `$ such that $`\text{Im}\iota (\pi )=\text{Im}\iota (f)`$ in $`H^{}(Y_{\overline{k}},_l)`$. Indeed, using projectors (and replacing $`H^{}(Y_{\overline{k}},_l)`$ with the corresponding subspace $`V)`$, the assertion is reduced to the case where $`A`$ is a full matrix algebra over a skew field. Then a matrix can be modified to a simple form by using actions of invertible matrices from both sides as well-known (and a conjugate of an idempotent is an idempotent).
(ii) For a morphism of $`k`$-varieties $`f:XY`$, the idempotent corresponding to the image of $`f^{}`$ is obtained by considering the graph of $`f`$ and applying the above argument to the disjoint union of $`X`$ and $`Y`$ if $`dimX=dimY`$. In general we can replace $`X`$ or $`Y`$ by a product with $`^m`$ (using the Künneth decomposition for $`^m`$), see . Note that $`D(X\times _kY\times _k^m)`$ can be reduced to $`D(X\times _kY)`$.
4.3. Proof of (0.4). By (4.1) and (4.2) there exists an idempotent $`\pi `$ of $`A^n(Y\times _kY)`$ such that $`\text{Im}\iota (\pi )`$ coincides with the intersection of the primitive part with the image of the Cech restriction or co-Cech Gysin morphism.
We show in general that for a projector $`\pi `$ of $`A^n(Y\times _kY)`$, the restriction of the pairing to $`\text{Im}\iota (\pi )`$ is nondegenerate if $`\text{Im}\iota (\pi )`$ is contained in $`H^j(Y_{\overline{k}},_l)^{\text{prim}}`$. For this it is enough to show
$$\text{Im}\iota (\pi )\text{Ker}\iota (\pi ^{})=0,$$
$`(\mathrm{4.3.1})`$
where $`\pi ^{}`$ is the composition of , $`{}_{}{}^{t}\pi `$, as in (4.1.3). We see that (4.3.1) follows from (4.1.3) if $`\pi `$ corresponds to a simple motive, because the intersection is defined as a motive (and the forgetful functor associating the underlying vector space commutes with Im, Ker and the intersection). In general, we consider a simple submotive of $`\text{Im}\pi `$. Let $`\pi _0`$ be the projector defining it. Then (4.3.1) holds for $`\pi _0`$, and we get a decomposition
$$H^j(Y_{\overline{k}},_l)^{\text{prim}}=\text{Im}\iota (\pi _0)(\text{Ker}\iota (\pi _0^{})H^j(Y_{\overline{k}},_l)^{\text{prim}}),$$
which is defined motivically, and is compatible with $`\text{Im}\iota (\pi )`$, i.e.
$$\text{Im}\iota (\pi )=\text{Im}\iota (\pi _0)(\text{Ker}\iota (\pi _0^{})\text{Im}\iota (\pi )).$$
Therefore, replacing $`H^j(Y_{\overline{k}},_l)^{\text{prim}}`$ with $`\text{Ker}\iota (\pi _0^{})H^j(Y_{\overline{k}},_l)^{\text{prim}}`$, we can proceed by induction. This completes the proof of (0.4).
5. Frobenius Action
5.1. Weil Conjecture. Let $`Y`$ be an equidimensional smooth projective $`k`$-variety, and put $`A=A^n(Y\times _kY)`$ where $`n=dimY`$. Let $`q=|k|`$, $`p=\text{char}k`$. We denote by $`g`$ the graph of the $`q`$-th power Frobenius. Then $`g`$ belongs to the center of $`A`$ (using the compatibility with the Galois action). To simplify the relation with the eigenvalues, let $`P(T)`$ denote the characteristic polynomial of the action of the Frobenius $`g`$ on $`H^{}(Y_{\overline{k}},_l)(:=_iH^i(Y_{\overline{k}},_l))`$ such that the eigenvalues of the Frobenius action are the roots of $`P(T)`$ (this normalization is different from the one used in , , etc.) Then $`P(T)`$ is a monic polynomial with integral coefficients and is independent of $`l(p)`$. The eigenvalues of the action of $`g`$ on $`H^i(Y_{\overline{k}},_l)`$ are algebraic integers whose image by any embedding of $`\overline{}`$ into $``$ has absolute value $`q^{i/2}`$, see .
Let $`E`$ be an algebraic number field. Consider the decomposition $`P(T)=_jP_j(T)^{m_j}`$ in $`E[T]`$ where the $`P_j(T)`$ are monic irreducible polynomials whose coefficients are algebraic integers in $`E`$. Let $`v`$ be a prime (i.e. a finite place) of $`E`$ over a given prime number $`l(p)`$. Let $`E_v`$ denote the completion of $`E`$ at $`v`$, and $`\overline{E}_v`$ denote an algebraic closure of $`E_v`$ (which is isomorphic to $`\overline{}_l)`$. Consider the ring of correspondences $`A_E`$ with $`E`$-coefficients. We have a natural inclusion
$$\iota :A_E\text{End}(H^{}(Y_{\overline{k}},E_v)),$$
and a natural surjection $`A_{}EA_E`$, but the injectivity of the last morphism is not clear unless the conjecture $`D(Y\times _kY)`$ holds.
As well-known, there exist $`R_j(T)E[T]`$ such that $`\pi _j:=R_j(g)A_E`$ is an idempotent, $`_j\pi _j=1`$ and the characteristic polynomial of the action of $`g`$ on $`\text{Im}\iota (\pi _j)`$ is $`P_j(T)^{m_j}`$. Note that $`\text{Im}\iota (\pi _j)`$ is contained in some cohomology group $`H^i(Y_{\overline{k}},E_v)`$ by , see .
5.2. Proposition. Let $`M_j`$ be the motive defined by $`\pi _j`$. Assume $`m_j=1`$. Then the endomorphism ring $`\text{End}(M_j)`$ is generated by $`g`$ over $`E`$, and is isomorphic to $`E[T]/(P_j(T))`$.
Proof. Let $`F`$ be a minimal Galois extension of $`E`$ containing all the roots of $`P_j[T]`$. Let $`d_j=\mathrm{deg}P_j`$. For a root $`\alpha `$ of $`P_j[T]`$ in $`F`$, let $`E[\alpha ]`$ denote the subfield of $`F`$ generated by $`\alpha `$, which is isomorphic to $`E[T]/(P_j(T))`$. By the same argument as above, we have a projector $`\pi _\alpha `$ in $`A_{E[\alpha ]}`$ such that
$$g\pi _\alpha =\alpha \pi _\alpha .$$
$`(\mathrm{5.2.1})`$
Take any $`\mathrm{\Gamma }A_E`$. There exists a polynomial $`h(T)E[T]`$ such that
$$\mathrm{\Gamma }\pi _\alpha \mathrm{\Delta }=h(\alpha )E[\alpha ],$$
$`(\mathrm{5.2.2})`$
where $`\mathrm{\Delta }`$ is the diagonal cycle of $`Y`$, and the left-hand side denotes the intersection number. Note that $`h`$ is independent of $`\alpha `$ using the action of $`\text{Gal}(F/E)`$ because, fixing one root $`\beta `$, we have for each root $`\alpha `$ an automorphism $`\sigma `$ of $`F/E`$ such that $`\alpha =\beta ^\sigma `$. By the Lefschetz trace formula, $`h(\alpha )`$ coincides with the eigenvalue of the action of $`\mathrm{\Gamma }`$ on the $`\alpha `$-eigenspace of $`g`$ up to a sign independent of $`\alpha `$. Here $`\alpha `$ denotes also the image of $`\alpha `$ by the embedding $`FF_v`$ where $`v`$ is a prime of $`F`$ over $`l`$. So the assertion follows.
5.3. Image of correspondences. Let $`X`$ be an equidimensional smooth projective $`k`$-variety, and $`\mathrm{\Gamma }\text{CH}^r(X\times _kY)_E`$. Assume that $`E`$ is a subfield of $``$, $`\mathrm{deg}P_j2`$ for any $`j`$, and the roots of $`P_j`$ are not real if $`\mathrm{deg}P_j=2`$. Let $`\beta `$ be an eigenvalue in $`\overline{E}_v`$ of the Frobenius action on $`\text{Im}\mathrm{\Gamma }_{}H^i(Y_{\overline{k}},E_v)^{\text{prim}}`$, and assume that its multiplicity as an eigenvalue of the Frobenius action on $`H^i(Y_{\overline{k}},E_v)^{\text{prim}}`$ is $`1`$ (i.e. the dimension of the generalized eigenspace in $`H^i(Y_{\overline{k}},E_v)^{\text{prim}}`$ is $`1`$). By definition $`\beta `$ is a root of some $`P_j(T)`$. Assume further that $`P_j(T)`$ is irreducible over $`E_v`$, and $`\mathrm{deg}P_j=2`$ (because the case $`\mathrm{deg}P_j=1`$ is trivial). Let $`\beta ^{}`$ be the other root of $`P_j(T)`$. Then $`\beta ^{}`$ is an eigenvalue of the Frobenius action on $`\text{Im}\mathrm{\Gamma }_{}H^i(Y_{\overline{k}},E_v)^{\text{prim}}`$, and hence $`\text{Im}\mathrm{\Gamma }_{}H^i(Y_{\overline{k}},E_v)^{\text{prim}}`$ contains $`\text{Im}\iota (\pi _j)`$ on which the canonical pairing is nondegenerate. Indeed, the pairing is compatible with the Frobenius action and has value in $`E_v(i)`$ on which the action of $`g`$ is a multiplication by $`q^i`$. So we are interested in the problem: When is $`P_j`$ irreducible over $`E_v`$?
Let $`F`$ be the field generated by all the roots $`\alpha `$ of $`P`$ in $``$. Since $`\alpha \overline{\alpha }=q^r`$ with $`r`$, the complex conjugation on $`F`$ (induced by that on $`/)`$ is in the center of $`\text{Gal}(F/)`$, because $`\alpha ^\sigma \overline{\alpha }^\sigma =q^r`$ and hence $`\overline{\alpha }^\sigma =\overline{\alpha ^\sigma }`$ for $`\sigma \text{Gal}(F/)`$. Let $`E`$ be the subfield of $`F`$ fixed by the complex conjugation. By the Tchebotarev density theorem, there are infinitely many primes $`v`$ whose density is $`1/2`$ and such that $`P_j`$ modulo $`v`$ is irreducible over the residue field (hence $`P_j`$ is irreducible over $`E_v`$) considering the conjugacy class of the above complex conjugation in $`\text{Gal}(F/E)`$. So we get the following.
5.4. Proposition. For prime numbers $`l`$ such that some $`v`$ as above is over $`l`$, Conjecture (0.1) is true if each eigenvalue of the Frobenius action on the intersections of the primitive part with $`\text{Im}\rho `$ and with $`\text{Im}\gamma `$ has multiplicity $`1`$ as an eigenvalue of the Frobenius action on the primitive cohomology $`H^j(Y^{(i+1)},_l)^{\text{prim}}`$ for $`j2`$ and $`i0`$.
5.5. Remark. If Conjecture (0.1) holds for a general hyperplane section of the generic fiber, then it is enough to consider the intersection with $`\text{Im}\rho `$ by Theorem (0.3). For a general $`l`$, Conjecture (0.1) in the case of (5.4) can be reduced to the conjecture $`D`$ (or $`B`$). Thus we can avoid the positivity in this simple case. However, it is not easy to avoid it even in the case where the multiplicity $`1`$ holds for each irreducible component of $`Y^{(i+1)}`$. Indeed, the ambiguity of the pairing on the motive $`M_j`$ in (5.2) is given by an automorphism $`h[g](=[T]/(P_j(T)))`$ where $`E=`$. If there is a morphism of $`M_j`$ to a direct sum of simple motives which are isomorphic to $`M_j`$ and indexed by $`r`$, and if the morphism to the $`r`$-th factor is given by a correspondence $`\mathrm{\Gamma }_r`$, then the pull-back of the pairing corresponds to the sum of $`h_r:=\mathrm{\Gamma }_r^{}\mathrm{\Gamma }_r[g]`$ up to a sign, and the problem is closely related to the standard conjecture of Hodge index type.
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# Spectroscopy of Blue Stragglers and Turnoff Stars in M67 (NGC 2682)1footnote 11footnote 1 Based in-part on observations obtained with the Hobby-Eberly Telescope, which is a joint project of the University of Texas at Austin, the Pennsylvania State University, Stanford University, Ludwig-Maximillians-Universität München, and Georg-August-Universität Göttingen.
## 1 Introduction
Blue straggler (BS) stars were first identified as an unusual subclass of stars in the cluster M3 (Sandage, 1953). Since that time, stragglers have been identified in most, if not all, globular clusters that have accurate broad-band photometry. Blue straggler candidates have also been identified in a large number of open clusters, e.g. Ahumada & Lapasset (1995).
Blue stragglers are commonly defined as those stars that are brighter and bluer (hotter) than the main sequence (MS) turnoff of the majority of the cluster stars, but along an apparent extension of the main sequence. At present, the leading explanation involves mass transfer in and/or merger of a binary star system, or the collision of stars (whether or not a binary is involved). These mechanisms can create a main sequence star with a mass greater than would be expected given the age of the cluster.
The above definition of blue straggler is overly restrictive if this is the correct explanation. Undoubtedly the same processes act on lower mass stars, although the remnant stars would not be easily identifiable if they were not brighter than the turnoff luminosity. The color restriction is also unrealistic since real scatter in the colors of blue stragglers in globular clusters has been observed, e.g. Fusi Pecci et al. (1992). Sills & Bailyn (1999) find that this color scatter results from a lack of large-scale mixing during the creation process, which means that the amount of helium in the core can vary from straggler to straggler.
Hydrodynamical studies (Lombardi et al., 1995; Sandquist et al., 1997) earlier predicted that large-scale mixing does not occur during stellar collisions. A collisionally-produced straggler is not expected to burn significant lithium initially even if high enough temperatures are reached because the timescale for the straggler to thermally adjust to a new equilibrium configuration is much shorter than the lithium destruction timescale Sills et al. (1997). Binary mass transfer is likely to lead to a more easily observable abundance pattern since the last of the transferred gas should have come from deep within the parent star (e.g, Sarna & de Greve, 1996).
In this paper we have attempted to examine this question by taking spectra of blue straggler stars in the open cluster M67 (NGC 2682). M67 makes an ideal target because it is relatively close ($`VM_V9.7`$; Dinescu et al., 1995), and also fairly old ($`4.0`$ Gyr; Dinescu et al., 1995). In addition, the density of stars in open clusters is lower than in nearly all globular clusters, making binary evolution and collisions involving binaries the most likely formation mechanisms (Leonard & Linnell, 1992). Globular cluster blue stragglers are not yet practically observable because of their distances. Although much attention has been spent on M67 blue stragglers, few observations are available for the coolest ones.
Few blue stragglers have abundance analyses. Mathys (1991) observed a sample of 11 blue stragglers in M67, although only 2 were examined in detail. Andrievsky (1998) observed 4 blue stragglers in Praesepe (NGC 2632) from the catalog of Ahumada & Lapasset (1995), although only one (HD 73666) is undeniably a straggler according to its photometry. In observing the blue stragglers of M67 we are most interested in species such as Li, C, N, and O that are processed by nuclear reactions at different temperatures (and therefore, depths) in a star. The reddest blue stragglers in the cluster have effective temperatures that can allow us to get abundances for all of these elements, and are close enough to the temperatures of turnoff stars that we can straightforwardly make comparisons with stars having (presumably) unmixed envelopes.
In §2 we describe the spectroscopic observations, and in §3 we present the abundance analysis for our sample. In §4 we discuss the constraints we can place on blue straggler formation mechanisms in this cluster.
## 2 Observations
Table 1 displays the program stars, basic photometry, and the signal-to-noise obtained. The identification numbers we use are from the proper motion study of Sanders (1977). We chose our blue straggler candidates from Ahumada & Lapasset (1995) and our MS members from the proper motion study of Zhao et al. (1993). All of the chosen M 67 stars are 99% probability proper motion members. The photometry for our program stars was taken from Fan et al. (1996), Montgomery et al. (1993), Chevalier & Ilovaisky (1991), Gilliland et al. (1991), and Sanders (1989). The $`BV`$ colors of Sanders and Gilliland et al. were adjusted by 0.006 and $`0.002`$, respectively. The $`V`$ magnitudes of Sanders, Gilliland et al. and Chevalier & Ilovaisky were adjusted by 0.017, $`0.009`$, 0.012, respectively. Several measurements made by Sanders were found to be in disagreement with the other sources, including ancillary measurements by Anupama et al. (1994), Murray et al. (1965), and Murray & Clements (1968). For this reason, the Sanders measurements for S 1183, S 821, S 975, and S 984 were not used to derive the average photometry. The final adopted values are given in Table 1, and the CMD of the cluster and the observed blue stragglers are shown in Figure 1. We observed one blue straggler candidate (Ahumada & Lapasset, 1995) in NGC 7789. Due to poor weather additional NGC 7789 candidates were not observed.
The majority of the observations were taken with the McDonald Observatory 2.7 m ”2d-coude” (Tull et al., 1995). The resultant spectra have $`R=30,000`$ and cover from 3800 Å to 10100 Å. Each star was observed on several different nights in order to look for radial velocity variations. Radial velocities were obtained from individual spectra, but the S/N quoted in Table 1 are those of the combined spectra. The combined spectra were also used for the abundance analysis. Additional observations were obtained with the Hobby-Eberly Telescope (Ramsey et al., 1994, HET) during commissioning of the UFOE (Harlow et al., 1996) spectrograph in early 1999 and in the first science period of 2000 with the UFOE on the HET. The HET spectra were used to obtain radial velocities and have lower resolution ($`R=11,000`$) than the 2.7 m spectra, and so are not included in the combined spectra. The HET spectra cover from 4600 Å to 9100 Å. In March 2000, the McDonald Observatory 2.1 m Cassegrain echelle was used to obtain a few additional observations of two of the program stars. Despite having resolution higher than the 2.7 m ”2d-coude” spectra, these were not included in the abundance analysis because the slit length did not allow proper sky subtraction, which would not interfere with the radial velocity study since M67 is well separated from the solar rest velocity. The 2.1 m spectra cover from 5400 Å to 6600 Å. Table 2 lists the heliocentric Julian dates of our observations. In addition to the program stars one additional M67 star (MMJ5554) was observed. This star is only $`3^{\prime \prime }`$ away from S 975 and its color makes it either a turnoff star or a possible blue straggler candidate. A set of bright, rapidly-rotating hot stars were observed at various air masses with both telescopes for sky line identification. Figure 2 exhibits a sample of the spectra obtained on the McDonald 2.1m.
## 3 Analysis
### 3.1 Radial and Rotational Velocities
The spectra were analyzed using the IRAF echelle package. Initial data reduction techniques included bias removal, flat field correction, scattered light removal, sky subtraction and wavelength calibration. To remove instrumental wavelength zero-point errors the spectra were shifted in velocity so that the telluric features fell at their rest velocities. The largest zero-point correction was 1.1 km s<sup>-1</sup>. Each program star spectrum was divided by the hot star spectrum with the most similar telluric line depths. The spectra were then cross-correlated against a synthetic spectrum to obtain the radial velocities, which were then corrected to the heliocentric rest frame. A single synthetic spectrum was used for all the stars ($`T_{eff}=`$ 6400 K, log g = 4.3, $`v_t`$ = 2.0, \[Fe/H\] = –0.05, and no instrumental or rotational broadening). The code used to create the synthetic spectrum is the latest version of MOOG (Sneden, 1973). To determine the radial velocities of the broad component of S 1082 we took an additional step to remove the narrow component. We fit the spectrum with a very high order spline (high enough to remove any broad component) and then divided this normalized spectrum into the original spectrum leaving only the blaze function and any broad line components. These cleaned broad-line spectra were then subjected to the analysis described above. Table 2 lists the heliocentric radial velocities and errors for each spectrum. The mean radial velocity of M67 is 33.6 $`\pm `$ 0.72 km s<sup>-1</sup> (Girard et al., 1989). We detect 3 binary periods not previously known for stars in our sample. S 821 and S 984 are found to have radial velocity variations with periods of 26.259 $`\pm `$ 0.002d and 1.465 $`\pm `$ 0.001, respectively. Figure 3 shows the phased radial velocity data from Mathieu et al. (1986) and this work for these periods. S 821 probably has an eccentricity of $`0.4\pm 0.1`$, while the data from S 984 is consistent with a circular orbit. S 1082 is found to be a double lined spectroscopic binary and will be discussed in §3.
The 2.7 m ”2d-coude” spectra were combined after shifting to the heliocentric rest velocity. New synthetic spectra were created with $`T_{eff}=`$ 6400 K, log g = 4.3, $`v_t`$ = 2.0, \[Fe/H\] = –0.05, a Gaussian $`R=30,000`$ profile, and various rotational velocities. The unbroadened spectrum was cross-correlated against the broadened spectra and the full width half maxima (FWHM) of the correlation peaks were measured. These FWHM were then compared to the FWHM of the correlation peaks created by the cross-correlation of the combined program star spectra and the unbroadened synthetic spectrum. From this comparison we determine the rotational velocity of the combined spectrum. We adopt a conservative lower limit (9 km s<sup>-1</sup>) to our ability to measure rotational velocities. Table 1 lists the rotational velocities found with this method.
We detect rotation in our spectra for only 3 of the blue straggler stars in our sample: S 975, S 997, and S 1082. All have relatively low rotation speeds. In addition, we marginally detect rotation in one of the main sequence stars (S 1271). The signal-to-noise ratios for the main sequence stars was comparable to those for the blue stragglers.
We observed one blue straggler candidate in NGC 7789, M 1251, but found it to be a cluster non-member based on its radial velocity.
### 3.2 Abundance Analysis
The spectral lines used in our analysis were taken from Shetrone (1996), Mathys (1991), and Edvardsson et al. (1993). The oscillator strengths were taken from these sources as well as the National Institute of Standards and Technology Atomic Spectra Database. Each line’s equivalent width was measured with the IRAF task splot. These lines, oscillator strengths, and equivalent widths are listed in Table 3. We employed the latest version of the LTE spectral analysis code MOOG (Sneden, 1973) and the Kurucz (1993) grid of ATLAS models in an iterative abundance analysis.
The initial estimate of the effective temperature for each star was determined from the dereddened $`BV`$ color (we assumed E$`(BV)=0.05`$). Effective temperatures were then determined from the relationship $`T_{eff}=1808(BV)^26103(BV)+8899`$ (Soderblom et al., 1993). Each effective temperature was then fine tuned by forcing the slope of abundances from Fe I lines versus excitation potential to be zero. This tuning of the temperature was restricted to the range of colors consistent with the errors in the photometry. The two largest corrections were found for S 1082 (152 K; likely binary blending of colors, see §3) and S 815 (42 K). The surface gravity was set by enforcing abundance equilibrium between the Fe I and Fe II lines. Based upon the errors in the Fe I and Fe II abundances we estimate the average uncertainty in the surface gravity to be 0.15 dex (1 sigma). The microturbulent velocity was determined by forcing the slope of Fe I line abundances versus equivalent width to be zero. These adjustments to the spectral model for each star were done iteratively until all the parameters ($`T_{eff}`$, log g, $`v_t`$) fulfilled the above requirements within the abundance uncertainty of each spectrum. The derived atmospheric parameters and abundances for each star are listed in Table 4.
Synthetic spectra were also employed to verify the abundances of the C I lines. Using the derived atmospheric parameters for each star a synthetic spectrum was generated and plotted over the observed spectrum and visually inspected. This method did not yield any significant abundance differences.
The three O I lines found at 7773Å are known to need NLTE corrections. These corrections can become very significant at high temperatures (T $`>`$ 6800 K). Over the narrow range of temperatures we considered the NLTE corrections should be relatively small. However, we have made some attempt to correct for NLTE effects using the lookup table employed by Gratton et al. (1999). Even if the magnitude of our corrections are incorrect, the relative oxygen abundances of the stars can be viewed with some confidence. In Figure 4 we display the combined equivalent widths of the O I lines plotted against their colors. The stars from Mathys (1991) are shown as squares. The large crosses are taken from the Varenne & Monier (1999) study of oxygen in the Hyades. Varenne & Monier (1999) found a constant abundance along the main sequence. Note the rapid rise in O I EW near $`BV=0.43`$ or $`T_{eff}=`$ 6840 K, caused by the onset of NLTE effects. The open and filled circles represent abundance analyses with constant oxygen abundance (\[O/Fe\] = 0.0) for several temperatures with the NLTE corrections of Gratton et al. (1999) and Faraggiana et al. (1988), respectively. A summary of our oxygen abundances is given in Table 4.
## 4 Discussion
Because the stragglers in M67 have been observed extensively, we have a better idea of their present status than the stragglers in any other cluster. It is therefore worth briefly summarizing what is known about the stars in our sample so that we can test whether the compositions reflect differences in straggler formation mechanisms.
### 4.1 Discussion of Individual Stars
#### 4.1.1 Blue Stragglers
S 975 (F 90): This star is known to be a binary ($`P=1221`$ d) with small eccentricity ($`e=0.088\pm 0.060`$) from the work of Latham & Milone (1996). (Leonard, 1996) indicates that the system is likely to have resulted from case C (asymptotic giant branch) mass transfer. Landsman et al. (1998) assert that there is strong evidence of a hot companion in their ultraviolet photometry, and conclude that the system is likely to be the result of Algol-type binary mass transfer, although they adopt a photometric temperature (Mathys, 1991) that is more than 300 K cooler than our derived value. Mathys states that the photometry could be contaminated by a nearby star, so that his temperature estimate is questionable. The contaminating star must be cooler than S 975 in order to produce the photometric temperature. MMJ 5554 ($`V=12.839,BV=0.532`$) is probably the cause since it is only $`3^{\prime \prime }`$ away from S 975, and its colors place it very near the bluest extension of the turnoff. Because our spectral temperature did not need to be significantly corrected from our $`BV`$ photometric temperature, it is unlikely that this star contaminated our spectrum. Judging from Figure 2 of Landsman et al. (1998), it probably did not contribute significantly to the UIT UV flux.
According to Landsman et al., “uncertainty of 200 K in the value of $`T_{eff}`$ corresponds to about a 0.5 mag uncertainty in the predicted UIT magnitude”, and “the UIT magnitude of F90 is probably uncertain by close to a factor of 2” because it was contaminated by the light of the bright blue straggler S 977. Based on this and our higher $`T_{eff}`$, the evidence for a hot companion to S 975 is considerably weaker, although there is still an unaccounted-for UV excess. The agreement between our photometric and spectroscopic temperatures indicates that the companion does not contribute significantly to the visual flux (consistent with being a white dwarf). However, we have been unable to detect the companion’s spectral signature, so we are unable to judge whether the spectrum is contaminated. The star does have an unusually large O I EW though (see §4.2.2), which may betray the presence of circumstellar material.
S 984 (F 134): Because S 984 lies within 0.75 mag of the fiducial line of the cluster in the color-magnitude diagram, it is possible that it is a non-interacting binary composed of two main sequence stars. The velocity dispersion in the radial velocity data of Mathieu et al. (1986) was marginally higher than average. After combining the velocity data from Mathieu et al. (1986) with our derived velocities, we conducted a periodicity analysis and found a likely period of 1.465 $`\pm `$ 0.001 days. The peak-to-peak variation is between about 3 and 5 km s<sup>-1</sup>, depending on whether we accept one outlying observation from Mathieu et al.
The period, radial velocity variation, and photometry do constrain the identity of S 984 though. If we hypothesize that S 984 is a binary composed of stars that have not interacted, it is possible to match the photometry with two main sequence stars having $`V12.7`$ and 13.4. This is the maximum brightness contrast that is possible if the primary is to be on the fiducial line of the cluster. However, this constrains the mass ratio to be $`0.9<q1.0`$. The orbital inclination is then constrained to be within 1.5 degrees of face-on. The lithium abundance, which is consistent with that of a single turnoff star, is low for a tidally-locked binary (as is implied if the radial velocity variations reflect orbital motions). So, there is a fairly strong (although circumstantial) case that S 984 is a true blue straggler, and not a photometric blend of two stars. There may still be a companion, but it is unlikely to contribute to the flux. If this is the case, S 984 is probably the result of a collisional merger, since binary mass transfer should result in a lack of surface lithium. For these reasons, we regard it as likely that the spectrum of S 984 is at most minimally contaminated by another star.
S 997 (F 124): This star is known to be in an eccentric ($`e=0.342\pm 0.082`$) binary of period 4913 d (Latham & Milone, 1996). The eccentricity of the system probably rules out a history as a mass transfer binary. The fact that the primary is a blue straggler casts doubt on a binary merger since it would either have required the system to have been in a rare triple system, or to have tidally captured a star after merger (Leonard, 1996). Either way, we do not have evidence to judge whether its spectrum is contaminated.
S 1082 (F 131; ES Cancrii): Pritchet & Glaspey (1991) and Mathys (1991) reported that the spectrum of S 1082 has a possible composite nature. X-ray emission has been detected from this star by Belloni et al. (1993). Landsman et al. (1998) find evidence of a hot subluminous companion in their ultraviolet photometry. However, Landsman et al. (1998) base their predicted UIT magnitudes on a temperature 120 K cooler than found in our analysis. Since a 200 K change in temperature corresponds to 0.5 m<sub>152</sub> magnitudes, their evidence for a hot subluminous companion is weakened, although not eliminated. The discovery of a variable secondary H$`\alpha `$ spectral feature by Mathys (1991) and van den Berg et al. (1999) suggests the presence of a hot rapidly rotating secondary or an outflowing wind. Further, the spectroscopy of van den Berg et al. (1999) appears to rule out magnetic activity as the source of the X-ray emission. This evidence supports the hypothesis that this is an Algol-type binary currently undergoing mass transfer.
In addition to being able to detect the second spectral component in H$`\alpha `$, the Na D lines, and the O I triplet as Mathys (1991) did, we also detected the secondary in H$`\beta `$, H$`\gamma `$, Paschen(8862), the Ca II IR triplet, lines of moderately strong neutral element species (e.g. Fe I, Ca I) and a few strong Fe II lines. The abundances derived for 1082 are 0.2 dex lower than the cluster mean. This can be explained if the second component produces 40–50$`\%`$ of the continuum flux but only a tiny ($``$10$`\%`$) contribution to the line flux. The rotational velocity implied for the second spectral feature is 90$`\pm `$20 km s<sup>-1</sup>, which explains the small contribution to the line flux.
For the first time, we are able to see radial velocity variations in the secondary (broad-lined) component amounting to at least 25 km s<sup>-1</sup>. Our phase coverage is not good enough to state more than a lower limit. Figure 2 exhibits the series of spectra we obtained of S 1082 on the McDonald 2.1m. The broad component can clearly be seen when compared with the spectra of S 984. In addition, the broad component of S 1082 moves redward over the 5 hours the star was observed. For the primary, we find a peak-to-peak amplitude of about 5.5 km s<sup>-1</sup>. From the data of (Mathieu et al., 1986), there is a peak-to-peak amplitude of about 7 km s<sup>-1</sup>.
Goranskij et al. (1992) found from photometry that the system has two partial eclipses per orbit of different depth, with a period of 1.0677978 d. Using this period (or half of it), we see no sign of a repeating pattern in the radial velocity data. Allowing the period to vary, we find a highest probability period of 1.87 $`\pm `$ 0.05 d, with less likely values of 0.7 d and 2.3 d (see Figure 3). However, none of these periods unveils a satisfactory pattern in the radial velocity data. We have two sequences of spectra taken on succeeding nights (centered around phase 0.5 of the photometric period) that indicate sudden ($`2`$ h) jumps in radial velocity for both components. Thus, the radial velocity data imply that we are seeing the signature of material outside of the photospheres of the two stars.
The spectral temperature derived for the narrow component is more than 100 K hotter than the photometry predicts. Using the color as the composite temperature for the system, our spectroscopic temperature suggests that the narrow line component may be 200 K hotter than the broad lined component. Further evidence for a temperature difference between the two components comes from the spectrum: the broad component’s high excitation lines of O I have smaller equivalent widths than those of the narrow component, while the low excitation Na D lines have higher equivalent widths. This implies a cooler temperature, although strong Fe II lines may indicate that a lower surface gravity also comes into play. The relatively large brightness (M$`{}_{V}{}^{}2.5`$) and hot temperature ($`T6800`$ K) of the secondary indicates that it has probably been severely disturbed by its interaction with the primary.
S 2204 (F 130): We observed this star initially believing it to be a main sequence star. Evidence indicates that it is very probably a faint blue straggler, and it is identified as such in the catalog of Ahumada & Lapasset (1995). The proper motion studies of Girard et al. (1989) and Sanders (1977) cite membership probabilities of 98% and 83% respectively, and radial velocity measurements by Mathieu et al. (1986) and us put it right on the cluster average. The photometry of Montgomery et al. (1993) puts it ($`BV=0.448`$) over 0.07 mag to the blue of the bluest stars at the turnoff, and over 0.1 mag from the fiducial line of the cluster. The quoted mean error in the photometry was 0.016. Significant contamination of the spectrum by a main sequence companion is not likely since this would either make the star brighter than the turnoff, or closer in color.
A hot, faint companion could also potentially create the observed colors. Although S 2204 was in the field of the ultraviolet observations of Landsman et al. (1998), it was not detected. Their detection of several white dwarf candidates thereby places an upper limit on the temperature of a possible white dwarf companion. Landsman et al. (1998) estimate their detection threshold for a white dwarf with a pure hydrogen atmosphere would be about 21,300 K. A white dwarf below the detection threshold would change the broadband color by a negligible amount.
For stars close to the turnoff, the possibility of delayed star formation (a very unlikely possibility for brighter blue stragglers) should be considered. For S 2204, a delay of more than 1 Gyr (about 25% of the cluster’s age) would be necessary according to the isochrone fits of Dinescu et al. (1995) if S 2204 is an undisturbed main sequence star. On this basis we discard the delayed star formation picture for S 2204.
With the elimination of these hypotheses, the most likely explanation of S 2204 is as a blue straggler of the merged-star or mass-transfer variety. Although there is currently not enough radial velocity information to more strongly test for the presence of a companion to S 2204, we argue that the primary would be identified as a blue straggler even without one, and that such a companion would not significantly modify the spectrum of the straggler.
#### 4.1.2 Turnoff Stars
S 821: We have detected binary motions for this object with a period $`P=26.259\pm 0.002`$ d and eccentricity $`e0.4\pm 0.1`$ (Figure 3). Broad-band photometry places it very close to the fiducial line of the cluster. This probably indicates that the spectrum is uncontaminated by the companion, although there is a chance that it is composed of two fainter, nearly equal mass stars.
### 4.2 Composition Analysis
The blue stragglers we have observed can be divided into two probable categories: mass transfer binaries (S 975, S 1082), and collision products (S 984, S 997). At present we do not have enough information to classify S 2204. We also have spectra for 5 stars near the main sequence turnoff (S 815, S 821, S 1183, and S 1271). Our goal is to compare the abundances to see if we can find chemical signatures of the blue straggler formation process.
#### 4.2.1 Lithium
If significant mixing of a star’s envelope has occurred, then the lithium abundance may show the results first since it is consumed at relatively low temperatures ($`T>\mathrm{\hspace{0.33em}2.5}\times 10^6K`$). Jones et al. (1999) present the most recent compilation of lithium abundances for M67 stars.
Our results are given in Table 4. Two of our stars (S 997 and S 2204) were previously observed by Garcia Lopez et al. (1988), who found upper limits. We derive upper limits that are 0.5 and 0.8 dex lower, respectively. With the exception of S 984, we are only able to derive upper limits for the stars in our sample. The values for the main sequence stars are consistent with the lower envelope of lithium abundances (see Jones et al., 1999, for references to other papers on M67 lithium abundances).
Hobbs & Mathieu (1991) previously observed the relatively cool blue stragglers S 1072 and S 1082, while Pritchet & Glaspey (1991) observed stragglers S 752, S 997, S 1082, S 1263, S 1267, S 1280, and S 1284. Both sets of authors found only upper limits for their samples. So, the lithium abundances do not as yet provide a dependable way of distinguishing among different formation mechanisms.
#### 4.2.2 CNO elements
The CNO elements can potentially constrain the amount of mixing that occurs in a star, since the CNO cycle starts to shuffle abundances at around $`10^7`$ K. Surface material is unlikely to reach this temperature in main sequence stars in the mass range present in M67, even if violent interactions occurred. However, if mass transfer were to occur on the giant branch, processed material might be dropped onto the surface of the straggler.
Our results for carbon and nitrogen are given in Table 4, while our oxygen results are found in Table 5. Our nitrogen abundances are very uncertain because of the very small equivalent widths found at this temperature, so we will not discuss them further except to note that the nitrogen abundances for the blue stragglers appear to be lower (0.2 – 0.3 dex) than the abundances found in the turnoff stars. The carbon abundances for all of the stars in our sample fall very close to the solar ratio. Comparing the oxygen abundances of the blue stragglers as a group to those of the main sequence stars, the blue stragglers (dropping S 975) also have an average abundance near that of the turnoff stars. However, several stragglers have anomalously high or low values.
The straggler S 975 is problematic to analyze. On one hand, the evidence points to mass transfer of the type that could produce an observable chemical signature. On the other hand, the composite nature of the system appears to be affecting our equivalent width measurements. The measured O abundance is rather high compared to all of the other stars in our sample. The expectation is that the mass transfer should lead to a depletion of oxygen if there was CNO processing. The anomalous abundance is possibly due to excess line flux contributed by a companion star or to inaccurate NLTE abundance corrections since S 975 lies very near the point where the NLTE corrections become large.
For the two stars (S 975 and S 1082) that overlap with the blue straggler sample of Mathys (1991), the agreement of the O I triplet equivalent widths is very good. We have attempted to test the hypothesis that the blue stragglers and main sequence stars observed all have the same oxygen abundance. As seen in Figure 4, the majority of the stars have EWs that are consistent. The exceptions are S 975 (which appears to have a companion that affects the line flux), S 1434 (B–V = 0.12, EW(O I) = 1071 Mathys (1991)) (which also has a high EW, possibly related to a post-main sequence companion inferred from infrared photometry; Peterson et al., 1984), and S 968 (B–V = 0.12, EW(O I) = 687 Mathys (1991)). S 968 is known to be an Am star, which may account for the difference in that case. However, S 752 (B–V = 0.29, EW(O I) = 754 Mathys (1991)) is also an Am star, and it does not appear to have an anomalous equivalent width. Finally, we find that S 1082 has a low O EW because of the continuum flux contribution of an nearly equal brightness secondary.
To aid in the interpretation of this figure we compare stars from other studies. The large crosses are taken from the Varenne & Monier (1999) study of oxygen along the main sequence of the Hyades. Varenne & Monier found a nearly constant abundance along the main sequence. Note the rapid rise in O I EW at $`BV=0.43`$ ($`(BV)_0=0.38`$) or $`T_{eff}=`$ 6840 K due to NLTE effects. The open circles and filled circles represent synthetic abundance analysis employing constant oxygen abundance (\[O/Fe\] = 0.0) for several temperature with the NLTE corrections of Gratton et al. (1999) and Faraggiana et al. (1988), respectively.
To the accuracy of our measurements, we see no evidence for differences between the abundances of the blue stragglers and the turnoff stars.
In an analysis of two hotter blue stragglers in M67, Mathys (1991) found total CNO abundances that were lower than those of giants in the cluster. Carbon and oxygen in particular were significantly depleted relative to iron. None of the blue stragglers in our sample show this pattern of depletion. Nor do our blue stragglers show differences in C and O abundances compared to our turnoff stars. The abundance results are thus consistent with a collisional formation. However, the low eccentricities of some of the binaries containing stragglers are difficult to explain in that case.
#### 4.2.3 Other Elements
Nuclear processing near the hydrogen burning shell of a giant star could produce anomalous abundances in other, more easily observable elements. Surface anomalies in globular cluster giant stars have been observed in a number of globular clusters (e.g. M3, M13, M15, M71, M92, 47 Tuc).
Gilroy & Brown (1991) showed that the carbon isotope ratio C<sup>12</sup>/C<sup>13</sup> in M67 giants is lower than can be predicted by standard evolutionary theory. However, the required extra mixing does not extend into regions where the oxygen abundance is modified, as is found in metal-poor globular clusters. In fact, the amount of mixing seen in disk giants (Lambert & Ries, 1981) is small in comparison to globular cluster giants (Shetrone, 1996). Na, an element which is burned at a relatively low temperature (Langer et al., 1997; Cavallo et al., 1998), has not been studied well in open cluster giants.
Here in particular our ability to discern which stragglers are most likely to be free from contamination is important. The best candidates (S 984 and S 2204) show fairly minor differences in abundance with the turnoff stars. We must also be aware that there may also be selection effects in that stragglers that form by some mechanisms may have unavoidably contaminated spectra. The straggler S 1082 shows abundances that set it off from the other stars in our sample — in particular Ni and Ba have lower than expected abundances. The measured iron abundance is very different from all of the other cluster stars. The temptation is to think that this may relate to the mass transfer that is inferred. But as mentioned above this star is part of a binary having a composite spectrum that is difficult to disentangle.
### 4.3 Interpretation
#### 4.3.1 Lithium
To date there has not been a detection of lithium in a blue straggler in an open cluster. There are several possible reasons for this, some of which have interesting consequences.
If two stars were to mix completely during a collision, depletions of a factor of 50 might be expected (Pritchet & Glaspey, 1991) because lithium would get diluted throughout the blue straggler. However, hydrodynamical studies of stellar collisions (Lombardi et al., 1995; Sandquist et al., 1997) indicate that the collision itself is unlikely to mix the stars substantially — the high entropy gas near the surfaces of the input stars remains at high entropy near the surface of the remnant (Lombardi et al., 1995). As a result, a collisionally produced straggler is not expected to burn significant lithium initially, even if high enough temperatures are reached (Sills et al., 1997).
For blue stragglers at or above the turnoff of the cluster that were formed via stellar collisions, the lithium abundance observed at the surface is therefore likely to be the same as the abundance at the surface of the input star with the higher surface entropy. For main sequence stars, this would be the more massive star. If the difference in mass of the input stars is relatively large, gas from the lower mass star is likely to be completely sequestered in the interior of the straggler where it would be unobservable.
For open clusters, it is well known that lithium abundance drops with decreasing effective temperature along the main sequence (e.g. Jones et al., 1999). For an open cluster as old as M67, two-body collisions that produce an identifiable blue straggler could involve a star of as low a mass as 0.6 M as the more massive star. According to standard stellar models, main sequence stars with masses less than about 0.8 M are likely to have burned enough of their lithium that it would be undetectable. Standard models are currently unable reproduce the lithium abundance patterns in the Hyades though. Extrapolating from Hyades measurements, stars with masses less than about 1.0 M are likely to be undetectable.
While collisions of 0.3 M and a 1.0 M stars could create blue stragglers with observable lithium, two factors work against this. First, such collisions are relatively rare thanks to the bias in the cluster luminosity function toward stars relatively near the turnoff (Montgomery et al., 1993). Second, significant main sequence Li depletion is possible after the blue straggler is formed. In the Hyades, a dip in lithium abundances is observed for stars with $`T_{eff}6600`$ K (the total range is about 300 K). It is believed that this feature may be due to meridional circulation or shear turbulence (for a recent review, see Talon & Charbonnel, 1998). For M67 no current main sequence members fall in this region, but some blue stragglers do. S 997 and S 2204 fall in the Li gap and may have depleted their Li even if they arrived on the blue straggler sequence with a significant supply. The age of the Hyades ($`625`$ Myr; Perryman et al., 1998) provides an upper limit for the lithium destruction timescale for these stars if this mechanism does act. The lifetime of a blue straggler falling in that range of effective temperatures would be about 3 Gyr, in which case there is less than 25% probability of finding lithium. Blue stragglers which form with more helium in their cores (because the lower mass input star was relatively evolved) would have higher probabilities of showing lithium, but shorter lifetimes as blue stragglers.
More massive blue stragglers ($`T_{eff}>\mathrm{\hspace{0.33em}7000}`$ K) would be more likely to retain surface lithium because both the straggler and the input stars lack extensive surface convection zones, and because the stragglers are hotter than the lithium gap. A possible problem is the mass loss associated with a collision. In hydrodynamical models of globular cluster star collisions by Lombardi et al. (1995) and Sandquist et al. (1997), this mass loss was less than approximately 6% of the total input mass, with the largest mass loss occuring for relatively rare head-on collisions. However, mass loss would tend to remove the most lithium-rich material (for example, see Table 4 of Lombardi et al. (1995)). For stars with masses between about 1.0 and 1.25 M (approximately the turnoff mass), the lithium-rich regions of the star add up to approximately 4-5% of the mass.
Blue stragglers formed via binary mass transfer are unlikely to show surface lithium (Sills et al., 1997). In principle, the brighter, more massive blue stragglers are the ones to observe to look for lithium evidence that different mechanisms are responsible for the formation of M67 blue stragglers. Unfortunately, the Li line is very temperature sensitive and becomes extremely weak in hot massive stars.
#### 4.3.2 CNO elements
Although they are less likely to show the effects of relatively low levels of mixing, CNO elements may prove to be more useful in distinguishing between stragglers that have formed by collisions and by binary mass transfer. Because the CNO elements undergo nuclear processing at higher temperatures than lithium, it is very unlikely that the surface abundances will be modified as the result of a collisional merger. Massive input stars will not have modified their surface CNO abundances because their convection zones are not extensive enough, while low mass input stars will in general become completely engulfed in a higher mass star due to their lower entropies.
On the other hand, if Algol-type binary mass transfer has occurred (on the red giant branch or asymptotic giant branch), the mass visible at the surface of the blue straggler is taken from deep inside the primary where it could have undergone nuclear processing. Spectroscopy of giant branch stars in M67 shows that first dredge-up modifies surface CN abundances in excess of predictions by standard evolution calculations (Gilroy & Brown, 1991; Brown, 1987). So almost any amount of mass transfer by a giant is likely to modify surface abundances of blue stragglers formed in this way. Other species like Na might also be affected (see Pinsonneault, 1997, and references therein). The straggler’s companion (the original primary) is likely to become either a helium or CO white dwarf as a result.
This is an improvement over lithium and other light elements in that it is a pattern change and not a destruction of the observable species. If the companion’s contribution to the spectrum can be removed, this could finally provide a means of determining the formation path a blue straggler took.
## 5 Conclusions
Our spectroscopic analysis of 4 relatively cool blue stragglers does not provide definitive evidence of composition modifications resulting from the formation process. Based on lithium abundance patterns as a function of effective temperature for main sequence stars in the Hyades and M67, we find that current upper limits for the lithium abundances of all spectroscopically observed blue stragglers are consistent with no mixing during formation, even though the upper limits fall below measured abundances for turnoff stars. Although the observable O lines are formed under NLTE conditions, the O I triplet equivalent widths for stars in our sample and those of Mathys (1991) are consistent with constant abundance. The three exceptions (S 968, S 975, and S 1434) may be caused by emission from or absorption by other sources in a binary system. Four of the five blue stragglers and all of the main sequence stars have projected rotational speeds of less than 20 km s<sup>-1</sup>, while the fifth straggler (S 975) has $`v\mathrm{sin}i50`$ km s<sup>-1</sup>.
Although we are unable to determine the existence of binary companions to some of the blue stragglers in our sample, we are able to provide arguments about the degree of contamination of the spectra of most of the stragglers. We confirm the spectroscopic detection of a binary companion to the straggler S 1082. From our spectra, we measure a projected rotational speed of $`90\pm 20`$ km s<sup>-1</sup> for the secondary, and find that its radial velocity amplitude is at least 25 km s<sup>-1</sup>. The lack of a pattern in the radial velocity data (and the variability on the order of hours) provides additional evidence that this system is currently undergoing mass transfer. The primary is found to be 100 – 200 K hotter than its companion. S 1082 may in the end provide us with definitive proof that binary mass transfer can produce blue stragglers.
We would like to thank the anonymous referee for a number of very useful comments that improved the discussion in this paper. E.L.S. would like to thank R. Taam for support (under NSF grant AST-9415423) and R. Taam and P. Etzel for helpful conversations during the course of this work.
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# On supermatrix idempotent operator semigroups
## 1 Introduction
Operator semigroups play an important role in mathematical physics viewed as a general theory of evolution systems . Its development covers many new fields , but one of vital for modern theoretical physics directions — supersymmetry and related mathematical structures — was not considered before in application to operator semigroup theory. The main difference between previous considerations is the fact that among building blocks (e.g. elements of corresponding matrices) there exist noninvertible objects (divisors of zero and nilpotents) which by themselves can form another semigroup. Therefore, we have to take that into account and investigate it properly, which can be called a $`semigroup\times semigroup`$ method.
Here we study continuous supermatrix representations of idempotent operator semigroups firstly introduced in for bands. Usually matrix semigroups are defined over a field $`𝕂`$ (on some nonsupersymmetric generalizations of $`𝕂`$-representations see ). But after discovery of supersymmetry the realistic unified particle theories began to be considered in superspace . So all variables and functions were defined not over a field $`𝕂`$, but over Grassmann-Banach superalgebras over $`𝕂`$ , becoming in general noninvertible, and therefore they should be considered by semigroup theory, which was claimed in , and some semigroups having nontrivial abstract properties were found . Also, it was shown that supermatrices of the special (antitriangle) shape can form various strange and sandwich semigroups not known before . Here we consider one-parametric semigroups (for general theory see ) of antitriangle supermatrices and corresponding superoperator semigroups. The first ones continuously represent idempotent semigroups and second ones lead to new superoperator semigroups with nontrivial properties.
Let $`\mathrm{\Lambda }`$ be a commutative Banach $`_2`$-graded superalgebra over a field $`𝕂`$ (where $`𝕂=,`$ $``$ or $`_p`$) with a decomposition into the direct sum: $`\mathrm{\Lambda }=\mathrm{\Lambda }_0\mathrm{\Lambda }_1`$. The elements $`a`$ from $`\mathrm{\Lambda }_0`$ and $`\mathrm{\Lambda }_1`$ are homogeneous and have the fixed even and odd parity defined as $`\left|a\right|\stackrel{def}{=}\left\{i\{0,1\}=_2\right|a\mathrm{\Lambda }_i\}`$. The even homomorphism $`𝔯_{body}:\mathrm{\Lambda }𝔹`$ is called a body map and the odd homomorphism $`𝔯_{soul}:\mathrm{\Lambda }𝕊`$ is called a soul map , where $`𝔹`$ and $`𝕊`$ are purely even and odd algebras over $`𝕂`$ and $`\mathrm{\Lambda }=𝔹𝕊`$. It can be thought that, if we have the Grassmann algebra $`\mathrm{\Lambda }`$ with generators $`\xi _i,\mathrm{},\xi _n`$ $`\xi _i\xi _j+\xi _j\xi _i=0,\mathrm{\hspace{0.17em}1}i,jn,`$ in particular $`\xi _i^2=0`$ ($`n`$ can be infinite, and only this case is nontrivial and interesting), then any even $`x`$ and odd $`\varkappa `$ elements have the expansions (which can be infinite)
$`x`$ $`=x_{body}+x_{soul}=x_{body}+x_{12}\xi _1\xi _2+x_{13}\xi _1\xi _3+\mathrm{}=`$
$`x_{body}+{\displaystyle \underset{1rn}{}}{\displaystyle \underset{1<i_1<\mathrm{}<i_{2r}n}{}}x_{i_1\mathrm{}i_{2r}}\xi _{i_1}\mathrm{}\xi _{i_{2r}}`$ (1)
$`\varkappa `$ $`=\varkappa _{soul}=\varkappa _1\xi _1+\varkappa _2\xi _2+\mathrm{}+x_{123}\xi _1\xi _2\xi _3+\mathrm{}=`$
$`{\displaystyle \underset{1rn}{}}{\displaystyle \underset{1<i_1<\mathrm{}<i_rn}{}}\varkappa _{i_1\mathrm{}i_{2r1}}\xi _{i_1}\mathrm{}\xi _{i_{2r1}}`$ (2)
From (1)-(2) it follows
###### Corollary 1.
The equations $`x^2=0`$ and $`x\varkappa =0`$ can have nonzero nontrivial solutions (appearing zero divisors and even nilpotents, while odd objects are always nilpotent).
###### Conjecture 2.
If zero divisors and nilpotents will be included in the following analysis as elements of matrices, then one can find new and unusual properties of corresponding semigroups.
For that we should consider general properties of supermatrices and introduce their additional reductions .
## 2 Supermatrices and their even-odd classification
Let us consider $`\left(p|q\right)`$-dimensional linear model superspace $`\mathrm{\Lambda }^{p|q}`$ over $`\mathrm{\Lambda }`$ (in the sense of ) as the even sector of the direct product $`\mathrm{\Lambda }^{p|q}=\mathrm{\Lambda }_0^p\times \mathrm{\Lambda }_1^q`$ . The even morphisms $`\text{Hom}_0(\mathrm{\Lambda }^{p|q},\mathrm{\Lambda }^{m|n})`$ between superlinear spaces $`\mathrm{\Lambda }^{p|q}\mathrm{\Lambda }^{m|n}`$ are described by means of $`\left(m+n\right)\times \left(p+q\right)`$-supermatrices (for some nontrivial properties see ). In what follows we will treat noninvertible morphisms on a par with invertible ones .
First we consider $`\left(1+1\right)\times \left(1+1\right)`$-supermatrices describing the elements from $`\text{Hom}_0(\mathrm{\Lambda }^{1|1},\mathrm{\Lambda }^{1|1})`$ in the standard $`\mathrm{\Lambda }^{1|1}`$ basis
$$M\left(\begin{array}{cc}a& \alpha \\ \beta & b\end{array}\right)\text{Mat}_\mathrm{\Lambda }\left(1|1\right)$$
(3)
where $`a,b\mathrm{\Lambda }_0,\alpha ,\beta \mathrm{\Lambda }_1,`$ $`\alpha ^2=\beta ^2=0`$ (in the following we use Latin letters for elements from $`\mathrm{\Lambda }_0`$ and Greek letters for ones from $`\mathrm{\Lambda }_1`$, and all odd elements are nilpotent of index 2).
The supertrace and Berezinian (superdeterminant) are defined by
$$\text{str}M=ab,$$
(4)
$$\text{Ber}M=\frac{a}{b}+\frac{\beta \alpha }{b^2}.$$
(5)
First term corresponds to triangle supermatrices, second term - to antitriangle ones. So we obviously have different two dual types of supermatrices .
###### Definition 3.
Even-reduced supermatrices are elements from $`\text{Mat}_\mathrm{\Lambda }\left(1|1\right)`$ of the form
$$M_{even}\left(\begin{array}{cc}a& \alpha \\ 0& b\end{array}\right)\text{RMat}_\mathrm{\Lambda }^{even}\left(1|1\right)\text{Mat}_\mathrm{\Lambda }\left(1|1\right).$$
(6)
Odd-reduced supermatrices are elements from $`\text{Mat}_\mathrm{\Lambda }\left(1|1\right)`$ of the form
$$M_{odd}\left(\begin{array}{cc}0& \alpha \\ \beta & b\end{array}\right)\text{RMat}_\mathrm{\Lambda }^{odd}\left(1|1\right)\text{Mat}_\mathrm{\Lambda }\left(1|1\right).$$
(7)
###### Conjecture 4.
The odd-reduced supermatrices have a nilpotent (but nonzero) Berezinian
$$\text{Ber}M_{odd}=\frac{\beta \alpha }{b^2}0,\left(\text{Ber}M_{odd}\right)^2=0.$$
(8)
###### Remark 5.
Indeed this property (8) prevented in the past the use of this type (odd-reduced) of supermatrices in physics. All previous applications (excluding ) were connected with triangle (even-reduced, similar to Borel) ones and first term in Berezinian $`\text{Ber}M=\frac{a}{b}`$ (5).
The even- and odd-reduced supermatrices are mutually dual in the sense of the Berezinian addition formula
$$\text{Ber}M=\text{Ber}M_{even}+\text{Ber}M_{odd}.$$
(9)
For sets of matrices we use corresponding bold symbols, and the set product is standard
$$𝐌𝐍\stackrel{def}{=}\left\{MN\right|M,N\text{Mat}_\mathrm{\Lambda }\left(1|1\right)\}.$$
The matrices from $`\text{Mat}\left(1|1\right)`$ form a linear semigroup of $`\left(1+1\right)\times \left(1+1\right)`$-supermatrices under the standard supermatrix multiplication $`𝔐\left(1|1\right)\stackrel{def}{=}\left\{𝐌\right|\}`$ . Obviously, the even-reduced matrices $`𝐌_{even}`$ form a semigroup $`𝔐_{even}\left(1|1\right)`$ which is a subsemigroup of $`𝔐\left(1|1\right)`$, because of $`𝐌_{even}𝐌_{even}𝐌_{even}`$ and the unity is in $`𝔐_{even}\left(1|1\right)`$. This trivial observation leads to general structure (Borel) theory for matrices: triangle matrices form corresponding substructures (subgroups and subsemigroups). It was believed before that in case of supermatrices the situation does not changed, because supermatrix multiplication is the same . But they did not take into account zero divisors and nilpotents appearing naturally and inevitably in supercase.
###### Conjecture 6.
Standard (lower/upper) triangle supermatrices are not the only substructures due to unusual properties of zero divisors and nilpotents appearing among elements (see (1)-(2) and Corollary 1).
It means that in such consideration we have additional (to triangle) class of subsemigroups. Then we can formulate the following general
###### Problem 1.
For a given $`n,m,p,q`$ to describe and classify all possible substructures (subgroups and subsemigroups) of $`\left(m+n\right)\times \left(p+q\right)`$-supermatrices.
First example of such new substructures are $`\mathrm{\Gamma }`$-matrices considered below.
###### Conjecture 7.
These new substructures lead to new corresponding superoperators which are represented by one-parameter substructures of supermatrices.
Therefore we first consider possible (not triangle) subsemigroups of supermatrices.
## 3 Odd-reduced supermatrix semigroups
In general, the odd-reduced matrices $`M_{odd}`$ do not form a semigroup, since their multiplication is not closed in general $`𝐌_{odd}𝐌_{odd}𝐌`$. Nevertheless, some subset of $`𝐌_{odd}`$ can form a semigroup . That can happen due to the existence of zero divisors in $`\mathrm{\Lambda }`$, and so we have $`𝐌_{odd}𝐌_{odd}𝐌_{odd}=𝐌_{odd}^{smg}\mathrm{}`$.
To find $`𝐌_{odd}^{smg}`$ we consider a $`\left(1+1\right)\times \left(1+1\right)`$ example. Let $`\alpha ,\beta \mathrm{\Gamma }_{set}`$, where $`\mathrm{\Gamma }_{set}\mathrm{\Lambda }_1`$. We denote $`\text{Ann}\alpha \stackrel{def}{=}\left\{\gamma \mathrm{\Lambda }_1\right|\gamma \alpha =0\}`$ and $`\text{Ann}\mathrm{\Gamma }_{set}=\underset{\alpha \mathrm{\Gamma }}{}\text{Ann}\alpha `$ (here the intersection is crucial). Then we define left and right $`\mathrm{\Gamma }`$-matrices
$$𝐌_{odd(L)}^\mathrm{\Gamma }\stackrel{def}{=}\left(\begin{array}{cc}0& \mathrm{\Gamma }_{set}\\ \text{Ann}\mathrm{\Gamma }_{set}& b\end{array}\right),$$
(10)
$$𝐌_{odd(R)}^\mathrm{\Gamma }\stackrel{def}{=}\left(\begin{array}{cc}0& \text{Ann}\mathrm{\Gamma }_{set}\\ \mathrm{\Gamma }_{set}& b\end{array}\right).$$
(11)
###### Proposition 8.
The $`\mathrm{\Gamma }`$-matrices $`𝐌_{odd(L,R)}^\mathrm{\Gamma }𝐌_{odd}`$ form subsemigroups of $`𝔐\left(1|1\right)`$ under the standard supermatrix multiplication, if $`b\mathrm{\Gamma }\mathrm{\Gamma }`$.
###### Definition 9.
$`\mathrm{\Gamma }`$-semigroups $`𝔐_{odd(L,R)}^\mathrm{\Gamma }\left(1|1\right)`$ are subsemigroups of $`𝔐\left(1|1\right)`$ formed by the $`\mathrm{\Gamma }`$-matrices $`𝐌_{odd(L,R)}^\mathrm{\Gamma }`$ under supermatrix multiplication.
###### Corollary 10.
The $`\mathrm{\Gamma }`$-matrices are additional to triangle supermatrices substructures which form semigroups.
Let us consider general square antitriangle $`\left(p+q\right)\times \left(p+q\right)`$-supermatrices (having even parity in notations of ) of the form
$$M_{odd}^{p|q}\stackrel{def}{=}\left(\begin{array}{cc}0_{p\times p}& \mathrm{\Gamma }_{p\times q}\\ \mathrm{\Delta }_{q\times p}& B_{q\times q}\end{array}\right),$$
(12)
where ordinary matrix $`B_{q\times q}`$ consists of even elements and matrices $`\mathrm{\Gamma }_{p\times q}`$ and $`\mathrm{\Delta }_{q\times p}`$ consist of odd elements (we drop their indices below). The berezinian of $`M_{odd}^{p|q}`$ can be obtained from the general formula by reduction and in case of invertible $`B`$ (which is implied here) is (cf. (8))
$$BerM_{odd}^{p|q}=\frac{det\left(\mathrm{\Gamma }B^1\mathrm{\Delta }\right)}{detB}.$$
(13)
###### Assertion 11.
A set of supermatrices $`𝐌_{odd}^{p|q}`$ form a semigroup $`𝔐_{odd}^\mathrm{\Gamma }\left(p|q\right)`$ of $`\mathrm{\Gamma }^{p|q}`$-matrices, if $`\mathrm{\Gamma }_{set}\mathrm{\Delta }_{set}=0`$, i.e. antidiagonal matrices are orthogonal, and $`\mathrm{\Gamma }_{set}𝐁\mathrm{\Gamma }_{set}`$, $`𝐁\mathrm{\Delta }_{set}\mathrm{\Delta }_{set}`$.
###### Proof.
Consider the product
$$M_{odd_1}^{p|q}M_{odd_2}^{p|q}=\left(\begin{array}{cc}\mathrm{\Gamma }_1\mathrm{\Delta }_2& \mathrm{\Gamma }_1B_2\\ B_1\mathrm{\Delta }_2& B_1B_2+\mathrm{\Delta }_1\mathrm{\Gamma }_2\end{array}\right)$$
(14)
and observe the condition of vanishing even-even block, which gives $`\mathrm{\Gamma }_1\mathrm{\Delta }_2=0`$, and other conditions follows obviously. ∎
From (14) it follows
###### Corollary 12.
Two $`\mathrm{\Gamma }^{p|q}`$-matrices satisfy the band relation $`M_1M_2=M_1`$, if $`\mathrm{\Gamma }_1B_2=\mathrm{\Gamma }_1,`$ $`B_1\mathrm{\Delta }_2=\mathrm{\Delta }_2,`$ $`B_1B_2+\mathrm{\Delta }_1\mathrm{\Gamma }_2=B_1`$.
###### Definition 13.
We call a set of $`\mathrm{\Gamma }^{p|q}`$-matrices satisfying additional condition $`\mathrm{\Delta }_{set}\mathrm{\Gamma }_{set}=0`$, a set of strong $`\mathrm{\Gamma }^{p|q}`$-matrices.
Strong $`\mathrm{\Gamma }^{p|q}`$-matrices have some extra nice features and all supermatrices considered below are of this class.
###### Corollary 14.
Idempotent strong $`\mathrm{\Gamma }^{p|q}`$-matrices are defined by relations $`\mathrm{\Gamma }B=\mathrm{\Gamma },`$ $`B\mathrm{\Delta }=\mathrm{\Delta },`$ $`B^2=B`$.
The product of $`n`$ strong $`\mathrm{\Gamma }^{p|q}`$-matrices $`M_i`$ has the following form
$$M_1M_2\mathrm{}M_n=\left(\begin{array}{cc}0& \mathrm{\Gamma }_1A_{n1}B_n\\ B_1A_{n1}\mathrm{\Delta }_n& B_1A_{n1}B_n\end{array}\right),$$
(15)
where $`A_{n1}=B_2B_3\mathrm{}B_{n1}`$, and its berezinian is
$$Ber\left(M_1M_2\mathrm{}M_n\right)=\frac{det\left(\mathrm{\Gamma }_1A_{n1}\mathrm{\Delta }_n\right)}{det\left(B_1A_{n1}B_n\right)}.$$
(16)
## 4 One-even-parameter supermatrix idempotent semigroups
Here we investigate one-even-parameter subsemigroups of $`\mathrm{\Gamma }`$-semigroups and as a particular example for clearness of statements consider $`𝔐_{odd}\left(1|1\right)`$, where all characteristic features taking place in general $`\left(p+q\right)\times \left(p+q\right)`$ as well can be seen. These formulas will be applied for establishing corresponding superoperator semigroup properties.
A simplest semigroup can be constructed from antidiagonal nilpotent supermatrices of the shape
$$Y_\alpha \left(t\right)\stackrel{def}{=}\left(\begin{array}{cc}0& \alpha t\\ \alpha & 0\end{array}\right).$$
(17)
where $`t\mathrm{\Lambda }^{1|0}`$ is an even parameter of the Grassmann algebra $`\mathrm{\Lambda }`$ which continuously ”numbers” elements $`Y_\alpha \left(t\right)`$ and $`\alpha \mathrm{\Lambda }^{0|1}`$ is a fixed odd element of $`\mathrm{\Lambda }`$ which ”numbers” the sets $`𝐘_\alpha =\underset{t}{}Y_\alpha \left(t\right)`$.
###### Definition 15.
The supermatrices $`Y_\alpha \left(t\right)`$ together with a null supermatrix $`Z\stackrel{def}{=}\left(\begin{array}{cc}0& 0\\ 0& 0\end{array}\right)`$ form a continuous null semigroup $`_\alpha \left(1|1\right)=\{𝐘_\alpha Z;\}`$ having the null multiplication
$$Y_\alpha \left(t\right)Y_\alpha \left(u\right)=Z.$$
(18)
###### Assertion 16.
For any fixed $`t\mathrm{\Lambda }^{1|0}`$ the set $`\{Y_\alpha \left(t\right),Z\}`$ is a 0-minimal ideal in $`_\alpha \left(1|1\right)`$.
###### Remark 17.
If we consider, for instance, a one-even-parameter odd-reduced supermatrix $`R_\alpha \left(t\right)=\left(\begin{array}{cc}0& \alpha \\ \alpha & t\end{array}\right)`$, then multiplication of $`R_\alpha \left(t\right)`$ is not closed since $`R_\alpha \left(t\right)R_\alpha \left(u\right)=\left(\begin{array}{cc}0& \alpha u\\ \alpha t& tu\end{array}\right)𝐑_\alpha =\underset{t}{}R_\alpha \left(t\right)`$. Any other possibility except ones considered below also do not give closure of multiplication.
Thus the only nontrivial closed systems of one-even-parameter odd-reduced (antitriangle) $`\left(1+1\right)\times \left(1+1\right)`$ supermatrices are $`𝐏_\alpha =\underset{t}{}P_\alpha \left(t\right)`$ where
$$P_\alpha \left(t\right)\stackrel{def}{=}\left(\begin{array}{cc}0& \alpha t\\ \alpha & 1\end{array}\right)$$
(19)
and $`𝐐_\alpha =\underset{t}{}Q_\alpha \left(u\right)`$ where
$$Q_\alpha \left(u\right)\stackrel{def}{=}\left(\begin{array}{cc}0& \alpha \\ \alpha u& 1\end{array}\right).$$
(20)
First, we establish multiplication properties of supermatrices $`P_\alpha \left(t\right)`$ and $`Q_\alpha \left(u\right)`$. Obviously, that they are idempotent.
###### Assertion 18.
Sets of idempotent supermatrices $`𝐏_\alpha `$ and $`𝐐_\alpha `$ form left zero and right zero semigroups respectively with multiplication
$`P_\alpha \left(t\right)P_\alpha \left(u\right)`$ $`=P_\alpha \left(t\right),`$ (21)
$`Q_\alpha \left(t\right)Q_\alpha \left(u\right)`$ $`=Q_\alpha \left(u\right).`$ (22)
if and only if $`\alpha ^2=0`$.
###### Proof.
It simply follows from supermatrix multiplication law and general previous considerations. ∎
###### Corollary 19.
The sets $`𝐏_\alpha `$ and $`𝐐_\alpha `$ are rectangular bands since
$`P_\alpha \left(t\right)P_\alpha \left(u\right)P_\alpha \left(t\right)`$ $`=P_\alpha \left(t\right),`$ (23)
$`P_\alpha \left(u\right)P_\alpha \left(t\right)P_\alpha \left(u\right)`$ $`=P_\alpha \left(u\right)`$ (24)
and
$`Q_\alpha \left(u\right)Q_\alpha \left(t\right)Q_\alpha \left(u\right)`$ $`=Q_\alpha \left(u\right),`$ (25)
$`Q_\alpha \left(t\right)Q_\alpha \left(u\right)Q_\alpha \left(t\right)`$ $`=Q_\alpha \left(t\right)`$ (26)
with components $`t=t_0+Ann\alpha `$ and $`u=u_0+Ann\alpha `$ correspondingly.
They are orthogonal in sense of
$$Q_\alpha \left(t\right)P_\alpha \left(u\right)=E_\alpha ,$$
(27)
where
$$E_\alpha \stackrel{def}{=}\left(\begin{array}{cc}0& \alpha \\ \alpha & 1\end{array}\right)$$
(28)
is a right “unity” and left “zero” in semigroup $`𝐏_\alpha `$, because
$$P_\alpha \left(t\right)E_\alpha =P_\alpha \left(t\right),E_\alpha P_\alpha \left(t\right)=E_\alpha $$
(29)
and a left “unity” and right “zero” in semigroup $`𝐐_\alpha `$, because
$$Q_\alpha \left(t\right)E_\alpha =E_\alpha ,E_\alpha Q_\alpha \left(t\right)=Q_\alpha \left(t\right).$$
(30)
It is important to note that
$$P_\alpha \left(t=1\right)=Q_\alpha \left(t=1\right)=E_\alpha ,$$
(31)
and so $`𝐏_\alpha 𝐐_\alpha =E_\alpha `$. Therefore, almost all properties of $`𝐏_\alpha `$ and $`𝐐_\alpha `$ are similar, and we will consider only one of them in what follows. For generalized Green’s relations and more detail properties of odd-reduced supermatrices see .
## 5 Odd-reduced supermatrix operator semigroups
Let us consider a semigroup $`𝒫`$ of superoperators $`𝖯\left(t\right)`$ (see for general theory ) represented by the one-even-parameter semigroup $`𝐏_\alpha `$ of odd-reduced supermatrices $`P_\alpha \left(t\right)`$ (19) which act on $`\left(1|1\right)`$-dimensional superspace $`^{1|1}`$ as follows $`P_\alpha \left(t\right)𝚇`$, where $`𝚇=\left(\begin{array}{c}x\\ \varkappa \end{array}\right)^{1|1}`$, where $`x`$ \- even coordinate, $`\varkappa `$ \- odd coordinate $`\left(\varkappa ^2=0\right)`$ having expansions (1) and (2) respectively (see Corollary 1). We have a representation $`\rho :𝒫𝐏_\alpha `$ with correspondence $`𝖯\left(t\right)P_\alpha \left(t\right)`$, but (as is usually made, e.g. ) we identify space of superoperators with the space of corresponding matrices (nevertheless, we use here operator notations for convenience).
###### Definition 20.
An odd-reduced “dynamical” system on $`^{1|1}`$ is defined by an odd-reduced supermatrix-valued function $`𝖯():_+𝔐_{odd}\left(1|1\right)`$ and “time evolution” of the state $`𝚇\left(0\right)^{1|1}`$given by the function $`𝚇\left(t\right):_+^{1|1}`$, where
$$𝚇\left(t\right)=𝖯\left(t\right)𝚇\left(0\right)$$
(32)
and can be called as orbit of $`𝚇\left(0\right)`$ under $`𝖯()`$.
###### Remark 21.
In general the definition, the continuity, the functional equation and most of conclusions below hold valid also for $`t^{1|0}`$ (as e.g. in \[5, p. 9\]) including “nilpotent time” directions (see Corollary 1).
From (21) it follows that
$$𝖯\left(t\right)𝖯\left(s\right)=𝖯\left(t\right),$$
(33)
and so superoperators $`𝖯\left(t\right)`$ are idempotent. Also they form a rectangular band, because of
$`𝖯\left(t\right)𝖯\left(s\right)𝖯\left(t\right)`$ $`=𝖯\left(t\right),`$ (34)
$`𝖯\left(s\right)𝖯\left(t\right)𝖯\left(s\right)`$ $`=𝖯\left(s\right).`$ (35)
We observe that
$$𝖯\left(0\right)=\left(\begin{array}{cc}0& 0\\ \alpha & 1\end{array}\right)𝖨=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),$$
(36)
as opposite to the standard case . A “generator” $`𝖠=𝖯^{}\left(t\right)`$ is
$$𝖠=\left(\begin{array}{cc}0& \alpha \\ 0& 0\end{array}\right),$$
(37)
and so the standard definition of generator
$$𝖠=\underset{t0}{lim}\frac{𝖯\left(t\right)𝖯\left(0\right)}{t}.$$
(38)
holds and for difference we have the standard relation
$$𝖯\left(t\right)𝖯\left(s\right)=𝖠\left(ts\right).$$
(39)
The following properties of the generator $`𝖠`$ take place
$`𝖯\left(t\right)𝖠`$ $`=𝖹,`$ (40)
$`\mathrm{𝖠𝖯}\left(t\right)`$ $`=𝖠,`$ (41)
where “zero operator” $`𝖹`$ is represented by the null supermatrix, $`𝖠^2=𝖹,`$ and therefore generator $`𝖠`$ is a nilpotent of second degree.
From (38) it follows that
$$𝖯\left(t\right)=𝖯\left(0\right)+𝖠t.$$
(42)
###### Definition 22.
We call operators which can be presented as a linear supermatrix function of $`t`$ a $`t`$-linear superoperators.
From (42) it follows that $`𝖯\left(t\right)`$ is a $`t`$-linear superoperator.
###### Proposition 23.
Superoperators $`𝖯\left(t\right)`$ cannot be presented as an exponent (as for the standard superoperator semigroups $`𝖳\left(t\right)=e^{𝖠t}`$ ).
###### Proof.
In our case
$$𝖳\left(t\right)=e^{𝖠t}=𝖨+𝖠t=\left(\begin{array}{cc}1& \alpha t\\ 0& 1\end{array}\right)𝐏_\alpha .$$
(43)
###### Remark 24.
Exponential superoperator $`𝖳\left(t\right)=e^{𝖠t}`$ is represented by even-reduced supermatrices $`𝖳():_+𝔐_{even}\left(1|1\right)`$ , but idempotent superoperator $`𝖯\left(t\right)`$ is represented by odd-reduced supermatrices $`𝖯():_+𝔐_{odd}\left(1|1\right)`$ (see Definition 3).
Nevertheless, the superoperator $`𝖯\left(t\right)`$ satisfies the same linear differential equation
$$𝖯^{}\left(t\right)=𝖠𝖯\left(t\right)$$
(44)
as the standard exponential superoperator $`𝖳\left(t\right)`$ (the initial value problem )
$$𝖳^{}\left(t\right)=𝖠𝖳\left(t\right).$$
(45)
That leads to the following
###### Corollary 25.
In case initial state does not equal unity $`𝖯\left(0\right)𝖨`$, there exists an additional class of solutions of the initial value problem (44)-(45) among odd-reduced (antidiagonal) idempotent $`t`$-linear (nonexponential) superoperators.
Let us compare behavior of superoperators $`𝖯\left(t\right)`$ and $`𝖳\left(t\right)`$. First of all, their generators coincide
$$𝖯^{}\left(0\right)=𝖳^{}\left(0\right)=𝖠.$$
(46)
But powers of $`𝖯\left(t\right)`$ and $`𝖳\left(t\right)`$ are different $`𝖯^n\left(t\right)=𝖯\left(t\right)`$ and $`𝖳^n\left(t\right)=𝖳\left(nt\right)`$. In their common actions the superoperator which is from the left transfers its properties to the right hand side as follows
$`𝖳^n\left(t\right)𝖯\left(t\right)`$ $`=𝖯\left(\left(n+1\right)t\right),`$ (47)
$`𝖯^n\left(t\right)𝖳\left(t\right)`$ $`=𝖯\left(t\right).`$ (48)
Their commutator is nonvanishing
$$\left[𝖳\left(t\right)𝖯\left(s\right)\right]=𝖯^{}\left(0\right)t=𝖳^{}\left(0\right)t=𝖠t,$$
(49)
which can be compared with the pure exponential commutator (for our case) $`\left[𝖳\left(t\right)𝖳\left(u\right)\right]=0`$ and idempotent commutator
$$\left[𝖯\left(t\right)𝖯\left(s\right)\right]=𝖯^{}\left(0\right)\left(ts\right)=𝖠\left(ts\right).$$
(50)
###### Assertion 26.
All superoperators $`𝖯\left(t\right)`$ and $`𝖳\left(t\right)`$ commute in case of “nilpotent time” and
$$tAnn\alpha .$$
(51)
###### Remark 27.
The uniqueness theorem \[5, p. 3\] holds only for $`𝖳\left(t\right)`$, because the nonvanishing commutator $`[𝖠,𝖯\left(t\right)]=𝖠0`$.
###### Corollary 28.
The superoperator $`𝖳\left(t\right)`$ is an inner inverse for $`𝖯\left(t\right)`$, because of
$$𝖯\left(t\right)𝖳\left(t\right)𝖯\left(t\right)=𝖯\left(t\right),$$
(52)
but it is not an outer inverse, because
$$𝖳\left(t\right)𝖯\left(t\right)𝖳\left(t\right)=𝖯\left(2t\right).$$
(53)
Let us try to find a (possibly noninvertible) operator $`𝖴`$ which connects exponential and idempotent superoperators $`𝖯\left(t\right)`$ and $`𝖳\left(t\right)`$.
###### Assertion 29.
The “semi-similarity” relation
$$𝖳\left(t\right)𝖴=\mathrm{𝖴𝖯}\left(t\right)$$
(54)
holds if
$$𝖴=\left(\begin{array}{cc}\sigma \alpha & \sigma \\ 0& \rho \alpha \end{array}\right)$$
(55)
which is noninvertible triangle and depends from two odd constants, and the “adjoint” relation
$$𝖴^{}𝖳\left(t\right)=𝖯\left(t\right)𝖴^{}$$
(56)
holds if
$$𝖴^{}=\left(\begin{array}{cc}0& \alpha vt\\ \alpha u& v\end{array}\right)$$
(57)
which is also noninvertible antitriangle and depends from two even constants and “time”.
Note that $`𝖴`$ is nilpotent of third degree, since $`𝖴^2=\sigma \rho 𝖠`$, but the “adjoint” superoperator is not nilpotent at all if $`v`$ is not nilpotent.
Both $`𝖠`$ and $`𝖹`$ behave as zeroes, but $`𝖸\left(t\right)`$ (see (17)) is a two-sided zero for $`𝖳\left(t\right)`$ only, since
$$𝖳\left(t\right)𝖸\left(t\right)=𝖸\left(t\right)𝖳\left(t\right)=𝖸\left(t\right),$$
(58)
but
$`𝖯\left(t\right)𝖸\left(t\right)`$ $`=𝖸\left(0\right),`$ (59)
$`𝖸\left(t\right)𝖯\left(t\right)`$ $`=𝖠t.`$ (60)
If we add $`𝖠`$ and $`𝖹`$ to superoperators $`𝖯\left(t\right)`$, then we obtain an extended odd-reduced noncommutative superoperator semigroup $`𝒫_{odd}=𝖯\left(t\right)𝖠𝖹`$ with the following Cayley table (for convenience we add $`𝖸\left(t\right)`$ and $`𝖳\left(t\right)`$ as well)
| $`12`$ | $`𝖯\left(t\right)`$ | $`𝖯\left(s\right)`$ | $`𝖠`$ | $`𝖹`$ | $`𝖸\left(t\right)`$ | $`𝖳\left(t\right)`$ | $`𝖳\left(s\right)`$ |
| --- | --- | --- | --- | --- | --- | --- | --- |
| $`𝖯\left(t\right)`$ | $`𝖯\left(t\right)`$ | $`𝖯\left(t\right)`$ | $`𝖹`$ | $`𝖹`$ | $`𝖯\left(t\right)`$ | $`𝖯\left(t\right)`$ | $`𝖯\left(t\right)`$ |
| $`𝖯\left(s\right)`$ | $`𝖯\left(s\right)`$ | $`𝖯\left(s\right)`$ | $`𝖹`$ | $`𝖹`$ | $`𝖯\left(s\right)`$ | $`𝖯\left(s\right)`$ | $`𝖯\left(s\right)`$ |
| $`𝖠`$ | $`𝖠`$ | $`𝖠`$ | $`𝖹`$ | $`𝖹`$ | $`𝖹`$ | $`𝖠`$ | $`𝖠`$ |
| $`𝖹`$ | $`𝖹`$ | $`𝖹`$ | $`𝖹`$ | $`𝖹`$ | $`𝖹`$ | $`𝖹`$ | $`𝖹`$ |
| $`𝖸\left(t\right)`$ | $`𝖠t`$ | $`𝖠s`$ | $`𝖹`$ | $`𝖹`$ | $`𝖹`$ | $`𝖸\left(t\right)`$ | $`𝖸\left(t\right)`$ |
| $`𝖳\left(t\right)`$ | $`𝖯\left(2t\right)`$ | $`𝖯\left(t+s\right)`$ | $`𝖠`$ | $`𝖹`$ | $`𝖸\left(t\right)`$ | $`𝖳\left(2t\right)`$ | $`𝖳\left(t+s\right)`$ |
| $`𝖳\left(s\right)`$ | $`𝖯\left(t+s\right)`$ | $`𝖯\left(2s\right)`$ | $`𝖠`$ | $`𝖹`$ | $`𝖸\left(t\right)`$ | $`𝖳\left(t+s\right)`$ | $`𝖳\left(2s\right)`$ |
(61a)
It is easily seen that associativity in the left upper square holds, and so the table (61a) is actually represents a semigroup of superoperators $`𝒫_{odd}`$ (under supermatrix multiplication).
The analogs of the “smoothing operator” $`𝖵\left(t\right)`$ are
$`𝖵_P\left(t\right)`$ $`={\displaystyle \underset{0}{\overset{t}{}}}𝖯\left(s\right)𝑑s={\displaystyle \frac{t}{2}}\left(𝖯\left(t\right)+𝖯\left(0\right)\right)=\left(\begin{array}{cc}0& \alpha {\displaystyle \frac{t^2}{2}}\\ \alpha t& t\end{array}\right),`$ (64)
$`𝖵_T\left(t\right)`$ $`={\displaystyle \underset{0}{\overset{t}{}}}𝖳\left(s\right)𝑑s={\displaystyle \frac{t}{2}}\left(𝖳\left(t\right)+𝖳\left(0\right)\right)=\left(\begin{array}{cc}t& \alpha {\displaystyle \frac{t^2}{2}}\\ 0& t\end{array}\right).`$ (67)
Let us consider the differential sequence of sets of superoperators $`𝖯\left(t\right)`$
$$𝖲_n\stackrel{}{}𝖲_{n1}\stackrel{}{}\mathrm{}𝖲_1\stackrel{}{}𝖲_0\stackrel{}{}𝖠\stackrel{}{}𝖹,$$
(68)
where $`=d/dt`$ and
$$𝖲_n=\underset{t}{}\frac{t^n}{n\left(n1\right)\mathrm{}1}𝖯\left(\frac{t}{n+1}\right),$$
(69)
and by definition
$`𝖲_0`$ $`={\displaystyle \underset{t}{}}𝖯\left(t\right),`$ (70)
$`𝖲_1`$ $`={\displaystyle \underset{t}{}}𝖵_P\left(t\right).`$ (71)
Now we construct an analog of the standard operator semigroup functional equation
$$𝖳\left(t+s\right)=𝖳\left(t\right)𝖳\left(s\right).$$
(72)
Using the multiplication law (33) and manifest representation (19). for the idempotent superoperators $`𝖯\left(t\right)`$ we can formulate
###### Definition 30.
The odd-reduced idempotent superoperators $`𝖯\left(t\right)`$ satisfy the following generalized functional equation
$$𝖯\left(t+s\right)=𝖯\left(t\right)𝖯\left(s\right)+𝖭(t,s),$$
(73)
where
$$𝖭(t,s)=𝖯^{}\left(t\right)s.$$
The presence of second term $`𝖭(t,s)`$ in the right hand side of the generalized functional equation (73) can be connected with nonautonomous and deterministic properties of systems describing by it . Indeed, from (32) it follows that
$`𝚇\left(t+s\right)`$ $`=𝖯\left(t+s\right)𝚇\left(0\right)=𝖯\left(t\right)𝖯\left(s\right)𝚇\left(0\right)+𝖯^{}\left(t\right)s𝚇\left(0\right)`$ (74)
$`=𝖯\left(t\right)𝚇\left(s\right)+𝖯^{}\left(t\right)s𝚇\left(0\right)𝖯\left(t\right)𝚇\left(s\right)`$
as opposite to the always implied relation for exponential superoperators $`𝖳\left(t\right)`$ (translational property )
$$𝖳\left(t\right)𝚇\left(s\right)=𝚇\left(t+s\right),$$
(75)
which follows from (72). Instead of (75), using the band property (33) we obtain
$$𝖯\left(t\right)𝚇\left(s\right)=𝚇\left(t\right),$$
(76)
which can be called the “moving time” property.
###### Problem 2.
Find a “dynamical system” with time evolution satisfying the “moving time” property (76) instead of the translational property (75).
###### Assertion 31.
For “nilpotent time” satisfying (51) the generalized functional equation (73) coincides with the standard functional equation (72), and therefore the idempotent operators $`𝖯\left(t\right)`$ describe autonomous and deterministic “dynamical” system and satisfy the translational property (75).
###### Proof.
Follows from (51) and (74). ∎
###### Problem 3.
Find all maps $`𝖯():_+𝔐\left(p|q\right)`$ satisfying the generalized functional equation (73).
We turn to this problem later, and now consider some features of the Cauchy problem for idempotent superoperators.
## 6 Cauchy problem
Let us consider an action (32) of superoperator $`𝖯\left(t\right)`$ in superspace $`^{1|1}`$ as $`𝚇\left(t\right)=𝖯\left(t\right)𝚇\left(0\right)`$, where the initial components are $`𝚇\left(0\right)=\left(\begin{array}{c}x_0\\ \varkappa _0\end{array}\right)`$. From (32) the evolution of the components has the form
$$\left(\begin{array}{c}x\left(t\right)\\ \varkappa \left(t\right)\end{array}\right)=\left(\begin{array}{c}\alpha \varkappa _0t\\ \alpha x_0+\varkappa _0\end{array}\right)$$
(77)
which shows that superoperator $`𝖯\left(t\right)`$ does not lead to time dependence of odd components. Then from (77) we see that
$$𝚇^{}\left(t\right)=\left(\begin{array}{c}\alpha \varkappa _0\\ 0\end{array}\right)=const.$$
(78)
This is in full agreement with an analog of the Cauchy problem for our case
$$𝚇^{}\left(t\right)=𝖠𝚇\left(t\right).$$
(79)
###### Assertion 32.
The solution of the Cauchy problem (79) is given by (32), but the idempotent superoperator $`𝖯\left(t\right)`$ can not be presented in exponential form as in the standard case , but only in the $`t`$-linear form $`𝖯\left(t\right)=𝖯\left(0\right)+𝖠te^{𝖠t},`$ as we have already shown in (42).
This allows us to formulate
###### Theorem 33.
In superspace the solution of the Cauchy initial problem with the same generator $`𝖠`$ is two-fold and is given by two different type of superoperators:
1. Exponential superoperator $`𝖳\left(t\right)`$ represented by the even-reduced supermatrices;
2. Idempotent $`t`$-linear superoperator $`𝖯\left(t\right)`$ represented by the odd-reduced supermatrices.
For comparison the standard solution of the Cauchy problem (79)
$$𝚇\left(t\right)=𝖳\left(t\right)𝚇\left(0\right)$$
in components is
$$\left(\begin{array}{c}x\left(t\right)\\ \varkappa \left(t\right)\end{array}\right)=\left(\begin{array}{c}x_0+\alpha \varkappa _0t\\ \varkappa _0\end{array}\right),$$
(80)
which shows that the time evolution of even coordinate is also in nilpotent even direction $`\alpha \varkappa _0`$ as in (77), but with addition of initial (possibly nonilpotent) $`x_0`$, while odd coordinate is (another) constant as well. That leads to
###### Assertion 34.
“Even” and “odd” evolutions coincide if even initial coordinate vanishes $`x_0=0`$ or common starting point is pure odd $`𝚇\left(0\right)=\left(\begin{array}{c}0\\ \varkappa _0\end{array}\right)`$.
A very much important formula is the condition of commutativity
$$[𝖠,𝖯\left(t\right)]𝚇\left(t\right)=𝖠𝚇\left(t\right)=\left(\begin{array}{c}\alpha \varkappa \left(t\right)\\ 0\end{array}\right)=0,$$
(81)
which satisfies, when $`\alpha \varkappa \left(t\right)=0`$, while in the standard case the commutator $`[𝖠,𝖳\left(t\right)]𝚇\left(t\right)=0`$, i.e. vanishes without any additional conditions .
## 7 Superanalog of resolvent for exponential and idempotent superoperators
For resolvents $`𝖱_P\left(z\right)`$ and $`𝖱_T\left(z\right)`$ we use analog the standard formula from in the form
$`𝖱_P\left(z\right)`$ $`={\displaystyle \underset{0}{\overset{\mathrm{}}{}}}e^{zt}𝖯\left(t\right)𝑑t,`$ (82)
$`𝖱_T\left(z\right)`$ $`={\displaystyle \underset{0}{\overset{\mathrm{}}{}}}e^{zt}𝖳\left(t\right)𝑑t.`$ (83)
Using the supermatrix representation (19) we obtain
$`𝖱_P\left(z\right)`$ $`=\left(\begin{array}{cc}0& {\displaystyle \frac{\alpha }{z^2}}\\ {\displaystyle \frac{\alpha }{z}}& {\displaystyle \frac{1}{z}}\end{array}\right),`$ (86)
$`𝖱_T\left(z\right)`$ $`=\left(\begin{array}{cc}{\displaystyle \frac{1}{z}}& {\displaystyle \frac{\alpha }{z^2}}\\ 0& {\displaystyle \frac{1}{z}}\end{array}\right).`$ (89)
We observe, that $`𝖱_T\left(z\right)`$ satisfies the standard resolvent relation
$$𝖱_T\left(z\right)𝖱_T\left(w\right)=\left(wz\right)𝖱_T\left(z\right)𝖱_T\left(w\right),$$
(90)
but its analog for $`𝖱_P\left(z\right)`$
$$𝖱_P\left(z\right)𝖱_P\left(w\right)=\left(wz\right)𝖱_P\left(z\right)𝖱_P\left(w\right)+\frac{wz}{zw^2}𝖠$$
(91)
has additional term proportional to the generator $`𝖠`$.
## 8 Properties of $`t`$-linear idempotent operators
Here we consider properties of general $`t`$-linear (super)operators of the form
$$𝖪\left(t\right)=𝖪_0+𝖪_1t,$$
(92)
where $`𝖪_0=𝖪\left(0\right)`$ and $`𝖪_1=𝖪^{}\left(0\right)`$ are constant (super)operators represented by $`\left(n\times n\right)`$ matrices or $`(p+q)\times `$.$`\left(p+q\right)`$ supermatrices with $`t`$ (“time”) independent entries. Obviously, that the generator of a general $`t`$-linear (super)operator is
$$𝖠_K=𝖪^{}\left(0\right)=𝖪_1.$$
(93)
We will find system of equations for $`𝖪_0`$ and $`𝖪_1`$ for some special cases appeared in above consideration.
###### Assertion 35.
If a $`t`$-linear (super)operator $`𝖪\left(t\right)`$ satisfies the band equation (33)
$$𝖪\left(t\right)𝖪\left(s\right)=𝖪\left(t\right),$$
(94)
then it is idempotent and the constant component (super)operators $`𝖪_0`$ and $`𝖪_1`$ satisfy the system of equations
$`𝖪_0^2`$ $`=𝖪_0,`$ (95)
$`𝖪_1^2`$ $`=𝖹,`$ (96)
$`𝖪_1𝖪_0`$ $`=𝖪_1,`$ (97)
$`𝖪_0𝖪_1`$ $`=𝖹,`$ (98)
from which it follows, that $`𝖪_0`$ is idempotent, $`𝖪_1`$ is nilpotent, and $`𝖪_1`$ is right divisor of zero and left zero for $`𝖪_0`$.
For non-supersymmetric operators we have
###### Corollary 36.
The components of $`t`$-linear operator $`𝖪\left(t\right)`$ have the following properties: idempotent matrix $`𝖪_0`$ is similar to an upper triangular matrix with $`1`$ on the main diagonal and nilpotent matrix $`𝖪_1`$ is similar to an upper triangular matrix with $`0`$ on the main diagonal .
Comparing with the previous particular super case (42) we have $`𝖪_0=𝖯\left(0\right)`$ and $`𝖪_1=𝖠=𝖯^{}\left(0\right)`$.
###### Remark 37.
In case of $`\left(p+q\right)\times \left(p+q\right)`$ supermatrices the triangularization properties of Corollary 36 do not hold valid due to presence divisors of zero and nilpotents among entries (see Corollary 1), and so the inner structure of the component supermatrices satisfying (95)-(98) can be much different from the standard non-supersymmetric case .
Let us consider the structure of $`t`$-linear operator $`𝖪\left(t\right)`$ satisfying the generalized functional equation (73).
###### Assertion 38.
If a $`t`$-linear (super)operator $`𝖪\left(t\right)`$ satisfies the generalized functional equation
$$𝖪\left(t+s\right)=𝖪\left(t\right)𝖪\left(s\right)+𝖪^{}\left(t\right)s,$$
(99)
then its component (super)operators $`𝖪_0`$ and $`𝖪_1`$ satisfy the system of equations
$`𝖪_0^2`$ $`=𝖪_0,`$ (100)
$`𝖪_1^2`$ $`=𝖹,`$ (101)
$`𝖪_1𝖪_0`$ $`=𝖪_1,`$ (102)
$`𝖪_0𝖪_1`$ $`=𝖹,`$ (103)
We observe that the systems (95)-(98) and (100)-(103) are fully identical. It is important to observe the connection of the above properties with the differential equation for $`t`$-linear (super)operator $`𝖪\left(t\right)`$
$$𝖪^{}\left(t\right)=𝖠_K𝖪\left(t\right).$$
(104)
Using (93) we obtain the equation for components
$`𝖪_1^2`$ $`=𝖹,`$ (105)
$`𝖪_1𝖪_0`$ $`=𝖪_1.`$ (106)
That leads to the following
###### Theorem 39.
For any $`t`$-linear (super)operator $`𝖪\left(t\right)=`$ $`𝖪_0+𝖪_1t`$ the next statements are equivalent:
1. $`𝖪\left(t\right)`$ is idempotent and satisfies the band equation (94);
2. $`𝖪\left(t\right)`$ satisfies the generalized functional equation (99);
3. $`𝖪\left(t\right)`$ satisfies the differential equation (104) and has idempotent time independent part $`𝖪_0^2=𝖪_0`$ which is orthogonal to its generator $`𝖪_0𝖠=𝖹`$.
## 9 General $`t`$-power-type idempotent operators
Let us consider idempotent (super)operators which depend from time by power-type function, and so they have the form
$$𝖪\left(t\right)=\underset{m=0}{\overset{n}{}}𝖪_mt^m,$$
(107)
where $`𝖪_m`$ are $`t`$-independent (super)operators represented by $`\left(n\times n\right)`$ matrices or $`(p+q)\times `$.$`\left(p+q\right)`$ supermatrices. This power-type dependence of is very much important for super case, when supermatrix elements take value in Grassmann algebra, and therefore can be nilpotent (see (1)–(2) and Corollary 1).
We now start from the band property $`𝖪\left(t\right)𝖪\left(s\right)=𝖪\left(t\right)`$ and then find analogs of the functional equation and differential equation for them. Expanding the band property (94) in component we obtain $`n`$-dimensional analog of (95)-(98) as
$`𝖪_0^2`$ $`=𝖪_0,`$ (108)
$`𝖪_i^2`$ $`=𝖹,\mathrm{\hspace{0.33em}1}in,`$ (109)
$`𝖪_i𝖪_0`$ $`=𝖪_i,\mathrm{\hspace{0.33em}1}in,`$ (110)
$`𝖪_0𝖪_i`$ $`=𝖹,\mathrm{\hspace{0.33em}1}in,`$ (111)
$`𝖪_i𝖪_j`$ $`=𝖹,\mathrm{\hspace{0.33em}1}i,jn,ij.`$ (112)
###### Proposition 40.
The $`n`$-generalized functional equation for any $`t`$-power-type idempotent (super)operators (107) has the form
$$𝖪\left(t+s\right)=𝖪\left(t\right)𝖪\left(s\right)+𝖭_n(t,s),$$
(113)
where
$$𝖭_n(t,s)=\underset{m=1}{\overset{n}{}}\underset{l=m}{\overset{n}{}}𝖪_l\frac{l\left(l1\right)\mathrm{}\left(lm+1\right)}{m!}s^mt^{lm}.$$
(114)
###### Proof.
For the difference using the band property (94) we have $`𝖭_n(t,s)=𝖪\left(t+s\right)𝖪\left(t\right)𝖪\left(s\right)=𝖪\left(t+s\right)𝖪\left(t\right)`$. Then we expand in Taylor series around $`t`$ and obtain $`𝖭_n(t,s)=\underset{m=1}{\overset{n}{}}𝖪^{\left(m\right)}\left(t\right){\displaystyle \frac{s^m}{m!}}`$, where $`𝖪^{\left(m\right)}\left(t\right)`$ denotes $`n`$-th derivative which is a finite series for the power-type $`𝖪\left(t\right)`$ (107). ∎
The differential equation for idempotent (super)operators coincide with the standard initial value problem only for $`t`$-linear operators. In case of the power-type operators (107) we have
###### Proposition 41.
The $`n`$-generalized differential equation for any $`t`$-power-type idempotent (super)operators (107) has the form
$$𝖪^{}\left(t\right)=𝖠_K𝖪\left(t\right)+𝖴_n\left(t\right),$$
(115)
where
$$𝖴_n\left(t\right)=\{\begin{array}{cc}0\hfill & n=1\hfill \\ \underset{m=2}{\overset{n}{}}m𝖪_mt^{m1}\hfill & n2\hfill \end{array}.$$
(116)
###### Proof.
To find the difference $`𝖴_n\left(t\right)`$ we use the expansion (107) and the band conditions for components (108)–(112). ∎
## 10 Conclusion
In general one-parametric semigroups and corresponding superoperator semigroups represented by antitriangle idempotent supermatrices and their generalization for any dimensions $`p,q,m,n`$ have many unusual and nontrivial properties . Here we considered only some of them related to their connection with functional and differential equations. It would be interesting to generalize the above constructions to higher dimensions and to study continuity properties of the introduced idempotent superoperators. These questions will be investigated elsewhere.
###### Acknowledgement.
The author is grateful to Jan Okniński for valuable remarks and kind hospitality at the Institute of Mathematics, Warsaw University, where this work was begun. Also fruitful discussions with W. Dudek, A. Kelarev, G. Kourinnoy, W. Marcinek and B. V. Novikov are greatly acknowledged.
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# Verifying Termination of General Logic Programs with Concrete Queries
## 1 Introduction
For a program in any computer language, in addition to having to be logically correct, it should be terminating. Due to the recursive nature of logic programming, however, a logic program may more likely be non-terminating than a procedural program. Termination of logic programs then becomes one of the most important topics in logic programming research. Because the problem is extremely hard (undecidable in general), it has been considered as a never-ending story; see for a comprehensive survey.
The goal of termination analysis is to establish a characterization of termination of a logic program and design algorithms for automatic verification. A lot of methods for termination analysis have been proposed in the last decade (e.g., see ). A majority of these existing methods are the norm- or level mapping-based approaches, which consist of inferring mode/type information, inferring norms/level mappings, inferring models/interargument relations, and verifying some well-founded conditions (constraints). For example, Ullman and Van Gelder and Plümer focused on establishing a decrease in term size of some recursive calls based on interargument relations; Apt, Bezem and Pedreschi , and Bossi, Cocco and Fabris provided characterizations of Prolog left-termination based on level mappings/norms and models; Verschaetse , Decorte, De Schreye and Fabris , and Martin, King and Soper exploited inferring norms/level mappings from mode and type information; De Schreye and Verschaetse , Brodsky and Sagiv , and Lindenstrauss and Sagiv discussed automatic inference of interargument/size relations; De Schreye, Verschaetse and Bruynooghe , and Mesnard addressed automatic verification of the well-founded constraints. Very recently, Decorte, De Schreye and Vandecasteele presented an elegant unified termination analysis that integrates all the above components to produce a set of constraints that, when solved, yields a termination proof.
It is easy to see that the above methods have among others the following features.
1. They are compile-time approaches in the sense that they make termination analysis only relying on some static information about the structure (of the source code) of a logic program, such as modes/types, norms (i.e. term sizes of atoms of clauses)/level mappings, models/interargument relations, and the like.
2. They are suitable for termination analysis with respect to (abstract) query patterns . A query pattern defines a class of concrete queries,<sup>1</sup><sup>1</sup>1The difference between an abstract query pattern and a concrete query is similar to that between a class and an object in object-oriented programming languages. such as ground queries, bounded queries, well-moded queries, etc.
Our observation shows that some dynamic information about the structure of a concrete infinite SLDNF-derivation, such as repetition of selected subgoals and clauses and recursive increase in term size, plays a crucial role in characterizing the termination. Such dynamic features are hard to capture unless we evaluate some related concrete queries. This suggests that methods of extracting and utilizing dynamic features for termination analysis should be exploited.
Another observation comes from real programming practices. Consider the following situation: Given a logic program $`P`$ and a query pattern $`Q`$, applying a termination analysis yields a conclusion that $`P`$ is not terminating w.r.t. $`Q`$. In most cases, this means that there are a handful of concrete queries of the pattern $`Q`$ evaluating which may lead to infinite SLDNF-derivations. In order to improve the program, users (programmers) most often want to figure out how the non-termination happens by posing a few typical concrete queries and evaluating them step by step while determining which derivations would most likely extend to infinite ones. Such a debugging process is both quite time consuming and tricky. Doing it automatically is of great significance. Obviously, the above mentioned termination analysis techniques cannot help with such job. This suggests that methods of termination analysis for concrete queries should be developed.
The above two observations motivated the research of this paper. In this paper, we introduce an effective method for termination analysis w.r.t. concrete queries. The basic idea is as follows: First, since non-termination is caused by an infinite (generalized) SLDNF-derivation, we directly make use of some essential structural characteristics of an infinite derivation (such as variants, expanded variants, etc.) to characterize the termination. Then, given a logic program and a set of concrete queries, we evaluate these queries while dynamically collecting and applying certain structural features to predict (based on the characterization) if we are on the track to an infinite derivation. Such a process of query evaluation is guaranteed to terminate by a necessary condition of an infinite derivation. Finally, we provide the user with either an answer Yes, meaning that the logic program is terminating w.r.t. the given set of queries, or a finite (generalized) SLDNF-derivation that would most likely lead to an infinite derivation. In the latter case, the user can improve the program following the guidance of the informative derivation.
Although the termination problem is undecidable in general, our method works effectively for a vast majority of general logic programs with non-floundering queries. In fact, the methodology used in this paper is partly borrowed from loop checking $``$ another research topic in logic programming, which focuses on detecting and eliminating infinite loops in SLD-trees (e.g., see ). Therefore, our work bridges termination analysis with loop checking, the two problems which have been studied separately in the past despite their close technical relation with each other .
The plan of the paper is as follows. In Section 2, we introduce a notion of a generalized SLDNF-tree, which is the basis of our method. Roughly speaking, a generalized SLDNF-tree is a set of SLDNF-trees augmented with an ancestor-descendant relation on their subgoals. In Section 3, we prove a necessary and sufficient condition for an infinite generalized SLDNF-derivation. In Section 4.1, we formally define the notion of termination, which is slightly different from that of De Schreye and Decorte . In Section 4.2, we develop an algorithm for automatically verifying termination of a general logic program with concrete queries and prove its properties. We will use some representative logic programs to illustrate the effectiveness of the algorithm. In Section 5, we mention some related work on termination analysis and on loop checking. We end in Section 6 with some concluding remarks and further work.
### 1.1 Preliminary
We present our notation and review some standard terminology of logic programs as described in .
Variables begin with a capital letter, and predicate, function and constant symbols with a lower case letter. Let $`A`$ be an atom/function. The size of $`A`$, denoted $`|A|`$, is the count of function symbols, variables and constants in $`A`$. We use $`rel(A)`$ to refer to the predicate/function symbol of $`A`$, and use $`A[i]`$ to refer to the $`i`$-th argument of $`A`$, $`A[i][j]`$ to refer to the $`j`$-th argument of the $`i`$-th argument, and so on. Let $`S`$ be a set or a list. We use $`|S|`$ to denote the number of elements in $`S`$.
###### Definition 1.1
Let $`A`$ be an atom with the list $`[X_1,\mathrm{},X_m]`$ of distinct variables. By variable renaming on $`A`$ we mean to substitute the variables of $`A`$ with another list $`[Y_1,\mathrm{},Y_m]`$ of distinct variables. Two atoms $`A`$ and $`B`$ are said to be variants if after variable renaming (on $`A`$ or $`B`$) they become the same.
For instance, let $`A=p(a,X,Y,X)`$ and $`B=p(a,Z,Y,Z)`$. By substituting $`[X,Y]`$ for $`[Z,Y]`$, $`B`$ becomes the same as $`A`$, so $`A`$ and $`B`$ are variants. However, $`A`$ and $`C=p(a,Z,Y,W)`$ are not variants because there is no variable substitution that makes them the same. Note that any atom $`A`$ is a variant of itself.
###### Definition 1.2
A (general) logic program is a finite set of clauses of the form
$`AL_1,\mathrm{},L_n`$
where $`A`$ is an atom and $`L_i`$s are literals. $`A`$ is called the head and $`L_1,\mathrm{},L_n`$ is called the body of the clause. If a general logic program has no clause with negative literals in its body, it is called a positive program.
###### Definition 1.3
A goal is a headless clause $`L_1,\mathrm{},L_n`$ where each literal $`L_i`$ is called a subgoal. $`L_1,\mathrm{},L_n`$ is called a (concrete) query. When $`n=0`$, the “$``$” symbol is omitted.
The initial goal, $`G_0=L_1,\mathrm{},L_n`$, is called a top goal. Without loss of generality, we shall assume throughout the paper that a top goal consists only of one atom (i.e. $`n=1`$ and $`L_1`$ is a positive literal).
###### Definition 1.4
A control strategy consists of two rules: one rule for selecting one goal from among a set of goals, and one rule for selecting one subgoal from the selected goal.
The second rule in a control strategy is usually called a selection or computation rule in the literature. Throughout the paper we use a fixed depth-first, left-most control strategy (as in Prolog). So the selected subgoal in each goal is the left-most subgoal.
Trees are commonly used to represent the search space of a top-down proof procedure. For convenience, a node in such a tree is represented by $`N_i:G_i`$ where $`N_i`$ is the name of the node and $`G_i`$ is a goal labeling the node. Assume no two nodes have the same name. Therefore, we can refer to nodes by their names.
## 2 Generalized SLDNF-Trees
Non-termination of general logic programs results from infinite derivations. In order to characterize infinite derivations more precisely, in this section we extend the standard SLDNF-trees to include some new features.
To characterize an infinite derivation we need first to define the ancestor-descendant relation on its selected subgoals. Informally, $`A`$ is an ancestor subgoal of $`B`$ if the proof of $`A`$ needs (or in other words goes via) the proof of $`B`$. For example, let $`M:A,A_1,\mathrm{},A_m`$ be a node in an SLDNF-tree, and $`N:B_1,\mathrm{},B_n,A_1,\mathrm{},A_m`$ be a child node of $`M`$ that is generated by resolving $`M`$ on the subgoal $`A`$ with a clause $`AB_1,\mathrm{},B_n`$. Then $`A`$ at $`M`$ is an ancestor subgoal of all $`B_i`$s at $`N`$. However, such relationship does not exist between $`A`$ at $`M`$ and any $`A_j`$ at $`N`$. It is easily seen that all $`B_i`$s at $`N`$ inherit the ancestor subgoals of $`A`$ at $`M`$.
The ancestor-descendant relation can be explicitly expressed using an ancestor list introduced in , which is a set of pairs $`(Node,Atom)`$ where $`Node`$ is the name of a node and $`Atom`$ is the selected subgoal at $`Node`$. The ancestor list of a subgoal $`L_j`$ is $`AL_{L_j}=\{(N_1,A_1),\mathrm{},(N_k,A_k)\}`$, showing that $`A_1`$ at node $`N_1`$, …, and $`A_k`$ at node $`N_k`$ are all ancestor subgoals of $`L_j`$. For instance, in the above example, if the ancestor list of the subgoal $`A`$ at node $`M`$ is $`AL_A`$ then the ancestor list of each $`B_i`$ at node $`N`$ is $`AL_{B_i}=\{(M,A)\}AL_A`$.
Augmenting SLDNF-trees with ancestor lists leads to the following definition of SLDNF-trees.
###### Definition 2.1 (SLDNF-trees)
Let $`P`$ be a general logic program, $`G_0=A_0`$ a top goal, and $`R`$ a depth-first, left-most control strategy. The SLDNF-tree $`T_{G_0}`$ for $`P\{G_0\}`$ via $`R`$ is defined as follows.
1. The root node is $`N_0:G_0`$ with the ancestor list $`AL_{A_0}=\{\}`$ for $`A_0`$.
2. Let $`N_i:L_1,\mathrm{},L_m`$ be a node in the tree selected by $`R`$. If $`m=0`$ then $`N_i`$ is a success leaf, marked by $`\mathrm{}_t`$. Otherwise, we distinguish between the following two cases.
1. If $`L_1`$ is a positive literal, then for each clause $`BB_1,\mathrm{},B_n`$ such that $`L_1`$ and $`B`$ are unifiable, $`N_i`$ has a child node
$`N_s:(B_1,\mathrm{},B_n,L_2,\mathrm{},L_m)\theta `$
where $`\theta `$ is an mgu (i.e. most general unifier) of $`L_1`$ and $`B`$, the ancestor list for each $`B_k\theta `$ is $`AL_{B_k\theta }=\{(N_i,L_1)\}AL_{L_1}`$, and the ancestor list for each $`L_k\theta `$ is $`AL_{L_k\theta }=AL_{L_k}`$. If there exists no clause whose head can unify with $`L_1`$ then $`N_i`$ has a single child node $``$ a failure leaf, marked by $`\mathrm{}_f`$.
2. If $`L_1=\neg A`$ is a ground negative literal, then build a partial SLDNF-tree $`T_A`$ for $`P\{A\}`$ via $`R`$ where $`A`$ inherits the ancestor list of $`L_1`$, until the first success leaf is generated. If $`T_A`$ has a success leaf then $`N_i`$ has a single child node $``$ a failure leaf, $`\mathrm{}_f`$. Otherwise, if all branches of $`T_A`$ end with a failure leaf then $`N_i`$ has a single child node
$`N_s:L_2,\mathrm{},L_m`$
where all $`L_k`$ inherit the ancestor lists of $`L_k`$ at node $`N_i`$.
Note that in this paper we do not discuss floundering $``$ a situation where a non-ground negative subgoal is selected by $`R`$ (see for discussion on such topic). In contrast to SLDNF-trees, an SLDNF-tree has the following two new features.
1. An ancestor list $`AL_{L_j}`$ is attached to each subgoal $`L_j`$. In particular, subgoals of a subsidiary SLTNF-tree $`T_A`$ built for solving a subgoal $`L_1=\neg A`$ inherit the ancestor list of $`L_1`$ (see item 2b). This is especially useful in identifying infinite derivations across SLTNF-trees (see Example 2.1). Note that a negative subgoal will never be an ancestor subgoal.
2. To handle a ground negative subgoal $`L_1=\neg A`$, only a partial subsidiary SLTNF-tree $`T_A`$ is generated by stopping at the first success leaf (see item 2b). The reason for this is that it is totally unnecessary to exhaust the remaining branches of $`T_A`$ because they would have no new answer for $`A`$. This can not only improve the efficiency of query evaluation, but also avoid some possible infinite derivations (see Example 2.2). In fact, Prolog achieves this by using cuts to skip the remaining branches of $`T_A`$ (e.g. see SICStus Prolog ).
For convenience, we use dotted edges $`\mathrm{`}\mathrm{`}\mathrm{}^{\prime \prime }`$ to connect parent and child SLDNF-trees, so that infinite derivations across SLDNF-trees can be clearly identified. Moreover, we refer to $`T_{G_0}`$, the top SLDNF-tree, along with all its descendant SLDNF-trees as a generalized SLDNF-tree for $`P\{G_0\}`$, denoted $`GT_{G_0}`$. Therefore, a path of a generalized SLDNF-tree may come across several SLDNF-trees through dotted edges. Any such a path starting at the root node $`N_0:G_0`$ is called a generalized SLDNF-derivation. A generalized SLDNF-derivation is successful (resp. failed) if it ends at a success leaf (resp. at a failure leaf).
Thus, there may occur two types of edges in a generalized SLDNF-tree, “$`\stackrel{C}{}`$” and $`\mathrm{`}\mathrm{`}\mathrm{}^{\prime \prime }`$. For convenience, we use $`\mathrm{`}\mathrm{`}^{\prime \prime }`$ to refer to either of them. We also use $`N_i:G_i\stackrel{C_1}{}\mathrm{}\stackrel{C_m}{}N_k:G_k`$ to represent a segment of a generalized SLDNF-derivation, which generates $`N_k:G_k`$ from $`N_i:G_i`$ by applying the set of clauses $`\{C_1,\mathrm{}C_m\}`$. Moreover, for any node $`N_i:G_i`$ we use $`L_i^1`$ to refer to the selected (i.e. left-most) subgoal in $`G_i`$.
###### Example 2.1
Let $`P_1`$ be a general logic program and $`G_0`$ a top goal, given by
| | $`P_1:`$ | $`p(X)\neg p(f(X))`$. $`C_{p_1}`$ |
| --- | --- | --- |
| | $`G_0:`$ | $`p(a).`$ |
The generalized SLDNF-tree $`GT_{p(a)}`$ for $`P_1\{G_0\}`$ is shown in Figure 1, where $`\mathrm{}`$ represents an infinite extension. We see that $`GT_{p(a)}`$ consists of one infinite generalized SLDNF-derivation.
###### Example 2.2
Consider the following general logic program and top goal.
| | $`P_2:`$ | $`p\neg q`$. $`C_{p_1}`$ |
| --- | --- | --- |
| | | $`q`$. $`C_{q_1}`$ |
| | | $`qq`$. $`C_{q_2}`$ |
| | $`G_0:`$ | $`p.`$ |
The generalized SLDNF-tree $`GT_p`$ for $`P_2\{G_0\}`$ is depicted in Figure 2 (a). For the purpose of comparison, the SLDNF-trees for $`P_2\{p\}`$ are shown in Figure 2 (b). Note that Figure 2 (a) is finite, whereas Figure 2 (b) is not.
We now formally define the ancestor-descendant relation.
###### Definition 2.2
Let $`N_i:G_i`$ and $`N_k:G_k`$ be two nodes in a generalized SLDNF-derivation, and $`A`$ and $`B`$ be the selected subgoals in $`G_i`$ and $`G_k`$, respectively. We say that $`A`$ is an ancestor subgoal of $`B`$, denoted $`A_{ANC}B`$, if $`A`$ is in the ancestor list $`AL_B`$ of $`B`$. When $`A`$ is an ancestor subgoal of $`B`$, we refer to $`B`$ as a descendant subgoal of $`A`$, $`N_i`$ as an ancestor node of $`N_k`$, and $`N_k`$ as a descendant node of $`N_i`$.
## 3 Characterizing an Infinite Generalized SLDNF-Derivation
In this section we establish a necessary and sufficient condition for an infinite generalized SLDNF-derivation.
In , a concept of expanded variants is introduced, which captures some key structural characteristics of certain subgoals in an infinite SLD-derivation. We observe that it applies to general logic programs as well. That is, infinite generalized SLDNF-derivations can be characterized based on expanded variants.
###### Definition 3.1
Let $`A`$ and $`A^{}`$ be two atoms or functions. $`A^{}`$ is said to be an expanded variant of $`A`$, denoted $`A^{}_{EV}A`$, if after variable renaming on $`A^{}`$ it becomes $`B`$ that is the same as $`A`$ except that there may be some terms at certain positions in $`A`$ each $`A[i]\mathrm{}[k]`$ of which grows in $`B`$ into a function $`B[i]\mathrm{}[k]=f(\mathrm{},A[i]\mathrm{}[k],\mathrm{})`$. Such terms like $`A[i]\mathrm{}[k]`$ in $`A`$ are then called growing terms w.r.t. $`A^{}`$.
As an illustration, let $`A=p(X,g(X))`$ and $`A^{}=p(Y,g(h(Y)))`$. By renaming $`Y`$ with $`X`$, $`A^{}`$ becomes $`B=p(X,g(h(X)))`$, which is the same as $`A`$ except that $`B[2][1]=h(A[2][1])`$. Therefore, $`A^{}`$ is an expanded variant of $`A`$ with a growing term $`A[2][1]`$. Here are a few more examples: $`p(X,Y)_{EV}p(Z,W)`$, $`p(f(a))_{EV}p(a)`$, $`p(g(a),f(g(h(X))))_{EV}p(a,f(h(Y)))`$, and $`p([X_1,X_2,X_3])_{EV}p([X_1,X_4])`$ (note that $`[X_1,X_2,X_3]=[X_1|[X_2,X_3]]`$).
It is immediate from Definition 3.1 that variants are expanded variants with the same size.
###### Theorem 3.1
Let $`D`$ be an infinite generalized SLDNF-derivation without infinitely large subgoals. Then there are infinitely many goals $`G_{g_1},G_{g_2},\mathrm{}`$ in $`D`$ such that for any $`j1`$, $`L_{g_j}^1_{ANC}L_{g_{j+1}}^1`$ and $`L_{g_j}^1`$ and $`L_{g_{j+1}}^1`$ are variants.
Proof. Let $`D`$ be of the form
For each derivation step $`N_i:G_i\stackrel{C}{}N_{i+1}:G_{i+1}`$, where $`L_i^1`$ is a positive subgoal and $`C=AB_1,\mathrm{}B_n`$ such that $`A\theta =L_i^1\theta `$ under an mgu $`\theta `$, we do the following:
1. If $`n=0`$, which means $`L_i^1`$ is proved at this step, mark node $`N_i`$ with #.
2. Otherwise, the proof of $`L_i^1`$ needs the proof of $`B_j\theta `$ $`(j=1,\mathrm{},n)`$. If all descendant nodes of $`N_i`$ in $`D`$ have been marked with #, which means that all $`B_j\theta `$ have been proved at some steps in $`D`$, mark node $`N_i`$ with #.
Note that the root node $`N_0`$ will never be marked by #, for otherwise $`G_0`$ would have been proved and $`D`$ should have ended at a success leaf. After the above marking process, let $`D`$ become
where all nodes except $`N_0,N_{i_1},N_{i_2},\mathrm{},N_{i_k},\mathrm{}`$ are marked with #. Since we use the depth-first, left-most control strategy, for any $`j0`$ the proof of $`L_{i_j}^1`$ needs the proof of $`L_{i_{j+1}}^1`$ (let $`i_0=0`$), for otherwise $`N_{i_j}`$ would have been marked with #. That is, $`L_{i_j}^1`$ is an ancestor subgoal of $`L_{i_{j+1}}^1`$. Moreover, $`D`$ must contain an infinite number of such nodes because if $`N_{i_k}:G_{i_k}`$ was the last one, which means that all nodes after $`N_{i_k}`$ were marked with #, then $`L_{i_k}^1`$ would be proved, so that $`N_{i_k}`$ should be marked with #, a contradiction.
The above proof shows that $`D`$ has an infinite number of selected subgoals $`L_{i_1}^1,L_{i_2}^1,\mathrm{}`$ such that $`L_{i_j}^1_{ANC}L_{i_{j+1}}^1`$ $`(j1)`$. Since all subgoals in $`D`$ are bounded in size, and any general logic program has only a finite number of clauses and predicate, function and constant symbols, there must be an infinite number of subgoals $`L_{g_1}^1,L_{g_2}^1,\mathrm{}`$ among the $`L_{i_j}^1`$s that are variants. This concludes the proof.
###### Theorem 3.2
Let $`D`$ be an infinite generalized SLDNF-derivation with infinitely large subgoals. Then there are infinitely many goals $`G_{g_1},G_{g_2},\mathrm{}`$ in $`D`$ such that for any $`j1`$, $`L_{g_j}^1_{ANC}L_{g_{j+1}}^1`$ and $`L_{g_{j+1}}^1_{EV}L_{g_j}^1`$ with $`|L_{g_{j+1}}^1|>|L_{g_j}^1|`$.
The following lemma is required to prove this theorem.
###### Lemma 3.3
Let $`S=\{C_1,\mathrm{},C_n\}`$ be a finite set of clauses. Let $`D`$ be an infinite generalized SLDNF-derivation of the form
where for any $`j1`$, $`\{C_{j_1},\mathrm{},C_{j_{n_j}}\}=S`$, $`L_{i_j}^1_{ANC}L_{i_{j+1}}^1`$, and $`L_{i_j}^1`$ is a variant of $`L_{i_{j+1}}^1`$ except for a few terms (at least one) at certain positions in $`L_{i_{j+1}}^1`$ that increase in size w.r.t. $`L_{i_j}^1`$. Then there are an infinite sequence of subgoals $`L_{g_1}^1,L_{g_2}^1\mathrm{}`$ among the $`L_{i_j}^1`$s such that for any $`k1`$, $`L_{g_{k+1}}^1_{EV}L_{g_k}^1`$ with $`|L_{g_{k+1}}^1|>|L_{g_k}^1|`$.
Proof. Since $`L_{i_{j+1}}^1`$ being an expanded variant of $`L_{i_j}^1`$ with $`|L_{i_{j+1}}^1|>|L_{i_j}^1|`$ is determined by those arguments of $`L_{i_j}^1`$ whose size increases w.r.t. $`L_{i_{j+1}}^1`$, for simplicity of presentation we ignore the remaining arguments of $`L_{i_j}^1`$ that are variants of the corresponding ones in $`L_{i_{j+1}}^1`$. Since the $`L_{i_j}^1`$s are generated by repeatedly applying the same set $`S`$ of clauses, the increase in their term size must be made in a fixed, regular way (assume that our programs contain no built-in’s such as $`assert(.)`$ and $`retract(.)`$). In order to facilitate the analysis of such regular increase in term size, with no loss in generality, let each $`L_{i_j}^1`$ be of the form $`p([T_1^j,\mathrm{},T_{m_j}^j])`$, which contains a single argument that is a list, such that the list of $`L_{i_{j+1}}^1`$ is derived from the list $`L=[T_1^j,\mathrm{},T_{m_j}^j]`$ of $`L_{i_j}^1`$ via a DELETE-ADD-SHUFFLE process which consists of deleting the first $`0n^{}m_j`$ elements from $`L`$, then adding $`n^+n^{}`$ elements to the front of the remaining part of $`L`$, and finally shuffling the list of elements obtained (in a fixed way).<sup>2</sup><sup>2</sup>2See Example 4.4 for an illustration.
Let $`\stackrel{}{X}`$ represent a sequence of elements, such as $`X_1,X_2,\mathrm{},X_m`$. We distinguish the following two cases based on the ADD operation.
1. The $`n^+`$ added elements are independent of the list $`L`$ of elements of $`L_{i_j}^1`$. Since the same set $`S`$ of clauses is applied, the set $`E_1`$ of elements added to $`L_{i_2}^1`$ from $`L_{i_1}^1`$ must be the same as or a variant of the set added to $`L_{i_3}^1`$ from $`L_{i_2}^1`$ that must be the same as or a variant of the set added to $`L_{i_4}^1`$ from $`L_{i_3}^1`$, and so on. Let $`E_2\{T_1^1,\mathrm{},T_{m_1}^1\}`$ be such that each $`TE_2`$ occurs in the derivation $`D`$ infinite times. That is, no $`TE_2`$ will be removed by the DELETE operation. (But all $`T\{T_1^1,\mathrm{},T_{m_1}^1\}E_2`$ will be deleted at some derivation steps in $`D`$ by the DELETE operation.) Then there must be an infinite sequence $`L_{f_1}^1,L_{f_2}^1,\mathrm{}`$ among the $`L_{i_j}^1`$s such that for any $`k1`$, all elements in $`L_{f_k}^1`$ come from $`E_1E_2`$ (or its variant). Since $`E_1E_2`$ contains only a finite number of elements, no matter what shuffling approach is used, there must be an infinite sequence $`L_{g_1}^1,L_{g_2}^1\mathrm{}`$ among the $`L_{f_j}^1`$s such that for any $`k1`$, let $`L_{g_k}^1=p([S_1,\mathrm{},S_n])`$, then after variable renaming $`L_{g_{k+1}}^1`$ will become like $`p([\stackrel{}{S_1},\mathrm{},\stackrel{}{S_n}])`$ such that each $`S_l`$ is in $`\stackrel{}{S_l}`$. Obviously, these $`S_l`$s in $`L_{g_k}^1`$ are growing terms w.r.t. $`L_{g_{k+1}}^1`$. That is, $`L_{g_{k+1}}^1_{EV}L_{g_k}^1`$ with $`|L_{g_{k+1}}^1|>|L_{g_k}^1|`$.
2. Some of the $`n^+`$ added elements depend on the list of $`L_{i_j}^1`$. For simplicity of presentation and without loss of generality assume that from $`L_{i_1}^1`$ to $`L_{i_2}^1`$ only one added element, say of the form $`f(A_1^1,g(A_2^1))`$, depends on the list $`[T_1^1,\mathrm{},T_{m_1}^1]`$ of $`L_{i_1}^1`$. Let $`E_1`$ be the set of elements added to $`L_{i_2}^1`$ that are independent of the list $`[T_1^1,\mathrm{},T_{m_1}^1]`$. That is, the set of added elements from $`L_{i_1}^1`$ to $`L_{i_2}^1`$ is $`E_1\{f(A_1^1,g(A_2^1))\}`$. Let $`E=E_1\{T_1^1,\mathrm{},T_{m_1}^1\}`$. Then $`E`$ can be considered to be the domain of the two arguments $`A_1^1`$, $`A_2^1`$ in $`f(A_1^1,g(A_2^1))`$, i.e. $`A_1^1,A_2^1E`$. Since all the $`L_{i_j}^1`$s in the derivation $`D`$ are generated recursively by applying the same set $`S`$ of clauses, for any $`j1`$ from $`L_{i_j}^1`$ to $`L_{i_{j+1}}^1`$ the set of added elements should be $`E_1\{f(A_1^j,g(A_2^j))\}`$ where $`A_1^j,A_2^jE\{f(A_1^k,g(A_2^k))|k<j\}`$. It is easy to see that any element of $`L_{i_j}^1`$ with an infinitely large size must be of the form $`f(A_1^{\mathrm{}},g(A_2^{\mathrm{}}))`$ where $`A_1^{\mathrm{}}`$ or $`A_2^{\mathrm{}}`$ or both are of the form $`f(A_1^{\mathrm{}},g(A_2^{\mathrm{}}))`$. We now consider the following two cases.
1. No $`L_{i_j}^1`$ in $`D`$ contains elements with an infinitely large size. Let $`N`$ be the largest size of an element $`f(\mathrm{\_},g(\mathrm{\_}))`$ in the $`L_{i_j}^1`$s and let $`E_2=\{f(A_1^k,g(A_2^k))|k1`$ such that $`|f(A_1^k,g(A_2^k))|N\}`$. Obviously $`E_2`$ is finite. So the elements of any $`L_{i_j}^1`$ in $`D`$ come from $`EE_2`$. Since $`EE_2`$ is finite, the number of the combinations of elements of $`EE_2`$ is finite. Since the $`L_{i_j}^1`$s in $`D`$ grow towards an infinitely large size, some combinations must repeat in $`D`$ infinite times. This suggests that no matter what shuffling approach is used, there must be an infinite sequence $`L_{g_1}^1,L_{g_2}^1\mathrm{}`$ among the $`L_{i_j}^1`$s such that for any $`k1`$, let $`L_{g_k}^1=p([S_1,\mathrm{},S_n])`$, then after variable renaming $`L_{g_{k+1}}^1`$ will become like $`p([\stackrel{}{S_1},\mathrm{},\stackrel{}{S_n}])`$ such that each $`S_l`$ is in $`\stackrel{}{S_l}`$. These $`S_l`$s in $`L_{g_k}^1`$ are growing terms w.r.t. $`L_{g_{k+1}}^1`$, thus $`L_{g_{k+1}}^1_{EV}L_{g_k}^1`$ with $`|L_{g_{k+1}}^1|>|L_{g_k}^1|`$.
2. As $`j\mathrm{}`$, some elements in $`L_{i_j}^1`$ grow towards an infinitely large size. Let $`T^{j+1}`$ be an element in the list of $`L_{i_{j+1}}^1`$ that (as $`j\mathrm{}`$) grows towards an element with an infinitely large size. Then there must be an element $`T^j`$ in the list of $`L_{i_j}^1`$ that grows towards an element with an infinitely large size via $`T^{j+1}`$ (otherwise, $`T^{j+1}`$ would not grow towards an infinitely large element since we apply the same set $`S`$ of clauses from $`L_{i_{j+1}}^1`$ to $`L_{i_{j+2}}^1`$ as from $`L_{i_j}^1`$ to $`L_{i_{j+1}}^1`$). This means that $`T^{j+1}`$ is the same as (or a variant of) $`T^j`$ or $`T^{j+1}=f(A_1^j,g(A_2^j))`$ such that $`T^j`$ (or its variant) is in $`A_1^j`$ or $`A_2^j`$. Obviously $`T^{j+1}`$ is an expanded variant of $`T^j`$. Generalizing such argument, for each infinitely large element $`T_l^{\mathrm{}}`$ in the list of $`L_i_{\mathrm{}}^1`$ we have an infinite sequence
with each $`T^j`$ in the list of $`L_{i_j}^1`$ such that $`T^1`$ grows towards $`T_l^{\mathrm{}}`$ via $`T^2`$ that grows towards $`T_l^{\mathrm{}}`$ via $`T^3`$, and so on. Obviously, for any $`k>j1`$ $`T^k`$ is an expanded variant of $`T^j`$. So in this case each $`T^j`$ is called a growing element w.r.t. $`T_l^{\mathrm{}}`$. Note that if there are more than one element in some $`L_{i_j}^1`$ that grow towards $`T_l^{\mathrm{}}`$ via $`T^{j+1}`$, such as in the case $`T^{j+1}=f(T_1^j,g(T_2^j))`$ with $`T_1^j,T_2^j`$ in $`L_{i_j}^1`$, only one $`T^j`$ of them is selected as the growing element w.r.t. $`T_l^{\mathrm{}}`$ based on the following criterion: $`T^j`$ must be an expanded variant of $`T^{j1}`$. If still more than one element meet such criterion, select an arbitrary one.
We now partition the list $`[T_1^j,\mathrm{},T_{m_j}^j]`$ of each $`L_{i_j}^1`$ into two parts: the sublist $`GE_{i_j}`$ of growing elements and the sublist $`NGE_{i_j}`$ of non-growing elements. That is, $`T[T_1^j,\mathrm{},T_{m_j}^j]`$ is in $`GE_{i_j}`$ if it is a growing element w.r.t. some $`T_l^{\mathrm{}}`$. Clearly, for each $`k>j1`$ $`|GE_{i_k}||GE_{i_j}|`$. Since for any $`V_1,\mathrm{},V_m`$ in $`GE_{i_j}`$ there are $`m`$ elements $`V_1^{},\mathrm{},V_m^{}`$ in $`GE_{i_k}`$ such that each $`V_l^{}`$ is an expanded variant of $`V_l`$, there must be an infinite sequence $`L_{f_1}^1,L_{f_2}^1,\mathrm{}`$ among the $`L_{i_j}^1`$s such that for any $`k>j1`$ $`GE_{f_k}`$ is an expanded variant of $`GE_{f_j}`$. That is, let $`GE_{f_j}=[V_1,\mathrm{},V_m]`$, then after variable renaming $`GE_{f_k}`$ becomes $`[\stackrel{}{V_1},\mathrm{},\stackrel{}{V_m}]`$ such that each $`\stackrel{}{V_l}`$ contains an element that is an expanded variant of $`V_l`$.
Now consider the elements of $`L_{f_j}^1`$s. Since all infinitely large elements have been covered by the growing elements, the size of any non-growing element is bounded by some constant, say $`N`$. Let $`E_2=\{f(A_1^k,g(A_2^k))|k1`$ such that $`|f(A_1^k,g(A_2^k))|N\}`$. So all non-growing elements of any $`L_{f_j}^1`$ come from $`EE_2`$. Since $`EE_2`$ is finite and the sequence $`L_{f_1}^1,L_{f_2}^1,\mathrm{}`$ is infinite, some combinations of elements of $`EE_2`$ must occur in infinitely many $`L_{f_j}^1`$s. This implies that no matter what shuffling approach is used, there must be an infinite sequence $`L_{g_1}^1,L_{g_2}^1\mathrm{}`$ among the $`L_{f_j}^1`$s such that for any $`k1`$, let $`L_{g_k}^1=p([S_1,\mathrm{},S_n])`$, then after variable renaming $`L_{g_{k+1}}^1`$ will become like $`p([\stackrel{}{S_1},\mathrm{},\stackrel{}{S_n}])`$ such that each $`\stackrel{}{S_l}`$ contains an element that is an expanded variant of $`S_l`$. That is, $`L_{g_{k+1}}^1_{EV}L_{g_k}^1`$ with $`|L_{g_{k+1}}^1|>|L_{g_k}^1|`$.
Proof of Theorem 3.2. By the proof of Theorem 3.1, $`D`$ contains an infinite number of selected subgoals $`L_1^1,L_2^1,\mathrm{}`$ such that $`L_j^1_{ANC}L_{j+1}^1`$ $`(j1)`$. Since any logic program has only a finite number of clauses, there must be a set of clauses in the program that are invoked an infinite number of times in $`D`$. Let $`S=\{C_1,\mathrm{},C_n\}`$ be the set of all different clauses that are used an infinite number of times in $`D`$. Then $`D`$ can be depicted as
where for any $`j1`$, $`L_{i_j}^1_{ANC}L_{i_{j+1}}^1`$ and $`\{C_{j_1},\mathrm{},C_{j_{n_j}}\}=S`$. Since any logic program has only a finite number of predicate, function and constant symbols and $`D`$ contains subgoals with infinitely large size, there must be an infinite sequence $`L_{f_1}^1,L_{f_2}^1\mathrm{}`$ among the $`L_{i_j}^1`$s such that for any $`l1`$, $`L_{f_l}^1`$ is a variant of $`L_{f_{l+1}}^1`$ except for a few terms in $`L_{f_{l+1}}^1`$ that increase in size. Hence by Lemma 3.3, there is an infinite sequence $`L_{g_1}^1,L_{g_2}^1\mathrm{}`$ among the $`L_{f_l}^1`$s such that for any $`k1`$ $`L_{g_{k+1}}^1_{EV}L_{g_k}^1`$ with $`|L_{g_{k+1}}^1|>|L_{g_k}^1|`$.
###### Theorem 3.4
$`D`$ is an infinite generalized SLDNF-derivation if and only if it is of the form
such that
1. For any $`j1`$, $`L_{g_j}^1_{ANC}L_{g_{j+1}}^1`$ and $`L_{g_{j+1}}^1_{EV}L_{g_j}^1`$.
2. For any $`j1`$ $`|L_{g_j}^1|=|L_{g_{j+1}}^1|`$, or for any $`j1`$ $`|L_{g_j}^1|<|L_{g_{j+1}}^1|`$.
3. For any $`j1`$, the set of clauses used to derive $`L_{g_{j+1}}^1`$ from $`L_{g_j}^1`$ is the same as that of deriving $`L_{g_{j+2}}^1`$ from $`L_{g_{j+1}}^1`$, i.e. $`\{C_{j_1},\mathrm{},C_{j_{n_j}}\}=\{C_{(j+1)_1},\mathrm{},C_{(j+1)_{n_{j+1}}}\}`$.
Proof. $`()`$ Straightforward.
$`()`$ By Theorems 3.1 and 3.2, $`D`$ is of the form
where for any $`j1`$, $`L_{i_j}^1_{ANC}L_{i_{j+1}}^1`$ and $`L_{i_{j+1}}^1_{EV}L_{i_j}^1`$. In particular, when all subgoals in $`D`$ are bounded in size, by Theorem 3.1 for any $`j1`$ $`|L_{i_j}^1|=|L_{i_{j+1}}^1|`$. Otherwise, by Theorem 3.2 for any $`j1`$ $`|L_{i_j}^1|<|L_{i_{j+1}}^1|`$.
Since any logic program has only a finite number of clauses, there must be a set $`S=\{C_1,\mathrm{},C_n\}`$ of clauses in the program that are invoked an infinite number of times in $`D`$. This means that there exists an infinite sequence of subgoals $`L_{g_1}^1,L_{g_2}^1,\mathrm{}`$ among the $`L_{i_j}^1`$s such that for any $`j1`$ $`L_{g_{j+1}}^1`$ is derived from $`L_{g_j}^1`$ by applying the set $`S`$ of clauses. That is, $`D`$ is of the form
such that the three conditions of this theorem hold.
Theorem 3.4 is the principal result of this paper. It captures two crucial characteristics of an infinite generalized SLDNF-derivation: repetition of selected subgoals and clauses, and recursive increase in term size. Repetition leads to variants, whereas recursive increase introduces growing terms. It is the characterization of these key (dynamic) features that allows us to design a mechanism for automatically testing termination of general logic programs.
## 4 Testing Termination of General Logic Programs
### 4.1 Definition of Termination
In , a generic definition of termination of logic programs is presented.
###### Definition 4.1
Let $`P`$ be a general logic program, $`S_Q`$ a set of queries and $`S_R`$ a set of selection rules. $`P`$ is terminating w.r.t. $`S_Q`$ and $`S_R`$ if for each query $`Q_i`$ in $`S_Q`$ and for each selection rule $`R_j`$ in $`S_R`$, all SLDNF-trees for $`P\{Q_i\}`$ via $`R_j`$ are finite.
Observe that the above definition considers finite SLDNF-trees for termination. That is, if $`P`$ is terminating w.r.t. $`Q_i`$ then all (complete) SLDNF-trees for $`P\{Q_i\}`$ must be finite. This does not seem to apply to Prolog where there exist cases in which, although $`P`$ is terminating w.r.t. $`Q_i`$ and $`R_j`$, some (complete) SLDNF-trees for $`P\{Q_i\}`$ are infinite. Example 2.2 gives such an illustration, where Prolog terminates with a positive answer.
In view of the above observation, we present the following definition based on a generalized SLDNF-tree.
###### Definition 4.2
Let $`P`$ be a general logic program, $`S_Q`$ a finite set of queries and $`R`$ a fixed depth-first, left-most control strategy. $`P`$ is terminating w.r.t. $`S_Q`$ and $`R`$ if for each query $`Q_i`$ in $`S_Q`$, the generalized SLDNF-tree for $`P\{Q_i\}`$ via $`R`$ is finite.
The above definition implies that $`P`$ is terminating w.r.t. $`S_Q`$ and $`R`$ if and only if there is no infinite generalized SLDNF-derivation in any generalized SLDNF-tree $`GT_{Q_i}`$. So the following result is immediate from Theorem 3.4.
###### Theorem 4.1
$`P`$ is terminating w.r.t. $`S_Q`$ and $`R`$ if and only if for each query $`Q_i`$ in $`S_Q`$ there is no infinite generalized SLDNF-derivation in $`GT_{Q_i}`$ of the form
that meets the three conditions of Theorem 3.4.
### 4.2 An Algorithm for Automatically Testing Termination
Theorem 4.1 provides a necessary and sufficient condition for termination of a general logic program. Obviously, such a condition cannot be directly used for automatic verification because it requires generating an infinite generalized SLDNF-derivation to see if the three conditions of Theorem 3.4 are satisfied.
As we mentioned before, the three conditions of Theorem 3.4 capture two most important structural features of an infinite generalized SLDNF-derivation. Therefore, we may well use these conditions to predict possible infinite generalized SLDNF-derivations based on some finite generalized SLDNF-derivations. Although the predictions may not always be guaranteed to be correct (since the termination problem is undecidable in general), it can be correct in a vast majority of cases. That is, if the three conditions of Theorem 3.4 are satisfied by some finite generalized SLDNF-derivation, the underlying general logic program is most likely non-terminating. This leads to the following definition.
###### Definition 4.3
Let $`P`$ be a general logic program, $`S_Q`$ a finite set of queries and $`R`$ a depth-first, left-most control strategy. Let $`d>1`$ be a depth bound. $`P`$ is said to be most-likely non-terminating w.r.t. $`S_Q`$ and $`R`$ if for some query $`Q_i`$ in $`S_Q`$, there is a generalized SLDNF-derivation of the form
| $`N_0:Q_i\mathrm{}`$ | $`N_{g_1}:G_{g_1}\stackrel{C_{1_1}}{}\mathrm{}\stackrel{C_{1_{n_1}}}{}`$ |
| --- | --- |
| | $`N_{g_2}:G_{g_2}\stackrel{C_{2_1}}{}\mathrm{}\stackrel{C_{2_{n_2}}}{}`$ |
| | |
| | $`N_{g_d}:G_{g_d}\stackrel{C_{d_1}}{}\mathrm{}\stackrel{C_{d_{n_d}}}{}N_{g_{d+1}}:G_{g_{d+1}}`$ |
such that
1. For any $`jd`$, $`L_{g_j}^1_{ANC}L_{g_{j+1}}^1`$ and $`L_{g_{j+1}}^1_{EV}L_{g_j}^1`$.
2. For any $`jd`$ $`|L_{g_j}^1|=|L_{g_{j+1}}^1|`$, or for any $`jd`$ $`|L_{g_j}^1|<|L_{g_{j+1}}^1|`$.
3. For any $`jd`$, the set of clauses used to derive $`L_{g_{j+1}}^1`$ from $`L_{g_j}^1`$ is the same as that of deriving $`L_{g_{j+2}}^1`$ from $`L_{g_{j+1}}^1`$, i.e. $`\{C_{j_1},\mathrm{},C_{j_{n_j}}\}=\{C_{(j+1)_1},\mathrm{},C_{(j+1)_{n_{j+1}}}\}`$.
###### Theorem 4.2
Let $`P`$, $`S_Q`$ and $`R`$ be as defined in Definition 4.3.
1. If $`P`$ is not terminating w.r.t. $`S_Q`$ and $`R`$ then it is most-likely non-terminating w.r.t. $`S_Q`$ and $`R`$.
2. If $`P`$ is not most-likely non-terminating w.r.t. $`S_Q`$ and $`R`$ then it is terminating w.r.t. $`S_Q`$ and $`R`$.
Proof: 1. If $`P`$ is not terminating w.r.t. $`S_Q`$ and $`R`$, by Definition 4.2 for some query $`Q_iS_Q`$ there exists an infinite generalized SLDNF-derivation in $`GT_{Q_i}`$. The result is then immediate from Theorem 3.4.
2. If $`P`$ is not most-likely non-terminating w.r.t. $`S_Q`$ and $`R`$ and, on the contrary, it is not terminating w.r.t. $`S_Q`$ and $`R`$, then by the first part of this theorem we reach a contradiction.
It is easily seen that the converse of the above theorem does not hold. The following algorithm is to determine most-likely non-termination.
###### Algorithm 4.1
Testing termination of a general logic program.
* Input: A general logic program $`P`$, a finite set of queries $`S_Q=\{Q_1,\mathrm{},Q_m\}`$, and a depth-first, left-most control strategy $`R`$.
* Output: Yes or a generalized SLDNF-derivation $`D`$.
* Method: Apply the following procedure.
| procedure $`Test(P,S_Q,R)`$ |
| --- |
| | begin |
| 1 | | For each query $`Q_iS_Q`$, construct the full generalized SLDNF-tree |
| | | $`GT_{Q_i}`$ for $`P\{Q_i\}`$ via $`R`$ unless a generalized SLDNF-derivation |
| | | $`D`$ is encountered that meets the three conditions of Definition 4.3, |
| | | in which case return $`D`$ and stop the procedure; |
| 2 | | Return Yes |
| | end |
###### Theorem 4.3
Algorithm 4.1 terminates. It returns Yes if and only if $`P`$ is not most-likely non-terminating w.r.t. $`S_Q`$ and $`R`$.
Proof: If for each query $`Q_iS_Q`$ the generalized SLDNF-tree $`GT_{Q_i}`$ for $`P\{Q_i\}`$ is finite, line 1 of Algorithm 4.1 will be completed in finite time, so that Algorithm 4.1 will terminate in finite time. Otherwise, let all generalized SLDNF-trees $`GT_{Q_i}`$ with $`i<m`$ be finite and $`GT_{Q_{i+1}}`$ be infinite. Let $`D`$ be the first infinite derivation in $`GT_{Q_{i+1}}`$. By Theorem 3.4, $`D`$ must be of the form
such that the three conditions of Theorem 3.4 hold. Obviously, such an infinite derivation will be detected at the node $`N_{g_{d+1}}:G_{g_{d+1}}`$, thus Algorithm 4.1 will stop here.
When Algorithm 4.1 ends with an answer Yes, all generalized SLDNF-trees for all queries in $`S_Q`$ must have been generated without encountering any derivation $`D`$ that meets the three conditions of Definition 4.3. This shows that $`P`$ is not most-likely non-terminating w.r.t. $`S_Q`$ and $`R`$. Conversely, if $`P`$ is not most-likely non-terminating w.r.t. $`S_Q`$ and $`R`$, Algorithm 4.1 will not stop at line 1. It will proceed to line 2 with an answer Yes returned.
###### Theorem 4.4
The following hold:
1. If Algorithm 4.1 returns Yes then $`P`$ is terminating w.r.t. $`S_Q`$ and $`R`$.
2. If $`P`$ is not terminating w.r.t. $`S_Q`$ and $`R`$ then Algorithm 4.1 will return a generalized SLDNF-derivation $`D`$ that meets the conditions of Definition 4.3.
Proof: 1. By Theorem 4.3 Algorithm 4.1 returning Yes implies $`P`$ is not most-likely non-terminating w.r.t. $`S_Q`$ and $`R`$. The result then follows from Theorem 4.2.
2. By Theorem 4.2, when $`P`$ is not terminating w.r.t. $`S_Q`$ and $`R`$, it is most-likely non-terminating w.r.t. $`S_Q`$ and $`R`$. So there exist generalized SLDNF-derivations in some generalized SLDNF-trees $`GT_{Q_i}`$ that meet the three conditions of Definition 4.3. Obviously, the fist such derivation $`D`$ will be captured at line 1 of Algorithm 4.1, which leads to an output $`D`$.
### 4.3 Examples
We use the following very representative examples to illustrate the effectiveness of our method. In the sequel, we choose the smallest depth bound $`d=2`$.
###### Example 4.1
Applying Algorithm 4.1 to the logic program $`P_1`$ of Example 2.1 with a query $`Q_1=p(a)`$ will return a generalized SLDNF-derivation $`D`$, which is the path from $`N_0`$ to $`N_4`$ in Figure 1. $`D`$ is informative enough to suggest that $`P_1`$ is not terminating w.r.t. $`Q_1`$.
###### Example 4.2
Applying Algorithm 4.1 to the logic program $`P_2`$ of Example 2.2 with a query $`Q_1=p`$ will return an answer Yes. That is, $`P_2`$ is terminating w.r.t. $`Q_1`$. However, for the query $`Q_2=q`$ applying Algorithm 4.1 will return the following generalized SLDNF-derivation
which strongly suggests that $`P_2`$ is not terminating w.r.t. $`Q_2`$.
###### Example 4.3
Consider the following widely used program:<sup>3</sup><sup>3</sup>3It represents a large class of well-known logic programs such as $`member`$, $`subset`$, $`merge`$, $`quick`$-$`sort`$, $`reverse`$, $`permutation`$, and so on.
| | $`P_3:`$ | $`append([],X,X)`$. $`C_{a_1}`$ |
| --- | --- | --- |
| | | $`append([X|Y],U,[X|Z])append(Y,U,Z)`$. $`C_{a_2}`$ |
Assume the following types of queries (borrowed from ):
| | $`Q_1=append([1,2],[3],L)`$, |
| --- | --- |
| | $`Q_2=append([1,2],[3],[4])`$, |
| | $`Q_3=append(L_1,L_2,[1,2])`$, |
| | $`Q_4=append(L_1,[1,2],L_3)`$, |
| | $`Q_5=append(L_1,L_2,L_3)`$, |
| | $`Q_6=append([X|Y],[],Y)`$, |
| | $`Q_7=append([X|Y],Y,[Z|Y])`$. |
The generalized SLDNF-trees $`GT_{Q_1}`$, $`GT_{Q_2}`$, and $`GT_{Q_3}`$ are all finite and contain no expanded variants in any derivations. So Algorithm 4.1 will return Yes when executing $`Test(P_3,\{Q_1,Q_2,Q_3\},R)`$. That is, $`P_3`$ is terminating w.r.t. the first three types of queries.
When evaluating $`Q_4`$, Algorithm 4.1 will return a generalized SLDNF-derivation as shown in Figure 3 (a). Note that all selected subgoals in the derivation are variants. Similar derivations will be returned when applying Algorithm 4.1 to $`Q_5`$ and $`Q_6`$. Applying Algorithm 4.1 to $`Q_7`$ will yield a generalized SLDNF-derivation as shown in Figure 3 (b). Note that the selected subgoal at node $`N_3`$ is an expanded variant of the subgoal at $`N_2`$ that is an expanded variant of the subgoal at $`N_1`$. That is, $`append(Y_2,[X_1|[X_2|Y_2]],Y_2)_{EV}append(Y_1,[X_1|Y_1],Y_1)_{EV}append(Y,Y,Y)`$.
It is clear that the derivations of Figures 3 (a) and (b) can be infinitely extended by repeatedly applying the clause $`C_{a_2}`$, thus $`P_3`$ is non-terminating w.r.t. the queries $`Q_4Q_7`$.
###### Example 4.4
The following program illustrates how a list of terms grows recursively through a DELETE-ADD-SHUFFLE process.
| | $`P_4:`$ | $`p([X_1,X_2|Y])q([X_1,X_2|Y],Z),reverse(Z,[],Z_1),p(Z_1)`$. $`C_{p_1}`$ |
| --- | --- | --- |
| | | $`q([X_1,X_2|Y],[X_3,f(X_1,X_2),X_2|Y]).`$ $`C_{q_1}`$ |
| | | $`reverse(Z,[],Z_1)`$ $`Z_1`$ is the reversed list of $`Z`$. $`C_{rev}`$ |
Given a subgoal $`p([X_1,X_2|Y])`$, applying $`C_{p_1},C_{q_1},C_{rev}`$ successively will
| | DELETE | $`X_1`$, thus yielding a list $`[X_2|Y]`$, |
| --- | --- | --- |
| | ADD | $`X_3`$ and $`f(X_1,X_2)`$, thus yielding a list $`[X_3,f(X_1,X_2),X_2|Y]`$, and |
| | SHUFFLE | the list $`[X_3,f(X_1,X_2),X_2|Y]`$ by reversing its components. |
Note that the addition of $`X_3`$ is independent of the original list $`[X_1,X_2|Y]`$, but $`f(X_1,X_2)`$ is generated based on the list. This means that given a query of the form $`p([T_1,\mathrm{},T_m])`$, a new variable $`X`$ and a function $`f(A_1^j,A_2^j)`$ will be added each cycle $`\{C_{p_1},C_{q_1},C_{rev}\}`$ is applied, where the domain of $`A_1^j`$ and $`A_2^j`$ is the closure of the function $`f(\mathrm{\_},\mathrm{\_})`$ over $`\{X,T_1,\mathrm{},T_m\}`$ (up to variable renaming). As an illustration, consider an arbitrary query $`Q_1=p([a,b])`$. Applying Algorithm 4.1 to $`Q_1`$ will return a generalized SLDNF-derivation as shown in Figure 4, where for the selected subgoals $`L_{12}^1,L_6^1,L_0^1`$ at nodes $`N_{12}`$, $`N_6`$ and $`N_0`$, we have $`L_0^1_{ANC}L_6^1_{ANC}L_{12}^1`$, $`L_{12}^1_{EV}L_6^1_{EV}L_0^1`$, and $`|L_{12}^1|>|L_6^1|>|L_0^1|`$. We see the following terms added due to the repeated applications of $`\{C_{p_1},C_{q_1},C_{rev}\}`$:
| | From | $`N_0`$ to $`N_3`$ | $`X_1,f(a,b),`$ |
| --- | --- | --- | --- |
| | | $`N_3`$ to $`N_6`$ | $`X_2,f(b,f(a,b)),`$ |
| | | $`N_6`$ to $`N_9`$ | $`X_3,f(X_1,f(a,b)),`$ |
| | | $`N_9`$ to $`N_{12}`$ | $`X_4,f(X_2,f(b,f(a,b))).`$ |
Apparently, the generalized SLDNF-derivation of Figure 4 can be infinitely extended. Thus $`P_4`$ is non-terminating w.r.t. $`Q_1`$ (and all queries of the form $`p([T_1,\mathrm{},T_m])`$).
###### Example 4.5 ()
Consider the following well-known game program:
| | $`P_5:`$ | $`win(X)move(X,Y).\neg win(Y)`$. $`C_{w_1}`$ |
| --- | --- | --- |
| | | $`move(a,b)`$ $`for`$ $`(a,b)𝒢`$ where $`𝒢`$ is an acyclic finite graph. $`C_{m_1}`$ |
Assume the following two types of queries:
| | $`Q_1=win(a)`$, |
| --- | --- |
| | $`Q_2=win(X)`$. |
Since $`𝒢`$ is an acyclic finite graph, no expanded variants occur in any generalized SLDNF-derivations. Therefore, Algorithm 4.1 will terminate for both $`Q_1`$ and $`Q_2`$ with an answer Yes. That is, $`P_5`$ is terminating w.r.t. $`\{Q_1,Q_2\}`$.
###### Example 4.6 ()
The following general logic program is used to compute the transitive closure of a graph.
| | $`P_6:`$ | $`r(X,Y,E,V)member([X,Y],E)`$. $`C_{r_1}`$ |
| --- | --- | --- |
| | | $`r(X,Z,E,V)member([X,Y],E),\neg member(Y,V),r(Y,Z,E,[Y|V])`$. $`C_{r_2}`$ |
| | | $`member(X,[X|T])`$. $`C_{m_1}`$ |
| | | $`member(X,[Y|T])member(X,T)`$. $`C_{m_1}`$ |
Queries over this program are of the form $`r(X,Y,e,[X])`$ where $`X`$, $`Y`$ are nodes and $`e`$ is a graph specified by a finite list of its edges denoted by $`[Node,Node]`$. Such a query is supposed to succeed when $`[X,Y]`$ is in the transitive closure of $`e`$. The last argument of $`r(X,Y,e,[X])`$ acts as an accumulator in which a list of nodes is maintained which should not be reused when looking for a path connecting $`X`$ with $`Y`$ in $`e`$ (to keep the search path acyclic). As an example, let $`e=\{[[a,b],[b,c],[c,a]]\}`$. We consider the following three types of queries:
| | $`Q_1=r(a,c,e,[a])`$, |
| --- | --- |
| | $`Q_2=r(a,Y,e,[a])`$, |
| | $`Q_3=r(X,Y,e,[X])`$. |
The generalized SLDNF-trees $`GT_{Q_1}`$, $`GT_{Q_2}`$, and $`GT_{Q_3}`$ are depicted in Figures 5, 6 and 7, respectively. Since there is no expanded variant in any generalized SLDNF-derivations, Algorithm 4.1 will return Yes when executing $`Test(P_6,\{Q_1,Q_2,Q_3\},R)`$. That is, $`P_6`$ is terminating w.r.t. these three types of queries.
It is interesting to observe that for each of the above logic programs, $`P_1P_6`$, it is terminating if and only if applying Algorithm 4.1 to it yields an answer Yes. In fact, this is true for all representative logic programs we have currently collected in the literature. However, due to the undecidability of the termination problem, it is unavoidable that there exist cases in which Algorithm 4.1 returns a generalized SLDNF-derivation $`D`$, but $`P`$ is a terminating logic program. We have created such a rarely used program.
###### Example 4.7
Consider the following logic program and top goal, where the function $`size(Z)`$ returns the number of elements in the list $`Z`$ (e.g. $`size([a,b])=2`$).
| | $`P_7:`$ | $`p([X|Y],N)size([X|Y])<N,p([X,X|Y],N)`$. $`C_{p_1}`$ |
| --- | --- | --- |
| | $`G_0:`$ | $`p([a],100).`$ |
The generalized SLDNF-tree $`GT_{G_0}`$ for $`P_7\{G_0\}`$ is shown in Figure 8. It is easy to see that for any $`i0`$, the subgoal $`L_{2i}^1`$ at $`N_{2i}`$ is an ancestor subgoal of the subgoal $`L_{2(i+1)}^1`$ at $`N_{2(i+1)}`$, while $`L_{2(i+1)}^1`$ is an expanded variant of $`L_{2i}^1`$ with $`|L_{2(i+1)}^1|>|L_{2i}^1|`$. $`P_7`$ is terminating w.r.t. the query $`p([a],100).`$ However, applying Algorithm 4.1 (with $`d=2`$) will return a generalized SLDNF-derivation $`D`$ that is the segment between $`N_0`$ and $`N_4`$ in Figure 8. Apparently, in order for Algorithm 4.1 to return Yes the depth bound $`d`$ should not be less than 100.
## 5 Related Work
Our work is related to both termination analysis and loop checking.
### 5.1 Work on Termination Analysis
Concerning termination analysis, we refer the reader to the survey of Decorte, De Schreye and Vandecasteele for a comprehensive bibliography.
There are two essential differences between existing termination analysis techniques and ours. The first difference is that theirs are static approaches, whereas ours is a dynamic one. Static approaches only make use of the syntactic structure of the source code of a logic program to establish some well-founded conditions/constraints that, when satisfied, yield a termination proof. Since non-termination is caused by an infinite generalized SLDNF-derivation, which contains some essential dynamic characteristics (such as expanded variants and the repeated application of some clauses) that are hard to capture in a static way, static approaches appear to be less precise than a dynamic one. For example, it is difficult to apply a static approach to prove the termination of program $`P_2`$ in Example 2.2 with respect to a query pattern $`p`$. Moreover, although some static approaches (e.g., see ) are automatizable, searching for an appropriate level mapping or computing some interargument relations could be very complex. For our approach, the major work is to identify expanded variants, which is easy to complete.
The second difference is that existing methods are suitable for termination analysis with respect to query patterns, whereas ours is for termination analysis with respect to concrete queries. The advantage of using query patterns is that if a logic program $`P`$ is shown to be terminating with respect to a query pattern $`Q`$, it is terminating with respect to all instances of $`Q`$ that could be an infinite set of concrete queries. However, if $`P`$ is shown to be not terminating with respect to $`Q`$, which usually means that $`P`$ is terminating with respect to some instances of $`Q`$ but is not with respect to the others, we cannot apply existing termination analysis methods to make such a further distinction. In contrast, our method can make termination analysis for each single concrete query posed by the user and provide explanations about how non-termination happens. This turns out to be very useful in real programming practices. Observe that in developing a software in any computer languages we always apply some typical cases (i.e. concrete parameters as inputs) to test for the correctness or termination of the underlying programs, with an assumption that if the software works well with these typical cases, it would work well with all cases of practical interests.
From the above discussion, it is easy to see that our method plays a complementary role with respect to existing termination analysis approaches (i.e. static versus dynamic and query patterns versus concrete queries).
### 5.2 Work on Loop Checking
Loop checking is a run-time approach towards termination. It locates nodes at which SLD-derivations step into a loop and prunes them from SLD-trees. Informally, an SLD-derivation
$`N_0:G_0N_1:G_1\mathrm{}N_i:G_i\mathrm{}N_k:G_k\mathrm{}`$
is said to step into a loop at a node $`N_k:G_k`$ if there is a node $`N_i:G_i`$ ($`0i<k`$) in the derivation such that $`G_i`$ and $`G_k`$ are sufficiently similar. Many mechanisms related to loop checking have been presented in the literature (e.g. see ). We mention here a few representative ones.
Bol, Apt and Klop introduced the Equality check and the Subsumption check. These loop checks can detect loops of the form
$`N_0:G_0N_1:G_1\mathrm{}N_i:G_i\mathrm{}N_k:G_k`$
where either $`G_k`$ is a variant or an instance of $`G_i`$ (for the Equality check), i.e. $`G_k=G_i\theta `$ under a substitution $`\theta `$, or $`G_i`$ is included in $`G_k`$ under a substitution $`\theta `$ (for the Subsumption check), i.e. $`G_kG_i\theta `$. However, they cannot handle infinite SLD-derivations of the form
$`N_0:p(X)N_1:p(f(X))\mathrm{}N_i:p(f(\mathrm{}f(X)\mathrm{}))\mathrm{}`$
Sahlin introduced the OS-check (see also ). It determines infinite loops based on two parameters: a depth bound $`d`$ and a size function size. Informally, OS-check says that an SLD-derivation may go into an infinite loop if it generates an oversized subgoal. A subgoal $`A`$ is said to be oversized if it has $`d`$ ancestor subgoals in the SLD-derivation that have the same predicate symbol as $`A`$ and whose size is smaller than or equal to $`A`$.
Bruynooghe, De schreye and Martens presented a framework for partial deduction with finite unfolding that, when applied to loop checking, is very similar to OS-check. That is, it mainly relies on term sizes of (selected) subgoals and a depth bound. See for a detailed comparison of these works.
OS-check (similarly the method of Bruynooghe, De schreye and Martens) is complete in the sense that it cuts all infinite loops. However, because it merely takes the number of repeated predicate symbols and the size of subgoals as its decision parameters, without referring to the informative internal structure of the subgoals, the underlying decision is fairly unreliable; i.e. many non-loop derivations may be pruned unless the depth bound $`d`$ is set sufficiently large.
Using expanded variants, in we proposed a series of loop checks, called VAF-checks (for Variant Atoms loop checks for logic programs with Functions). These loop checks are complete and much more reliable than OS-check. However, they cannot deal with infinite recursions through negation like that in Figure 1.
The work of the current paper can partly be viewed as an extension of from identifying infinite SLD-derivations to identifying infinite generalized SLDNF-derivations. It is worth noting that termination analysis is merely concerned with the characterization and identification of infinite derivations, but loop checking is also concerned about how to prune infinite derivations. The latter work heavily relies on the semantics of a logic program, especially when an infinite recursion through negation occurs. Bol discussed loop checking for locally stratified logic programs under the perfect model semantics .
## 6 Conclusions
We have presented a method of verifying termination of general logic programs with respect to concrete queries. A necessary and sufficient condition is established and an algorithm for automatic testing is developed. Unlike existing termination analysis approaches, our method does not need to search for a model or a level mapping, nor does it need to compute an interargument relation based on additional mode or type information. Instead, it detects infinite derivations by directly evaluating the set of queries of interest. As a result, some key dynamic features of a logic program can be extracted and employed to predict its termination. Such idea partly comes from loop checking. Therefore, the work of this paper bridges termination analysis with loop checking, the two problems which have been studied separately in the past despite their close technical relation with each other.
It is worth mentioning that the practical purpose of termination analysis is to assist users to write terminating programs. Our method exactly serves for this purpose. When Algorithm 4.1 outputs Yes, the logic program is terminating; otherwise it provides users with a generalized SLDNF-derivation of the form as shown in Figures 3 or 4. Such a derivation may most likely lead to an infinite derivation, thus users can improve their programs following the informative guidance. (In this sense, our method is quite like a spelling mechanism used in a word processing system, which always indicates most likely incorrect spellings.)
Due to the undecidability of the termination problem, there exist cases in which a logic program is terminating but Algorithm 4.1 would not say Yes unless the depth bound $`d`$ is set sufficiently large (see Example 4.7). Although $`d=2`$ works well for a vast majority of logic programs (see Examples 4.1 \- 4.6), how to choose the depth bound in a general case then presents an interesting open problem.
Tabled logic programming is receiving increasing attention in the community of logic programming (e.g. see ). Verbaeten, De Schreye and K. Sagonas recently exploited termination proofs for positive logic programs with tabling. For future research, we are going to extend the work of the current paper to deal with general logic programs with tabling.
## Acknowledgements
We would like to thank Danny De Schreye for his constructive comments on our work and valuable suggestions for its improvement.
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# Nuclear Spin Relaxation for Higher Spin
## Abstract
We study the relaxation of a spin $`I`$ that is weakly coupled to a quantum mechanical environment. Starting from the microscopic description, we derive a system of coupled relaxation equations within the adiabatic approximation. These are valid for arbitrary $`I`$ and also for a general stationary non–equilibrium state of the environment. In the case of equilibrium, the stationary solution of the equations becomes the correct Boltzmannian equilibrium distribution for given spin $`I`$. The relaxation towards the stationary solution is characterized by a set of relaxation times, the longest of which can be shorter, by a factor of up to $`2I`$, than the relaxation time in the corresponding Bloch equations calculated in the standard perturbative way.
Nuclear magnetic resonance is a well–established method for testing electronic properties in solids . In recent years, it became possible to apply this technique not only in three dimensions, but also to a two–dimensional electron system, the quantum Hall ferromagnet that is realized in semiconductor heterostructures in a strong magnetic field. The experimental work lead to the unexpected conclusion that a new kind of low–energy states, Skyrmions, can be formed and can determine the nuclear relaxation processes in these systems when one Landau sub–level of one spin direction is filled .
As a theoretical description of spin relaxation, Bloch’s equations have been successfully used for about fifty years now. While these phenomenological equations are applicable in a wide range of cases, their microscopic derivation reveals two main restrictions. First, as was already discussed in the original work , the derivation becomes strictly valid, if either the spin is $`I=1/2`$ or the temperature of the bath is large compared to the resonance frequency. But the spin in the system under study can be $`I=3/2`$ (for <sup>69</sup>Ga, <sup>71</sup>Ga, and <sup>75</sup>As; cf. Ref. ), or higher in the case of magnetic impurities. Further, the progress in the experimental techniques now lets a regime of temperatures and magnetic fields come into reach, in which the temperature of the bath is of the same order as the nuclear resonance energy (nuclear Zeeman energy). The second restriction in the derivation of the phenomenological equations demands that the environment (bath), the quantum mechanical degrees of freedom causing the spin relaxation, be in thermodynamic equilibrium. But in the case of the quantum Hall ferromagnet, the nuclear spins are coupled to a two–dimensional electron gas in which the electron–electron interaction plays the key role, since all single particle states are degenerate into a single Landau level due to the strong magnetic field; such a system, dominated by Skyrmion–states, is not necessarily in equilibrium. Thus, it appears worthwhile to reconsider the microscopic derivation of phenomenological equations for the spin relaxation in order to investigate whether there is a significant difference between the general case and a case in which the Bloch equations are valid.
It is the purpose of this work to relax the two above restrictions by investigating the general case of an arbitrary spin $`I`$ and also an arbitrary stationary state of the environment responsible for the relaxation. We refrain from studying a specific mechanism and consider instead the general case of a magnetic moment coupled to a bath of other quantum degrees of freedom. This magnetic moment can be a nuclear spin, or also a magnetic impurity. In the following, we use the terms “nuclear spin” for the magnetic moment and “electrons” for the bath - usually the latter is called “lattice”. Then, the contribution of the nuclear spin to the Hamiltonian is
$$H=\gamma \widehat{\stackrel{}{I}}(\stackrel{}{B}_0+\widehat{\stackrel{}{B}}).$$
(1)
Here, the magnetic moment, $`\gamma \widehat{\stackrel{}{I}}`$ (where $`\widehat{\stackrel{}{I}}`$ is the spin), couples to an effective magnetic field. It is well known that for $`I>1/2`$, there is an additional term causing relaxation, the electric quadrupole moment of the nucleus coupled to an inhomogeneous external electric field. Here, we have omitted this term in $`H`$, since the model Eq. (1) already suffices for the demonstration of our method; an inclusion of a quadrupolar coupling is straightforward. The effective magnetic field acting on the nuclear spin $`\widehat{\stackrel{}{I}}`$ in Eq. (1) contains an operator $`\widehat{\stackrel{}{B}}`$ of the electronic system. One can picture $`\widehat{\stackrel{}{B}}`$ as being proportional to the electrons’ spin. Its longitudinal component, $`\widehat{B}_z`$, modifies the eigenvalues of the nuclear spin system, while the transverse components cause transitions between eigenstates. There is also a fixed part of the magnetic field in $`z`$–direction, $`\stackrel{}{B}_0=B_0\widehat{e}_z`$, which acts as an external field. The coupling beween nuclear spin and electrons is supposed to be weak in the sense that we can use the adiabatic approximation as discussed below. We do not need to make any assumptions about the electronic subsystem’s Hamiltonian or the electronic subsystem’s state. This Hamiltonian may contain electron–electron interactions and the subsystem may be in an arbitrary stationary state, equilibrium or non–equilibrium.
We want to derive kinetic equations for the expectation value of the spin vector $`\widehat{\stackrel{}{I}}`$ ($`\widehat{\stackrel{}{I}}{}_{}{}^{2}=I(I+1)`$). To this end, we need to take all components of the density matrix into account, not just only the spin vector. Then, as will be seen below, it is advantageous to use the spherical tensor operators $`\widehat{T}_{LM}`$ as a complete basis in the space of operators acting on the state of a spin $`I`$. The $`\widehat{T}_{LM}`$ are the irreducible tensor operators in the spherical coordinate representation . For actual calculations, the following definition using the spherical harmonics $`Y_{LM}`$ proves very helpful:
$$\widehat{T}_{LM}:=𝒩_{IL}(\widehat{\stackrel{}{I}}\stackrel{}{})^Lr^LY_{LM}(\widehat{r}).$$
(2)
$`\widehat{T}_{LM}`$ is a polynomial in the components of the spin operator. It is independent of the auxiliary variable $`\stackrel{}{r}`$, since on the r.h.s., there are $`L`$ derivatives acting on a polynomial of order $`L`$ ($`\widehat{r}=\stackrel{}{r}/r`$ denotes the unit vector; we use the conventions of Ref. for the spherical harmonics $`Y_{LM}`$). For a spin $`I`$, the $`(2I+1)^2`$ operators $`\widehat{T}_{LM}`$ with $`L=0\mathrm{}2I`$ and $`M=L\mathrm{}L`$ form a complete system of operators acting in the spinspace. The normalization, $`𝒩_{IL}=2^L\sqrt{4\pi (2IL)!/(2I+L+1)!}/L!`$, is choosen such that $`Sp\left\{\widehat{T}_{LM}^{}\widehat{T}_{L^{}M^{}}\right\}=\delta _{LL^{}}\delta _{MM^{}},`$ and we have $`\widehat{T}_{LM}^{}=(1)^M\widehat{T}_{L,M}`$. Specific expressions for the operators $`\widehat{T}_{L,M}`$ (for $`L=0\mathrm{}4`$) can be found in Table IV of Ref. .
After having established the basic notation, we proceed now to describe the derivation of the kinetic equation for the average $`T_{LM}(t)=\widehat{T}_{LM}`$, to find the stationary solution and to study finally the relaxation towards the stationary solution. Here and below, the brackets $`\mathrm{}`$ stand for the state of the combined system of spin and electrons. We use the framework of the Keldysh method in order to derive the kinetic equation. As in our earlier work on the electron spin relaxation , we shall employ the adiabatic approximation; that means that in the equation of motion, the effect of the coupling between spin and electrons is neglected beyond the first order in the spin’s eigenenergies and also neglected beyond the second order in the relaxation times, cf. Ref. . The unperturbed motion of the spin is a precession with the frequency $`\omega _0=\gamma B_0`$. Then, we have up to the second order of the perturbation theory in the coupling $`\widehat{V}=\gamma \widehat{\stackrel{}{I}}\widehat{\stackrel{}{B}}`$:
$`\left(i_tM\omega _0\right)T_{LM}(t)=`$ (3)
$`[\widehat{T}_{LM}(t),\widehat{V}(t)]`$ (4)
$`i{\displaystyle _{\mathrm{}}^t}𝑑t_1[[\widehat{T}_{LM}(t),\widehat{V}(t)],\widehat{V}(t_1)]+𝒪(\widehat{V}^3).`$ (5)
In the spirit of the adiabatic approximation, we now consider the terms of the perturbation series in higher than second order as giving rise to an additional (weak) time dependence of the spin operators due to the relaxation process, and we decouple the expectation values of spin and electron operators. The first order term, $`[\widehat{T}_{LM},\widehat{\stackrel{}{I}}](t)\widehat{\stackrel{}{B}}(t)`$, describes the Knight shift, the shift of the nuclear resonance frequency due to the coupling to the electrons. Since we want to focus on the relaxation, we disregard this corrections of the spin’s eigenfrequency in the following. Then, we get
$`\left(i_tM\omega _0\right)T_{LM}(t)=`$ (6)
$`i{\displaystyle \underset{m;m^{}=0,\pm 1}{}}([[\widehat{T}_{LM},\widehat{I}_m],\widehat{I}_m^{}](t)C_{m;m^{}}`$ (7)
$`\text{ }+\{[\widehat{T}_{LM},\widehat{I}_m],\widehat{I}_m^{}\}(t)R_{m;m^{}})`$ (8)
In deriving Eq. (8), we took the unperturbed time dependence of the operator $`\widehat{\stackrel{}{I}}(t_1)`$ into account, but neglected the difference $`tt_1`$ in its time dependence due to the relaxation. In Eq. (8), equal–time commutators ($`[,]`$) and anticommutators ($`\{,\}`$) can be evaluated, and the expectation values can again be expressed by $`T_{LM}(t)`$ as we shall see below. Thus, it is the use of the spherical tensor operators that makes it possible to derive closed coupled equations for the relaxation of the spin. Most crucial are the averages of the electronic subsystem that enter Eq. (8). Since we do not assume thermodynamic equilibrium for the electronic subsystem, we get both, a correlation function and also a response function which are independent and which we denote by $`C_{m;m^{}}`$ and $`R_{m;m^{}}`$, respectively. They are given by
$`C_{m;m^{}}`$ $`=`$ $`{\displaystyle \frac{\gamma ^2}{2}}{\displaystyle \underset{\mathrm{}}{\overset{t}{}}}𝑑t_1e^{im^{}\omega _0(tt_1)}\{\widehat{B}_m(t),\widehat{B}_m^{}(t_1)\}`$ (9)
$`R_{m;m^{}}`$ $`=`$ $`{\displaystyle \frac{\gamma ^2}{2}}{\displaystyle \underset{\mathrm{}}{\overset{t}{}}}𝑑t_1e^{im^{}\omega _0(tt_1)}[\widehat{B}_m(t),\widehat{B}_m^{}(t_1)]`$ (10)
Our convention for the vector components of $`\widehat{\stackrel{}{I}}`$ and $`\widehat{\stackrel{}{B}}`$ is: $`\widehat{I}_{\pm 1}=\widehat{I}_x\pm i\widehat{I}_y`$, $`\widehat{I}_0=\widehat{I}_z`$ and $`\widehat{B}_{\pm 1}=(\widehat{B}_x\pm i\widehat{B}_y)/2`$, $`\widehat{B}_0=\widehat{B}_z`$.
Commutators and anticommutators in Eq. (8) can be calculated either directly from the definition Eq. (2), or with the aid of the theory of spherical tensors (cf. Ref. ). The commutators express the behavior of the $`\widehat{T}_{LM}`$ under rotations:
$`[\widehat{T}_{LM},\widehat{I}_z]`$ $`=`$ $`M\widehat{T}_{LM}`$ (12)
$`[\widehat{T}_{LM},\widehat{I}_\pm ]`$ $`=`$ $`\sqrt{L(L+1)M(M\pm 1)}\widehat{T}_{L,M\pm 1}.`$ (13)
The anticommutators can also be expanded in $`\widehat{T}_{LM}`$; specializing the general result containing Racah coefficients (see Eq. (17) of ref. ) to our case, we get for $`m=0,\pm 1`$:
$$\{\widehat{T}_{LM},\widehat{I}_m\}=a_{L+1;M}^{(m)}\widehat{T}_{L+1,M+m}+b_{L;M}^{(m)}\widehat{T}_{L1,M+m}$$
(14)
with the coefficients
$`a_{L;M}^{(0)}=b_{L;M}^{(0)}=\sqrt{L^2M^2}c_L,`$ (15)
$`a_{L;M}^{(\pm 1)}=b_{L;M\pm 1}^{(1)}=\sqrt{(L+1\pm M)(L\pm M)}c_L,`$ (16)
where
$$c_L=\sqrt{\frac{(2I+1)^2L^2}{(2L)^21}}.$$
(17)
Inserting relations (Nuclear Spin Relaxation for Higher Spin) and (14) into Eq. (8) solves our task of deriving a closed set of relaxation equations for arbitrary nuclear spin $`I`$. The equations are linear in the expectation values of the spherical tensor operators $`\widehat{T}_{LM}`$. The properties of the electronic system enter the equations parametrically in the form of correlation functions and response functions. Since the spin vector is given by the $`\widehat{T}_{1M}`$, a relaxation equation for $`\widehat{\stackrel{}{I}}`$ can be extracted from the system of equations (8). In the case of $`I=1/2`$, this is particularly simple, since $`\widehat{T}_{LM}=0`$ for $`L>1`$ (and $`\widehat{T}_{00}=1/\sqrt{2}`$ is constant) and we recover the Bloch equations, see below.
In this general form, the equations (8) are still not very transparent. Therefore, we now make additional assumptions regarding the correlation and response functions entering Eq. (8). The term $`\widehat{I}_z\widehat{B}_z`$ in the perturbation $`\widehat{V}`$ (corresponding to $`m,m^{}=0`$ in Eq. (8)) just changes the spin’s resonance frequency $`\omega _0`$. Since we already omitted the Knight shift, the first order correction term to this frequency, we now also omit consistently the second order contributions resulting from $`m=0`$ or $`m^{}=0`$ in Eq. (8). Next, the terms with $`m=m^{}`$ in Eq. (8) are neglected too. In the electron system, they would correspond to a total change of the $`z`$–component of the spin by $`2`$; if the electrons’ state is a strict eigenstate to $`z`$–component of the spin, such expectation values vanish. Under these assumptions, both correlation and response functions, $`C_{m;m^{}}`$ and $`R_{m;m^{}}`$, now become “diagonal” in the index $`m`$, $`C_{m;m^{}}=\delta _{m;m^{}}C_m`$ and $`R_{m;m^{}}=\delta _{m;m^{}}R_m`$. Now, from the general structure of Eq. (8), it is obvious that the equations do not couple $`T_{LM}(t)`$ for different $`M`$. So finally, we arrive at our general result, valid for arbitrary spin $`I`$ and for a non–equilibrium state of the bath
$`\left(_t+iM\omega _0\right)T_{LM}(t)=`$ (18)
$`{\displaystyle \frac{1}{2\tau _+}}\left[L(L+1)M^2\right]T_{LM}(t)`$ (19)
$`{\displaystyle \frac{1}{2\tau _{}}}\{L\sqrt{(L+1)^2M^2}c_{L+1}T_{L+1,M}(t)`$ (20)
$`\text{ }(L+1)\sqrt{L^2M^2}c_LT_{L1,M}(t)\}.`$ (21)
We have defined two different times by $`1/\tau _+:=2(C_1+C_1)`$ and $`1/\tau _{}:=2(R_1R_1)`$. These times are independent, as long as we assume a non–equilibrium state of the electronic system. We have $`\tau _+0`$ and the ratio $`\tau _+/\tau _{}`$ can be expressed as $`\tau _+/\tau _{}=(1q)/(1+q)`$ with positive $`q`$. Thus, $`|\tau _+/\tau _{}|1`$. If the electrons are in equilibrium at temperature $`T`$, the fluctuation–dissipation theorem results in $`q=\mathrm{exp}(\omega _0/T)`$.
We want to stress that these relaxation equations are valid (under the assumptions stated above), no matter which specific relaxation mechanism one wants to consider. The special mechanism enters the equations in the form of two time scales, $`\tau _+`$ and $`\tau _{}`$. In the case of thermodynamic equilibrium in the bath, their ratio, $`\tau _+/\tau _{}`$, is fixed by the temperature, and only the time $`\tau _+`$ – which also depends on temperature – is specific for the relaxation mechanism.
Relaxation equations generally serve two purposes. First, they determine a stationary state. Second, they describe the relaxation towards this state. We now want to discuss both of these points.
Stationary solution
The study of the stationary solution of Eq. (21) serves as an important test of our procedure, since, in the case of equilibrium in the bath, the result is obvious. With the ansatz $`T_{LM}(t)=\delta _{M;0}T_L`$, one derives a recursion relation
$$c_{L+1}T_{L+1}=\frac{\tau _{}}{\tau _+}T_L+c_LT_{L1},L=1\mathrm{}2I$$
(22)
where $`T_0=(2I+1)^{1/2}`$, $`T_{2I+1}=0`$. The solution is
$$I_z=\frac{2}{3}I(I+1)(1q)_{21}/_{11}$$
(23)
where $`_{ij}`$ are elements of the matrix $`=M_1M_2\mathrm{}M_{2I}`$ with
$$M_L=\left(\genfrac{}{}{0pt}{}{1+q(1q)^2c_{L+1}^2}{10}\right).$$
(24)
Evaluation of this matrix product yields – for a general non–equilibrium state of the bath – the following stationary distribution for the $`z`$–component of the spin vector
$$I_z=\underset{m=I}{\overset{I}{}}mq^m/\underset{m=I}{\overset{I}{}}q^m.$$
(25)
If we assume now the electronic system is in equilibrium at temperature $`T`$, we find that $`I_z`$ is given by the Brillouin function, the well–known correct equilibrium distribution for a spin $`I`$.
Relaxation
The general solution of the coupled relaxation equations Eq. (21) is a superposition of exponentially decaying terms. We are mostly interested in the relaxation of the spin components $`\widehat{I}_z\widehat{T}_{10}`$ and $`(\widehat{I}_x+i\widehat{I}_y)\widehat{T}_{11}`$, which are described by the part $`M=0`$ and $`M=1`$ of the full system of equations. For $`M=0`$ and $`M=1`$, there are $`2I`$ different relaxation times, for $`M=2`$ there are $`2I1`$ times; in total we have $`I(2I+3)`$ times. Here, each case of $`I`$ needs to be discussed separately.
We start with spin $`1/2`$. In this case, the Eq. (21) gives immediately the standard Bloch equations :
$`_tI_z(t)={\displaystyle \frac{1}{\tau _+}}\left(I_z(t){\displaystyle \frac{\tau _+}{2\tau _{}}}\right)`$ (26)
$`\left(_t+i\omega _0\right)I_+(t)={\displaystyle \frac{1}{2\tau _+}}I_+(t)`$ (27)
More interesting, because of its experimental relevance , is the case of spin $`3/2`$. The relaxation equations for $`M=0`$ couple the expectation values of three operators: $`\widehat{T}_{10}`$, $`\widehat{T}_{20}`$, and $`\widehat{T}_{30}`$. These equations are easily solved. In Fig. 1, we show the flow of the expectation values towards the stationary state at the origin in the case of $`\tau _+/\tau _{}=0.8`$. The inverse time scales, $`\lambda `$, which determine the relaxation are shown for both the longitudinal ($`\lambda _L`$, $`M=0`$) and the transverse equations ($`\lambda _T`$, $`M=1`$) in units of $`1/\tau _+`$ as functions of the temperature in the case of equilibrium in the bath in Fig. 2. Between the high–temperature limit $`T>>\omega _0`$ and the low–temperature limit $`T<<\omega _0`$, both the largest longitudinal and the largest transverse relaxation time decrease by a factor of three as compared to the relaxation time in Bloch approximation calculated in the usual perturbative way.
For general $`I`$, we determine both the longitudinal and the transverse relaxation times in the following limiting cases. For $`T>>\omega _0`$ ($`\tau _+/\tau _{}0`$), the longitudinal relaxation time is $`\tau _+`$ and the transverse $`2\tau _+`$ as in the Bloch equations. In the opposite case, $`T<<\omega _0`$ ($`\tau _+/\tau _{}1`$), all the operators $`\widehat{T}_{LM}`$ become equally important in Eq. (21). Then, the $`2I`$ relaxation times are as follows ($`n=0,1,\mathrm{}2I1`$): In the longitudinal case, $`M=0`$, we get $`\tau _+/[2(n+1)In(n+1)]`$ and the longest relaxation time ($`n=0`$) is always twofold degenerate; in the transverse case, $`M=1`$, we get $`\tau _+/[(2n+1)In^2]`$ and here, the longest relaxation time ($`n=0`$) is always non–degenerate. These explicite expressions are a conjecture based on an explicite evaluation of the relaxation equations for spin $`I`$ up to $`7/2`$.
Starting from the microscopic description, we have derived relaxation equations valid for an arbitrary spin $`I`$ coupled to a quantum mechanical environment in a non–equilibrium state. The solution of these coupled equations shows that, compared to the Bloch equations, there is an additional temperature dependence in the relaxation time which can decrease the relaxation time by a factor of up to $`2I`$.
We would like to thank the Deutsche Forschungsgemeinschaft for its support (Ap 47/3-1, 436 RUS 17/73/00). Yu. A.B. wishes to acknowledge support from grants RFFI-00-02-17292, IR-97-0076, and INTAS 99-01146, and wishes to thank the PTB, where this work was performed, for its hospitality.
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# Dynamics of Entangled Polymeric Fluids in Two-roll Mill studied via Dynamic Light Scattering and Two-color flow Birefringence. I. Steady flow
## I Introduction
The majority of prior experimental studies of the dynamics of polymeric fluids were based on the simple shear flow. The main advantage is that the complete kinematics of these homogeneous flow-fields is known a priory, and experimental measurements of velocity-gradients are not necessary. On the other hand, the majority of the polymer processing applications employ flows that are inhomogeneous and often involve a mixture of simple shear and pure extension. This have provided one motivation for recent experimental as well as theoretical studies of the behavior of the polymeric fluids in “mixed” or purely extensional flows (also called the “strong” flows, which in two-dimensions would require the magnitude of the strain-rate to exceed the vorticity). Generally, such strong flows appear only as a local part of some otherwise globally inhomogeneous and “weak” flows. Examples of such flows include the stagnation zone in a cross-slot or cross-jet device; the converging zone of a contracting channel; and the stagnation region in two- or four-roll mills. The velocity-field in these local region for any polymeric fluid will be very different from that of a Newtonian fluid of similar density and viscosity, because the flows are globally inhomogeneous and the fluid’s viscoelastic nature. It is therefore insufficient to measure only the birefringence (or stress in a rheological experiment), without simultaneously measuring the velocity-gradient field by some technique.
Different indirect techniques, e.g., Particle Imaging Velocimetry and Laser Doppler Velocimetry, hot wire anemometers, forced Rayleigh scattering, holographic grating methods, etc. have been employed in the past for the measurement of velocity-gradients in time varying flows, but each exhibits some characteristic drawbacks. Provided that the flow is laminar and nonchaotic, the dynamic light scattering (DLS) technique is one extremely valuable method for the direct measurement of velocity-gradients in strong flows. The high degree of spatial resolution of the laser light as well as the very high temporal resolution of the present day correlators has made it possible to use DLS for precise determination of velocity-gradients for both steady and time-dependent flows even in very small-scale flow systems with velocity-gradients that change relatively in time. Further, unlike other methods mentioned above, DLS is nonintrusive in nature.
The most frequently used indirect rheo-optical probe of polymer conformation in two-dimensional flows has been the flow birefringence technique. Since it provides an average measure across the flow in third direction, it is only approximately local in that sense. Early studies using traditional flow birefringence technique were primarily restricted to steady flows. On the other hand, recent studies for both dilute and concentrated polymer solutions in time-dependent flows have used two-color flow birefringence (TCFB) and/or phase-modulated birefringence techniques to simultaneously measure the retardance (and hence the birefringence, $`\mathrm{\Delta }n`$) and the orientation angle, $`\chi `$, of the principle axis of the refractive-index tensor with respect to the axes fixed on the flow-cell, in a single experiment.
The theoretical understanding of the dynamics of entangled polymer solution and melts is far from complete. Using de Gennes’ original idea of “reptation”, Doi and Edwards (DE) proposed a tube model, which views the polymer chain as a sequence of segments confined in a tube-like region. The tube is defined by the topological constraints on the lateral motion and orientation imposed by neighboring chains. The segmental stretching can relax on a timescale similar to that for an unconstrained chain (dilute systems), i.e., the Rouse time $`\tau _R`$, and the segmental orientations can only relax via reptational dynamics i.e., through longitudinal diffusion to escape the tube which requires a much longer timescale $`\tau _d`$, called the reptation or disengagement time. These two time scales are related via $`N_e`$, the number of entanglement points per chain: $`\tau _d=3N_e\tau _R`$. The DE model is expected to be valid for highly entangled ($`N_e>>1`$) polymeric samples. Thus the original DE model neglects the chain-stretching effects by assuming that the “snap-back” of the stretched chain to its equilibrium length is instantaneous (i.e., $`\tau _R<<\tau _d`$) and develops constitutive equations which incorporate segmental orientational dynamics only. This simplified model has been quite successful in many rheometric flows such as steady simple shear and oscillatory shear, where chain-stretching is expected to be insignificant. On the other hand, serious inherent limitations of DE theory have long been recognized for the case of strong, extension-dominated flows and also in transient shear flows, where chain-stretching is important. Over the past decade, this model has seen several improvements to include chain-stretching in the uniform, non-uniform and finitely-extensible form. The predictions of the resulting so-called Doi-Edwards-Marrucci-Grizzuti (DEMG) model has been studied in detail in a recent set of papers by Mead et al..
The aim of the present paper is to understand the dynamics of entangled, high molecular weight polymeric fluids subjected to the strong, extension-dominated flows (i.e., for strong two-dimensional flows in the sense of Astarita). To this end, we present two-color flow birefringence and dynamic light scattering results for a series of entangled polystyrene solutions in steady, near-homogeneous, planar flows created by a co-rotating two-roll mill. A detailed comparison of the experimental results for birefringence, orientation angle and generalized extensional viscosity (defined below, and obtained via stress-optical relation) for the three entangled samples to the predictions from the numerical simulation of the DEMG model is carried out using the measured flow data as input to the model. The effect of molecular weight and number of entanglements per chain of the polymers on the results are demonstrated.
In sec. II we provide the necessary details related to our experiments, namely, the samples we used (sec. II.A); the linear viscoelastic measurements (sec. II.B); and the experimental apparatus (sec. II.C). The two-color flow birefringence and the dynamic light scattering techniques and apparatus involved are described in brief along with the two-roll mill flow-cell in sec. II.C. Section III deals with the results of the steady-state flow experiments (sec. III.A) using the DLS method and comparison of the steady flow TCFB results with the DEMG model (sec. III.B). Finally, sec. IV contains a summary of our findings and conclusion.
## II Experimental details
### A Materials
The Newtonian fluid sample used was a suspension of spherical polystyrene particles (polyballs) of diameter $`0.11\mu `$m (Polysciences Inc., U.S.A.) in glycerol (Sigma Chemical Co., U.S.A., ACS Reagent) at a concentration of 150 ppm. As shown in Table I, three non-Newtonian fluids, named PS81, PS82 and PS2, were made using two standard, fairly monodisperse polystyrene samples (Tosoh Co., Japan), namely $`M_w=8.42\times 10^6`$, $`M_w/M_n=1.17`$, Lot. No. TS-31 for one and $`M_w=2.89\times 10^6`$, $`M_w/M_n=1.09`$, Lot. No. TS-10 for the other. Here $`M_w`$ and $`M_n`$ specify the weight-average and the number-average molecular weights respectively. The solvents were prepared by adding polystyrene oligomer with $`M_w=6000`$ or $`2500`$ and a broad molecular weight distribution to toluene (Aldrich Chemical Co., Inc., U.S.A., ACS spectrophotometric grade) at mixing ratios by weight (toluene: 2500 PS = 43:57 for PS81 and PS82, and 48:52 for PS2), and allowing the mixture to dissolve at room temperature ($`=20^{}`$C) for a few days with occasional slow mixing using a magnetic stirrer. Measured amounts of high molecular weight polystyrene samples were then added to these oligomer solvents to achieve the desired concentrations, $`c`$, shown in Table I. The complete solution was then thoroughly mixed in a glass bottle for several hours by continuously rolling the bottle on its side in a Rollacell (New Brunswick Scientific Co., Inc., U.S.A.). In order to reduce the background scattering in the light scattering experiments, the samples were filtered using Milipore filters with $`5\mu `$m pore size for the polymer solutions and $`0.45\mu `$m pore size for the polyball suspension.
The flow-cell was pre-cleaned with acetone (Fisher Scientific, U.S.A., ACS spectranalyzed), oven-dried at $`50^{}`$C and brought back to the room temperature everytime before the fluid sample was transferred into the cell via an air-tight siphon system. To minimize the loss of toluene in the precess of fluid transfer, the use of an inactive gas in the siphon system was necessary. For this purpose, we used compressed nitrogen gas to maintain a minimum pressure differential in the siphon to fill up the cell volume ($`100`$ mL) in a couple of hours.
### B Linear viscoelastic measurements
The linear viscoelastic shear flow properties of the three polymer solutions were measured, at a temperature of $`20^{}`$C, using a Rheometrics Mechanical Spectrometer (RMS-800) in a cone-and-plate geometry with a cone diameter of 40 mm and a gap angle of $`4^{}`$. The measured modulii of complex dynamic viscosities $`|\eta ^{}|`$ of the three liquids are plotted against the angular frequency $`\omega `$ in Fig. 1. We note that the viscosity of PS2 is about an order of magnitude higher than the other two samples; it is a much less mobile liquid. The sample PS81 is a couple of times more viscous than PS82. The typical nature of the three curves seen in Fig. 1 is very similar, namely, $`|\eta ^{}|`$ is almost constant at low frequency $`\omega 3\times 10^2`$ rad s<sup>-1</sup> and beyond this $`\omega `$ it shows a “shear” thinning behavior, following an empirical power law $`|\eta ^{}|=K\omega ^n`$, over the entire experimental frequency range, with the correlation coefficient $`R^2`$ being about $`0.99`$ or higher in all three cases. The power law constant $`K`$ and the exponent $`n`$, obtained from a linear regression of the data are given in Table I. We note that although the values of $`K`$ are significantly different in these samples, the slope of the complex dynamic viscosity versus angular frequency curves in a double-logarithmic plot in Fig. 1 is about same for two solutions with $`N_e13`$ and larger than the slope with $`N_e7`$.
From the low frequency Newtonian plateau of these curves, the zero “shear-rate” viscosity $`\eta _0`$ is extracted for these samples and is shown in Table I. The number of entanglement points per chain, $`N_e`$, can be defined by analogy with the classical theory of rubber elasticity,
$$G_N^0=\left(\frac{cRT}{M_e}\right)\left(\frac{cRT}{M_w}\right)N_e,$$
(1)
where $`c`$ is the polymer concentration in g cm<sup>-3</sup>, $`M_e=M_w/N_e`$ is the molecular weight between entanglements, $`R`$ is the gas constant and $`T`$ is the absolute temperature. We have used the approximate expression for the plateau modulus,
$$G_N^0=3.44\times 10^6\times c^{2.4}\text{dyn/cm}^2(c0.1\text{g cm}^3),$$
(2)
obtained experimentally by Osaki et al. for polystyrene solutions in a good solvent. The values obtained using this equation agrees reasonably well with measured values. Thus Eqns. (1) and (2) leads to
$$N_e=1.41\times 10^4\times M_w\times c^{1.4}$$
(3)
for our experiments and the calculated values of $`N_e`$ for the polystyrene samples are shown in Table I.
The dynamic modulii of PS81, PS82 and PS2 are measured on the same rheometer RMS-800 with the same pair of cone and plate as for the data in Fig. 1 and are shown in Fig. 2. Both, the elastic modulus ($`G^{}`$) and the viscous modulus ($`G^{\prime \prime }`$) show a strong frequency dependence for all three cases. The values of the modulii $`G^{}`$ and $`G^{\prime \prime }`$ for PS81 is higher than that of PS82, both having an order of magnitude lower values than in the case of PS2. Each of these pair of complex modulii shows a pattern of $`G^{\prime \prime }`$ dominating over $`G^{}`$ in the low frequency “terminal zone” and the values of $`G^{}`$ taking over that of $`G^{\prime \prime }`$ in the higher frequency “plateau region”, typical of polymeric fluids. In the low frequency flow region, the storage modulus $`G^{}`$ approaches a quadratic dependence on $`\omega `$ and the loss modulus $`G^{\prime \prime }`$ shows a linear dependence, again a typical behavior of polymers with a narrow molecular weight distribution. At higher frequencies $`G^{}`$ and $`G^{\prime \prime }`$ tend to converge, which may be due to the formation of transient entanglement network. The value of $`\omega `$ at which the storage modulus falls below the loss modulus is highest in PS2 and lowest in PS81. Assuming that this crossing point frequency corresponds to the inverse of the disengagement time $`\tau _d`$, we have estimated the values of $`\tau _d`$ to be $`67.50`$ s, $`54.18`$ s and $`16.80`$ s respectively for PS81, PS82 and PS2. Using the prediction $`\tau _R=\tau _d/(3N_e)`$ from reptation theory, we have calculated the Rouse time, $`\tau _R`$, for each sample as reported in Table I. To check these values, we can use the correlation of Menezes and Graessley
$$\tau _R=\left[\frac{6\left(M_c\right)_0^{a1}}{\pi ^2\rho RT}\right]\frac{\eta _0}{M_w^{(a2)}\nu ^{(a1)b+1}},$$
(4)
where for polystyrene we use $`(M_c)_0=33,000`$, $`a=3.4`$, $`b=1.3`$, density $`\rho =1.07`$ g cm<sup>-3</sup> and volume fraction $`\nu =\rho c`$. The calculated Rouse times are $`1.95`$ s, $`3.84`$ s and $`0.87`$ s for PS81, PS82 and PS2 respectively. It is unclear to us why the values calculated from Eqn. (4) are not in closer agreement with those obtained here experimentally, though they are in the same ballpark. In what follows, we use the experimentally determined values listed in Table I.
### C Experimental set-up
A schematic diagram of our experimental system is shown in Fig. 5. It is composed of two main parts: the TCFB optical arrangement that is used to measure the optical anisotropy of the fluid and the DLS set-up that is used to measure the flow parameters of the fluid. In the following, we will briefly touch upon some specific details of the apparatus relevant to the TCFB and DLS studies reported in this paper. Further details of the data analysis technique used for TCFB and DLS have been described elsewhere.
#### 1 Two-color flow birefringence
The TCFB technique simultaneously measures both retardance and orientation $`\chi `$ of the principal axis of the refractive-index tensor $`\underset{¯}{\underset{¯}{𝐧}}`$ relative to the axes \[$`(x,y)`$ in Figs. 3 and 4\] fixed in the flow device. The light source is an Argon-ion laser (Spectra Physics Model 2020) operating at $`300`$ mW, which emits two intense wavelengths at $`\lambda _B=4880`$ Å(blue) and $`\lambda _G=5145`$ Å(green). As shown in Fig. 5, with the use of several optical elements, each of these beams (independently polarized at $`45^{}`$ relative to each other) are made to pass along an identical optical path through the sample in a two-roll mill which is placed between crossed polarizers. Since the flow is non-homogeneous, to ensure the collinearity of the two beams (with identical optical properties, e.g., beamwidth, Gaussian beam-profile etc.) passing through the same element of fluid in the flow-cell is very crucial in this set-up. Each of the measured intensities (at the photodetectors) depends upon the current degree of optical anisotropy of the sample in the plane of flow \[$`(x,y)`$ plane in Fig. 3\], — i.e., the birefringence, $`\mathrm{\Delta }n=n_{}n_{}`$ — and $`\chi `$, which provides a measure of the average orientation of Kuhn segments in the polymeric liquid. Further details of the optical system are described in Ref. and will thus not be repeated here. The measured overall error due to nonidealities of the optical components indicate that the maximum extinction ratios detectable in this set up with the flow device (loaded with polymer solution), for both colors, are typically $`𝒪(2\times 10^5)`$.
#### 2 The two-roll mill and the flow
The two-roll mill is a flow device consisting of two cylinders driven simultaneously by a single DC stepping motor (Superior Electric, U.S.A., SLO-SYN motor type MO62-FD-09). The cylinders rotate at the same angular velocity and in the same direction to generate a flow between the roller pair, that can be approximated at the central stagnation point by a linear (or homogeneous), planar velocity-field of the form
$$v=\stackrel{}{𝐯}\stackrel{}{r}.$$
(5)
Here, $`\stackrel{}{r}`$ is the position vector defined in the flow plane w.r.t. the $`(x,y)`$ frame, shown in Figs. 3 and 4. The velocity-gradient tensor $`\stackrel{}{𝐯}`$ in this coordinate system is given by
$$\stackrel{}{𝐯}=\dot{\gamma }\left[\begin{array}{cc}ϵ& 1\\ \lambda & ϵ\end{array}\right],0<\lambda 1.$$
(6)
If we assume further that the flow is symmetric about the central $`(x,z)`$ plane between the rollers, as in the case for a Newtonian fluid at zero Reynolds number, then
$$ϵ0(\text{symmetric flow}).$$
(7)
In that case, Eqn. (6) can be completely characterized by two scaler parameters, namely, the magnitude of the velocity-gradient $`\dot{\gamma }`$ $`|\stackrel{}{𝐯}|`$ (also called the “shear-rate”) and the flow-type parameter $`\lambda `$. The parameter $`\lambda `$ is a measure of the ratio of strain-rate to vorticity, and is defined as,
$$\frac{\underset{¯}{\underset{¯}{𝐄}}}{\underset{¯}{\underset{¯}{𝛀}}}=\frac{1+\lambda }{1\lambda }.$$
(8)
Thus, $`\lambda =0`$ corresponds to a simple shear flow, $`\lambda =1`$ to a purely extensional (also called hyperbolic) flow, and intermediate values of $`\lambda `$ represent a “mixed” shear and elongational flow-type (also called a “strong” or extension dominated flow, since the strain-rate exceeds the vorticity). In sharp contrast to other extensional flow devices, polymer molecules at the stagnation region of a two-roll mill have very long residence time and hence are subjected to very large total strains. Consequently, they can, in principle, achieve the maximum change in conformation that is consistent with a particular strain-rate. For a Newtonian fluid, at very low Reynolds numbers, where the symmetric flow assumption is valid, the creeping flow solutions for a two-roll mill in an unbounded fluid approximately relate the values of the corresponding velocity-gradient $`\dot{\gamma }_N`$ and the flow-type parameter $`\lambda _N`$ in the stagnation region with the geometry of the cell (i.e., the roller radius $`R`$ and the gapwidth $`2h`$), via
$`\lambda _N=\left({\displaystyle \frac{4\mathrm{coth}K_f}{K_f}}1\right)^1,\dot{\gamma }_N={\displaystyle \frac{A\omega }{\lambda _N}}`$ (9)
$`\text{with}A={\displaystyle \frac{K_f}{K_f+\mathrm{sinh}K_f\mathrm{cosh}K_f}},`$ (10)
where $`\omega `$ is the angular velocity of the rollers and the parameter $`K_f`$ is given by
$$1+\frac{h}{R}=\mathrm{cosh}K_f.$$
(11)
The symmetry axes of the flow-field, the principal optical axes of the solution, as well as the relative orientation of the blue polarizer and analyzer for the two-roll mill set is shown in Fig. 3. It may be noted that the Newtonian value of the flow-type parameter, $`\lambda _N`$, is also related to the acute angle of crossing, $`2\phi `$, of the streamlines passing through the stagnation point, according to
$$\lambda _N=\mathrm{tan}^2\phi .$$
(12)
The present two-roller can, in general, be used with a pair of rollers chosen from a set of eight such pairs of different diameters covering $`0\lambda 0.25`$. For this study, the dimensions of the rollers used are specified in Table II, which corresponds to $`\lambda =\lambda _N=0.1501`$ for a Newtonian fluid. The subscripts “th” and “exp” in Table II refer to the theoretical and experimental values, respectively. Eqns. (10) and (11) allow us to estimate the theoretical value of $`(\dot{\gamma }_N/\omega )_{\text{th}}`$ for this set of rollers, as given in Table II.
The stepping motor is interfaced to a computer (Hewlett Packard, model 9133 for TCFB or PC-486 for DLS) via a GPIB (National Instruments, U.S.A., model NI-488.2) switch board as shown in Fig. 5. The motor has a fast response time of $`𝒪(10`$ ms) and a maximum acceleration of 100 000 steps s<sup>-2</sup>. Gears with reduction ratios $`5:1`$ or $`20:1`$ are used between the motor and the flow device to cover a wide range of velocity-gradients. The accuracy involved in repositioning of the flow cell with respect to the incident beam, even when it is dismounted for the purpose of replacing the solution, is always better than $`0.0025`$ cm. in $`x`$, $`y`$ and $`z`$ directions (as defined by three translation stages used to control the respective movements) and $`0.001`$ rad in the azimuthal orientation $`\varphi `$ of the flow-cell (defined by the micrometer used to control the orientation) \[Fig. 4\]. The flow-cell was thermostated within $`\pm 0.2^{}`$C via a temperature-regulated waterbath circulator.
#### 3 Dynamic light scattering
The optical setup for the dynamic light scattering experiment was built around the TCFB apparatus, as can be seen from Fig. 5, by mounting the necessary optical accessories on a triangular optical rail placed on the rotating arm of a massive goniometer that defines the scattering angle, $`\theta `$, and hence the scattering vector $`\stackrel{}{q}`$ \[$`q|\stackrel{}{q}|=\frac{4\pi n}{\lambda _B}\mathrm{sin}(\theta /2)`$ as shown in Fig. 4, where $`n`$ is the refractive-index of the solution\]. The green beam of the laser is blocked using a beam stop near P<sub>G</sub>, in Fig. 5. The incident blue beam scattered by the sample is polarized and collimated to project an image of the scattering volume at the photomultiplier tube (PMT) \[Hamamatsu, model R647-04\]. A blue line filter, F<sub>B</sub>, obstructs any spurious light from reaching the detector. The preamplified and discriminated PMT pulses are fed to a 72-channel correlator (Brookhaven Instruments, U.S.A., model BI-2030) to construct time-resolved intensity autocorrelation functions. The commercial correlator software has been modified and used in PC-486 computer to control the flow experiment, as well as to analyze the correlation functions.
Provided that the seed particles are isotropic scatterers, and the time scale associated with the velocity-gradient \[$`t_{\dot{\gamma }}(q\dot{\gamma }L)^1`$, where $`L`$ is the laser beamwidth\] is much smaller than the time scale for diffusive motion \[$`t_D(Dq^2)^1`$\], the homodyne intensity correlation function for a general linear flow in the beam coordinates $`(x^{},y^{},z^{})`$ \[Fig. 4\] is given by
$$F_2(\stackrel{}{q},t)=\beta ^{}\left|d^3\stackrel{}{r}^{}I(\stackrel{}{r}^{})\mathrm{exp}\{i\stackrel{}{q}^{}\stackrel{}{𝐯}^{}\stackrel{}{r}^{}t\}\right|^2.$$
(13)
Here, $`\beta ^{}`$ is the coherence factor defined by the optical geometry. The velocity-gradient tensor $`\stackrel{}{𝐯}`$ of Eqn. (6) can be expressed in the beam coordinates as
$$\stackrel{}{𝐯}^{}=𝐐\stackrel{}{𝐯}𝐐^T\text{ where }𝐐=\left[\begin{array}{cc}\mathrm{cos}\varphi & \mathrm{sin}\varphi \\ \mathrm{sin}\varphi & \mathrm{cos}\varphi \end{array}\right].$$
(14)
In Fig. 4, the beam coordinates are chosen such that the scattering vector $`\stackrel{}{q}^{}`$ is orthogonal to the $`y^{}`$ axis, i.e.,
$$\stackrel{}{q}^{}\stackrel{}{𝐯}^{}\stackrel{}{r}^{}=q_xf(\varphi )x^{}+q_xg(\varphi )y^{},$$
(15)
where $`q_x=q\mathrm{cos}(\theta /2)`$, $`f(\varphi )=\mathrm{cos}\varphi (ϵ\mathrm{cos}\varphi \dot{\gamma }\lambda \mathrm{sin}\varphi )\mathrm{sin}\varphi (\dot{\gamma }\mathrm{cos}\varphi +ϵ\mathrm{sin}\varphi )`$ and $`g(\varphi )=\mathrm{sin}\varphi (ϵ\mathrm{cos}\varphi \dot{\gamma }\lambda \mathrm{sin}\varphi )+\mathrm{cos}\varphi (\dot{\gamma }\mathrm{cos}\varphi +ϵ\mathrm{sin}\varphi )`$.
Using a Gaussian intensity profile, Eqns. (14) and (15) in Eqn. (13), we get
$$F_2(\stackrel{}{q},t)=\beta \mathrm{exp}\left\{\frac{1}{2}L^2q^2t^2\mathrm{cos}^2\left(\frac{\theta }{2}\right)h(\varphi )\right\},$$
(16)
where $`\beta \beta ^{}L^6\mathrm{csc}^2\theta `$ and
$$h(\varphi )=f^2(\varphi )+g^2(\varphi ).$$
(17)
The values of $`h(\varphi )`$ are evaluated in Table III for three different azimuthal orientations of the two-roll mill. Thus, with the assumption of a symmetric flow \[Eqn. 7\], $`\dot{\gamma }`$ can be obtained from the decay rate $`h(\varphi )`$ of the measured correlation function \[Eqn. (16)\] at one orientation, $`\varphi =0`$, of the flow-cell and $`\lambda `$ from an additional measurement at a second orientation, $`\varphi =90^{}`$ (see, Table III). The correlation functions were accumulated only after the flow had attained its steady value for the corresponding motor speed. To improve the signal to noise ratio, each correlation function was obtained by averaging over many repetitions of the experiment, ranging from a minimum of 200 repetitions at the high roller speeds, up to a maximum of 800 repetitions at the low roller speeds. Using a simulated annealing Monte Carlo fitting procedure, these average correlation functions were then fitted to Eqn. (16). In order to verify that the experimentally obtained correlation functions are very close to exponential in nature so that the above procedure followed to extract the velocity-gradient components from the decay time is indeed justified, we required a correlation coefficient $`R^2`$, specifying the quality of the fit, exceeding $`0.99`$.
Contrary to the basic assumption of optical isotropy of scatterers used to derive Eqn. (16), Wang et al had shown that this equation could also be used for the scattering from (intrinsically anisotropic) polymer molecules. In that case, the pre-exponential factor $`\beta `$ is shown to be directly related to the components of the intrinsic molecular polarizability tensor.
## III Results
### A DLS steady flow results
#### 1 Flow characterization
The DLS technique provides us the opportunity to map the flow-field in a large region between the rollers. In order to look for changes in the steady flow-field due to viscoelasticity in the presence of polymers compared to that seen with a Newtonian fluid, we first characterize the flow using a Newtonian fluid. This provides a comparison with predictions from the creeping flow theory for an unbounded two-roll mill. The measured velocity-gradient $`\dot{\gamma }`$ along the $`x`$ axis through the stagnation point of the two-roll mill orientated at $`\varphi =0^{}`$ for an angular velocity $`\omega =1.57`$ rad/s is plotted along with two sets of data obtained by Wang et al, in terms of $`\dot{\gamma }/\omega `$ in Fig. 6. The filled symbols denote the points where the geometric construction of the flow-cell did not allow measurements, and those points are therefore obtained by the symmetric reflection of the hollow measured points at negative $`x`$ values. As can be clearly seen from this figure, different sets of experimental results obtained at widely separated times with different roller rotation rates overlap perfectly. In a small region surrounding the stagnation point the velocity-gradient were approximately constant, $`\dot{\gamma }/\omega `$ = $`2.8`$, certifying that the flow-field in this region is nearly homogeneous. Also, this value compares very well with the theoretically expected value of $`\dot{\gamma }_N/\omega `$ = $`2.7`$ (Table. II). From here onwards, we shall use $`\dot{\gamma }_N`$ and $`\lambda _N`$ in the text to specify the values of $`\dot{\gamma }`$ and $`\lambda `$ \[Eqns. (10) and (11)\] at the stagnation point of a two-roll mill filled with a Newtonian fluid, and the imposed motor speed will be measured in terms of the corresponding Newtonian strain-rate $`\dot{\gamma }_N\sqrt{\lambda }_N`$ .
In Fig. 7, we show the theoretical predictions for $`\dot{\gamma }_N`$ and $`\dot{\gamma }_N\lambda _N`$ \[Eqn. (10), (11) and Table II\] versus $`\dot{\gamma }_N\sqrt{\lambda }_N`$ in the form of straight lines. The experiments were carried out for 14 different Newtonian strain-rates, and for both parallel and perpendicular orientations of the sample cell. To extract the experimental values of $`\dot{\gamma }`$ and $`\dot{\gamma }\lambda `$ from the measured decay rates of the correlation function of Eqn. (16) at $`\varphi =0^{}`$ and $`90^{}`$ respectively, we required a value for the beamwidth $`L`$. With an initial guess of $`L=29`$ $`\mu `$m, that was reported for a previous set of experiments from this laboratory, we used a simulated annealing Monte Carlo fitting technique to find via the fit shown in Fig. 7 that $`L=32`$ $`\mu `$m yields the best match of the theoretical straight lines for both $`\dot{\gamma }`$ and $`\dot{\gamma }\lambda `$ over the entire range of experimental strain-rates. Compared to the previous experiments, we have performed the present experiments including much higher motor speeds. It is clearly seen in the figure that the corresponding data agrees well with the creeping flow solution. The inset shows that the extracted value of the flow-type parameter $`\lambda `$ also maintains a constant value as expected.
In a similar manner, the flow parameters are measured for the three polymeric samples, PS81, PS82, and PS2, subjected to steady-state flow conditions. The results are plotted in Figs. 8, 9, and 10 versus the non-dimensionalized deformation rate, $`Wi_R=`$ $`\dot{\gamma }\sqrt{\lambda }\tau _R`$ , called the Weissenberg number (based on the Rouse time, $`\tau _R`$, and measured values of $`\dot{\gamma }`$ and $`\lambda `$). The Weissenberg number is defined as the ratio of a characteristic relaxation time for the polymeric fluid to a characteristic time for deformation. It specifies the ability of a flow to generate departures of the polymer configuration from its static equilibrium value. The same data are also plotted against $`\dot{\gamma }_N\sqrt{\lambda }_N`$ , in the inset of these figures, for a comparison with the Newtonian values shown by the straight lines. As expected, there is a substantial reduction in both $`\dot{\gamma }`$ and $`\dot{\gamma }\lambda `$ from the corresponding Newtonian value. The deviation from the Newtonian values increases with the increasing roller rotation rate. Somewhat unexpectedly, however, both $`\dot{\gamma }`$ and $`\dot{\gamma }\lambda `$ still appear to increase approximately linearly with the Newtonian strain-rate. It is also noteworthy that this increase is very similar for the two samples, PS81 and PS2, which have the similar number of entanglements ($`N_e13`$) per chain (see, Table I). On the other hand, PS82, with $`N_e7`$ shows a different slope for its almost linear increment with increasing $`\dot{\gamma }_N\sqrt{\lambda }_N`$ . When plotted against the Weissenberg number (Figs. 8 and 9), the qualitative feature of the curves for $`\dot{\gamma }`$ and $`\dot{\gamma }\lambda `$ appear to be quite different than that against $`\dot{\gamma }_N\sqrt{\lambda }_N`$ : they become much more linear and owing to an order of magnitude smaller $`\tau _R`$, the increment for PS2 with the increasing rate of deformation is much steeper than the other two samples. Also, $`\dot{\gamma }`$ versus $`Wi_R`$ data for the two samples with similar molecular weight but different $`N_e`$ are similar up to $`Wi_R3`$. The quantitative similarity in deformation rate dependence for the samples with similar $`N_e`$ is more apparent for the flow-type parameter $`\lambda `$, as can be seen in Fig. 10. The qualitative nature of this dependence for all three samples appears to be similar, namely, it reaches an almost constant value ($`0.1`$ for the samples with $`N_e13`$ and $`0.05`$ for PS82) at high deformation rates $`Wi_R0.4`$ ($`\dot{\gamma }_N\sqrt{\lambda }_N`$ $`1`$) and at low rates of deformation exceeds the Newtonian values. There is an intermediate transition regime between these these two different behaviors. The parameter $`\lambda `$ attaining values higher than the Newtonian value at the lowest $`Wi_R`$ studied, seems surprising since in the limit of extremely low deformation rate ($`Wi_R0`$), one should expect even a Non-Newtonian fluid to exhibit Newtonian behavior. We are still away from this limit in that the Weissenberg number corresponding to the lowest motor speeds for the steady flow experiments are $`0.05`$ (PS2), $`0.1`$ (PS81) and $`0.2`$ (PS82). As mentioned earlier, each data point in $`\dot{\gamma }`$ and $`\dot{\gamma }\lambda `$ curves are obtained by averaging over several repeated experiments and then the $`\lambda `$ values are obtained by point to point division of the two. Following a similar procedure, we have calculated the error-bars on the $`\lambda `$ values from the standard deviations of the repeated experiments for $`\dot{\gamma }`$ and $`\dot{\gamma }\lambda `$ and plotted in Fig. 11. As expected, the error-bars increase for lower $`Wi_R`$, because of the division of smaller values of the corresponding $`\dot{\gamma }\lambda `$ by $`\dot{\gamma }`$ . Even within the limits of these error-bars we can clearly see that $`\lambda `$ initially exceeds $`\lambda _N`$ . We have confirmed that this finding is not an experimental artifact by two other experimental means employed in this study, namely, the measurements of birefringence and flow parameters in transient flow conditions, and is described in section III.A.3 of this paper.
#### 2 Flow symmetry
As we have mentioned before, the assumption that the flow in the two-roll mill is symmetric about the plane passing exactly midway between the co-rotating rollers $`(x,z)`$ plane in Figs. 3 and 4) is valid for the creeping flow of a Newtonian fluid. For a viscoelastic fluid, one might expect the symmetry of the flow to change compared to the Newtonian case. Using the measured values of $`\dot{\gamma }`$ and $`\dot{\gamma }\lambda `$ (Figs. 8 & 9) at several different steady Newtonian strain-rates, we have calculated the expected values of $`\dot{\gamma }\sqrt{(1+\lambda ^2)}`$ . Then, by repeating exactly same steady flow conditions that were used to measure $`\dot{\gamma }`$ and $`\lambda `$, we have extracted the velocity-gradient component $`\dot{\gamma }\sqrt{(1+\lambda ^2)}`$ from the decay rates of the correlations functions measured directly with the two-roll mill oriented at $`\varphi =45^{}`$ (see, Table III). The result is plotted versus the corresponding calculated values at the same Newtonian strain-rates with the connected hollow symbols in Fig. 12 for the three viscoelastic samples. The hollow squares represent the measured values of $`\dot{\gamma }\sqrt{(1+\lambda ^2)}`$ for the Newtonian fluid at $`\varphi =45^{}`$. Both these and the theoretical values, i.e., $`\dot{\gamma }_N\sqrt{(1+\lambda _N^2)}`$ (shown by the solid line) at several strain-rates are plotted in Fig. 12 against the expected values of the same parameter that are calculated using the measured $`\dot{\gamma }`$ and $`\dot{\gamma }\lambda `$ data from Fig. 7. If the symmetry of the flow is maintained, then Eqn. (16) should be strictly valid for all orientations, $`\varphi `$, of the flow-cell. This means that each curve in Fig. 12 should be linear. As can be clearly seen from the figure, this is best followed for the Newtonian fluid indicating the fact that the flow is symmetric in this case, as expected. On the other hand, the data for the entangled fluids indicate that the flow symmetry is maintained only approximately, i.e., within the experimental accuracy of $`10\%`$, for the first five or six data points (i.e., for $`\dot{\gamma }_N\sqrt{\lambda }_N`$ $`<0.5`$ s<sup>-1</sup> ). The data for PS81 and PS2 (both with $`N_e13`$) deviate more from the linearity with the increase of $`\dot{\gamma }_N\sqrt{\lambda }_N`$ , but for PS82 ($`N_e7`$) the experimental data are fairly close to a straight line over the entire range of strain-rates studied. At high rates of deformation, the measured $`\dot{\gamma }\sqrt{(1+\lambda ^2)}`$ is higher than the expected value for PS81 and PS82, but is lower in the case of PS2.
In order to quantify how flow symmetry changes with the deformation rate, we now proceed to evaluate the parameter $`ϵ`$ of Eqn. (6). To do this, we can first extract the values for $`h(\varphi =0^{})`$, $`h(\varphi =45^{})`$ and $`h(\varphi =90^{})`$ from the experimentally measured decay rates of the correlation function at the three aforementioned orientations of the two-roll mill. When these values are used in conjunction with Table III, we get three equations involving the three unknowns $`\dot{\gamma }`$ , $`\lambda `$ and $`ϵ`$. We have used the Newton-Raphson’s method to numerically solve these equations for the case of a Newtonian fluid as well as for a representative case of the viscoelastic fluids of our study, namely, for PS2. The results are shown in Fig. 13. Our experiments with a Newtonian fluid shows that $`ϵ`$ is zero, i.e., the flow is perfectly symmetric, for strain-rates upto about 2 s<sup>-1</sup> . Then, as the deformation rate is increased the value of $`ϵ`$ increases almost linearly, though its value remains extremely small ($`10^5`$) even at the highest strain-rates studied, justifying the fact that, for a Newtonian fluid the flow-symmetry is maintained. This also explains why the first seven data points ($`\dot{\gamma }_N\sqrt{\lambda }_N`$ $`<2`$ s<sup>-1</sup> ) for the Newtonian fluid in Fig. 12 lies on the theoretical solid line and then slightly deviates from the line for the higher values of $`\dot{\gamma }_N\sqrt{\lambda }_N`$ . For a non-Newtonian fluid, on the other hand, the flow-symmetry is maintained for only up to $`\dot{\gamma }_N\sqrt{\lambda }_N`$ $`<0.5`$ s<sup>-1</sup> . Our transient flow experiments with these polymer solutions are carried out in this range of deformation rates (which corresponds to $`Wi_R<1`$) and will be reported in a subsequent paper.
#### 3 Steady flow and flow-type parameter
Let us now look into the effect of the polymer solutions on the steady-state values of the flow-type parameter $`\lambda `$ at different angular velocities of the rollers. We have briefly referred to this point in Fig. 10. Since, we have also performed extensive DLS and TCFB experiments on these samples subjected to transient flows, namely, the startup flows from rest, for several constant Weissenberg numbers ($`Wi_R<1`$), it would be worthwhile to check how the asymptotic values of $`\lambda `$ extracted from each of those different experimental techniques vary with $`Wi_R`$ and also how they compare with the steady flow results of Fig. 10. From the transient DLS experiments, we can directly get the asymptotic values of $`\lambda `$ but for the TCFB experiments, we should follow an indirect method to calculate $`\lambda `$. This can be done because for $`Wi_R<1`$ the flow retains its symmetry, as we have shown before. The steady-state orientation of the principal optic axis at appreciable strain-rates is expected to approach the outflow axis of the flow, or, in other words, $`\chi \chi ^{}\phi `$ \[see, Fig. 3\]. Hence, $`\lambda `$ may be calculated from the asymptotic orientation angle, $`\chi ^{}`$, of the flow birefringence data for the startup flows, via a relation similar to Eqn. 12, i.e.,
$$\lambda \mathrm{tan}^2\chi ^{}.$$
(18)
We have presented the results in Figs. 14, 15 and 16 for the samples PS81, PS82, and PS2 respectively and compared the results with the theoretical value of $`\lambda _N=0.1501`$. For each of these three polymeric samples, the overall qualitative match of the data as well as the quantitative match between the $`\lambda `$ values calculated using three different experimental procedures, is quite satisfactory. The fair quantitative agreement between the asymptotic values of $`\lambda `$ obtained from the transient flow experiments and the steady flow data for $`\lambda `$ confirms that the startup experiments with the entangled samples have almost reached their steady values in the total predetermined evolution time of $`t_e=30`$ s. In agreement with the steady-state result, the transient experiments also show $`\lambda `$ values exceeding $`\lambda _N`$ at the lowest rates of deformation studied, which again are higher than $`Wi_R=0`$, where we should expect that $`\lambda =\lambda _N`$. In fact, for PS81, at the three lowest rates of deformations studied, the extracted value of $`\lambda `$ from the transient experiments do show a trend to decrease after reaching a maximum.
### B The DEMG model comparison of the steady flow TCFB results
Probably one of the most used theoretical models to describe the dynamics of entangled polymers is the Doi-Edwards (DE) tube (reptation) model. The primary drawback of this model, as is clearly evident by its failure to predict the polymer dynamics at high deformation rates $`\dot{\gamma }\sqrt{\lambda }`$ $`>1/\tau _R`$, is the fact that the primitive chain is assumed to be inextensible. Subsequently, in an effort to improve its predictions, the Doi-Edwards-Marrucci-Grizzuti (DEMG) model was developed to incorporate chain-stretching into the original DE tube model. As noted earlier, the DEMG model contains two widely separated time scales, one for the relaxation of orientation ($`\tau _d`$) and a much shorter one for the relaxation of the stretch of polymer segments ($`\tau _R`$). Thus, as the strain-rate $`\dot{\gamma }\sqrt{\lambda }`$ is increased, segmental orientation takes place first when $`\dot{\gamma }\sqrt{\lambda }`$ $`1/\tau _d`$ or $`Wi_d=`$ $`\dot{\gamma }\sqrt{\lambda }\tau _d`$ $`=𝒪(1)`$ and segmental stretching then “starts” at a much larger strain-rate when $`Wi_R=`$ $`\dot{\gamma }\sqrt{\lambda }\tau _R`$ $`=𝒪(1)`$. Also, it is expected that segmental orientation will take place earlier in dilute systems than in entangled systems. In agreement with the aforesaid expectations, a number of numerical studies have suggested that the DEMG model has more success over the DE model when compared to experimental results for startup of a simple shear flow at high deformation rates ($`Wi_R1`$). Following this, in a recent numerical study, the DEMG model was developed further by “adding” a non-linear finitely extensible spring force.
For an entangled solution, birefringence measurements directly relate to changes in the conformation of the polymer molecules. As mentioned earlier, the complete characterization of the flow-induced anisotropy of a polymeric fluid requires determination of the birefringence, $`\mathrm{\Delta }n`$, as well as the orientation angle, $`\chi `$, of the principal axes of the refractive index tensor. In this section, we present the steady-state two-color flow birefringence results and compare them with the predictions of the DEMG model using the flow parameters measured via DLS as input. We assume that the velocity-gradient tensor in $`75\mu `$m diameter birefringence measurement zone surrounding the stagnation point of the co-rotating two-roll mill can be approximated by Eqn. (6) and (7). The values of the model parameters $`N_e`$, $`\tau _R`$, and $`n_t`$ calculated for the three solutions are given in the Table I. Here $`n_t`$ is the number of Kuhn statistical subunits in the chain based on the assumption of 10 monomers per subunit and is calculated from the molecular weight of the polymer. The steady-state experimental results for the dimensionless birefringence, $`\mathrm{\Delta }n/(CG_N^0)`$, for the three entangled samples are presented in Fig. 17 as functions of $`Wi_R`$. The birefringence, $`\mathrm{\Delta }n`$, is scaled by the birefringence $`CG_N^0`$ (obtained by using the stress-optical law) that would be present at a stress level equal to the plateau modulus, $`G_N^0`$. Here, $`C`$ is the stress-optical coefficient which in the free-jointed chain model is defined as
$$C=\frac{2\pi }{45}\left[\frac{(\overline{n}+2)^2}{\overline{n}}\right]\alpha _{}\alpha _{},$$
(19)
where $`\alpha _{}`$ and $`\alpha _{}`$ are the components of the intrinsic molecular polarizability tensor $`\underset{¯}{\underset{¯}{\alpha }}(t)`$ along and perpendicular to the Kuhn statistical segment, and $`\overline{n}`$ is the bulk refractive-index of the medium. This is justified only up to a modest level ($`50\%`$) of fractional chain-extension, so that non-Gaussian effects can be neglected. To scale the experimental birefringence, an approximate expression given in Eqn. (2) for the plateau modulus, $`G_N^0`$, and the literature value of $`C=5\times 10^{10}`$ cm<sup>2</sup>/dyn for polystyrene solutions is used. The model predictions corresponding to the experimental curves in Fig. 17 are displayed in Fig. 18 (both on the same scale). Each individual point on the predicted curves of Fig. 18 is the steady-state value obtained from separate runs of DEMG model numerical simulations using the measured values of $`\dot{\gamma }`$ and $`\lambda `$ for each corresponding point on the experimental curves (Fig. 17). The model clearly predicts the existence of three flow regimes in the steady-state, namely, the linear viscoelastic regime seen at low rates of deformation $`Wi_d<𝒪(1)`$, the non-linear viscoelastic regime seen at the intermediate deformation rates lying between $`Wi_d>𝒪(1)`$ and $`Wi_R<𝒪(1)`$, and the highly non-linear viscoelastic regime seen at high deformation rates $`Wi_R>𝒪(1)`$. In each of these flow regimes, the effect of three parameters, namely, $`n_t`$, $`N_e`$, and $`\lambda `$ on the nature of deformation rate dependence of the birefringence, orientation angle, stretch and viscosity, as predicted by the DEMG model, has been discussed in detail in a recent publication by Mead et al..
In the linear viscoelastic regime, the slope of each experimental curve on the double-logarithmic plots (Fig. 17) is indeed unity, similar to the predicted curves (Fig. 18). According to the predictions of Ref. , in this flow regime, the relative positions of the birefringence curves from top to bottom for these entangled samples is primarily determined by the values of $`N_e`$ and $`\lambda `$ (curves with lower $`N_e`$ and higher $`\lambda `$ will fall lower), while $`n_t`$ has no effect. As clearly seen from Figs. 17 and 18, the lower value of $`N_e`$ for PS82 keeps the $`\mathrm{\Delta }n/(CG_N^0)`$ curve lower than those of the other two samples. The lower value of $`\lambda `$ for PS82 in this regime (Fig. 10) would predict otherwise. The slight difference between the curves for PS2 and PS81 (both having $`N_e13`$) seen both in the experimental and predicted plots, is due to the slight difference in the $`\lambda `$ values between the two samples in this regime, as can be clearly seen from Fig. 10. In this flow regime, although the initial low levels of segmental orientation, as predicted in Fig. 20, is indeed seen for the experimental curves in Fig. 19, the rate of decrease of $`\chi `$ with increasing $`Wi_R`$, as seen in the experiments for all samples, is much faster than predicted. Also, the orientation angles for PS81 at a fixed value of $`Wi_R`$ in this regime lies in between those for the other two liquids which again is in accord with the theoretical expectations (taking into account the effects of the corresponding values of $`n_t`$, $`N_e`$ and measured $`\lambda `$).
When the second flow regime of non-linear viscoelasticity is approached, $`Wi_d=𝒪(1)`$, the scaled birefringence curves are predicted to depart from linearity and begin to approach a plateau, as seen in the DEMG results of Fig. 18. The width of the plateau is directly proportional to $`N_e`$, and is therefore expected to be about half as wide for PS82 ($`N_e7`$) as for PS2 and PS81 ($`N_e13`$). The experimental curves in this regime do show non-linearity for all three samples, but the collapse of the curves and an existence of a narrow plateau is seen only for PS81 and PS2 and not for PS82. The dynamics in this regime are controlled by strong segmental orientation of the tube without stretching. Referring to Fig. 19, we see that the monotonic decrease of the orientation angle with increasing $`Wi_R`$, seen in experiments for all samples is in qualitative agrement with the predictions of the DEMG model, resulting from the gradual unwinding and straightening out of the tube. The relative positions of the $`\chi `$ versus $`Wi_R`$ curves from top to bottom in Fig. 20 for the three samples are the result of a competition between the effects of $`N_e`$ and $`\lambda `$ them. Having a lower value of $`\lambda `$ (see, Fig. 10) prompts the orientational angle for PS82 to fall below those for the other two samples, but as can be seen in Figs. 19 and 20, the effect of $`N_e`$ is stronger in this case too, since the curves with a lower value of $`N_e`$ (that for PS82, here) are predicted to show a lower degree of orientation at any fixed $`Wi_R`$ in this flow regime. Similarly, PS2 exhibits a higher degree of orientation than for PS81 at a fixed $`Wi_R`$, since $`\lambda `$ is lower for PS2 in this intermediate regime of flow. As $`Wi_R1`$, the tube is expected to straighten out to its full length and become orientated in the direction of the outflow axis \[Fig. 3\] with the orientation angle reaching its asymptotic value, $`\chi ^{}`$, given in Eqn. (18). Using the asymptotic values of $`\lambda 0.1`$ for PS2 and PS81 and $`\lambda 0.05`$ for PS82 from Fig. 10, we obtain $`\chi ^{}17.55^{}`$ for PS2 and PS81 and $`\chi ^{}12.60^{}`$ for PS82 which are in agrement with the results in Fig. 20. Interestingly, the theoretical predictions for the orientation angles for PS2 and PS81 in Fig. 20 show an “undershoot” before reaching $`\chi ^{}`$. For the highest rates of deformation, $`Wi_R1.5`$, studied for PS82, the predicted orientation angle has not reached its asymptotic value. The experimental curves of $`\chi `$, as seen in Fig. 19, qualitatively follow the behavior shown by the predicted curves except for that they do not show an undershoot behavior. In addition, $`\chi ^{}17.5^{}`$ and $`\chi ^{}16^{}`$ for PS2 and PS81, respectively and although the final value of $`\chi `$ at the highest $`Wi_R`$ studied for PS82 is $`17.4^{}`$, it does not seem to have saturated; and for $`0.3<Wi_R1.5`$ the $`\chi `$ values for PS81 and PS82 overlap.
The third flow region of highly non-linear viscoelasticity ($`Wi_R1`$) is predicted to show three distinct signatures. Firstly, the tube orientation should be complete as indicated by $`\chi =\chi ^{}`$. We have discussed about the behavior of the predicted and experimental orientation angles for all polystyrene samples in the preceding paragraph. Secondly, the onset of chain-stretching is predicted to take place with a clear evidence of the dimensionless birefringence exceeding its plateau value of unity as shown in Fig. 18. This is evident in the experiments too, as displayed in Fig. 17. Thirdly, the onset of chain-stretching is also marked by the collapse of all the birefringence curves to a single universal curve at the end of the plateau i.e., they should follow the scaling behavior $`\mathrm{\Delta }n/(CG_N^0)=f(Wi_R)`$ in this region, as is clearly shown by the predicted traces in Fig. 18. This feature is well demonstrated in the experiments (Fig. 17) by PS81 and PS2 but not by PS82, which fails to show a well-defined plateau too. At high rates of deformation ($`Wi_R>1`$), the chain-stretching dynamics prompts the different birefringence curves to reach different asymptotes proportional to $`n_t/N_e`$. The theoretical curves in Fig. 18 show these features: the increase in the birefringence with $`Wi_R`$ in this regime is faster for PS82 ($`n_t=8420`$ and $`N_e7`$) than for PS81 ($`n_t=8420`$ and $`N_e13`$), which is in turn faster than for PS2 ($`n_t=2890`$ and $`N_e13`$). However, none of the $`\mathrm{\Delta }n/(CG_N^0)`$ curves are predicted to reach their asymptotic values ($`648`$ for PS81, $`1203`$ for PS82 and $`222`$ for PS2) for the highest $`Wi_R`$ studied, though the curves for PS81 and PS82 indeed show a tendency towards saturation in Fig. 18. At the highest rates of deformations, the experimental birefringence levels (Fig. 17) for all samples are significantly below their saturation values, and they do not show a tendency to saturate. Contrary to the prediction, the experimental $`\mathrm{\Delta }n/(CG_N^0)`$ curve for PS81 falls below that for PS2. We note that the relative level of birefringence for the different samples as depicted by the individual curves in Figs. 17 and 18 are strongly dependent on the choice of the value of the relaxation modulus, $`G_N^0`$ \[Eqn. (2)\]. However, we believe that the qualitative as well as quantitative agreement between the experimental and predicted birefringence curves well justifies the choice of value for $`G_N^0`$.
Following Mead et al., we define the generalized viscosity function (which is the conventional viscosity appropriately generalized for mixed shear and extensional flow) as follows:
$$\eta (\dot{\gamma },\lambda )=\frac{\sigma _{xy}}{\dot{\gamma }(1+\lambda )}.$$
(20)
This reduces to the usual definition of the shear viscosity and the planar extensional viscosity in the limits of $`\lambda =0`$ and $`\lambda =1`$, respectively. The scaling factor $`(1+\lambda )`$ makes the viscosity independent of the flow-type $`\lambda `$ at low strain-rates. In the principal axes of the rate of deformation tensor, which are at $`45^{}`$ from the coordinate axes used in Figs. 3 and 4, we get $`\eta (\dot{\gamma },\lambda )=\frac{\sigma _{xx}^{}\sigma _{yy}^{}}{\dot{\gamma }(1+\lambda )}`$. Using the stress-optical relationship, which states that the birefringence tensor is proportional and coaxial to the stress tensor, one can relate the birefringence, which is the difference of the principal values of the refractive-index tensor, $`\underset{¯}{\underset{¯}{𝐧}}`$, in the plane of the flow (i.e., $`\mathrm{\Delta }n=n_{xx}^{}n_{yy}^{}`$), to the stress as $`\sigma _{xy}=\frac{\mathrm{\Delta }n}{2C}\mathrm{sin}(2\chi )`$, so that Eqn. (20) can be written as
$$\eta (\dot{\gamma },\lambda )=\frac{\mathrm{\Delta }n\mathrm{sin}(2\chi )}{2C\dot{\gamma }(1+\lambda )}.$$
(21)
We present the predicted and the experimental steady-state viscosity function versus the measured $`Wi_R`$ for the three polystyrene solutions in Figs. 21 and 22, respectively. The predicted values of $`\eta `$ are calculated from Eqn. (21) by using the predicted birefringence and extinction angle shown in Figs. 18 and 20, and by scaling the viscosity with the zero shear viscosity,
$$\eta _0=\frac{G_N^0\tau _d\pi ^2}{45}.$$
(22)
Similarly the experimental plots shown in Fig. 22 are calculated with the use of the measured $`\mathrm{\Delta }n`$ (Fig. 17) and $`\chi `$ (Fig. 19) and then by normalizing by the experimental $`\eta _0`$ given in Table I. In sharp contrast to the observation that dilute solutions, (e.g., Boger fluids) show strain-rate thickening at high strain-rates, for entangled solutions the viscosity function is predicted to show different behaviors in the three earlier defined flow regimes. In the first flow regime of linear viscoelasticity, the viscosity is predicted to be nearly constant and equal to its zero shear-rate value, $`\eta _0`$, (Fig. 21). This is approximately followed by the experimental curves (Fig. 22). In the intermediate flow regime \[between $`Wi_d>𝒪(1)`$ and $`Wi_R<𝒪(1)`$\], where the stress is generated primarily via orientation of the tube segments, the model predicts a strong strain-rate thinning behavior with the degree of thinning being an increasing function of $`N_e`$ and a decreasing function of $`\lambda `$, as shown in Fig. 21. In the experiments, we see that $`N_e`$ has a stronger effect than $`\lambda `$ on the degree of thinning. The predicted effect of $`\lambda `$ on the relative positions of $`\eta /\eta _0`$ versus $`Wi_R`$ curves for PS81 and PS2 from top to bottom is not supported in the experiments.
The presence of an “undershoot” in the predicted orientation angle for PS2 (Fig. 20) in this flow regime clearly manifests itself in the predicted viscosity plot (Fig. 21) too. This does not seem to be the case with PS82, because for PS82 the small undershoot present in the predicted orientation angle does not show up in the predicted viscosity. Although, the experimental viscosity function shows a thinning behavior with the degree of thinning being much weaker in PS82 than the other two solutions, the model predicts much stronger viscosity thinning behavior than observed in the experiments (except for PS81 where they are similar). Also, the expected “undershoot” in the viscosity for PS2 is absent in the experiment.
In the third regime of flow \[$`Wi_R=𝒪(1)`$\] the model predicts a sharp upturn in the viscosity function as chain-stretching is predicted to take place. As the rate of deformation is increased, the viscosity is also expected to increase and finally reach an asymptote for $`Wi_R>>1`$. Surprisingly, however, except for a small upturn shown by the last three data points for PS2 in Fig. 22, these predicted features are completely absent in the experimental data. The cause for this lies in the fact that even though the experimental birefringence values of these samples (Fig. 17) rise above the plateau value, indicating chain-stretching, the rise is much weaker than predicted (Fig. 18). Also, the degree of orientation of the tube segments in the experiments (Fig. 19) fails to show the sharpness predicted by the model (Fig. 20) in this flow regime.
The above two effects compete with each other to nullify any signature of chain-stretching that could have been otherwise seen in the generalized viscosity function. We have noted before that the present model has improved the original DEMG model by using finitely extensible freely jointed chains instead of infinitely extensible Gaussian chains. In the limit of the high value of $`n_t`$ (or $`M_w`$), the original version of the model is retrieved and at high deformation rates the chains stretch tremendously to a show a nearly singular viscosity or birefringence. This effect is apparent in the model predictions shown in Figs. 21 and 18, where both PS81 and PS82 (higher $`n_t`$ or $`M_w`$) show a steeper increase in $`\eta `$ and $`\mathrm{\Delta }n/(CG_N^0)`$ at $`Wi_R1`$) compared to PS2, but is absent in the experiments (Figs. 22 and 17). Unfortunately, the highest dimensionless rate of deformation or $`Wi_R`$ that we could reach in the experiments is only about 3 for PS81 and about 1 for the other two solutions. This deficiency is caused by the fact that the viscoelastic modification relative to the Newtonian flow was much stronger than we had initially expected. The measured values of $`\dot{\gamma }`$ and $`\lambda `$ were thus reduced compared to their Newtonian values, so that the originally calculated theoretical $`(Wi_R)_N^{\text{max}}(\dot{\gamma }_N\sqrt{\lambda }_N\tau _R)^{\text{max}}=𝒪(10)`$ corresponding to the maximum motor speed chosen for the experiments ultimately produced only $`(Wi_R)^{\text{max}}(\dot{\gamma }\sqrt{\lambda }\tau _R)^{\text{max}}=𝒪(1)`$. The absence of chain-stretching signatures in the viscosity may also have to do with the choice of sample compositions which gives rise to low values for the number of entanglements per chain. It would be interesting to carry out further experiments on samples with higher $`N_e`$ in this highly non-linear viscoelastic flow regime at high enough $`Wi_R`$, particularly to look into the chain-stretching effects much more carefully.
Let us now turn our attention into more quantitative comparisons between the steady-state results and the predictions of the DEMG model. We will first consider the case of the orientation angle $`\chi `$ for three samples as shown in Figs. 23, 24, and 25. Since there is no other parameter involved for scaling (e.g., in the case of the birefringence), direct comparisons between the experimental data and model predictions are possible in this case. Although the experimental data follows a similar qualitative decreasing trend in the orientation angle with increasing $`Wi_R`$, as predicted, the quantitative mismatch with the prediction is quite obvious in these figures. In the limit $`Wi_R0`$, it is expected that $`\chi `$ should be $`45^{}`$, corresponding to alignment of the principal optical axis with the principal axis of the rate of strain tensor. The experimental accuracy of both $`\mathrm{\Delta }n`$ and $`\chi `$ at the lowest rates of deformation is limited by the sensitivity of the apparatus in determining a minimum birefringence and its associated orientation. Also, the contribution from the residual glass birefringence is important at these low levels of polymer anisotropy. At low and intermediate $`Wi_R`$ the DEMG model underpredicts the orientation but at higher $`Wi_R`$ it overpredicts the same, and there is a crossing point between these two behaviors. By comparing the results on the orientation angle with those obtained from the DEMG model numerical simulations with identical parameters but with a constant flow-type parameter $`\lambda =\lambda _N=0.1501`$, we have confirmed that the opposite curvatures seen in the predicted curves below and above the crossing point depends on the way $`\lambda `$ changes with $`Wi_R`$. The fact that the DEMG model overpredicts the tendency of the flow to orient polymer chains towards the outflow axis at intermediate and high strain-rates has been observed in earlier experiments and it was speculated that “tube-dilation”, which is not incorporated in the DEMG model, may be responsible for such an effect. In simple terms, this means that the effective dilation of the tube radius is stronger for large rates of deformation, leading to a reduction of $`N_e`$ and hence the lower orientation angle seen in the experiments. Weather or not tube-dilation is present in the system, for $`Wi_R𝒪(1)`$, we expect that the polymer molecules will become oriented along the outflow axis of the flow-field. As we have noted earlier, and as can be seen clearly from these figures, the asymptotic value of $`\chi `$ is approximately reached in experiments for PS81 and PS2, but not for PS82.
The qualitative as well as quantitative match between the experimental and predicted birefringence curves for these entangled samples, shown in Figs. 26, 27 and 28, is quite satisfactory. This firstly points to the fact that the choice of the plateau modulus, as described earlier in this section, works very well for these entangled systems. Given the proper values of $`G_N^0`$, we can clearly see that DEMG model predictions for birefringence are fairly close to the experimental measurements, at least for low and intermediate $`Wi_R`$. We note that the plateau expected in the transition region between the dynamics dominated by segmental orientation and segmental stretch, is quite narrow, because of the small separation between $`\tau _d`$ and $`\tau _R`$ ($`\tau _d/\tau _R=3N_e=21`$ for PS82, and $`39`$ for PS81 & PS2). The plateau is smeared out in the experimental curves, since tube-dilation reduces $`N_e`$ and hence the separation $`3N_e`$ between these two time scales. This supports the tube-dilation idea. In the nonlinear viscoelastic regime, the chain-stretching predicted by the model is much stronger than what has been seen experimentally. In the limit $`Wi_R\mathrm{}`$, the model predicts a saturation of birefringence $`n_t/N_e`$ corresponding to a maximum chain-extension ratio $`\sqrt{n_t/N_e}`$. Comparing the maximum birefringence shown in Figs. 26, 27, and 28 with the above numbers, we see that the maximum chain-extension predicted by DEMG model at the highest $`Wi_R`$ studied are $`87.8\%`$, $`63.2\%`$ and $`28.5\%`$ for PS81, PS82 and PS2, respectively but the experimentally observed maximum chain-extension for these samples is only $`15\%`$.
The quantitative agreement between the experimental and the predicted $`\eta /\eta _0`$ for small rates of deformation in Figs. 29, 30, and 31 directly points to a proper choice of the calculated and the experimental values of $`\eta _0`$. In the intermediate flow regime the DEMG model predicts a strong strain-rate thinning behavior for the generalized extensional viscosity for entangled solutions. This also distinguishes the entangled solution from dilute ones, for which a monotonic increase in $`\eta `$ with the increase of $`\dot{\gamma }`$ is expected. In the experiments, we see that the solutions with $`N_e13`$ show much stronger thinning than the solutions with a lower value of $`N_e`$, consistent with the prediction. It has been shown in earlier experiments that the original DE model overpredicts the strength of shear-thinning in simple shear flows and the DEMG model also shows a similar behavior in extension-dominated flows. Similarly, we note that the excessive shear-thinning predicted by DEMG model for PS82 and PS2 (shown in Figs. 30 and 31, respectively) is due to the fact that the predicted orientation angle and the birefringence are too small in this intermediate range of $`Wi_R`$ (see, Figs. 24, 25, 27 and 28). The tendency of DEMG model to over-orient the polymer chain away from the principle axis of the rate of strain tensor, is speculated to be due to the fact that the model does not incorporate the tube-dilation effect which may be present in the real systems. For PS81, as can be seen from Fig. 29, the experimental curve for the viscosity closely follows the predictions of the DEMG model in the low and intermediate range of $`Wi_R`$ primarily due to the fact that, in this region, $`\mathrm{\Delta }n`$ is too small but $`\chi `$ is too large so as to compensate each other. It can be seen from these figures that as $`Wi_R`$ approaches $`𝒪(1)`$, the model predicts an upturn of the viscosity followed by a sharp increase owing to the prediction of a strong chain-stretching phenomenon. Apart from a small upturn seen in Fig. 31, the model fails to account for the observed behavior of the generalized extensional viscosity in this non-linear viscoelastic regime.
#### 1 Cox-Merz superposition
The empirical Cox-Merz relation states that the steady shear viscosity is equal to the modulus of the complex dynamic viscosity evaluated at the the angular frequency equal to the shear-rate. This is observed to be valid mostly in the linear viscoelastic regime, and also in the case of a simple shear flow ($`\lambda =0`$), primarily because chain-stretching is not significant in either of these cases. Thus, for extensional flows ($`\lambda >0`$) departure from Cox-Merz “rule” constitutes another means to study chain-stretching effects present in a given flow system. Following a suggestion of Mead et al. in Ref. , we have compared the measured generalized extensional viscosity, $`\eta `$, for the steady-state flow generated in a two-roll mill as a function of the measured velocity-gradient (or “shear-rate”) $`\dot{\gamma }`$ for the entangled solutions with the corresponding dynamic linear viscoelastic measurement, i.e., $`|\eta ^{}(\omega )|`$, in the form of “Cox-Merz plots” shown in Figs. 29, 30, and 31. Both parameters are scaled with the zero-shear viscosity $`\eta _0`$ (see, Table I). The match between the two sets of data at low deformation rates for all samples and that the values for these parameters are in the same ballpark provides an increased confidence on the choice of the corresponding $`\eta _0`$ values. At low and intermediate range of flow deformations, the values of these two sets of parameters are in a closer agreement for the samples with $`N_e13`$ than that for PS82 ($`N_e7`$). Apart from this, the overall dissimilarities between the deformation rate dependence of these two sets of results is obvious from these plots. At high deformation rates, chain-stretching effects become important, and, though it is not as strong as demanded by the DEMG model seen in Figs. 29, 30, and 31 in the previous section, its effect is strong enough to show an increased departure of $`\eta `$ from the Cox-Merz superposition. This departure is quite expected, as we have noted above, since the chain-stretching effects are not accounted for in this empirical rule. The departure seen in the intermediate range of deformation rates is because of the fact that even in this regime, the flow experienced by the polymers in the stagnation region of a two-roll mill is very different from a shear-flow. In short, these figures provide a quantitative picture of the relative importance of chain-stretching to orientation effects.
## IV Summary and Conclusions
The recently developed dynamic light scattering section of the two-color flow birefringence experimental setup is used to characterize the steady-state flow-fields for a Newtonian fluid as well as for viscoelastic, entangled polymeric fluids in a two-roll mill, by measuring the velocity-gradient and the flow-type parameter for these fluids. Assuming Newtonian, creeping flow symmetry of the flow-field, the flow parameters for a Newtonian fluid, obtained directly from the decay rate of the autocorrelation functions measured at the stagnation region of the flow, are shown to be very well described by the Frazer’s creeping flow solution.
The measured flow parameters for the polymeric fluids show clear departures from their corresponding Newtonian values, providing direct evidence of flow modification due to the conformational dynamics of the polymer molecules. Our results suggest that this departure may have a stronger dependence on the entanglement density of the polymer chains than on the molecular weight and concentration.
Within the limit of experimental error, the flow-field at the stagnation region of the two-roller is verified to retain its symmetry, for all rates of deformation studied with the Newtonian fluid and for the strain-rates $`\dot{\gamma }_N\sqrt{\lambda }_N`$ $`<0.5`$ s<sup>-1</sup> ($`Wi_R<1`$) with the polystyrene samples. The flow-type parameter for the polystyrene fluids, extracted from the steady-state DLS experiments, was found to exceed its Newtonian value at very low rates of deformation. This result was verified to be consistent from three different experimental means.
The dynamics of the entangled polymers under study was coupled with the changes in the flow-field. As we have noted above, the DLS experiments has probed the effect of the viscoelasticity on the velocity field compared to its Newtonian form. On the other hand, the polymer response induced by the flow-fields generated by a two-roll mill is studied using two-color flow birefringence experiments, which were in turn compared with the birefringence predictions from the DEMG model using the measured flow data as input to the model.
Similar to the model predictions, the experimental birefringence results clearly illustrate the existence of three steady-state flow regimes. In the first two regimes of low and moderate rates of deformation, where the dynamics is dominated by chain segment re-orientation and reptational diffusion, the model predictions are qualitatively, and some times, quantitatively, reproduced in the experiments. However, the model is clearly inadequate in describing the polymer dynamics at sufficiently high rates of deformation, where chain-stretching is important. Our data do show signatures of chain-stretching: the experimental birefringence exceeds the plateau value, the Weissenberg number scaling behavior shown at the onset of chain-stretching by at least two of the three polymer samples studied, and the departure of the generalized extensional viscosity data from the empirical “Cox-Merz superposition”. The comparisons of the experimental birefringence as well as the generalized extensional viscosity data unambiguously indicates that the DEMG model overpredicts chain-stretching in these entangled samples. The relative values of birefringence, orientation angle, and viscosity for different entangled polymeric samples are shown to be determined by the competition between the effects of the corresponding flow-type parameter, molecular weight (via $`n_t`$) and the number entanglements per chain. Our experiments have demonstrated that the effect of $`N_e`$ on these parameters is much stronger than that of the $`M_w`$ and $`\lambda `$. We note that the model fails in describing the “smearing” of the expected plateau region in the birefringence curves. In particular, it overpredicts the tendency of the flow to rotate the polymer chains toward the outflow axis, at high Weissenberg numbers, thereby predicting excessive thinning of the generalized viscosity in this regime of flow for at least two of the polystyrene samples of our study. It was suggested in an earlier paper that one of the prime causes for these deficiencies may be the fact that a real system experiences a conformation dependent decrease in the entanglement density (or a dilation in the tube diameter) as well as convective constraint-release, causing a decrease in the time scale for reptation with increasing flow strength that is not included in the present form of the DEMG model. This calls for further efforts to improve reptation based constitutive models to obtain a better match of its predictions to experiments on entangled polymeric systems.
###### Acknowledgements.
We thank Johan Remmelgas for helpful discussions, critical reading of the manuscript and his help in the flow-symmetry calculation. We acknowledge James P. Oberhauser for the help with the DEMG model numerical simulation code.
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# Electrostatic Attraction of Coupled Wigner Crystals: Finite Temperature Effects
## I Introduction
Electrostatic interactions play an important role in a system of charged macroions in an aqueous solution of neutralizing counterions. The macroions may be charged membranes, stiff polyelectrolytes such as DNA, or charged colloidal particles. Recently, there has been a great interest in understanding the attraction arising from correlations between highly-charged macroions as evidenced in experiments and in simulations. This attraction cannot be explained by the standard Poisson-Boltzmann (PB) treatment, even for an idealized system of two highly charged planar surfaces, with counterions distributed between them, since PB, being a mean-field theory, neglects correlations. Indeed, it has been proven recently that PB theory predicts only repulsions between two likely-charged objects. Recall that the PB solution for a single charged surface with charge density $`en`$ – where $`e`$ is the elementary charge and $`n`$ the areal density – immersed in a solution of neutralizing counterions of valence $`Z`$, predicts a length scale $`\lambda _{GC}=1/(2\pi l_BZn`$) (where $`l_B\frac{e^2}{ϵk_BT}7`$Å is the Bjerrum length below which electrostatics dominates the thermal energy in an aqueous solution of dielectric constant $`ϵ=80`$ ($`H_2O`$), $`k_B`$ is the Boltzmann constant, and $`T`$ is the temperature.) Physically, this Gouy-Chapman length $`\lambda _{GC}`$ defines a sheath near the charged surface within which most of the counterions are confined. For a moderately charged surface of $`n1/100`$Å<sup>-2</sup>, $`\lambda _{GC}`$ is of the order of few angstroms, and for highly charged surfaces and multivalent counterions $`Z>1`$, we have $`\lambda _{GC}<l_B`$, signaling the breakdown of PB theory. In this limit, fluctuations and correlations about the mean-field potential become so large that the solution to the PB equation is no longer valid.
To account for the attraction arising from correlations, two distinct approaches have been proposed. The first approach, based on charge fluctuations, treats the “condensed” counterion fluctuations in the Gaussian approximation. This theory predicts a long-ranged attraction which vanishes as $`T0`$. In the other approach based on “structural” correlations first proposed by Rouzina and Bloomfield , the attraction comes from the ground state configuration of the “condensed” counterions. Indeed, at low temperature, the “condensed” counterions crystallize on the charged surface to form a 2D Wigner crystal. When brought together, the counterions of two Wigner crystals correlate themselves to minimize the electrostatic energy. These staggered Wigner crystals attract each other via a short-ranged force that is strongest at $`T=0`$. Although the physical origin of the attraction is clear in each approach, the relationship between them remains somewhat obscure and their results in the $`T0`$ limit are somewhat contradictory. Therefore, it is desirable to formulate a unified approach which captures the physics of both mechanisms.
To this end, we attempt in this paper to develop a detailed physical picture of the electrostatic interaction at finite temperature between two planar Wigner crystals in the strong Coulomb coupling limit. Since correlation effects are essentially two-dimensional, we consider a model system composed of two uniformly charged planes a distance $`d`$ apart, each having a charge density $`en`$. Confined on the surfaces are negative point-like mobile charges of magnitude $`e`$. In order to understand correlation effects that are not captured by PB theory, we assume that the charges form a system of interacting Wigner crystals (see Fig. 1). In particular, we compute the electrostatic attraction between the two layers by explicitly taking into account both correlated fluctuations and “structural” correlations. (By “structural” correlations, we mean the residual ground state spatial correlations which remain at finite temperature.) In the former case, we obtain a long-ranged force ($`1/d^3`$), which agrees with the result based on Debye-Hückel (Gaussian) approximation. For the latter, a simple expression for the short-ranged force is derived, which shows that thermal fluctuations reduce its range, and which in the $`T0`$ limit agrees with the known exponentially decaying result. Furthermore, we argue that at zero temperature, there must also be a long-ranged force derived from the quantum fluctuations of the plasmons, in addition to this zero-temperature exponentially decaying force. At very low temperatures in the quantum regime, we obtain a new length scale $`\lambda _L`$ (to be defined below), where the attraction scales like $`d^2`$ when $`d<\lambda _L`$.
It should be pointed out that we have assumed a uniform charge distribution on the surface of the charged plates in our model for electrostatic attraction, mediated by the “condensed” counterions. This assumption of a uniform neutralizing background may not be a good approximation to real experimental settings, since charges on macroion surfaces are discrete. For monovalent counterions, they tend to bind to the charges on the surface and form dipolar molecules. Therefore, the ground state for this system may not be a Wigner crystal, which relies on mutual repulsion among charges for its stability, and short-ranged effects are likely to be important. However, for polyvalent counterions, a Wigner crystal is likely to form since each counterion does not bind to a particular charge on the surface, and a uniform background may be more appropriate. The detailed structure of the ground state as determined by short-ranged effects and valences will be the subject for another study.
Another point worth mentioning concerns the ordering of 2D solids which exhibit quasi-long-range-order (QLRO). It is well-known that a true long-range order is impossible for 2D systems with continuous symmetries. For a 2D solid, which may be described by continuum elasticity theory with nonzero long-wavelength elastic constants, the Fourier components of the density function $`n(𝐫)=_𝐆n_𝐆(𝐫)e^{i𝐆𝐫}`$ average (thermally) out to zero for a nonzero reciprocal lattice vector $`𝐆`$, i.e. $`n_𝐆(𝐫)=e^{i𝐆𝐮(𝐫)}=0,`$ where $`𝐮(𝐫)`$ are the displacements of the particles from their equilibrium positions, while the correlation function decays algebraically to zero: $`n_𝐆(𝐫)n_𝐆^{}(\mathrm{𝟎})r^{\eta _𝐆(T)}`$ with $`\eta _𝐆(T)=\frac{k_BTG^2(3\mu +\lambda )}{4\pi \mu (2\mu +\lambda )}`$, where $`\mu `$ and $`\lambda `$ are Lamé elastic constants. This slow power-law decay of the correlation function is very different from the exponential decay one would expect in a liquid. Hence the term QLRO. For a single 2D Wigner crystal, QLRO implies that the thermal average of the electrostatic potential at a distance $`d`$ above the plane is zero at any non-zero temperature, in contrast to a perfectly ordered lattice ($`T=0`$) where the electrostatic potential decays exponentially with $`d`$. This may lead to the conclusion that at finite temperatures the short-ranged force between two coupled Wigner crystals should likewise be zero. As we show below, this is not the case because the susceptibility, which measures the linear response of a 2D lattice to an external potential, nevertheless diverges at the reciprocal lattice vectors as in 3D solids.
This paper is organized as follows. In Sec. II, we derive an effective Hamiltonian which describes two interacting planar Wigner crystals starting from the zero temperature ground state. The total pressure is then decomposed into a long-ranged and a short-ranged component, which are evaluated in Sec. II A and II B, respectively, and a detailed discussion of our results is presented in Sec. II C. In Sec. III, we present an argument for a long-ranged attractive force arising from the zero-point fluctuations at zero temperature. In addition, we use the Bose-Einstein distribution to calculate the attractive long-ranged pressure in the quantum regimes.
## II Effective Hamiltonian and Pressure
We start with the Hamiltonian for two interacting Wigner crystals: $`=_0+_{int}`$. Here, $`_0`$ is the elastic Hamiltonian for two isolated Wigner crystals
$$\beta _0=\frac{1}{2}\underset{i}{}\frac{d^2𝐪}{(2\pi )^2}\mathrm{\Pi }_{\alpha \beta }(𝐪)u_\alpha ^{(i)}(𝐪)u_\beta ^{(i)}(𝐪),$$
(1)
where $`\beta ^1=k_BT`$, $`𝐮^{(i)}(𝐪)`$ is the Fourier transform of the in-plane displacement field of the charges in the $`i`$th layer ($`i=A`$ or $`B`$), $`\mathrm{\Pi }_{\alpha \beta }(𝐪)=\left[\frac{2\pi l_Bn^2}{q}P_{\alpha \beta }^L+\mu P_{\alpha \beta }^T\right]q^2`$ is the dynamical matrix, $`\mu 0.245n^{3/2}l_B`$ is the shear modulus in units of $`k_BT`$, $`P_{\alpha \beta }^L=q_\alpha q_\beta /q^2`$ and $`P_{\alpha \beta }^T=\delta _{\alpha \beta }q_\alpha q_\beta /q^2`$ are longitudinal and transverse projection operator, respectively. Here, Greek indices indicate Cartesian components. $`_{int}`$ is the electrostatic interaction between the two layers:
$$\beta _{int}=l_Bd^2𝐱d^2𝐱^{}\frac{[\rho _A(𝐱)n][\rho _B(𝐱^{})n]}{\sqrt{(𝐱𝐱^{})^2+d^2}},$$
(2)
where $`\rho _i(𝐱)`$ is the number density of charges in the $`i`$th layer. In order to capture the long-wavelength coupling as well as discrete lattice effects which are essential for our discussions on the short-ranged force, we employ a method, similar to that in Ref., which allows us to derive an effective Hamiltonian that is valid in the elastic regime where the density fluctuations are slowly varying in space, i.e. $`𝐮^{(i)}(𝐱)\mathrm{\hspace{0.17em}1}`$, but $`|𝐮^A(𝐱)𝐮^B(𝐱)|`$ need not be small compared to the lattice constant $`a`$.
Let us introduce a slowly varying field for each layer:
$$\varphi _\alpha ^{(i)}(𝐱)=x_\alpha u_\alpha ^{(i)}[\stackrel{}{\varphi }^{(i)}(𝐱)],$$
(3)
where the displacement field $`𝐮^{(i)}(𝐱)`$ is defined in such a way that it has no Fourier components outside of the Brillouin Zone (BZ). Then, the density $`\rho _i(𝐱)`$ can be written as:
$$\rho _i(𝐱)=\underset{l}{}\delta ^2(𝐑_l\stackrel{}{\varphi }^{(i)}(𝐱))det[_\alpha \varphi _\beta ^{(i)}(𝐱)],$$
(4)
where $`𝐑_l`$ are the equilibrium positions of the charges, i.e. the underlying lattice sites. Using the Fourier representation of the delta function and solving $`\varphi _\alpha ^{(i)}(𝐱)`$ iteratively in terms of the displacement field, we obtain a decomposition of the density for the $`i`$th layer into a slowly and a rapidly spatially varying pieces:
$$\rho _i(𝐱)nn𝐮^{(i)}(𝐱)+\underset{𝐆0}{}ne^{i𝐆[𝐱+𝐮^{(i)}(𝐱)]},$$
(5)
where $`𝐆`$ is a reciprocal lattice vector. Note that we have neglected terms that are products of the slowly and the rapidly varying terms. Physically, the first term represents density fluctuations for wavelengths greater than the lattice constant, and the second term represents the underlying lattice, modified by thermal fluctuations. Using the density decomposition (5), $`_{int}`$ may be written as
$`\beta _{int}`$ $`=`$ $`{\displaystyle \frac{d^2𝐪}{(2\pi )^2}\frac{2\pi l_B}{q}e^{qd}d^2𝐱d^2𝐱^{}e^{i𝐪(𝐱𝐱^{})}\left(n𝐮^A(𝐱)\underset{𝐆0}{}ne^{i𝐆[𝐱+𝐮^A(𝐱)]}\right)}`$ (7)
$`\times \left(n𝐮^B(𝐱^{}){\displaystyle \underset{𝐆^{}0}{}}ne^{i𝐆^{}[𝐱^{}+𝐜+𝐮^B(𝐱^{})]}\right),`$
where $`𝐜`$ is the relative displacement vector between two lattices of the different plane and we have used the fact that $`\frac{1}{\sqrt{x^2+d^2}}=\frac{d^2𝐪}{(2\pi )^2}e^{i𝐪𝐱}\frac{2\pi }{q}e^{qd}.`$ Again neglecting the products of slowly and rapidly varying terms, which give vanishingly small contributions when integrating over all space, $`_{int}`$ separates into two pieces: a long-wavelength term
$$\beta _{int}^L=\frac{d^2𝐪}{(2\pi )^2}\frac{2\pi l_Bn^2}{q}e^{qd}q_\alpha q_\beta u_\alpha ^A(𝐪)u_\beta ^B(𝐪),$$
(8)
and a short-wavelength term
$`\beta _{int}^S=+{\displaystyle \underset{𝐆0}{}}{\displaystyle \underset{𝐆^{}0}{}}{\displaystyle \frac{d^2𝐪}{(2\pi )^2}}`$ $`{\displaystyle \frac{2\pi l_Bn^2}{q}}e^{qd}`$ (9)
$`\times `$ $`{\displaystyle d^2𝐱d^2𝐱^{}e^{i𝐪(𝐱𝐱^{})}e^{i𝐆[𝐱+𝐮^A(𝐱)]}e^{i𝐆^{}[𝐱^{}+𝐜+𝐮^B(𝐱^{})]}}.`$ (10)
In order to obtain a tractable analytical treatment, we approximate this expression by splitting the sum over $`𝐆^{}`$ into two parts. The dominant part, with $`𝐆^{}=𝐆`$ is
$$\beta _{int}^S=\underset{𝐆0}{}\frac{d^2𝐪}{(2\pi )^2}\frac{2\pi l_Bn^2}{q}e^{qd}d^2𝐱d^2𝐱^{}e^{i(𝐪+𝐆)(𝐱𝐱^{})}e^{i𝐆[𝐮^A(𝐱)𝐮^B(𝐱^{})]},$$
(11)
where we have used $`e^{i𝐆𝐜}=1`$. The second part (those terms with $`𝐆^{}𝐆`$) contains extra phase factors which tend to average to zero in the elastic limit. As a first approximation, we neglect such terms. Finally, Eq. (11) can be systematically expanded using a gradient expansion:
$$\beta _{int}^S=\underset{𝐆0}{}\mathrm{\Delta }_𝐆(d)d^2𝐱\mathrm{cos}\left\{𝐆\left[𝐮^A(𝐱)𝐮^B(𝐱)\right]\right\}+O(_\alpha u_\beta ^{(i)}_\gamma u_\tau ^{(j)}),$$
(12)
where $`\mathrm{\Delta }_𝐆(d)=\frac{4\pi l_Bn^2}{G}e^{Gd}.`$ Putting Equations (1), (8), and (12) together, we obtain an effective Hamiltonian for the coupled planar Wigner crystals:
$`\beta _e=\beta _0+{\displaystyle \frac{d^2𝐪}{(2\pi )^2}\frac{2\pi l_Bn^2}{q}}`$ $`e^{qd}`$ $`q_\alpha q_\beta u_\alpha ^A(𝐪)u_\beta ^B(𝐪)`$ (13)
$``$ $`{\displaystyle \underset{𝐆0}{}}\mathrm{\Delta }_𝐆(d){\displaystyle d^2𝐱\mathrm{cos}\left\{𝐆\left[𝐮^A(𝐱)𝐮^B(𝐱)\right]\right\}}.`$ (14)
The second term in Eq. (14) comes from the long-wavelength couplings while the third term reflects the periodicity of the underlying lattice structure. This particular structure in the effective Hamiltonian, as will be demonstrated below, leads to a total force which is comprised of two pieces – an exponentially decaying (short-ranged) force and a long-ranged power-law force:
$$\mathrm{\Pi }(d)=\frac{1}{A_0}\frac{_{int}}{d}__e=\frac{1}{A_0}\frac{_{int}^S}{d}__e\frac{1}{A_0}\frac{_{int}^L}{d}__e=\mathrm{\Pi }_{SR}(d)+\mathrm{\Pi }_{LR}(d),$$
(15)
where $`A_0`$ is the area of the plane. It is important to emphasize that both forces are present simultaneously, although each force dominates at a different spatial scale – the long-ranged force dominates at large separations while the short-ranged force at small separations.
To calculate various expectation values in Eq. (15), it is convenient to transform the displacement fields into in-phase and out-of-phase displacement fields by $`𝐮^+(𝐱)=𝐮^A(𝐱)+𝐮^B(𝐱)`$ and $`𝐮^{}(𝐱)=𝐮^A(𝐱)𝐮^B(𝐱)`$, respectively, so that the effective Hamiltonian (14) separates into two independent parts: $`_e=_++_{}`$ with
$$\beta _+=\frac{1}{2}\frac{d^2𝐪}{(2\pi )^2}\mathrm{\Pi }_{\alpha \beta }^+(𝐪)u_\alpha ^+(𝐪)u_\beta ^+(𝐪),$$
(16)
and
$$\beta _{}=\frac{1}{2}\frac{d^2𝐪}{(2\pi )^2}\mathrm{\Pi }_{\alpha \beta }^{}(𝐪)u_\alpha ^{}(𝐪)u_\beta ^{}(𝐪)\underset{𝐆0}{}\mathrm{\Delta }_𝐆(d)d^2𝐱\mathrm{cos}[𝐆𝐮^{}(𝐱)],$$
(17)
where $`\mathrm{\Pi }_{\alpha \beta }^\pm (𝐪)=\frac{1}{2}\left[\frac{2\pi l_Bn^2}{q}(1\pm e^{qd})P_{\alpha \beta }^L+\mu P_{\alpha \beta }^T\right]q^2`$. Furthermore, at low temperature, where $`𝐮^{}(𝐱)`$ is small compared to the lattice constant $`a`$, the cosine term in Eq. (17) can be expanded up to second order in $`|𝐮^{}(𝐱)|`$ to obtain the “mass” terms. Within a harmonic approximation $`_{}`$, up to an additive constant, may be written as
$$\beta _{}\frac{1}{2}\frac{d^2𝐪}{(2\pi )^2}\mathrm{\Pi }_{\alpha \beta }^{}(𝐪)u_\alpha ^{}(𝐪)u_\beta ^{}(𝐪)+\frac{1}{2}\frac{d^2𝐪}{(2\pi )^2}\left[m_L^2P_{\alpha \beta }^L+m_T^2P_{\alpha \beta }^T\right]u_\alpha ^{}(𝐪)u_\beta ^{}(𝐪).$$
(18)
Here, $`m_{L,T}^2=4\pi l_Bn^2_{𝐆0}Ge^{Gd}=4\pi l_Bn^2\mathrm{\Delta }_0(d)`$. Note that the mass terms vanish exponentially with $`d`$ as also found in Ref.. The fact that the transverse $`m_T`$ and longitudinal “mass” $`m_L`$ are degenerate is related to the underlying triangular structure of the lattices. These “masses” are associated with the finite energy required to uniformly shear the two Wigner crystals, and thus give rise to a gap in the dispersion relations of the out-of-phase modes. In the next two subsections, we derive expressions for the long-ranged and the short-ranged pressure as given in Eq. (15) within the harmonic approximation.
### A Long-ranged Pressure
The long-ranged power-law force comes from the correlated long-wavelength density fluctuations (the plasmon modes). The shear modes do not contribute to this interaction since $`_\alpha P_{\alpha \beta }^Tu_\beta ^{(i)}(𝐱)=0.`$ Using Eqs. (8) and (15), we obtain an expression for the long-ranged force:
$$\beta \mathrm{\Pi }_{LR}(d)=\frac{2\pi l_B}{A_0}\frac{d^2𝐪}{(2\pi )^2}e^{qd}\delta \rho _A(𝐪)\delta \rho _B(𝐪),$$
(19)
where $`\delta \rho _i(𝐱)=n𝐮^{(i)}(𝐱)`$ is the long-wavelength density fluctuation to the lowest order. Making use of the equipartition theorem to evaluate
$`\delta \rho _A(𝐪)\delta \rho _B(𝐪)`$ $``$ $`{\displaystyle D𝐮^\pm (𝐪)\delta \rho _A(𝐪)\delta \rho _B(𝐪)e^{\beta [_++_{}]}}`$ (20)
$`=`$ $`A_0{\displaystyle \frac{n^2q^2}{4}}\left[{\displaystyle \frac{1}{\pi l_Bn^2q(1e^{qd})+m_L^2}}{\displaystyle \frac{1}{\pi l_Bn^2q(1+e^{qd})}}\right],`$ (21)
and substituting this into Eq. (19), we have the result
$$\mathrm{\Pi }_{LR}(d)=\frac{k_BT}{d^3}\alpha (\mathrm{\Delta }_0d),$$
(22)
where
$$\alpha (x)\frac{\zeta (3)}{8\pi }+\frac{x}{\pi }\left[\text{Ci}(\mathrm{\hspace{0.17em}2}\sqrt{x})\mathrm{cos}(\mathrm{\hspace{0.17em}2}\sqrt{x})+\text{Si}(\mathrm{\hspace{0.17em}2}\sqrt{x})\mathrm{sin}(\mathrm{\hspace{0.17em}2}\sqrt{x})\right],$$
(23)
$`\zeta `$ is the Riemann zeta function, and Ci$`(x)`$ and Si$`(x)`$ are the cosine and sine integral functions, respectively. In the large distance limit, the second term in Eq. (23) is exponentially suppressed and can be neglected, yielding $`\alpha =\frac{\zeta (3)}{8\pi }`$. Therefore, for large $`d`$ we have
$$\mathrm{\Pi }_{LR}(d)=\frac{\zeta (3)}{8\pi }\frac{k_BT}{d^{\mathrm{\hspace{0.17em}3}}}.$$
(24)
This is the well-known result from the Debye-Hückel approximation. Note also that the amplitude $`\frac{\zeta (3)}{8\pi }0.048`$ is universal for this interaction, induced by the long wavelength fluctuations.
### B Short-ranged Pressure
The short-ranged force which decays exponentially owes its existence to the “structural” correlations. It survives even at non-zero temperature, in contrast to the conclusion drawn from a single 2D Wigner crystal, as discussed in the Introduction. However, we expect on physical grounds the short-ranged force to be weakened by thermal fluctuations. To compute its temperature dependence explicitly, we start with the expression for this force derived from Eqs. (11) and (15):
$$\beta \mathrm{\Pi }_{SR}(d)=\mathrm{\hspace{0.17em}2}\pi l_Bn^2\underset{𝐆0}{}e^{\frac{G^2}{2}|𝐮^{}(0)|^2}f_𝐆(d),$$
(25)
where $`f_𝐆(d)=\frac{d^2𝐪}{(2\pi )^2}𝒮(𝐪𝐆)e^{qd}`$, $`𝒮(𝐪𝐆)=d^2𝐫e^{i(𝐪𝐆)𝐫}e^{\frac{G^2}{8}[B^+(𝐫)B^{}(𝐫)]}`$, and $`B^\pm (𝐫)=[𝐮^\pm (𝐫)𝐮^\pm (\mathrm{𝟎})]^2`$. Note that Eq. (25) is exact, provided all the averages are evaluated exactly. For a system of coupled perfect Wigner crystals at zero temperature, $`f_𝐆(d)=e^{Gd}`$. At finite temperature, but below the melting temperature $`T_m`$, we note that $`B^\pm (r)`$ varies very slowly in space, so that $`f_𝐆(d)`$ can be approximated by its zero temperature value: $`f_𝐆(d)e^{Gd}`$. Hence, we obtain
$$\beta \mathrm{\Pi }_{SR}(d)\mathrm{\hspace{0.17em}2}\pi l_Bn^2\underset{𝐆0}{}e^{Gd}e^{i𝐆[𝐮^A(\mathrm{𝟎})𝐮^B(\mathrm{𝟎})]}__e.$$
(26)
The thermal average of the displacement fields in Eq. (26) resembles a “Debye-Waller” factor and indicate the degree to which the short-ranged force is depressed by thermal fluctuations from its zero temperature maximum value. Because of the cosine term present in Eq. (14), this “Debye-Waller” factor is in general not zero, unlike the case of a single 2D Wigner crystal. However, if the system has melted into a Coulomb fluid, this cosine term, which comes from the lattice structure, would have to be modified.
The required expectation value in Eq. (26) only involves $`_{}`$. Within the harmonic approximation, the mean-square out-of-phase displacement field can be evaluated
$$|𝐮^{}(𝐱)|^2\frac{\lambda _D}{2\pi nd}\mathrm{ln}\left[\frac{d}{4\mathrm{\Delta }_0(d)a^2}\right]+\frac{1}{2\pi \mu }\mathrm{ln}\left[\frac{\mu }{8\pi l_Bn^2\mathrm{\Delta }_0(d)a^2}\right]\frac{G_0d}{2\pi }\left[\frac{\lambda _D}{nd}+\frac{1}{\mu }\right],$$
(27)
where $`\lambda _D=1/(2\pi l_Bn)`$, $`a`$ is the lattice constant, $`\mu 0.245n^{3/2}l_B`$ is the shear modulus of an isolated Wigner crystal in units of $`k_BT`$, and in the last line, we have approximated $`\mathrm{\Delta }_0(d)`$ by the first nonzero reciprocal lattice vector contribution: $`\mathrm{\Delta }_0(d)G_0e^{G_0d}`$. Note also that the logarithmic dependence on the “mass” ($`=4\pi n^2l_B\mathrm{\Delta }_0(d)`$) is a characteristic of 2D solids. Inserting Eq. (27) into Eq. (26), we obtain an expression for the short-ranged pressure at finite temperatures
$$\beta \mathrm{\Pi }_{SR}(d)\mathrm{\hspace{0.17em}2}\pi l_Bn^2e^{(1+\xi /2)G_0d}.$$
(28)
Here, the parameter $`\xi `$ defined by
$$\xi =\frac{G_0^2}{2\pi }\left(\frac{\lambda _D}{nd}+\frac{1}{\mu }\right),$$
(29)
characterizes the relative strengths of thermal fluctuations and the electrostatic energy of a Wigner crystal, i.e. $`\xi \frac{k_BTa}{e^2}`$. Thus, the sole effect of thermal fluctuations on the short-ranged force is to reduce its range: $`G_0G_0\left(\mathrm{\hspace{0.17em}1}+\frac{\xi }{2}\right)`$.
### C Discussion of Results
In summary, we have shown that the total pressure can be decomposed into a long-ranged $`\mathrm{\Pi }_{LR}`$ and a short-ranged pressure $`\mathrm{\Pi }_{SR}`$. Each force is computed at low temperatures, where the harmonic approximation is expected to be valid. The result for the total force is
$$\beta \mathrm{\Pi }(d)\mathrm{\hspace{0.17em}2}\pi l_Bn^2e^{(1+\xi /2)G_0d}\frac{\alpha (\mathrm{\Delta }_0d)}{d^3},$$
(30)
where $`\xi =\frac{G_0^2}{2\pi }\left[\frac{\lambda _D}{nd}+\frac{1}{\mu }\right]`$ and $`\alpha (\mathrm{\Delta }_0d)=\frac{\zeta (3)}{8\pi }`$ for large $`d`$. In Fig. 2, we have plotted $`\mathrm{\Pi }_{SR}`$ and $`\mathrm{\Pi }_{LR}`$ for two values of the coupling constant, $`\mathrm{\Gamma }\frac{l_B}{a}=150`$ and $`50`$. Not surprisingly, they show that $`\mathrm{\Pi }_{SR}`$ dominates for small $`d`$, and $`\mathrm{\Pi }_{LR}`$ for large $`d`$. However, it is interesting to note that even for high values of $`\mathrm{\Gamma }`$, $`\mathrm{\Pi }_{LR}`$ dominates as soon as $`da`$.
According to Eq. (28), the magnitude of $`\mathrm{\Pi }_{SR}`$ tends to decrease exponentially with temperature, as illustrated in Fig. 3. This strong decrease with increasing temperature is consistent with the Brownian dynamics simulations of Grønbech-Jensen et al.. The shortening of its range may be attributed to the generic nature of strong fluctuations in 2D systems, and can also be understood by the following scaling argument. Referring back to $`_{}`$ in Eq. (17), one can show that the anomalous dimension of the operator $`\mathrm{cos}[𝐆_0𝐮^{}(𝐱)]`$ is $`[`$ Length $`]^\xi `$ and correspondingly the dimension of $`\mathrm{\Delta }_{𝐆_0}(d)`$ is $`[`$ Length $`]^{\xi 2}`$. Since $`\mathrm{\Delta }_{𝐆_0}(d)`$ is the only relevant length scale in $`_{}`$, we must have $`e^{i𝐆_0𝐮^{}(\mathrm{𝟎})}\mathrm{\Delta }_{𝐆_0}^{\frac{\xi }{2\xi }}`$. Therefore, the short-ranged pressure scales like
$$\mathrm{\Pi }_{SR}(d)\mathrm{\Delta }_{𝐆_0}(d)\times e^{i𝐆_0𝐮^{}(\mathrm{𝟎})}\mathrm{\Delta }_{𝐆_0}\times \mathrm{\Delta }_{𝐆_0}^{\frac{\xi }{2\xi }}e^{G_0d\left(\frac{2}{2\xi }\right)}.$$
(31)
In the low temperature limit ($`\xi 1`$), we see that the range of $`\mathrm{\Pi }_{SR}`$ is $`G_0\left(\mathrm{\hspace{0.17em}1}+\frac{\xi }{2}\right)`$ as in Eq. (28). This scaling argument also suggests that at higher temperatures thermal fluctuations may have interesting nonperturbative effects. At zero temperature $`\xi =0`$, so $`\mathrm{\Pi }_{SR}`$ in Eq. (28) reproduces the known result of exponentially decaying attractive force.
The long-ranged pressure for large $`d`$ in Eq. (24) agrees exactly, including the prefactor, with the Debye-Hückel approximation. This is hardly surprising since the existence of long-wavelength plasmons (average density fluctuations) is independent of local structure, and they are present for solids and fluids alike. Thus, the asymptotic long-ranged power-law force must manifest itself even after QLRO is lost via a 2D melting transition driven by dislocations. Therefore, our formulation captures the essential physics of the attraction not only arising from the ground state “structural” correlations, but also from the high temperature charge-fluctuations.
## III Quantum Contributions to the Long-ranged Attraction
According to the classical calculations above, correlation effects give rise to a “structural” short-ranged and a long-ranged attractive force. Recall that the long-ranged force vanishes as $`T0`$, and that the short-ranged force is strongest at zero temperature but vanishes exponentially with distance. This observation suggests that for sufficiently large separations correlated attractions at finite temperatures are stronger than those arising from the perfectly correlated zero temperature ground state. However, we have pointed out in Ref. that zero-point fluctuations of the plasmons lead to an attractive long-ranged interaction, which exhibits an unusual fractional-power-law decay ($`d^{7/2}`$), in contrast to the zero-temperature van der Waals interaction ($`d^4`$). Hence, in the $`T0`$ limit, this “zero-point attraction” dominates the short-ranged “structural” force at large separations. Furthermore, we expect that quantum fluctuations persist at finite temperature, and in this section, we compute their temperature dependence.
Within the harmonic approximation to the effective Hamiltonian, the dispersion relations for the plasmons can be readily obtained:
$`\omega _1^2(q)`$ $`=`$ $`{\displaystyle \frac{8\pi e^2n}{mϵ}}\mathrm{\Delta }_0(d)+{\displaystyle \frac{2\pi e^2n}{mϵ}}q(\mathrm{\hspace{0.17em}1}e^{qd});`$ (32)
$`\omega _2^2(q)`$ $`=`$ $`{\displaystyle \frac{2\pi e^2n}{mϵ}}q(\mathrm{\hspace{0.17em}1}+e^{qd}),`$ (33)
where $`m`$ is the mass of the charges and $`\mathrm{\Delta }_0(d)e^{Gd}`$ is proportional to the energy gap (the “mass” term) for the out-of-phase mode. The plasmon modes are related to the correlated charge-density fluctuations in the two layers. At any finite temperature, the free energy of the low-lying plasmon excitations is given by
$$(d)/A_0=\frac{\mathrm{}}{2}\underset{i=1,2}{}\frac{d^2𝐪}{(2\pi )^2}\omega _i(𝐪)+k_BT\underset{i=1,2}{}\frac{d^2𝐪}{(2\pi )^2}\mathrm{ln}\left[1e^{\beta \mathrm{}\omega _i(𝐪)}\right],$$
(34)
where $`A_0`$ is the area of the plane. Since the energy gap $`\mathrm{\Delta }_0`$ is exponentially damped for large distances, its contribution to the free energy may be neglected in the large distance limit, where the long-ranged force is expected to be dominant.
The first term in Eq. (34) arising from the zero-point fluctuations has been computed in Ref. and gives the $`d^{7/2}`$ power-law mentioned above. An additional contribution to the pressure at finite temperature arises from the second term in Eq. (34),
$`\beta \mathrm{\Pi }_{LR}(d)={\displaystyle \frac{\mathrm{}\mathrm{\Lambda }}{4\pi d^{\mathrm{\hspace{0.17em}7}/2}}}{\displaystyle _0^{\mathrm{}}}dxx^{5/2}\{{\displaystyle \frac{1}{\mathrm{exp}[\eta \sqrt{x(1e^x)}]1}}{\displaystyle \frac{e^x}{\sqrt{1e^x}}}`$ (35)
$`{\displaystyle \frac{1}{\mathrm{exp}[\eta \sqrt{x(1+e^x)}]1}}`$ $`{\displaystyle \frac{e^x}{\sqrt{1+e^x}}}\},`$ (36)
where $`\mathrm{\Lambda }=\sqrt{\frac{2\pi e^2n}{mϵ}}`$ and $`\eta =\beta \mathrm{}\mathrm{\Lambda }/\sqrt{d}`$. We can evaluate this expression in two limits:
In the low-temperature limit $`\eta 1`$, Eq. (36) can be systematically expanded in powers of $`\eta ^1`$. The lowest order term is given by $`\mathrm{\Pi }_{LR}(d)=\overline{\alpha }\frac{k_BT}{\lambda _Ld^2}`$, where $`\lambda _La_B\frac{l_B}{2\lambda _D},`$ $`a_Bϵ\mathrm{}^2/(me^2)`$ is the effective Bohr radius, $`\overline{\alpha }\frac{1}{4\pi }_0^{\mathrm{}}𝑑x\frac{x^2}{e^x1}=\zeta (3)/(2\pi )`$, and $`\zeta `$ is the Riemann zeta function. We observe that the low temperature condition $`\eta >1`$ is equivalent to the short distance limit $`d<\lambda _L`$.
In the high temperature limit $`\eta 1`$ or the large distance limit $`d>\lambda _L`$, we expand the exponential in the denominator of Eq. (36) to obtain $`\mathrm{\Pi }_{LR}(d)=\alpha \frac{k_BT}{d^3},`$ where $`\alpha =\zeta (3)/(8\pi )`$. This result agrees with the classical calculation in Sec. II A as it should. Therefore, we have the following regimes for correlated attraction from plasmon fluctuations at finite temperature
$$\mathrm{\Pi }_{LR}(d)\{\begin{array}{cc}k_BT/d^3,\hfill & \text{ for }\lambda _L<d\text{,}\hfill \\ k_BT/(\lambda _Ld^2),\hfill & \text{ for }\lambda _L>d.\hfill \end{array}$$
(37)
We note that $`\lambda _L`$, in contrast to $`\lambda _D`$, increases with decreasing temperature, indicating, as one might expect, that quantum fluctuations are important at low temperatures. Furthermore, since $`\mathrm{\Pi }_{LR}(d)0`$ as $`T0`$, the attractive interaction as $`T0`$ is governed by zero-point fluctuations as emphasized above. In the strong Coulomb coupling limit $`l_B/\lambda _D100`$, we get $`\lambda _L3`$Å for $`ϵ100`$ and $`a_B1/20`$Å. Finally, it should be emphasized that quantum contributions to the long-ranged attraction are unlikely to be relevant for macroions. Our motivation here stems from the desire to understand the charge-fluctuation-induced attraction between coupled layers in a complete picture. However, our results may be relevant for electrons in bilayer semiconductor systems. Indeed, there are recent theoretical efforts devoted to this subject .
## IV Discussion and Conclusion
In this paper, we have studied analytically the electrostatic attraction between two planar Wigner crystals in the strong Coulomb coupling limit. We show that the total attractive pressure can be separated into a long-ranged and short-ranged component. The long-ranged pressure arises from correlated fluctuations and the short-ranged pressure from the ground state “structural” correlations. We also compute the very low temperature behavior of the fluctuation-induced attraction, where long-wavelength plasmon excitation must be described by Bose-Einstein statistics. The results are summarized in Fig. 4, showing different regimes for the charge-fluctuation-induced long-ranged attraction, including the high temperature results in Ref. and the characteristic decay length $`l_{SR}`$ for the short-ranged force. For small $`d`$, the short-ranged force is always dominant, but the decay length shrinks with increasing temperature. The crossover from the short-ranged to long-ranged dominant regimes occurs about $`da`$. Thus, for large $`da`$ only the long-ranged force is operative, which crosses over from $`d^{7/2}`$ at zero temperature to the finite temperature distance dependence of $`d^2`$ if $`d<\lambda _L`$ and $`d^3`$ if $`d>\lambda _L`$. This provides a unified description to the electrostatic attraction between two coupled Wigner crystals.
In addition, our formulation may offer further insights into the nature of the counterion-mediated attraction at short distances. As discussed in Sec. II B, the reason that the short-ranged force in Eq. (26) does not vanish is because of the cosine term in $`_{}`$, which represents the underlying lattice structures, and our results indicate that the strength of the short-ranged force decreases exponentially with temperature. However, at higher temperatures the expression for $`\mathrm{\Pi }_{SR}`$ in Eq. (28) is no longer valid, since the harmonic approximation breaks down. Indeed, the scaling argument leading to Eq. (31) suggests that if the full cosine term is retained, $`\mathrm{\Pi }_{SR}`$ may exhibit nonperturbative behaviors as $`\xi 2^{}`$.
To discuss qualitatively what happens at higher temperatures, we assume that $`\mathrm{\Delta }_𝐆(d)`$ is sufficiently small and the system of interacting Wigner crystals is below its melting temperature $`T_m`$. Then, the charges between the two layers may unlock via a Kosterlitz-Thouless (KT) type of transition, determined by the relevancy of the cosine term in $`_{}`$, at $`\xi =2`$ . (An order of magnitude estimate for the coupling constant is $`\mathrm{\Gamma }13`$.) In the locked phase, $`\xi 2`$, the periodic symmetry in $`_{}`$ is spontaneously broken, and the resulting state is well captured by the harmonic approximation. On the other hand, when $`\xi >2`$ the fluctuations are so large that the ground state becomes nondegenerate (gapless), i.e. the layers are decoupled. To compute $`\mathrm{\Pi }_{SR}`$ in the unlocked phase, $`_{int}^S`$ given in Eq. (12) can be treated as a perturbation in evaluating the “Debye-Waller” factor in Eq. (26). To the lowest order, we obtain
$$\mathrm{\Pi }_{SR}(d)\frac{k_BT}{\lambda _D^2a}\left(\frac{\xi 1}{\xi 2}\right)e^{2G_0d}.$$
(38)
We first note that this expression diverges as $`\xi 2^+`$, indicating the breakdown of the perturbation theory as the temperature is lowered. Furthermore, in contrast to Eq. (28), the range of $`\mathrm{\Pi }_{SR}`$ remains constant and the amplitude acquires a temperature dependence of $`1/T`$ (for large $`\xi 2`$), reminiscent of a high temperature expansion.
However, the above picture may be modified if the charges have melted into a Coulomb fluid via a dislocation-mediated melting transition before $`\xi 2^{}`$. If this is the case, further analysis is necessary to obtain a more complete picture of the high temperature phase. Although the spatial correlations in a system of coupled 2D Coulomb fluids are expected to be somewhat different from 2D Wigner crystals, the solid phase results above suggest a qualitative lower limit of $`\mathrm{\Gamma }13`$ at which $`\mathrm{\Pi }_{SR}`$ crosses over from low temperature to high temperature behavior. It may be of interest to note that in Ref. , an estimate for the upper limit of $`\mathrm{\Gamma }`$ at which the Poisson-Boltzmann equation breaks down is of the order of $`\mathrm{\Gamma }3`$. To describe the melting of coupled 2D Wigner crystals, excitations of dislocations must be introduced into the effective Hamiltonian Eq. (14) similar to what is done in Ref. . These considerations may help to establish an analytical theory of the attraction arising from counterion correlations, not captured by the Poisson-Boltzmann theory. The present formulation is a first step in that direction.
## V acknowledgment
We would like to thank Ramin Golestanian, T. C. Lubensky, A.W.W. Ludwig, and S. Safran for stimulating and helpful discussions. AL and PP acknowledge support from NSF grants MRL-DMR-9632716, DMR-9624091, and DMR-9708646. DL acknowledges support from Israel Science Foundation grant 211/97. HF acknowledges support from NSF Grant No. DMR-9870681.
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# The Semiclassical Approach to Small 𝑥 Physicsaafootnote aTalk given at the DIS-2000 conference, Liverpool, England, April 2000
## 1 Introduction
The semiclassical approach to deep inelastic scattering (DIS) at small $`x`$ exploits the target rest frame point of view $`^\mathrm{?}`$. In this frame the virtual photon interacts via its partonic fluctuation with the target. The target itself is considered as localized soft color field. In the small $`x`$ limit the partons are fast, and therefore their interaction with the target can be treated in an eikonal approximation $`^\mathrm{?}`$.
This picture of DIS allows a combined description of both inclusive and diffractive events. In particular in the case of diffractive scattering, like diffractive dissociation and the diffractive production of high-$`p_{}`$ jets, several interesting results can be obtained, even without explict numerical calculations.<sup>b</sup><sup>b</sup>bFor a comprehensive review of these topics see Ref. $`^\mathrm{?}`$ and references therein.
Here we focus on the most recent developement in the semiclassical approach, namely the determination of the gluon density at next-to-leading order (NLO) $`^\mathrm{?}`$, which can serve as input in the evolution equation. The gluon density is expressed in terms of a (non-perturbative) Wilson loop and can be evaluated in any model of the target color field.
## 2 The Gluon Distribution
To extract the gluon density it is convenient to use a ‘scalar photon’ (denoted by $`\chi `$) coupled directly to the gluon field $`^\mathrm{?}`$. The gluon density is then derived by matching the semiclassical and the parton model approach. To leading order this means that we have to equate the cross section for the transition $`\chi g`$ in an external field with the cross section of the process $`\chi gg`$ as given in the parton model. The result reads
$$xg^{(0)}(x,\mu ^2)=\frac{1}{12\pi ^2\alpha _s}d^2x_{}|\frac{}{y_{}}W_x_{}^𝒜(y_{})|_{y_{}=0}|^2,$$
(1)
where $`W^𝒜`$ indicates a Wilson loop in the adjoint representation describing the eikonalised interaction of a gluon in the external color field of the target. The gluon distribution $`xg^{(0)}(x,\mu ^2)`$ is a constant, and measures the averaged local field strength of the target.
At NLO, we write the gluon density as
$$xg(x,\mu ^2)=xg^{(0)}(x,\mu ^2)+xg^{(1)}(x,\mu ^2),$$
(2)
with $`xg^{(1)}(x,\mu ^2)`$ denoting the (scheme dependent) NLO correction. To extract this correction, the cross section for the transition $`\chi gg`$ in an external field has to be equated with the parton model cross section of the process $`\chi ggg`$. Without providing any details of the calculation, we quote here only the final result of the distribution in the $`\overline{\text{MS}}`$ scheme at NLO $`^\mathrm{?}`$,
$$xg^{(1)}(x,\mu ^2)=\frac{1}{\pi ^3}\left(\mathrm{ln}\frac{1}{x}\right)_{r^2(\mu )}^{\mathrm{}}\frac{dy_{}^2}{y_{}^4}\left\{d^2x{}_{}{}^{}\text{tr}W_x_{}^𝒜(y_{})\right\}.$$
(3)
The scheme dependence enters through the short-distance cutoff $`r^2(\mu )`$. The NLO gluon density shows a $`\mathrm{ln}(1/x)`$ enhancement at small $`x`$, and is sensitive to the large-distance structure of the target.
Using the model of a large hadron to describe the color field of the target, allows a comparison of our result with the one of Mueller $`^{\mathrm{?},\mathrm{?}}`$. We find agreement for both the integrated distribution in (3) and the unintegrated one $`^\mathrm{?}`$ not shown here. However, in Refs. $`^{\mathrm{?},\mathrm{?}}`$, where the main focus is on parton saturation, the scale dependence has not been discussed. Therefore, we provide for the first time a quantitative relation between the short-distance cutoff in Eq. (3) and the scale of the gluon distribution, which can only be achieved by matching the semiclassical approach with a treatment in the parton model.
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# Schemata Evolution and Building Blocks
## 1 Introduction
One of the most commonly asked questions about genetic algorithms (GAs) is: under what circumstances do GAs work well? Obviously an answer to this question would help immeasurably in knowing to which problems one can apply a GA and expect a high level of performance. However, to answer this question one has to answer a more fundamental question: how do GAs work? For example, in a typical optimization problem how does the GA arrive at a good solution? It is clear that in very complex problems this is not achieved via a random search in the state space. The search is structured. However, the question remains as to what is the nature of this structure. To put this question another way, if we think of individual string bits as “degrees of freedom”, the GA does not exhaustively search through the different combinations of individual bits, i.e. a search in the entire state space. Rather it searches through a restricted space spanned by different combinations of “effective degrees of freedom” (EDOF), which are combinations of the more fundamental “microscopic” bit degrees of freedom.
What exact form these EDOF take depends of course on the particular landscape under consideration. Hence, one might despair as to whether it was possible to say anything that applied to more than a specific case. However, it is not meaningless to try to understand if they exhibit generic properties, independent of the landscape, or at least properties that are common to a large class of possible landscapes. The building block hypothesis and the schema theorem , attempt to identify such generic features and as such have played an important role in GA theory, if one accepts that one of the principal goals of a theory is to provide a framework within which one can gain a qualitative understanding of the behavior of a system. The basic gist of the building block hypothesis is that short, low-order, highly fit schemata play a preeminent role in the evolution of a GA; i.e. that the relevant EDOF for a GA are short, low-order, highly fit schemata. The schema theorem tries to lend a more quantitative aspect to the hypothesis by showing that such schemata are indeed favored. This fact is deduced via an analysis of the destructive effects of crossover. However, as is well known, the schema theorem is an inequality and is such because it does not say anything precise about schema reconstruction. To understand better the interplay between schema destruction, schema reconstruction and schema length one requires an evolution equation that is exact, and where schemata are the fundamental objects considered.
Various exact evolution equations have been derived previously: wrote down exact equations for two-bit problems. Later these equations were extended to three and four-bit problems citewhitley. These equations allowed for an explicit analysis of string gains and losses. also presented an algorithm for generating evolution equations for larger problems that was equivalent to an earlier equation of Bridges and Goldberg . Although exact these equations are extremely unwieldy and it is difficult to infer general conclusions from their analysis. Another related approach is that of Vose and collaborators , , that treats GA evolution as a Markov chain. One of the chief drawbacks of all the above, with respect to an analysis of the schema theorem and the block hypothesis, is that the former are evolution equations for strings whereas the latter refer to schemata. Evidently an evolution equation that is amenable to interpretation and analysis that treats schemata as fundamental objects would be preferable. Such an equation has been derived recently <sup>1</sup><sup>1</sup>1After the completion of this work we became aware of the related work of Altenberg . , for the case of proportional selection and 1-point crossover. The chief aim of this paper is to analyze the schema theorem and the building block hypothesis in the light of this equation.
Crucially, we will be able to quantify the effect of schema reconstruction relative to that of schema destruction. Traditionally, crossover as a source of schema disruption has been emphasized , . This idea is at the heart of the schema theorem and the building block hypothesis. There has been some work towards a more positive point of view of crossover vis a vis reconstruction , but mainly in the light of the exploratory nature of crossover. Here, we will see exactly under what conditions schema reconstruction dominates destruction.
In analyzing the consequences of the evolution equation we will especially emphasize two ideas: effective fitness and EDOF. With respect to the former we will show that: if one thinks intuitively of fitness as representing the ability of a schema to propagate then effective fitness is a more relevant concept than the conventional idea of fitness. We will formulate a schema theorem in terms of the effective fitness showing that schemata with high effective fitness receive an exponentially increasing number of trials as a function of time. The second key idea, already mentioned, is that of EDOF. Generically one can think of a schema as an EDOF. However, schemata offer for every string a decomposition into $`2^N`$ different elements of a space with $`3^N`$ members. Not every decomposition will be useful. In fact, typically, only a small subset. So what do we mean by useful? EDOF, if they are to have any utility whatsoever, should not be very strongly coupled. This is a notion that is intimately associated with how epistasis is distributed in the problem. This type of thinking is common to many fields and generally is associated with the idea of finding a basis for a highly non-linear problem wherein it decomposes into a set of fairly independent sub-problems. An important feature of complex systems is that the EDOF are “scale” dependent. This scale dependence very often takes the form of a time dependence wherein the EDOF are different at different stages of evolution. This complicates life greatly in that if we find a useful decomposition of a problem at time $`t`$ we have no guarantee that it will remain a useful decomposition indefinitely into the future.
## 2 Coarse Graining and Schemata
As is well known to any scientist or engineer a good model of a system is one that captures the relevant features and ignores irrelevant details. The deemphasis of irrelevant details we can think of as a “coarse graining”. Of course, a great difficulty is that often what is relevant versus irrelevant depends on what one wants to say about the system, i.e. what level of description one requires. It also, more often than not, depends on time. One of the most obvious examples of this is evolution: the primitive constituents of life, amino acids, DNA, RNA etc, which represent the “microscopic” degrees of freedom, over time combined to form progressively more and more complicated EDOF such as cells, sponges, people etc. This evolution in time is intimately linked to an evolution in “scale” and a corresponding evolution in complexity.
In a GA specifying all the bits of a string gives us the most fine grained, microscopic description possible. For strings of size $`N`$ and a population of size $`n`$ there are $`Nn`$ degrees of freedom and, for a binary alphabet, $`O(n2^N)`$ possible population states. Consider the different classes of fitness maps that may be defined: first, $`f_G:GR^+`$, where $`G`$ denotes the space of genotypes (string states) and $`f_G`$ is the fitness function that assigns a number to a given genotype; second, $`f_Q:QR^+`$, where $`Q`$ is the space of phenotypes. These mappings may be explicitly time dependent. In fact, this will normally be the case when the “environment” is time dependent. They may also be injective or not, although the map $`f_Q`$ will usually be injective. If $`f_G`$ is many-to-one then there exist “synonymous” genotypes, i.e. the mapping is degenerate. If we assume there exists a map $`\varphi :GQ`$ between genotype and phenotype then we have $`f_G=f_Q\varphi `$, i.e. the composite map induces a fitness function on $`G`$. A schema, $`\xi `$, consists of $`N_2N`$ defined bits. The defining length of the schema, $`l`$, is the distance between its two extremal defining bits. The space of all schemata, $`S`$, may be partitioned according to schema order; i.e. $`S=_{N_2}S_{N_2}`$, where $`S_{N_2}`$ is the space of schemata of order $`N_2`$. The mapping $`g:GS`$ between strings and schemata is many-to-one. The degree of degeneracy of the map, $`g_{N_2}:GS_{N_2}`$, is $`2^{NN_2}`$. Except for the trivial case of a $`0`$-schema, maximal degeneracy occurs when $`N_2=1`$ where half of $`S`$ is mapped onto one schema. The fitness of a schema is the map $`f_S:SR^+`$, which is related to $`f_G`$ via the composite map $`f_Sg=f_G`$. Explicitly,
$`\overline{f}(\xi ,t)={\displaystyle \frac{\underset{C_i\xi }{}f(C_i,t)n(C_i,t)}{\underset{C_i\xi }{}n(C_i,t)}}`$ (1)
where $`f(C_i,t)`$ is the fitness of string $`C_i`$ at time $`t`$, $`n(C_i,t)`$ is the expected number of strings of type $`C_i`$ at time $`t`$ and the sums are over all strings in the population that contain $`\xi `$.
As mentioned, the total number of schemata for a binary alphabet is $`3^N`$. Why go to an even bigger space than the state space itself? One answer to this question is related to the idea of coarse graining. In defining a schema we average over all strings that contain the given schema. In such a sum we are summing over all possible values for the string bits $`C_i\xi `$ present in the population. A schema thus represents a coarse grained degree of freedom because we are forfeiting explicit information about the out of schema string bits. Clearly the lower the order of the schema the higher the degree of coarse graining, the maximal coarse graining being associated with the maximally degenerate schema where $`N_2=1`$.
A schema of order $`N_2`$ has only $`N_2`$ degrees of freedom and $`2^{N_2}`$ possible states. Given that one of the fundamental characteristics of complex systems is the existence of a large number of degrees of freedom and an exponentially large state space any methodology that purports to reduce the number of EDOF will prove very useful. To see this in the context of an explicit example let us say that we wish to calculate the average fitness in a GA evolving according to proportional selection with strings of size $`N`$, where $`N`$ is a multiple of $`2`$. The evolution equation for the expected number of strings of type $`C_i`$, $`n(C_i,t)`$, is
$`n(C_i,t+1)={\displaystyle \frac{f(C_i,t)}{\overline{f}(t)}}n(C_i,t).`$ (2)
The average population fitness, $`\overline{f}(t)`$, for the case of a non-time dependent landscape obeys the equation
$`\overline{f}(t+1)={\displaystyle \underset{C_i}{}}{\displaystyle \frac{f^2(C_i)}{\overline{f}(t)}}P(C_i,t)`$ (3)
where $`P(C_i,t)=n(C_i,t)/n`$, $`n`$ being the population size which we regard as being constant. As proportional selection is a stochastic process, for small population sizes one will typically see large fluctuations, i.e. in any given experiment one may well see large deviations between the results of (2) and (3) and the corresponding experimental quantities. However, taking averages over repeated experiments the results converge to those of the above equations. In fact, in the infinite population limit $`P(C_i,t)`$ will converge to the probability of finding string $`C_i`$ at time $`t`$. The string fitness maps every string state to $`R^+`$. If the population is large then many strings will be represented and hence many terms in $`\overline{f}(t)`$ will be non-zero. Thus, to calculate the evolution of $`\overline{f}(t)`$ one needs to solve $`2^N`$ coupled equations. Let us instead take the following approach: we will average over odd string positions in the population leaving strings, $`C_i^{}`$, (or rather now schemata) of $`N/2`$ definite bits that satisfy
$`n(C_i^{},t+1)={\displaystyle \frac{f(C_i^{},t)}{\overline{f}(t)}}n(C_i^{},t)`$ (4)
where now the fitness of $`C_i^{}`$ depends on time even if $`f(C_i)`$ didn’t. To calculate $`\overline{f}(t)`$ one now only needs to solve $`2^{N/2}`$ coupled equations. One can repeat this process, each step of coarse graining reducing the number of EDOF by one half, until we reach the situation where a $`1`$-schema has been reached. This cannot be further coarse grained of course, except by passing to the trivial situation wherein all string bits are summed over, i.e. a $`0`$-schema. At this level the evolution equation for the effective string of size one ($`1`$-schema $`\alpha `$) is
$`n(\alpha ,t+1)={\displaystyle \frac{f(\alpha ,t)}{\overline{f}(t)}}n(\alpha ,t)`$ (5)
Now, $`\alpha `$ can only take two values, $`1`$ and $`0`$ say, hence $`\overline{f}(t)=[f(1,t)n(1,t)+f(0,t)n(0,t)]/n`$. Thus, the problem of finding the average fitness has been reduced to that of solving a problem with one degree of freedom and two possible states!
So what’s the catch? The principal, and more fundamental, problem is that the genetic operators, principally reproduction and crossover, are defined at the microscopic level. In other words, as can be seen in (1), to assign a fitness to a schema one has to sum over the different strings in the population that contain it. Thus, to calculate quantities associated with the coarse grained degrees of freedom one must consider the microscopic degrees of freedom. One might be tempted therefore to think that even though there is an apparent reduction in the number of EDOF the net gain is canceled out by the fact that one has to return to the microscopic degrees of freedom in order to calculate their evolution under the genetic operators. If one wished to calculate the dynamics exactly then the above would be true. However, returning to the idea of emphasizing the relevant degrees of freedom it may well be that in the averaging process certain ones are more important than others, therefore allowing one to neglect, or treat as a perturbation, the effect of the irrelevant ones. In particular, near a fixed point of the dynamics one might well expect to see a simplification.
A second problem is that if we wish to ask a question about a particular string and we only have access to schemata of order $`N_2<N`$ then the question will be impossible to answer. In other words, if we are going to accept a coarse grained description then we can only ask questions about coarse grained variables. This in no way will affect the calculation of population variables such as average population fitness, standard deviation about the average fitness etc. Neither should it affect the ability of the GA to find an optimum as a fixed point of the dynamics as this can be represented in terms of optimal schemata.
In the above we discussed a particular coarse graining which led to a certain, definite set of schemata of order $`N/2`$, $`N/4,\mathrm{},1`$ associated with averaging over the odd bits of each successive coarse grained string. Generally there are very many different coarse grainings possible, $`3^N1`$ for a given string. Which are useful and which aren’t? This depends on the fitness landscape under consideration. What one wishes to do is to choose a coarse graining that gives rise to EDOF that are relatively weakly coupled. Finding such a coarse graining may well of course be very difficult. The coarse graining by factors of $`2`$ above is a proposal for an algorithm to calculate GA evolution. Whether this particular coarse graining would be useful, as mentioned, depends on the fitness landscape. Although the method might seem somewhat artificial it is important to emphasize that such methods, based on the idea of the renormalization group (see for example Goldenfeld (1989) for a review) have proved to be extremely effective in many areas of physics and applied mathematics and have yielded very good results on canonical optimization problems such as the Traveling Salesman problem.
Later we will emphasize that $`1`$-schemata are very useful coarse grained variables, as being of size $`1`$ they are immune to the effects of crossover. In terms of $`1`$-schemata the average fitness in the population is
$`\overline{f}(t)={\displaystyle \underset{i=\alpha }{\overset{N}{}}}f(\alpha ,t)P(\alpha ,t)`$ (6)
where the sum is over the $`N`$ possible $`1`$-schemata, $`f(\alpha ,t)`$ is the fitness of the $`1`$-schema $`\alpha `$ at time $`t`$ and $`P(\alpha ,t)=n(\alpha ,t)/n`$ is the expected proportion of strings present in the population that contain the $`1`$-schemata $`\alpha `$.
## 3 Schema Equation
In this section we will review the derivation of the schema evolution equation of , . Given that the microscopic degrees of freedom are strings however, we will first derive an equation for strings evolving under the effects of the three genetic operators: proportional selection, crossover and mutation. We will throughout only consider simple one-point crossover. The analysis can be repeated, with analogous results, for the case of $`n`$-point crossover.
We will consider the change in the expected number, $`n(\xi ,t)`$, of strings that contain a particular schema $`\xi `$, of order $`N_2`$ and length $`lN_2`$, as a function of time (generation). If mutation is carried out after crossover one finds that the expected relative proportion of $`C_i`$ in the population, $`P(C_i,t)=n(C_i,t)/n`$, satisfies
$`P(C_i,t+1)=𝒫(C_iC_i)P_c(C_i,t)+{\displaystyle \underset{C_jC_i}{}}𝒫(C_jC_i)P_c(C_j,t)`$ (7)
where the effective mutation coefficients are: $`𝒫(C_iC_i)=_{k=1}^N(1p_m(k))`$, which is the probability that string $`i`$ remains unmutated, and $`𝒫(C_jC_i)`$, the probability that string $`j`$ is mutated into string $`i`$ given by
$`𝒫(C_jC_i)={\displaystyle \underset{k\left\{C_jC_i\right\}}{}}p_m(k){\displaystyle \underset{k\left\{C_jC_i\right\}_c}{}}(1p_m(k))`$ (8)
where $`p_m(k)`$ is the mutation probability of bit $`k`$. For simplicity we assume it to be constant, though the equations are essentially unchanged if we also include a dependence on time. $`\left\{C_jC_i\right\}`$ is the set of bits that differ between $`C_j`$ and $`C_i`$ and $`\left\{C_jC_i\right\}_c`$, the complement of this set, is the set of bits that are the same. In the limit where $`p_m`$ is uniform, $`𝒫(C_iC_i)=(1p_m)^N`$ and $`𝒫(C_jC_i)=p_m^{d^H(i,j)}(1p_m)^{Nd^H(i,j)}`$, where $`d^H(i,j)`$ is the Hamming distance between the strings $`C_i`$ and $`C_j`$. The quantity $`P_c(C_i,t)`$ is the expected proportion of strings of type $`C_i`$ in the population after selection and crossover. Explicitly
$`P_c(C_i,t)=P^{}(C_i,t){\displaystyle \frac{p_c}{N1}}{\displaystyle \underset{C_jC_i}{}}{\displaystyle \underset{k=1}{\overset{N1}{}}}𝒞_{C_iC_j}^{(1)}(k)P^{}(C_i,t)P^{}(C_j,t)`$ (9)
$`+{\displaystyle \frac{p_c}{N1}}{\displaystyle \underset{C_jC_i}{}}{\displaystyle \underset{C_lC_i}{}}{\displaystyle \underset{k=1}{\overset{N1}{}}}𝒞_{C_jC_l}^{(2)}(k)P^{}(C_j,t)P^{}(C_l,t)`$
where $`p_c`$ is the crossover probability, $`k`$ is the crossover point, and the coefficients $`𝒞_{C_iC_j}^{(1)}(k)`$ and $`𝒞_{C_jC_l}^{(2)}(k)`$, represent the probabilities that, given that $`C_i`$ was one of the parents, it is destroyed by the crossover process, and the probability that given that neither parent was $`C_i`$ it is created by the crossover process. Explicitly
$`𝒞_{C_iC_j}^{(1)}(k)=\theta (d_L^H(i,j))\theta (d_R^H(i,j))`$ (10)
and
$`𝒞_{C_jC_l}^{(2)}(k)={\displaystyle \frac{1}{2}}[\delta (d_L^H(i,j))\delta (d_R^H(i,l))+\delta (d_R^H(i,j))\delta (d_L^H(i,l))]`$ (11)
where $`d_R^H(i,j)`$ is the Hamming distance between the right halves of the strings $`C_i`$ and $`C_j`$, “right” being defined relative to the crossover point, with the other quantities in (10) and (11) being similarly defined. $`\theta (x)=1`$ for $`x>0`$ and is $`0`$ for $`x=0`$, whilst $`\delta (x)=0x0`$ and $`\delta (0)=1`$. Finally, $`P^{}(C_i,t)=(f(C_i,t)/\overline{f}(t))P(C_i,t)`$ gives the expected proportion of strings $`C_i`$ after the selection step. Note that $`𝒞_{C_iC_j}^{(1)}(k)`$ and $`𝒞_{C_jC_l}^{(2)}(k)`$ are properties of the crossover process itself and therefore population independent. The equation (7) yields an exact expression for the expectation values, $`n(C_i,t)`$, and in the limit $`n\mathrm{}`$ yields the correct probability distribution governing the GA evolution.
The evolution equation we have derived takes into account exactly the effects of destruction and reconstruction of strings and, at least at the formal level, has the same content as other exact formulations of GA dynamics . It should also be formally equivalent to the equation of Bridges and Goldberg . Before passing to the case of schemata it is interesting to put the equation into a simpler form. To see this, consider first the destruction term. The matrix (10) restricts the sum to those $`C_j`$ that differ from $`C_i`$ in at least one bit both to the left and to the right of the crossover point. One can convert the sum over $`C_j`$ into an unrestricted sum by subtracting off those $`C_j`$ that have $`d_L^H(i,j)=0`$ and/or $`d_R^H(i,j)=0`$. Similarly one may write the reconstruction term as
$`{\displaystyle \underset{C_jC_i^L}{}}{\displaystyle \underset{C_lC_i^R}{}}P^{}(C_j,t)P^{}(C_l,t)`$ (12)
where $`C_i^L`$ is the part of $`C_i`$ to the left of the crossover point and correspondingly for $`C_i^R`$. However, by definition
$`\overline{f}(C_i^L,t)={\displaystyle \frac{\underset{C_jC_i^L}{}f(C_j,t)n(C_j,t)}{_{C_jC_i^L}n(C_j,t)}}`$ (13)
where $`_{C_jC_i^L}n(C_j,t)`$ is the total number of strings in the population that contain $`C_i^L`$. The final form of the string equation without mutation thus becomes
$`P(C_i,t+1)=P^{}(C_i,t){\displaystyle \frac{p_c}{N1}}{\displaystyle \underset{k=1}{\overset{N1}{}}}(P^{}(C_i,t)P^{}(C_i^L,t)P^{}(C_i^R,t))`$ (14)
with
$`P^{}(C_i^L,t)={\displaystyle \underset{C_jC_i^L}{}}P^{}(C_j,t)`$ (15)
and similarly for $`P^{}(C_i^R,t)`$. It is important to note here that in this form the evolution equation shows that crossover explicitly introduces the idea of a schema and the consequent notion of a coarse graining. $`C_i^L`$ and $`C_i^R`$ are schemata of order and length $`k`$ and $`Nk`$ respectively.
The analogous equation for schema evolution can be found by summing equation (7) over all strings that contain the schema of interest $`\xi `$. The result is
$`P(\xi ,t+1)=𝒫(\xi \xi )P_c(\xi ,t)+{\displaystyle \underset{\xi /_i}{}}𝒫(\xi /_i\xi )P_c(\xi /_i,t)`$ (16)
where
$`P_c(\xi ,t)=P^{}(\xi ,t){\displaystyle \frac{p_c}{N1}}{\displaystyle \underset{k=1}{\overset{l1}{}}}\left(P^{}(\xi ,t)P^{}(\xi _L,t)P^{}(\xi _R,t)\right)`$ (17)
and the sum in (16) is over all schemata $`\xi /_i`$ that differ by at least one bit from $`\xi `$ in one of the $`N_2`$ defining bits of $`\xi `$. All other quantities are the schema analogs of quantities defined in (7). The effective mutation coefficients $`𝒫(\xi \xi )`$ and $`𝒫(\xi /_i\xi )`$ are
$`𝒫(\xi \xi )=(1p_m)^{N_2}\mathrm{and}𝒫(\xi /_i\xi )=p_m^{d^H(\xi ,\xi /_i)}(1p_m)^{N_2d^H(\xi ,\xi /_i)}`$ (18)
where $`d^H(\xi ,\xi /_i)`$ is the Hamming distance between the schemata $`\xi `$ and $`\xi /_i`$.
A very interesting feature of the evolution equations we have presented is their form invariance under a coarse graining. Starting with the string equation any coarse graining to schemata of order $`N_2<N`$ yields an equation identical in form to that of its predecessor.
## 4 Effective Fitness
Having derived the schema evolution equation, before turning to an analysis of its many features, we wish to digress on the notion of fitness. The main intuitive idea behind fitness is that fitter parents have more offspring. In equation (16), neglecting for the moment mutation and crossover, taking the limit of a continuous time evolution one finds
$`P(\xi ,t)=P(\xi ,0)\mathrm{e}^{_0^ts_\xi 𝑑t^{}}`$ (19)
where $`s_\xi =\frac{\overline{f}(\xi ,t)}{\overline{f}(t)}1`$ is the selective advantage of the schema $`\xi `$. If $`s_\xi >0`$ the expected number of $`\xi `$ grows, whilst if $`s_\xi <0`$ it decreases. However, consider the following two simple cases. First, consider the effect of mutation without crossover in the context of a model that consists of $`2`$-schemata, $`11`$, $`01`$, $`10`$, $`00`$, where each schema can mutate to the two adjacent ones when the states $`11`$, $`10`$, $`00`$, $`01`$ are placed clockwise on a circle. For example, $`11`$ can mutate to $`10`$ or $`01`$ but not to $`00`$. This is evidently the limit where two-bit mutations are completely negligible compared to one-bit mutations. We assume a simple degenerate fitness landscape: $`f(11)=f(01)=f(10)=2`$, $`f(00)=1`$. In a random population, $`P(11)=\mathrm{}=P(00)=\frac{1}{4}`$. If there is uniform probability $`p_m`$ for each schema to mutate to an adjacent one then the evolution equation that describes this system is
$`P(i,t+1)=(12p_m)P^{}(i,t)+p_m(P^{}(i1,t)+P^{}(i+1,t))`$ (20)
For $`p_m=0`$ the steady state population is $`P(11)=P(01)=P(10)=1/3`$, $`P(00)=0`$. Thus we see the synonym symmetry of the landscape associated with the degeneracy of the states $`11`$, $`10`$ and $`01`$ is unbroken. However, for $`p_m>0`$, the schemata distribution at $`t=1`$, starting from a random distribution at $`t=0`$, is $`P(11)=2/7`$, $`P(01)=P(10)=(2p_m)/7`$, $`P(00)=(1+2p_m)/7`$. Thus, we see that there is an induced breaking of the landscape synonym symmetry due to the effects of mutation. In other words the population is induced to flow along what in the fitness landscape is a flat direction.
As a second example consider the $`2`$-schemata problem now with crossover but neglecting mutation, and with a fitness landscape where $`f(01)=f(10)=0`$ and $`f(11)=f(00)=1`$. The steady state solution of the schema evolution equation is
$`P(11)=P(00)={\displaystyle \frac{1}{2}}\left(1{\displaystyle \frac{p_c}{2}}{\displaystyle \frac{(l1)}{(N1)}}\right)P(01)=P(10)={\displaystyle \frac{p_c}{4}}{\displaystyle \frac{(l1)}{(N1)}}`$ (21)
For $`l=N`$ and $`p_c=1`$ we see that half the steady state population is composed of strings that have zero fitness!
Although the above examples are artificial they serve to make the point that the genetic operators can radically change the “effective” landscape in which the population evolves. The actual “bare” landscape associated purely with selection in the above offers very little intuition as to the true population evolution. Real populations can flow rapidly along flat directions and strings may be present even if they have zero fitness. To take this into account we propose using an “effective” fitness function , defined via
$`P(\xi ,t+1)={\displaystyle \frac{f_{\mathrm{eff}}(\xi ,t)}{\overline{f}(t)}}P(\xi ,t)`$ (22)
comparing with equation (16) one finds
$`f_{\mathrm{eff}}(\xi ,t)=𝒫(\xi \xi )\overline{f}(\xi ,t)+{\displaystyle \underset{\xi /_i}{}}𝒫(\xi /_i\xi ){\displaystyle \frac{P(\xi /_i,t)}{P(\xi ,t)}}\overline{f}(\xi /_i,t)`$
$`{\displaystyle \frac{p_c}{N1}}𝒫(\xi \xi )\overline{f}(t){\displaystyle \underset{k=1}{\overset{N1}{}}}\left({\displaystyle \frac{P^{}(\xi ,t)P^{}(\xi _L,t)P^{}(\xi _R,t)}{P(\xi ,t)}}\right)`$
$`{\displaystyle \frac{p_c}{N1}}{\displaystyle \underset{\xi /_i}{}}𝒫(\xi /_i\xi )\overline{f}(t){\displaystyle \underset{k=1}{\overset{N1}{}}}\left({\displaystyle \frac{P^{}(\xi /_i,t)P^{}(\xi /_{i_L},t)P^{}(\xi /_{i_R},t)}{P(\xi ,t)}}\right)`$ (23)
In the limit $`p_m0`$, $`p_c0`$ we see that $`f_{\mathrm{eff}}(\xi ,t)\overline{f}(\xi ,t)`$. The above also leads to the idea of an effective selection coefficient, $`s_{\mathrm{eff}}=f_{\mathrm{eff}}(\xi ,t)/\overline{f}(t)1`$, that measures directly selective pressure. If we think of $`s_{\mathrm{eff}}`$ as being approximately constant in the vicinity of time $`t_0`$, then $`s_{\mathrm{eff}}(t_0)`$ gives us the exponential rate of increase or decrease of growth of the schema $`\xi `$ at time $`t_0`$. In the limit of a continuous time evolution the solution of (22) is
$`P(\xi ,t)=P(\xi ,0)\mathrm{e}^{_0^ts_{\mathrm{eff}}𝑑t^{}}`$ (24)
In the case of the toy examples above: for mutations without crossover
$`f_{\mathrm{eff}}(i,t)=f_i+{\displaystyle \frac{p_m}{P(i,t)}}(f_{i1}P(i1,t)+f_{i+1}P(i+1,t)2f_iP(i,t))`$ (25)
At $`t=0`$, $`f_{\mathrm{eff}}(11,0)=2`$, $`f_{\mathrm{eff}}(01,0)=f_{\mathrm{eff}}(10,0)=2p_m`$ and $`f_{\mathrm{eff}}(00,0)=1+2p_m`$. Thus we see that the effective fitness function provides a selective pressure by selecting among the degenerate schemata those that have a higher probability to produce fit descendents.
Of course, the definition of effective fitness is not unique. Another natural definition follows from the split into those terms of the evolution equation that are linear in $`P(\xi ,t)`$ and those “source” terms that are independent of it. For instance, in the case of selection and crossover we have
$`P(\xi ,t+1)={\displaystyle \frac{f_{\mathrm{eff}}^{}(\xi ,t)}{\overline{f}(t)}}P(\xi ,t)+j(t)`$ (26)
where $`f_{\mathrm{eff}}^{}(\xi ,t)=(1p_c\frac{(l1)}{N1})\frac{\overline{f}(\xi ,t)}{\overline{f}(t)}`$ and $`j(t)=\frac{p_c}{N1}_{k=1}^{l1}P^{}(\xi _L,t)P^{}(\xi _R,t)`$. The corresponding effective selection coefficient is $`s_{\mathrm{eff}}^{}=((1p_c\frac{(l1)}{N1})\frac{\overline{f}(\xi ,t)}{\overline{f}(t)}1)`$. In the limit of a continuous time evolution (26) may be formally integrated to yield
$`P(\xi ,t)=\mathrm{e}^{_0^ts_{\mathrm{eff}}^{}(t^{})𝑑t^{}}P(\xi ,0)+\mathrm{e}^{_0^ts_{\mathrm{eff}}^{}(t^{})𝑑t^{}}{\displaystyle _0^t}j(t^{})\mathrm{e}^{_0^t^{}s_{\mathrm{eff}}^{}(t^{\prime \prime })𝑑t^{\prime \prime }}𝑑t^{}`$ (27)
## 5 Schema Theorem and Building Blocks
We now turn to a discussion of the schema theorem and the building block hypothesis. The standard “schema theorem” , or fundamental theorem of GAs, states that for a schema, $`\xi `$, of length $`l`$ evolving according to proportional selection and $`1`$-point crossover
$`P(\xi ,t+1)P^{}(\xi ,t)\left(1p_c\left({\displaystyle \frac{l1}{N1}}\right)N_2p\right),`$ (28)
and has the interpretation that schemata of higher than average fitness will be allocated exponentially more trials over time. The conventional schema theorem only provides us with a lower bound for the expected number of schemata due to the fact that it does not explicitly account for schema reconstruction. Equation (16), however, exactly takes into account the effect of schema reconstruction due to both mutation and crossover. Together with the definition of effective fitness in equation (23) of the previous section it allows one to state a new schema theorem:
Schema Theorem
$`P(\xi ,t+1)={\displaystyle \frac{f_{\mathrm{eff}}(\xi ,t)}{\overline{f}(t)}}P(\xi ,t)`$ (29)
The interpretation of this equation is clear and analogous to the old schema theorem: schemata of higher than average effective fitness will be allocated an “exponentially” increasing number of trials over time. We put the word exponentially in quotes as the real exponent, $`^ts_{\mathrm{eff}}𝑑t^{}`$, is not, except for very simple cases such as a flat fitness landscape, of the form $`\alpha t`$, where $`\alpha `$ is a constant. The illustrative examples of the last section show that there is potentially a strong difference between the standard selection based fitness and effective fitness as the latter takes into account the effect of all genetic operators. The fact that strings with zero selective fitness can receive an exponentially increasing number of trials shows quite clearly that effective fitness is a more relevant concept. In this sense our schema theorem does not just state the obvious — that fit schemata that are preserved by the crossover operator will prosper. Once again this emphasizes the role of the destructive effect of crossover. The novel element here is seeing exactly how important is schema reconstruction. In fact generically it is the dominant contribution.
The schema evolution equation we have derived possesses many interesting features one of the most interesting being the way that it relates evolution in time to different levels of coarse graining. To see this we first return to the string evolution equation (7). Up to this point we have presented our results in almost the most general way possible — for any type of landscape and taking into account both crossover and mutation. Throughout the rest of the paper we will concentrate more on the effect of crossover and thus neglect mutation. The reason for this is that we will mainly be concerned with the importance of schema length vis a vis the building block hypothesis. Mutation being a strictly local operator will not play a major role in this discussion. Note that this equation is written entirely in terms of the fundamental degrees of freedom — the strings. In passing to the form (14) we have performed a coarse graining by summing over all strings that contain $`C_i^L`$ irrespective of what lies to the right of the crossover point, and similarly for strings containing $`C_i^R`$. The implication is that the very nature of crossover imposes on us the idea of coarse graining, and more specifically the idea of a schema, given that $`C_i^L`$ and $`C_i^R`$ define schemata of order and size $`k`$ and $`Nk`$ respectively. In order to solve the equation (14) we need to know $`P(C_i^L,t)`$ and $`P(C_i^R,t)`$. However, these in turn obey evolution equations of the form
$`P(C_i^L,t+1)=P^{}(C_i^L,t){\displaystyle \frac{p_c}{N1}}{\displaystyle \underset{m=1}{\overset{k1}{}}}(P^{}(C_i^L,t)P^{}(C_i^{LL},t)P^{}(C_i^{LR},t))`$ (30)
where $`C_i^{LL}`$ and $`C_i^{LR}`$ are the left and right parts of $`C_i^L`$, left and right being defined relative to the crossover point $`m`$, where $`m<k`$. Now, $`C_i^{LL}`$ and $`C_i^{LR}`$ as schemata are more coarse grained than $`C_i^L`$, i.e. they are of lower order. Clearly this pattern of behavior continues, i.e. in order to calculate $`P(C_i,t+1)`$ one requires $`P(C_i^L,t)`$ and $`P(C_i^R,t)`$ which in their turn require $`P(C_i^{LL},t1)`$, $`P(C_i^{LR},t1)`$, $`P(C_i^{RL},t1)`$ and $`P(C_i^{RR},t1)`$ etc. For each step back in time we pass to more coarse grained degrees of freedom. $`C_i`$ thought of as a schema is of higher order than $`C_i^L`$ or $`C_i^R`$, which in their turn are of higher order than $`C_i^{LL}`$, $`C_i^{LR}`$, $`C_i^{RL}`$ and $`C_i^{RR}`$. So where does this process stop? The maximally coarse grained EDOF are $`1`$-schemata. It is not possible to cut a $`1`$-schemata and hence crossover is explicitly neutral, i.e. $`1`$-schemata obey the equation
$`P(i,t+1)=P^{}(i,t)`$ (31)
As a simple example consider a $`4`$-bit string $`ijkl`$. The hierarchical structure of one possible ancestral tree can be written as
$$\begin{array}{c}\text{ }\text{t+1}ijkl\text{ }\hfill \\ \text{ }\text{t}ijk,lij,kli,jkl\text{ }\hfill \\ \text{ }\text{t-1}ij,ki,jki,jk,ljk,lj,kl\text{ }\hfill \\ \text{ }\text{t-2}i,jj,kj,kk,l\text{ }\hfill \end{array}$$
This tree shows only the effect of the reconstruction term in the schema equation over the space of 3 generations. Of course there are many other processes that contribute to the appearance of $`ijkl`$ at time $`t+1`$ that involve various combinations of schemata destruction and reconstruction. As far as pure schemata reconstruction is concerned however we see that $`1`$-schemata play a privileged role as they represent the ultimate building blocks. For an $`N`$-bit string the maximum number of time steps before all ancestors are $`1`$-schemata is $`N1`$.
All the above equally applies to a generic schema, $`\xi `$, composed of schemata, $`\xi _L`$ and $`\xi _R`$ which in their turn are composed of the schemata $`\xi _{LL}`$, $`\xi _{RL}`$, $`\xi _{LR}`$ and $`\xi _{RR}`$ etc. It should be clear that the idea of building blocks is manifest in the very structure of our evolution equations. $`\xi _{LL}`$, $`\xi _{RL}`$, $`\xi _{LR}`$ and $`\xi _{RR}`$ are building blocks for $`\xi _L`$ and $`\xi _R`$ which in their turn are building blocks for $`\xi `$. The ultimate building blocks are of course the $`1`$-schemata. In the above example of a $`4`$-bit string or schema the four building blocks of order one, $`i`$, $`j`$, $`k`$ and $`l`$ combine to form building blocks of order two $`ij`$ and $`kl`$ which in turn combine with the building blocks of order one to form building blocks of order three, $`ijk`$ and $`jkl`$. The building blocks of order three combine with the blocks of order one and the blocks of order two combine together to give blocks of order four, and so the process continues.
In terms of the effective fitness, $`f_{\mathrm{eff}}^{}(\xi ,t)`$ introduced previously
$`P(\xi ,t)=\mathrm{e}^{_0^ts_{\mathrm{eff}}^{}(t^{})𝑑t^{}}P(\xi ,0)+{\displaystyle \frac{p_c\mathrm{e}^{_0^ts_{\mathrm{eff}}^{}(t^{})𝑑t^{}}}{N1}}{\displaystyle \underset{k=1}{\overset{N1}{}}}{\displaystyle _0^t}P^{}(\xi _L,t^{})P^{}(\xi _R,t^{})\mathrm{e}^{_0^t^{}s_{\mathrm{eff}}^{}(t^{\prime \prime })𝑑t^{\prime \prime }}𝑑t^{}`$ (32)
Up to now we have been able to analyze a general landscape. To arrive at more explicit, analytic formulae in an arbitrary landscape is prohibitively difficult. We will therefore temporarily restrict our attention to some more restrictive but simpler cases. We start with the case of a flat fitness landscape. In this case $`s_{\mathrm{eff}}=p_c(l1)/(N1)`$, hence
$`P(\xi ,t)=\mathrm{e}^{p_c\frac{(l1)}{(N1}t}P(\xi ,0)+\mathrm{e}^{p_c\frac{(l1)}{(N1)}t}{\displaystyle \frac{p_c}{N1}}{\displaystyle \underset{k=1}{\overset{N1}{}}}{\displaystyle _0^t}P^{}(\xi _L,t^{})P^{}(\xi _R,t^{})\mathrm{e}^{p_c\frac{(l1)}{(N1)}t^{}}𝑑t^{}`$ (33)
Notice that dependence on the initial condition, $`P(\xi ,0)`$, is exponentially damped unless $`\xi `$ happens to be a $`1`$-schema, the solution of the $`1`$-schemata equation being
$`P(i,t)=P(i,0)`$ (34)
An immediate consequence is that when considering the source term describing reconstruction the only non-zero terms that need to be taken into account are those which arise from $`1`$-schemata, as any higher order term will always have an accompanying exponential damping factor. Thus we see that the fixed point distribution for a GA with crossover evolving in a flat fitness landscape is
$`P^{}(\xi )=Lt_t\mathrm{}P(\xi ,t)={\displaystyle \underset{i=1}{\overset{N_2}{}}}P(i,0)`$ (35)
which is basically Geiringer’s Theorem in the context of schema distributions and simple crossover. We see here that the theorem appears in an extremely simple way as a consequence of the solution of the evolution equation.
Note that this fixed point distribution arises purely from the effects of reconstruction, the absence of which leads to a pure exponential damping and the unphysical behavior $`P(\xi )0`$. We can also see from the above that a version of Geiringer’s theorem will also hold in a more general non-flat landscape where selection is only very weak, where what we mean by weak is that $`\frac{\overline{f}(\xi ,t)}{\overline{f}(t)}(1+ϵ)`$ and $`ϵ<\frac{p_c\frac{(l1)}{(N1)}}{1p_c\frac{(l1)}{(N1)}}l>1`$. Under such circumstances once again anything other than a $`1`$-schema will be associated with an exponential damping factor. A distinction between the two cases however is that for a flat fitness landscape the fixed point is fixed by the initial proportions of the various $`1`$-schemata as there is no competition between them. Here, however, due to the non-trivial landscape certain $`1`$-schemata are preferred over others. A concrete example of such a landscape would be $`f_i=1+\alpha _i`$ where $`_i|\alpha _i|\frac{ϵ}{(2+ϵ)}`$ and $`f_i`$ is the fitness of the ith bit. Note there is no need to restrict to a linear fitness function here, arbitrary epistasis is allowed as long as it does not lead to large fitness deviations away from the mean. In this case $`\frac{\overline{f}(\xi ,t)}{\overline{f}(t)}<1+ϵ`$.
So what is the analog of the building block hypothesis here? Our schema theorem states that schemata of above average effective fitness will be allocated “exponentially” more trials over time. In the way the evolution equation is structured we see that the effective fitness in terms of the effects of crossover consists of a destruction term and a reconstruction term. Inherent in the structure of the reconstruction term is a form of the building block hypothesis — that higher order schemata are built from fit, shorter, lower order schemata. If $`P^{}(\xi ,t)>P^{}(\xi _L,t)P^{}(\xi _R,t)`$ then the effects of destruction will outweigh those of reconstruction, whilst if $`P^{}(\xi ,t)<P^{}(\xi _L,t)P^{}(\xi _R,t)`$ reconstruction will dominate. The content of this inequality is that reconstruction will dominate destruction if the probability to select the parts of a schema is greater than the probability to select the whole schema. Once again this is a general conclusion valid for any landscape. To give a more analytic slant we restrict to the case of two-schemata in a flat fitness landscape wherein one finds
$`s_{\mathrm{eff}}=p_c\left({\displaystyle \frac{l1}{N1}}\right)+p_c\left({\displaystyle \frac{l1}{N1}}\right){\displaystyle \frac{P(i,0)P(j,0)}{P(ij,t)}}`$ (36)
Thus we see that the effect of reconstruction is greater than that of destruction if $`i`$ and $`j`$ are negatively correlated. Notice that if reconstruction is more important then the contribution from the latter is maximized by maximizing the schema length, $`l`$. In other words large, rather than small, schemata are favored!
In general the fitness landscape itself induces correlations between $`\xi _L`$ and $`\xi _R`$. In this case there is a competition between the (anti-) correlating effect of the landscape and the mixing effect of crossover. Selection itself more often than not induces an anti-correlation between fit schemata parts, rather than a positive correlation. Indeed, in the neutral case of a non-epistatic landscape one has $`1+\frac{2N_2}{N}\delta f_\xi <(1+\frac{2N_L}{N}\delta f_{\xi _L})(1+\frac{2N_R}{N}\delta f_{\xi _R})`$ where $`\delta f_\xi `$, $`\delta f_{\xi _L}`$ and $`\delta f_{\xi _R}`$ are the fitness deviations of the schemata $`\xi `$, $`\xi _L`$ and $`\xi _R`$ from an average fitness which we have normalized to one half. Thus we see that selection induces an anti-correlation when $`\delta f_{\xi _L},\delta f_{\xi _R}>0`$ and hence in an uncorrelated initial population, $`P^{}(\xi ,t)<P^{}(\xi _L,t)P^{}(\xi _R,t)`$. This means that crossover plays an important role in allowing both parts of a successful schema to appear in the same individual. The effect of crossover is to weaken but not cancel completely the anti-correlations induced by selection and thus make it easier to find the whole schema. Indeed, it is possible to show that for a non-epistatic landscape that the contribution to population fitness from all schemata of length $`l`$, starting with a random initial population at time $`t`$, is independent of $`l`$ at time $`t+1`$ and is an increasing function of $`l`$ for large $`l`$ at time $`t+2`$ .
More complicated landscapes one has to examine on a case by case basis. Of course, it is always possible to invent a landscape where there is a fitness advantage associated with bits that are close together. However, it is equally easy to find one where there is a fitness advantage for bits that are widely separated. The non-epsitatic landscapes above are neutral in this respect and therefore any results about the nature of schemata and building blocks are a reflection of the geometric effect of crossover and not associated with bit-bit correlations induced by the landscape itself. We now have ample experimental evidence that this is the case as well for “generic” landscapes with epistasis such as the Kaufmann $`Nk`$ models. This evidence will be published elsewhere. A particularly interesting example of epistasis is deception as it has played an important role in the theory of GAs . The very nature of deception is such that the bits of a schema are less selected than the whole and hence we can see from the schema equation that in this circumstance destruction will outweigh reconstruction. However, this will only be totally deceptive if all possible schema reconstruction channels $`\xi _L+\xi _R\xi `$ are deceptive. For a schema of order $`N_2`$ there are $`N_21`$ such channels. Thus, for $`N_2`$ large it will typically be quite unlikely that all channels will be deceptive. If there exist non-deceptive channels then it is probable that the population will evolve in those directions. In fact, as the example of a two-schema shows, for every deceptive channel there is a non-deceptive one. One may explicitly see this from
$`P(11,t+1)=P^{}(11,t)p_c\left({\displaystyle \frac{l1}{N1}}\right)(P^{}(11,t)P^{}(00,t)P^{}(01,t)P^{}(10,t))`$
$`P(01,t+1)=P^{}(01,t)p_c\left({\displaystyle \frac{l1}{N1}}\right)(P^{}(01,t)P^{}(10,t)P^{}(11,t)P^{}(00,t))`$ (37)
Here $`11`$-channel deception, i.e. $`P^{}(11,t)P^{}(00,t)>P^{}(01,t)P^{}(10,t)`$, implies that the $`01`$-channel is non-deceptive. However, this is not much consolation if the $`11`$-schemata happens to be the optimum. If we start with a random population then $`11`$-channel deception is equivalent to the statement $`f_{\mathrm{eff}}(11)<f(11)`$. For something as simple as the two-schemata problem there is only a single $`11`$-channel. For the $`4`$-bit schemata $`ijkl`$ we see that there are six reconstruction channels in total. There are various ways to end at a totally deceptive problem. For instance, the three channels $`ijk+lijkl`$, $`ij+klijkl`$ and $`i+jklijkl`$ might all be deceptive. Alternatively all the $`1`$-schemata $``$ $`2`$-schemata channels might be deceptive. Generically the deviation of the effective fitness from the selective fitness will offer a reasonable measure of deception.
## 6 Conclusions
In this paper we have analyzed the Schema Theorem and the Building Block Hypothesis based on an exact evolution equation for GAs. At the level of the microscopic degrees of freedom, the strings, we established that the action of crossover by its very nature introduces the notion of a schema, the probability to reconstruct a given string being dependent on the probabilities for finding the right and left parts of the string relative to the crossover point in the two parents. These probabilities involve a coarse graining, i.e. an averaging over all strings that contain the constituent parts of the string, and hence represent schema probabilities. We saw that the same equation, after a suitable coarse graining, also described the evolution of any arbitrary schema.
One might enquire as to what advantages a formulation based on schemata, as has been presented here, has over other existing formulations, such as the Vose Markov chain model. Indeed, the value of schemata and the Schema theorem in understanding GA evolution has been seriously questioned There are many possible answers to this question: first a pragmatic one — that all “things” are made out of “building blocks”, whether they be tables, giraffes or computer programmes. Having an exact, amenable description of complex systems from the microscopic point of view is a vain hope. Complex systems and complex behaviour can much better be understood in terms of EDOF. EDOF, almost by definition, are much fewer in number than the microscopic degrees of freedom and hence, in principle, would offer a computationally simpler picture. However, the number of ways of combining the microscopic degrees of freedom into EDOF is very large, hence one might think that such a description is even more costly than one based on the microscopic degrees of freedom such as the Vose model. This would be true if in analysing the GA one had to search through all the possible “coarse grainings” available. For a given landscape, however, a preferred coarse graining will often suggest itself. Secondly, we believe strongly that approximation schemes for solving GA evolution equations will be much more forthcoming via a formulation in terms of schemata wherein one may appeal to all the intuition and machinery of the renormalization group.
We introduced the notion of effective fitness showing through explicit examples that it was a more relevant concept than pure selective fitness in governing the reproductive success of a schema. Based on this concept of effective fitness and our evolution equation we introduced a new schema theorem that showed that schemata of high effective fitness received an exponentially increasing number of trials as a function of time. We then went on to discuss the building block hypothesis. One of the more remarkable features of our equation is that it implicitly contains a version of the latter in that the structure of the reconstruction term relates in an ancestral tree the relation between a given schema and its more coarse grained ancestors as a function of time. This ancestral tree terminates at $`1`$-schemata, which are in some sense the ultimate building blocks as they cannot be destroyed by crossover. We also showed that generically there is no preference for short, low-order schemata. In fact if schema reconstruction dominates the opposite is true, typically large schemata will be favored. Only in deceptive problems does it generally seem that short schemata will be favored, and then only in totally deceptive problems as the system will tend to seek out the non-deceptive channels if they exist.
There are many points of departure from the present work to future research. On the theoretical side it will be very interesting to see if other exact results besides Geiringer’s theorem follow very simply from our evolution equation. A fundamental issue is trying to find approximation schemes within which the equations can be solved, as for a general landscape an exact solution will be impossible. In this respect, as mentioned, techniques familiar from statistical mechanics such as the renormalization group might well prove very useful. In fact the very structure of our evolution equation is very similar to that of a renormalization group equation, a theme we shall return to in a future publication. It is of course necessary to verify the equations numerically. Some work in this direction has already been done and further work has confirmed its qualitative conclusions . In this respect one has to tread carefully, as the interplay between selection and crossover can be very subtle as the work on Royal Road functions has shown. Although very simple we favor preliminary analytic analyses based on non-epistatic landscapes where one knows that there is no intrinsic inter-bit linkage due to the fitness landscape and therefore one can study the geometric effects of crossover in a more uncluttered environment.
### Acknowledgements
This work was partially supported through DGAPA-UNAM grant number IN105197. CRS wishes to thank Drs. David Goldberg and Michael Vose for enlightening discussions. We also thank the entire Complex Systems group at the ICN under NNCP (http://luthien.nuclecu.unam.mx/nncp) for maintaining a stimulating research atmosphere.
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# A New Look on the Electromagnetic Duality. Suggestions and Developments
## 1 Introduction
The recent developments in superstring and brane theories show that the concept of duality plays and probably will play more and more fundamental role in theoretical and mathematical physics. In today’s view duality has a very general sense and, shortly speaking, it is considered and understood as availability of (at least) two different ways to describe the same physical phenomena, thus, the set of all descriptions of the same physical situation factors over the equivalent ones and, in this way, the essential differences and peculiarities of the inequivalent descriptions are explicitly accounted for. This general approach requires a strong insight and a detailed knowledge of modern quantum and nonquantum field theories.
The roots of duality in field theory, as well as the roots of the whole field theory, can be found in classical electrodynamics, i.e. in Maxwell equations. Therefore, a detailed and thorough understanding of the duality nature of classical electrodynamics seems to be a very important first step in getting well acquainted with the duality notion. At first sight the electromagnetic duality seems quite simple and not much promising theoretical tool, but a closer look reveals the opposite. In fact, as we shall see later, a careful study of the duality brings us to the conclusion that the adequate mathematical model object of the electromagnetic field should have a more complicated structure than just the couple $`(𝑬,𝑩)`$, or $`𝐅_{\mu \nu }`$. As it is well known , one of the most frequently used ways of introducing duality makes use of complex electro-magnetic vectors $`(𝑬+i𝑩)`$ and transformation of the kind $`(𝑬+i𝑩)e^{i\phi }(𝐄+i𝑩)`$, but such a formal step seems closer to a peculiarity notice than to reveal an important feature of the theory, while the real presence of the $`^2`$-complex structure inside the electromagnetic theory seems to stay not fully understood as its structure property. Following such a line of consideration many questions stay unanswered, e.g. why complex vector fields on real manifolds; why $`i𝑩`$ and not $`i𝑬`$; why in the relativistic formulation $`i=\sqrt{1}`$ disappears and, instead, a complex structure (through the Hodge $``$) in the 6-dimensional module of 2-forms on the Minkowski space appears, etc.
As we shall see, an adequate (nonrelativistic or relativistic) formulation of classical electrodynamics really needs complex structure, but not complex vectors on real manifolds. The duality transformations turn to be those linear isomorphisms of the standard complex structure $``$ in $`^2`$, which are also isometries of the standard metric in $`^2`$. In our approach duality transformations appear as closer related to the invariance of the conserved quantities of the theory than to the invariance properties of the equations. And this should be expected since we can treat them as natural extension of the usual space-time isometries, which, as we know, are basic tools in formulating and computing the conservative characteristics of any physical field.
The first man who noticed the duality properties of electromagnetic equations was Heaviside . A further development of Heaviside’s notice was given later by Larmor . A more detailed study of electromagnetic duality was made by Rainich in the frame of General Relativity. A comparatively complete presentation may be found in the extensive paper of Misner and Wheeler also in the frame of General Relativity in connection with their attempt to geometrize classical physics and to give topological interpretation of charges. Electromagnetic duality has always been in sight of those trying to introduce magnetic charges and currents in the theory, ,. In one may find duality considerations in the frame of nonlinear electrodynamics of continuous media. A formal generalization for p-forms is given in . A modern consideration of electromagnetic duality, directed to superstring and brane theories may be found in ; the possible nonabelian generalizations are considered in .
In this paper we pursue three main purposes. First, we are going to give a brief nonrelativistic and relativistic reviews of what is usually called electromagnetic duality in vacuum and in presence of electric and magnetic charges, without referring to other physical theories. Second, we shall present a new understanding of this duality, leading to a new and, in our view, more adequate mathematical nature of the (nonrelativistic and relativistic) model objects plus complex structure. Finally, we shall represent classical electrodynamics entirely in terms of the new nonrelativistic and relativistic mathematical model objects.
## 2 Electromagnetic Duality
2.1 Nonrelativistic consideration. We consider first the pure field Maxwell equations
$$\mathrm{rot}𝑬+\frac{1}{c}\frac{𝑩}{t}=0,\mathrm{div}𝑩=0,$$
(1)
$$\mathrm{rot}𝑩\frac{1}{c}\frac{𝑬}{t}=0,\mathrm{div}𝑬=0.$$
(2)
First we note, that because of the linearity of these equations if $`(𝑬_i,𝑩_i),i=1,2,\mathrm{}`$ are a collection of solutions, then every couple of linear combinations of the form
$$𝑬=a_i𝑬_i,𝑩=b_i𝑩_i$$
(3)
(sum over the repeated $`i=1,2,\mathrm{}`$) with arbitrary constants $`(a_i,b_i)`$ gives a new solution.
We also note, that in the static case the pure field Maxwell equations reduce to:
$$𝐝𝑬=0,𝐝𝑩=0,𝐝𝑩=0,𝐝𝑬=0,$$
(4)
where $``$ is the euclidean Hodge $``$-operator, $`𝒅`$ is the exterior derivative and the vector fields $`𝑬`$ and $`𝑩`$ are identified with the corresponding 1-forms through the euclidean metric $`g`$. The substitution $`𝑬𝑬,𝑩𝑩`$, because of the relation $`^2=\pm id`$, turns the first equation of (4) into the fourth one and vice versa, the fourth - into the first one; also, the second equation is turned into the third one, and vice versa, the third - into the second one. Hence, in this special case we can talk about $``$-symmetry of the equations.
The important observation made by Heaviside , and later considered by Larmor , is that the substitution
$$𝑬𝑩,𝑩𝑬$$
(5)
transforms the first couple (1) of the pure field Maxwell equations into the second couple (2), and, vice versa, the second couple (2) is transformed into the first one (1). This symmetry transformation (5) of the pure field Maxwell equations is called special duality transformation, or SD-transformation. It clearly shows that the electric and magnetic components of the pure electromagnetic field are interchangeable and the interchange (5) transforms solution into solution. In the transformed solution the magnetic component is the former electric component and vice versa, i.e. the electric component may be considered as magnetic if needed, and then the magnetic component should be considered as electric. This feature of the pure electromagnetic field reveals its dual nature.
It is important to note that the SD-transformation (5) does not change the energy density $`8\pi 𝐰=𝑬^2+𝑩^2`$, the Poynting vector $`4\pi 𝑺=c(𝑬\times 𝑩)`$ , and the (nonlinear) Poynting relation
$$\frac{}{t}\frac{𝑬^2+𝑩^2}{8\pi }=\mathrm{div}𝑺.$$
Hence, from energy-momentum point of view two dual, in the sense of (5), solutions are indistinguishable.
Note that the substitution (5) may be considered as a transformation of the following kind:
$$(𝑬,𝑩)\begin{array}{cc}0& 1\\ 1& 0\end{array}=(𝑩,𝑬).$$
(6)
The following question now arises naturally: do there exist constants $`(a,b,m,n)`$, such that the linear combinations
$$𝑬^{}=a𝑬+m𝑩,𝑩^{}=b𝑬+n𝑩,$$
(7)
or in a matrix form
$$(𝑬^{},𝑩^{})=(𝑬,𝑩)\begin{array}{cc}a& b\\ m& n\end{array}=(a𝑬+m𝑩,b𝑬+n𝑩),$$
(8)
form again a vacuum solution? Substituting $`𝑬^{}`$ and $`𝑩^{\mathbf{}}`$ into Maxwell’s vacuum equations we see that the answer to this question is affirmative iff $`m=b,n=a`$, i.e. iff the corresponding matrix $`S`$ is of the form
$$S=\begin{array}{cc}a& b\\ b& a\end{array}.$$
(9)
The new solution will have now energy density $`𝐰^{}`$ and momentum density $`𝑺^{}`$ as follows:
$$𝐰^{}=\frac{1}{8\pi }\left(𝑬_{}^{}{}_{}{}^{2}+𝑩_{}^{}{}_{}{}^{2}\right)=\frac{1}{8\pi }(a^2+b^2)\left(𝑬^2+𝑩^2\right),$$
$$𝑺^{}=(a^2+b^2)\frac{c}{4\pi }𝑬\times 𝑩.$$
Obviously, the new and the old solutions will have the same energy and momentum if $`a^2+b^2=1`$, i.e. if the matrix $`S`$ is unimodular. In this case we may put $`a=\mathrm{cos}\alpha `$ and $`b=\mathrm{sin}\alpha `$, where $`\alpha =const`$, so transformation (8) becomes
$$\stackrel{~}{𝑬}=𝑬\mathrm{cos}\alpha 𝑩\mathrm{sin}\alpha ,\stackrel{~}{𝑩}=𝑬\mathrm{sin}\alpha +𝑩\mathrm{cos}\alpha .$$
(10)
Transformation (10) is known as electromagnetic duality transformation, or D-transformation. It has been a subject of many detailed studies in various aspects and contexts ,-. It also has greatly influenced some modern developments in non-Abelian Gauge theories, as well as some recent general views on duality in field theory, esp. in superstring and brane theories (classical and quantum). In the next section we shall study this transformation from a new point of view, following the idea that $`(𝑬,𝑩)`$ are two vector components of one mathematical object having some more complicated nature.
From physical point of view a basic feature of the D-transformation (10) is, that the difference between the electric and magnetic fields becomes non-essential: we may superpose the electric and the magnetic vectors, i.e. vector-components, of a general electromagnetic field to obtain new solutions. From mathematical point of view we see that Maxwell’s equations in vacuum, besides the usual linearity (3) mentioned above, admit also ”cross”-linearity, i.e. linear combinations of $`𝑬`$ and $`𝑩`$ of a definite kind define new solutions: with every solution $`(𝑬,𝑩)`$ of Maxwell’s vacuum equations a 2-parameter family of solutions can be associated by means of linear transformations given by matrices of the kind (9). If these matrices are unimodular, i.e. when $`a^2+b^2=1`$, then all solutions of the family have the same energy and momentum. In other words, the space of all solutions to the pure field Maxwell equations factors over the action (8) of the group of linear maps $`S:^2^2`$ represented by matrices of the kind (9).
It is well known that matrices $`S`$ of the kind (9) do not change the canonical complex structure $``$ in $`^2`$: we recall that if the canonical basis of $`^2`$ is denoted by $`(\epsilon ^1,\epsilon ^2)`$ then $``$ is defined by $`(\epsilon ^1)=\epsilon ^2`$, $`(\epsilon ^2)=\epsilon ^1`$, so if $`S`$ is given by (9) we have: $`S..S^1=`$. Hence, the electromagnetic D-transformations (10) coincide with the unimodular symmetries of the canonical complex structure $``$ in $`^2`$. This important in our view remark clearly points out that the canonical complex structure $``$ in $`^2`$ should be an essential element of classical electromagnetic theory, so we should not forget about it and in no way neglect it. Moreover, in my opinion, we must find an appropriate way to introduce $``$ explicitly in the theory, and if this is done we should expect a covariance with respect to general nondegenerate transformations (7), not just covariance with respect to transformations (9), which are usually considered as symmetries .
In presence of electric current $`𝒋_e`$ and electric charges $`\rho _e`$ Maxwell equations (1)-(2) lose this D-symmetry. In order to retain it magnetic charges with density $`\rho _m`$ and magnetic currents $`𝒋_m=\rho _m𝒗`$ are usually introduced, and of course, the Lorentz force is correspondingly modified. The new system of equations looks as follows:
$$\mathrm{rot}𝑬+\frac{1}{c}\frac{𝑩}{t}=\frac{4\pi }{c}𝒋_m,\mathrm{div}𝑩=4\pi \rho _m,$$
(11)
$$\mathrm{rot}𝑩\frac{1}{c}\frac{𝑬}{t}=\frac{4\pi }{c}𝒋_e,\mathrm{div}𝑬=4\pi \rho _e,$$
(12)
$$\mu _𝒗𝒗=\rho _e𝑬+\frac{1}{c}(𝒋_e\times 𝑩)+\rho _m𝑩\frac{1}{c}(𝒋_m\times 𝑬),$$
(13)
where $`\mu `$ is the mass density of the particles and it is assumed that particles do not vanish and do not arise. Now the whole, or extended, D-transformation looks in the following way:
$$\stackrel{~}{𝑬}=𝑬\mathrm{cos}\alpha 𝑩\mathrm{sin}\alpha ,\stackrel{~}{𝒋}_e=\stackrel{~}{q}_e𝒗=𝒋_e\mathrm{cos}\alpha 𝒋_m\mathrm{sin}\alpha =(\rho _e\mathrm{cos}\alpha \rho _m\mathrm{sin}\alpha )𝒗$$
(14)
$$\stackrel{~}{𝑩}=𝑬\mathrm{sin}\alpha +𝑩\mathrm{cos}\alpha ,\stackrel{~}{𝒋}_m=\stackrel{~}{q}_m𝒗=𝒋_e\mathrm{sin}\alpha +𝒋_m\mathrm{cos}\alpha =(\rho _e\mathrm{sin}\alpha +\rho _m\mathrm{cos}\alpha )𝒗.$$
(15)
Hence, the 1-parameter family of transformations (14-15) is a symmetry of the system (11-13).
The corresponding considerations concerning one particle carrying electric and magnetic charges may be found in . We shall omit this simple case.
Finally we note that D-transformations change the two well known invariants: $`I_1=(𝑩^2𝑬^2)`$ and $`I_2=2𝑬.𝑩`$ in the following way:
$$\stackrel{~}{I_1}=\stackrel{~}{𝑩}^2\stackrel{~}{𝑬}^2=(𝑩^2𝑬^2)\mathrm{cos}2\alpha +2𝑬.𝑩\mathrm{sin}2\alpha =I_1\mathrm{cos}2\alpha +I_2\mathrm{sin}2\alpha ,$$
(16)
$$\stackrel{~}{I_2}=2\stackrel{~}{𝑬}.\stackrel{~}{𝑩}=(𝑬^2𝑩^2)\mathrm{sin}2\alpha +2𝑬.𝑩\mathrm{cos}2\alpha =I_1\mathrm{sin}2\alpha +I_2\mathrm{cos}2\alpha .$$
(17)
It follows immediately that
$$\stackrel{~}{I_1}^2+\stackrel{~}{I_2}^2=I_1^2+I_2^2,$$
i.e. the sum of the squared invariants is a D-invariant.
2.2 Relativistic consideration. Recall Maxwell’s pure field equations in relativistic form
$$\mathrm{𝐝𝐅}=0,𝐝𝐅=0.$$
(18)
The Hodge $``$-operator is defined by the relation
$$\alpha \beta =\eta (\alpha ,\beta )\sqrt{|det(\eta _{\mu \nu })|}dxdydzd\xi ,$$
where $`\alpha `$ and $`\beta `$ are a $`p`$-forms on the Minkowski spacetime $`M`$, $`\xi =ct`$ is the time coordinate, and the Minkowski metric $`\eta `$ has signiture $`(,,,+)`$.
We note first 2 simple symmetries of (18).
$`1^o`$. The transformation $`𝐅𝐅`$ keeps the system (18) the same. This follows from the property $`(𝐅)=𝐅`$ of the $``$-operator when restricted on 2-forms.
$`2^o`$. Any conformal change of the Minkowski metric $`\eta f^2\eta `$, where $`f`$ is everywhere different from zero function on the Minkowski space $`M`$, keeps the restriction of the Hodge $``$ to 2-forms on $`M`$ the same, so (18) is conformally invariant.
Recalling the explicit form of $`𝐅`$ and $`𝐅`$ in canonical coordinates
$`𝐅=𝑩_3dxdy𝑩_2dxdz+𝑩_1dydz+𝑬_1dxd\xi +𝑬_2dyd\xi +𝑬_3dzd\xi `$ (19)
$`𝐅=𝑬_3dxdy𝑬_2dxdz+𝑬_1dydz𝑩_1dydz𝑩_2dyd\xi 𝑩_3dzd\xi `$ (20)
we see that $``$ replaces $`𝑬`$ with $`𝑩`$ and $`𝑩`$ with $`𝑬`$, i.e. the action of $``$ gives the SD-transformation. On the other hand an extended SD-transformation may be introduced by
$$(𝐅,𝐅)\begin{array}{cc}0& 1\\ 1& 0\end{array}=(𝐅,𝐅).$$
We note the difference: the $``$-operator transforms a $`p`$-form $`\beta `$ into a $`(4p)`$-form $`\beta `$, while the above SD-transformation transforms a couple of forms to another couple of forms.
As in the nonrelativistic case, this SD-transformation is readily extended to the full D-transformation
$$(𝐅,𝐅)(,\stackrel{~}{})=(𝐅,𝐅)\begin{array}{cc}\mathrm{cos}\alpha & \mathrm{sin}\alpha \\ \mathrm{sin}\alpha & \mathrm{cos}\alpha \end{array},$$
i.e.
$$=𝐅\mathrm{cos}\alpha +𝐅\mathrm{sin}\alpha ,\stackrel{~}{}=𝐅\mathrm{sin}\alpha +𝐅\mathrm{cos}\alpha .$$
(21)
The above transformations transform solutions to solutions. Moreover, in contrast to the nonrelativistic case, where linear combinations of $`𝑬`$ and $`𝑩`$ of special kind transform solutions to solutions, here every linear combination of $`𝐅`$ and $`𝐅`$, i.e. a transformation of the kind
$$_g=a𝐅+b𝐅,\stackrel{~}{}_g=m𝐅+n𝐅,$$
(22)
where the subscribe $`g`$ means ”general”, with arbitrary constants $`(a,b,m,n)`$ gives again a solution. As we shall see this is because some of the special properties of $`S`$ are now hidden in the $``$-operator through the pseudometric $`\eta `$, and the components $`𝐅`$ and $`𝐅`$ are interrelated.
It is important to note that transformation (21) keeps the energy-momentum tensor
$$𝐐_\mu ^\nu =\frac{1}{4\pi }[\frac{1}{4}𝐅_{\alpha \beta }𝐅^{\alpha \beta }\delta _\mu ^\nu 𝐅_{\mu \sigma }𝐅^{\nu \sigma }]=\frac{1}{8\pi }[𝐅_{\mu \sigma }𝐅^{\nu \sigma }(𝐅)_{\mu \sigma }(𝐅)^{\nu \sigma }],$$
(23)
and its divergence
$$_\nu 𝐐_\mu ^\nu =\frac{1}{4\pi }[𝐅_{\mu \nu }(_\sigma 𝐅^{\sigma \nu })+(𝐅)_{\mu \nu }(_\sigma (𝐅)^{\sigma \nu })]$$
(24)
the same: $`𝐐_\mu ^\nu (𝐅)=Q_\mu ^\nu ()`$, $`_\nu 𝐐_\mu ^\nu (𝐅)=_\nu 𝐐_\mu ^\nu ()`$.
We see that the relativistic formulation of Maxwell theory naturally admits general $`^2`$-covariance as far as transformations (22) are implied to act on two 2-forms of the kind $`(𝐅,𝐅)`$. As for the D-transformations, they are closely related to the symmetries of the energy-momentum quantities and relations.
The two quantities
$$(4\pi )^2Q_{\mu \nu }Q^{\mu \nu }=I_1^2+I_2^2,(4\pi )^2Q_{\mu \sigma }Q^{\nu \sigma }=\frac{1}{4}(I_1^2+I_2^2)\delta _\mu ^\nu $$
also enjoy the D-invariance. We note also that the eigen values of $`𝐐_\mu ^\nu `$ are D-invariant, as it should be, while the eigen values of $`𝐅`$ and $`𝐅`$, given respectively by
$$\lambda _{1,2}=\pm \sqrt{\frac{1}{2}I_1+\frac{1}{2}\sqrt{I_1^2+I_2^2}},\lambda _{3,4}=\pm \sqrt{\frac{1}{2}I_1\frac{1}{2}\sqrt{I_1^2+I_2^2}},$$
$$\lambda _{1,2}^{}=\pm \sqrt{\frac{1}{2}I_1+\frac{1}{2}\sqrt{I_1^2+I_2^2}},\lambda _{3,4}^{}=\pm \sqrt{\frac{1}{2}I_1\frac{1}{2}\sqrt{I_1^2+I_2^2}}$$
are not D-invariant. Only when $`I_1=I_2=0`$, the so called null field case, the eigen values of $`𝐅`$ and $`𝐅`$ are D-invariant since in this case they are zero.
The relativistic generalization of equations (11)-(13) is
$$_\nu 𝐅^{\nu \mu }=4\pi j_e^\mu ,$$
(25)
$$_\nu (𝐅)^{\nu \mu }=4\pi J_m^\mu ,$$
(26)
$$\mu c^2u^\nu _\nu u_\mu =𝐅_{\mu \nu }j_e^\nu (𝐅)_{\mu \nu }J_m^\nu ,$$
(27)
where $`\mu `$ is the invariant mass density, $`j_e^\mu =\rho _eu^\mu `$, $`J_m^\mu =(𝑱_m,J^4)=(𝒋_m,\rho _m)`$, and $`u^\mu `$ is the 4-velocity vector field. The generalized Lorentz force (13) is obtained from (27) after raising the index $`\mu `$ through $`\eta ^{\mu \nu }`$.
The above system (25)-(27) enjoys the following symmetry transformation:
$$𝐅𝐅;j_eJ_m;J_mj_e.$$
This invariance is a particular case of the more general invariance transformation given by relations (21) plus
$$j_e^{}=j_e\mathrm{cos}\alpha J_m\mathrm{sin}\alpha ,J_m^{}=j_e\mathrm{sin}\alpha +J_m\mathrm{cos}\alpha ,$$
which is readily checked. It should be noted that these invariances make use of the property $`^2=id`$ of the restricted to 2-forms Hodge $``$-operator, and of the constancy of the phase angle $`\alpha `$ in the D-transformation.
Finally we’d like to note the different physical sense of equation (27) compare to equations (25)-(26): equation (27) equalizes changes of energy-momentum densities, it is a direct differential form of the energy-momentum balance relation between the field and the particles, carrying mass and electric and magnetic charge; equations (25)-(26) are relativistic forms of the differential equivalents of the time changes of the flows through 2-surfaces of $`𝑬`$ and $`𝑩`$, moreover, these two equations identify quantities of different physical nature: it is hard to believe that differentiating field functions (the left hand side) we could obtain currents of charged mass particles (the right hand side).
## 3 The Suggestion and Developments
3.1 Nonrelativistic formulation. We summurize some of the D-features of the field description through Maxwell equations.
1. The D-invariance of the field equations is a mathematical representation of the dual electro-magnetic $`(𝑬,𝑩)`$-nature of the field. This dual nature is explicitly seen in the nonrelativistic form of Maxwell’s equations: $`𝑬`$ and $`𝑩`$ depend on each other but they can be always distinguished from each other.
2. All energy-momentum quantities and relations are D-invariant.
3. The D-transformation is represented by a rotation in a 2-dimensional vector space and acts through superposing the electric and magnetic components of the field and the corresponding currents.
4. The rotation angle $`\alpha `$ does not depend on the space-time coordinates.
5. The electro-magnetic duality transformation keeps and emphasises the field’s united nature.
The suggestion coming from these notices is that the electromagnetic field, considered as one physical object, has two physically distinguishable interrelated vector components, $`(𝑬,𝑩)`$, so the adequate mathematical model-object must have two vector components and must admit 2-dimensional linear transformations of its components, in particular, the 2-dimensional rotations should be closely related to the invariance properties of the energy-momentum characteristics of the field. But every 2-dimensional linear transformation requires a ”room where to act”, i.e. a 2-dimensional real vector space has to be explicitly pointed out. This 2-dimensional space has always been implicitly present inside the electromagnetic field theory, but has never been introduced explicitly. Let’s introduce it.
The electromagnetic field is mathematically represented (nonrelativisticaly) by an $`^2`$-valued differential 1-form $`\omega `$, such that in the canonical basis $`(\epsilon ^1,\epsilon ^2)`$ in $`^2`$ the 1-form $`\omega `$ looks as follows
$$\omega =𝑬\epsilon ^1+𝑩\epsilon ^2.$$
(28)
Remark. In (28), as well as later on, we identify the vector fields and 1-forms on $`^3`$ through the euclidean metric and we write, e.g. $`(𝑬𝑩)=𝑬\times 𝑩`$. Also, we identify $`(^2)^{}`$ with $`^2`$ through the euclidean metric.
Now we have to present equations (11)-(13) correspondingly, i.e. in terms of $`^2`$-valued objects.
We begin with the electric $`\rho _e`$ and magnetic $`\rho _m`$ densities, considering them as components of an $`^2`$ valued function $`𝒬`$, i.e.
$$𝒬=\rho _e\epsilon ^1+\rho _m\epsilon ^2.$$
(29)
The two currents $`𝒋_e`$ and $`𝒋_m`$, considered as 1-forms, become components of an $`^2`$-valued 1-form $`𝒥`$ as follows
$$𝒥=𝒋_e\epsilon ^1+𝒋_m\epsilon ^2.$$
(30)
As we mentioned earlier (p.4), the above assumption (28) requires a general covariance with respect to transformations in $`^2`$, so, the complex structure $``$ has to be introduced explicitly in the equations. In order to do this we recall that the linear map $`:^2^2`$ induces a map $`_{}:\omega _{}(\omega )=𝑬(\epsilon ^1)+𝑩(\epsilon ^2)=𝑩\epsilon ^1+𝑬\epsilon ^2`$. We recall also that every operator $`𝒟`$ in the set of differential forms is naturally extended to vector-valued differential forms according to the rule $`𝒟𝒟\times id`$, and $`id`$ is usually omitted. Having in mind the identification of vector fields and 1-forms through the euclidean metric we introduce now $``$ in Maxwell’s equations (11)-(13) through $`\omega `$ in the following way:
$$𝐝\omega \frac{1}{c}\frac{}{t}_{}(\omega )=\frac{4\pi }{c}_{}(𝒥),\delta \omega =4\pi 𝒬,$$
(31)
where $`\delta =𝐝`$ is the codifferential. Two other equivalent forms of (31) are given as follows:
$$𝐝\omega \frac{1}{c}\frac{}{t}_{}(\omega )=\frac{4\pi }{c}_{}(𝒥),\delta \omega =4\pi 𝒬,$$
$$𝐝_{}(\omega )+\frac{1}{c}\frac{}{t}\omega =\frac{4\pi }{c}𝒥,\delta \omega =4\pi 𝒬.$$
In order to verify the equivalence of (31) to Maxwell equations (11)-(13) we compute the marked operations. For the left-hand side of the first (31) equation we obtain
$$𝐝\omega \frac{1}{c}\frac{}{t}_{}(\omega )=(\mathrm{rot}𝑬+\frac{1}{c}\frac{𝑩}{t})\epsilon ^1+(\mathrm{rot}𝑩\frac{1}{c}\frac{𝑬}{t})\epsilon ^2,$$
and the right-hand side is
$$\frac{4\pi }{c}_{}(𝒥)=\frac{4\pi }{c}𝒋_m\epsilon ^1+\frac{4\pi }{c}𝒋_e\epsilon ^2$$
The second equation $`\delta \omega =4\pi 𝒬`$ is, obviously, equivalent to
$$\mathrm{div}𝑬\epsilon ^1+\mathrm{div}𝑩\epsilon ^2=4\pi \rho _e\epsilon ^1+4\pi \rho _m\epsilon ^2$$
since $`\delta =\mathrm{div}`$. Hence, (31) coincides with (11)-(13).
We shall emphasize once again that according to our general assumption (28) the field $`\omega `$ will have different representations in the different bases of $`^2`$. Changing the basis $`(\epsilon ^1,\epsilon ^2)`$ to any other basis $`\epsilon ^1^{}=\phi (\epsilon ^1),\epsilon ^2^{}=\phi (\epsilon ^2)`$, means, of course, that in equations (31) the field $`\omega `$ changes to $`\phi _{}\omega `$ and the complex structure $``$ changes to $`\phi \phi ^1`$. In some sense this means that we have two fields now: $`\omega `$ and $``$, but $``$ is given beforehand and it is not determined by equations (31). So, in the new basis the $``$-dependent equations of (31) will look like
$$𝐝\phi _{}\omega \frac{1}{c}\frac{}{t}(\phi \phi ^1)_{}(\phi _{}\omega )=\frac{4\pi }{c}(\phi \phi ^1)_{}(\phi _{}𝒥).$$
If $`\phi `$ is a symmetry of $`:\phi \phi ^1=`$, then we transform just $`\omega `$ to $`\phi _{}\omega `$.
In order to write down the Poynting energy-momentum balance relation we recall the product of vector-valued differential forms. Let $`\mathrm{\Phi }=\mathrm{\Phi }^ae_a`$ and $`\mathrm{\Psi }=\mathrm{\Psi }^bk_b`$ are two differential forms on some manifold with values in the vector spaces $`V_1`$ and $`V_2`$ with bases $`\{e_a\},a=1,\mathrm{},n`$ and $`\{k_b\},b=1,\mathrm{},m`$, respectively. Let $`f:V_1\times V_2W`$ is a bilinear map valued in a third vector space $`W`$. Then a new differential form, denoted by $`f(\mathrm{\Phi },\mathrm{\Psi })`$, on the same manifold and valued in $`W`$ is defined by
$$f(\mathrm{\Phi },\mathrm{\Psi })=\mathrm{\Phi }^a\mathrm{\Psi }^bf(e_a,k_b).$$
Clearly, if the original forms are $`p`$ and $`q`$ respectively, then the product is a $`(p+q)`$-form.
Assume now that $`V_1=V_2=^2`$ and the bilinear map is the exterior product:$`:^2\times ^2\mathrm{\Lambda }^2(^2)`$. Let’s compute the expression $`(\omega ,𝐝\omega )`$.
$$(\omega ,𝐝\omega )=(𝑬\epsilon ^1+𝑩\epsilon ^2,𝐝𝑬\epsilon ^1+𝐝𝑩\epsilon ^2)=(𝑬𝐝𝑩𝑩𝐝𝑬)\epsilon ^1\epsilon ^2=$$
$$=𝐝(𝑬𝑩)\epsilon ^1\epsilon ^2=𝐝((𝑬𝑩))\epsilon ^1\epsilon ^2=\delta (𝑬\times 𝑩)\epsilon ^1\epsilon ^2=$$
$$=\mathrm{div}(𝑬\times 𝑩)\epsilon ^1\epsilon ^2=\mathrm{div}(𝑬\times 𝑩)dxdydz\epsilon ^1\epsilon ^2.$$
Following the same rules we obtain
$$(\omega ,\frac{1}{c}\frac{}{t}_{}\omega )=\frac{1}{c}\frac{}{t}\frac{𝑬^2+𝑩^2}{2}dxdydz\epsilon ^1\epsilon ^2,$$
and
$$\frac{4\pi }{c}(\omega ,_{}𝒥)=\frac{4\pi }{c}(𝑬.𝒋_e𝑩.𝒋_m)dxdydz\epsilon ^1\epsilon ^2.$$
Hence, the generalized Poynting energy-momentum balance relation is given by
$$(\omega ,𝐝\omega \frac{1}{c}\frac{}{t}_{}\omega )=\frac{4\pi }{c}(\omega ,_{}𝒥).$$
(32)
As for the generalized Lorentz force $`\stackrel{}{}`$, staying on the right-hand side of eq.(13), it is presented by
$$\stackrel{}{}\epsilon ^1\epsilon ^2=[\frac{1}{c}(𝒋_e\times 𝑩𝒋_m\times 𝑬)+\rho _e𝑬+\rho _m𝑩]\epsilon ^1\epsilon ^2=\frac{1}{c}(\omega ,𝒥)+(\omega ,_{}𝒬).$$
(33)
Since the orthonormal 2-form $`\epsilon ^1\epsilon ^2`$ is invariant with respect to rotations (and even with respect to unimodular transformations in $`^2`$) we have the duality invariance of the above energy-momentum quantities and relations.
Hence, in our approach we have achieved a full covariance of the equations, given in the form (31). Indeed, the covariance with respect to arbitrary transformations in $`^3`$ is obvious, so we show now the covariance of (31) with respect to nonsingular linear transformations $`\phi :^2^2`$. Let $`\omega `$ satisfies (31), so we have to show that $`\phi _{}(\omega )=𝑬\phi (\epsilon ^1)+𝑩\phi (\epsilon ^2)`$ also satisfies (31). We apply $`\phi `$ from the left on (31) and obtain
$$\phi (𝐝\omega \frac{1}{c}\frac{}{t}_{}\omega )=(𝐝\phi _{}\omega \frac{1}{c}\frac{}{t}\phi _{}_{}\omega )=(𝐝\phi _{}\omega \frac{1}{c}\frac{}{t}\phi _{}_{}\phi _{}^1\phi _{}\omega )=$$
$$=(𝐝\phi _{}\omega \frac{1}{c}\frac{}{t}(\phi \phi ^1)_{}\phi _{}\omega )=(𝐝\frac{1}{c}\frac{}{t}_{}^{})\phi _{}\omega ,$$
where $`^{}=\phi \phi ^1`$. For the right-hand side of (31) we obtain
$$(\phi )_{}𝒥=(\phi \phi ^1)_{}\phi _{}𝒥=_{}^{}\phi _{}𝒥,$$
so, our assertion is proved.
Note the following simple forms of the energy density
$$\frac{1}{8\pi }(\omega ,_{}\omega )=\frac{𝑬^2+𝑩^2}{8\pi }\epsilon ^1\epsilon ^2,$$
and of the Poynting vector,
$$\frac{c}{8\pi }(\omega ,\omega )=\frac{c}{4\pi }𝑬\times 𝑩\epsilon ^1\epsilon ^2,$$
the D-invariance is obvious. As for the general $`^2`$ covariance of the second equation of (31) it is obvious.
Resuming, we may say that pursuing the correspondence: one physical object - one mathematical model-object, we came to the idea to introduce the $`^2`$-valued 1-form $`\omega `$ as the mathematical model-field. This, in turn, set the problem for general $`^2`$ covariance of the equations and this problem was solved through introducing explicitly the canonical complex structure $``$ in the dynamical equations (31) of the theory. The duality transformation appears now as an invariance property of the energy-momentum quantities and relations.
3.2 Relativistic formulation. As it was mentioned in the previous section in the relativistic formulation of Maxwell equations we have got already a general $`^2`$-covariance, but the two components, subject to the general $`^2`$-linear transformation, are not independent, in fact they are $`(𝐅,𝐅)`$. We look now at the situation from another point of view.
First we note that we shall follow the main idea of the nonrelativistic formulation, namely, to consider as a mathematical-model field some $`^2`$-valued differential form, being closely connected to the canonical complex structure $``$ in $`^2`$. But in contrast to the nonrelativistic case here we consider an $`^2`$-valued 2-form on the Minkowski space-time $`M=(^4,\eta )`$. In general such a 2-form $`\mathrm{\Omega }`$ looks as $`\mathrm{\Omega }=𝐅_1\epsilon ^1+𝐅_2\epsilon ^2=𝐅_a\epsilon ^a`$. Let now we are given two linear maps:
$$\mathrm{\Phi }:\mathrm{\Lambda }^2(M)\mathrm{\Lambda }^2(M),$$
$$\phi :^2^2.$$
These maps induce a map $`(\mathrm{\Phi },\phi ):\mathrm{\Lambda }^2(M,^2)\mathrm{\Lambda }^2(M,^2)`$ by the rule:
$$(\mathrm{\Phi },\phi )(\mathrm{\Omega })=(\mathrm{\Phi },\phi )(𝐅_a\epsilon ^a)=\mathrm{\Phi }(𝐅_a)\phi (\epsilon ^a).$$
It is natural to ask now is it possible the joint action of these two maps to keep $`\mathrm{\Omega }`$ unchanged, i.e. to have
$$(\mathrm{\Phi },\phi )(\mathrm{\Omega })=\mathrm{\Omega }.$$
In such a case the form $`\mathrm{\Omega }`$ is called $`(\mathrm{\Phi },\phi )`$-equivariant. If $`\phi `$ is a linear isomorphism and we identify $`\mathrm{\Phi }`$ with $`(\mathrm{\Phi },id)`$ and $`\phi `$ with $`(id,\phi )`$, we can equivalently write
$$\mathrm{\Phi }(\mathrm{\Omega })=\phi ^1(\mathrm{\Omega }).$$
If we specialize now: $`\phi =`$ we readily find that the $`(\mathrm{\Phi },)`$-equivariant forms $`\mathrm{\Omega }`$ must satisfy
$$(\mathrm{\Phi },)(\mathrm{\Omega })=\mathrm{\Phi }(𝐅_2)\epsilon ^1+\mathrm{\Phi }(𝐅_1)\epsilon ^2=𝐅_1\epsilon ^1+𝐅_2\epsilon ^2,$$
Hence, we must have $`\mathrm{\Phi }(𝐅_1)=𝐅_2`$ and $`\mathrm{\Phi }(𝐅_2)=𝐅_1`$, i.e. $`\mathrm{\Phi }^2=id`$. In other words, the property $`^2=id`$ is carried over to $`\mathrm{\Phi }`$: $`\mathrm{\Phi }^2=id`$. Since the Hodge $``$, restricted to 2-forms in Minkowski space, satisfies this condition, and according to expressions (19)-(20) in standard coordinates its action coincides with the special duality transformation, it is a natural candidate for $`\mathrm{\Phi }`$. Hence, working with $`(,)`$-equivariant 2-forms on Minkowski space, we can replace the action of $``$ with the action of the Hodge $``$-operator. And that’s why in relativistic electrodynamics we have general $`^2`$ covariance if we work with forms $`\mathrm{\Omega }`$ of the kind $`\mathrm{\Omega }=𝐅\epsilon ^1+𝐅\epsilon ^2`$. In the nonrelativistic formulation this is not possible to be done since we work there with 1-forms on $`^3`$ and no map $`\mathrm{\Phi }:\mathrm{\Lambda }^1(^3)\mathrm{\Lambda }^1(^3)`$ with the property $`\mathrm{\Phi }^2=id`$ exists, and we have to introduce the complex structure through $`^2`$ only.
Having in view these considerations our basic assumption for the algebraic nature of the mathematical-model object must read:
The electromagnetic field is (relativistically) represented by a $`(,)`$-equivariant 2-form $`\mathrm{\Omega }`$ on the Minkowski space-time:
$$\mathrm{\Omega }=𝐅\epsilon ^1+𝐅\epsilon ^2.$$
(34)
We note two algebraic properties of the forms of the kind (34). First, all such $`p`$-forms, where the basis is fixed, form a linear space. Further, if $`\mathrm{\Omega }`$ is of the form (34) with respect to some basis $`(\epsilon ^1,\epsilon ^2)`$, i.e. if $`(,)(\mathrm{\Omega })=\mathrm{\Omega }`$, this does not mean that it will have the same form with respect to any other basis $`(e_1,e_2)`$ of $`^2`$. But if these two bases are transformed into each other through a symmetry $`S`$ of $``$, i.e. If $`SS^1=`$, we shall prove that $`\mathrm{\Omega }`$ will have the same form (38) with respect to $`(e_1,e_2)`$. We shall use the notation: $`\mathrm{\Omega }_\epsilon 𝐅\epsilon ^1+𝐅\epsilon ^2`$. Now let $`\epsilon ^1=S(e^1)e^1=S^1(\epsilon ^1),\epsilon ^2=S(e^2)e^2=S^1(\epsilon ^2);SS^1=`$. The last relation is equivalent to any of the following two relations: $`S^1S=,S^1=S^1`$. We have also $`(\epsilon ^1)=\epsilon ^2`$ and $`(\epsilon ^2)=\epsilon ^1`$. So, we have to prove the relation: $`(,)(\mathrm{\Omega }_e)=\mathrm{\Omega }_e`$, where $`\mathrm{\Omega }_e=(id,S^1)(\mathrm{\Omega }_\epsilon )`$. First we prove $`(e^1)=e^2`$ and $`(e^2)=e^1`$:
$$SS^1(\epsilon ^1)=(\epsilon ^1)=\epsilon ^2=S(e^1)S^1(\epsilon ^2)=e^2=(e^1)(e^2)=e^1.$$
Further,
$$(,)\mathrm{\Omega }_e=(,S^1)\mathrm{\Omega }_\epsilon =(,S^1)\mathrm{\Omega }_\epsilon =(,S^1)(𝐅\epsilon ^2𝐅\epsilon ^1)=$$
$$=𝐅S^1(\epsilon ^2)𝐅S^1(\epsilon ^1)=𝐅e^2+𝐅e^1=\mathrm{\Omega }_e.$$
This means that the set of all $`(,)`$-equivariant forms is partitioned to nonoverlaping subclasses with respect to the action of the automorphisms of the complex structure $``$.
The relativistic pure field Maxwell equations (18), expressed through the $`(,)`$-equivariant 2-form $`\mathrm{\Omega }`$ have, obviously, general $`^2`$ covariance and are equivalent to
$$𝐝\mathrm{\Omega }=0.$$
(35)
In presence of electric and magnetic charges, making use of the definitions in the previous section, we introduce the generalized relativistic current $`𝒥_r`$ as follows
$$𝒥_r=j_e\epsilon ^1+J_m\epsilon ^2.$$
Hence, equations (25)-(26) acquire the form
$$𝐝\mathrm{\Omega }=4\pi 𝒥_r.$$
(36)
The generalized Lorentz force is given by
$$(𝒥_r,\mathrm{\Omega })=(𝐅_{\mu \sigma }j_e^\sigma (𝐅)_{\mu \sigma }J_m^\sigma )dx^\mu \epsilon ^1\epsilon ^2.$$
(37)
The divergence of the energy-momentum tensor $`𝐐_\mu ^\nu `$ is given by
$$(\delta \mathrm{\Omega },\mathrm{\Omega })=\frac{1}{4\pi }[𝐅_{\mu \nu }_\sigma 𝐅^{\sigma \nu }+(𝐅)_{\mu \nu }_\sigma (𝐅)^{\sigma \nu }]dx^\mu \epsilon ^1\epsilon ^2.$$
(38)
The energy-momentum tensor $`𝐐`$, considered as a symmetric 2-form $`𝐐_{\mu \nu }=𝐐_{\nu \mu }`$, is given by
$$(𝐐\epsilon ^1\epsilon ^2)(X,Y)=\frac{1}{8\pi }(i_X\mathrm{\Omega },i_Y_{}\mathrm{\Omega }),$$
(39)
where $`X`$ and $`Y`$ are 2 arbitrary vector fields, and $`i_X`$ is the inner product by the vector field $`X`$. Indeed,
$$i_X\mathrm{\Omega }=X^\mu 𝐅_{\mu \nu }dx^\nu \epsilon ^1+X^\mu (𝐅)_{\mu \nu }dx^\nu \epsilon ^2,$$
$$i_Y_{}\mathrm{\Omega }=\left[Y^\mu 𝐅_{\mu \nu }dx^\nu \right]\epsilon ^2\left[Y^\mu (𝐅)_{\mu \nu }dx^\nu \right]\epsilon ^1.$$
$$(i_X\mathrm{\Omega },i_Y_{}\mathrm{\Omega })=X^\mu Y^\nu [𝐅_{\mu \sigma }𝐅_\nu ^\sigma +(𝐅)_{\mu \sigma }(𝐅)_\nu ^\sigma ]dxdydzd\xi \epsilon ^1\epsilon ^2.$$
So, we obtain
$$\frac{1}{8\pi }(i_X\mathrm{\Omega },i_Y_{}\mathrm{\Omega })=\frac{1}{8\pi }X^\mu Y^\nu [𝐅_{\mu \sigma }F_\nu ^\sigma +(𝐅)_{\mu \sigma }(𝐅)_\nu ^\sigma ]\epsilon ^1\epsilon ^2.$$
The presence of the 2-form $`\epsilon ^1\epsilon ^2`$ introduces invariance with respect to unimodular transformations in $`^2`$.
Another definition of the energy-momentum tensor, making no use of the complex structure $``$, is through the canonical inner product $`g`$ in $`^2`$ as a bilinear map instead of $``$. Indeed, it is easy to see that the right-hand side of the relation
$$𝐐_{\mu \nu }X^\mu Y^\nu =\frac{1}{2}g(i(X)\mathrm{\Omega },i(Y)\mathrm{\Omega })$$
is equal to
$$\frac{1}{2}X^\mu Y^\nu [𝐅_{\mu \sigma }𝐅_\nu ^\sigma +(𝐅)_{\mu \sigma }(𝐅)_\nu ^\sigma ].$$
Resuming, we note the main differences with respect to the nonrelativistic formulation. First, the mathematical model-object is a 2-form $`\mathrm{\Omega }`$ on Minkowski space with values in $`^2`$, second, $`\mathrm{\Omega }`$ is $`(,)`$-equivariant. As for the usual duality transformations, they appear as particular $`^2`$-invariance properties of the conserved quantities and of the corresponding conservation relations.
The general conclusion of this section is that the $`^2`$ valued nonrelativistic 1-form $`\omega `$ and the relativistic 2-form $`\mathrm{\Omega }`$ seem to be natural and adequate mathematical model-objects of electromagnetic fields, while the duality transformations characterize the invariance properties of the conversed quantities and the corresponding conservation relations.
## 4 Conclusion
Here we are going to mention those points of the paper which from our point of view seem most important.
Classical electrodynamics works mainly with two concepts: charge and field. The charge carriers (called also field sources) are considered as point-like (or structureless) objects. The field is considered as generated by static or moving charges, and it is not defined at the points of its own source, hence, the point charges acquire topological sense. Passing to continuous charge distributions we write down currents on the right hand sides of Maxwell equations and forget about the topological nature of charges. Moreover, this identification of characteristics of the field represented by the left hand sides of the equations with characteristics of mass objects carrying electric (and possiply magnetic) charges, seems not well enough motivated from theoretical point of view. It would be more natural to write down equations which identify quantities of the same nuture, e.g. some energy-momentum balance relation between the field and the particles (recall our remark at the end of Sec.2).
The duality properties of the solutions reveal the internal structure of the field as having two vector components, which are
-differentially interrelated (through the equations), but
-algebraically distinguished.
Moreover, the adequate understanding of the duality properties requires explicitly introduced complex structure in the equations.
This resulted in making use of $`^2`$-valued differential forms, $`\omega `$ and $`\mathrm{\Omega }`$, as mathematical model objects, and corresponding complex structures $``$ and the relativistic Hodge $``$ restricted to 2-forms. The equations admit a full $`^2`$ covariance, while the duality properties appear as invariance properties of the conservative quantities and conservation relations. In fact, the action of the D-transformations in the linear space of vacuum solutions separates classes of solutions with the same energy-momentum.
Finally, we may expect that recognizing the structure of the field as a double vector-component one through its duality properties may open new ways of considerations and may generate new ideas and developments towards an appropriate nonlinearization of classical electrodynamics.
REFERENCES
1. Olive, D.,hep-th/9508089; Harvey, J.,hep-th/9603086; Gomez, C.,hep-th/9510023;
Verlinde, E.,hep-th/9506011
2. Heaviside, O., ”Electromagnetic Theory” (London, 1893); Electrical Papers, London, 1 (1892); Phyl.Trans.Roy.Soc., 183A, 423 (1893)
3. Larmor, J., Collected Papers. London, 1928
4. Rainich, G., Trans.Amer.Math.Soc., 27, 106 (1925)
5. Misner, C., Wheeler, J. Ann. of Phys., 2, 525 (1957)
6. ”The Dirac Monopole”, (Mir, Moscow, 1970)
7. Strajev, V., Tomilchik, L., ”Electrodynamics with Magnetic Charge”, (Nauka i Tehnika, Minsk, 1975)
8. Gibbons, G., Rasheed, D., hep-th/9506035
9. Chruscinski, D., hep-th/9906227
10. Chan Hong-Mo, hep-th/9503106; hep-th/9512173
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# Topological invariant variables in QCD
## 1. Introduction
The anomalous decay processes in both electrodynamics and chromodynamics are described by the spatial integral of the product of the magnetic and electric field strength tensors. In non-Abelian field theory, this integral is known as the time derivative of the winding number functional.
One of the main differences between QCD and QED is the fact of the non-invariance of the QCD winding number functional with respect to gauge transformations with stationary matrices (considered as maps of the coordinate space into the color group one ). This is a result of the nontrivial topological structure of the group of stationary gauge transformations. The condition that the functional for the degree of the map is a finite number (normalization) determines the ’winding number’ class of functions on which the gauge transformations act. These functions decrease at spatial infinity as O$`(1/r)`$. This means that the dynamic gluon fields behave similar to the fields of point charges (monopoles).
This difference between the classes of functions for dynamic fields in QED and QCD is the principal peculiarity of the gauge theory of strong interactions. The present paper is devoted to a discussion of the dynamical consequences of the ’winding number’ class of functions in QCD.
The present approach is based on the fact that the winding number class of functions contains the zero-mode of the Gauss law constraint . This zero-mode is identified with the winding number as a collective variable .
We derive an equivalent unconstrained system compatible with the simplest canonical quantization scheme in the form of Feynman’s path integral. From this path integral representation of the generating functional a unitary perturbation theory for non-Abelian gauge theories can be obtained (see ) which is similar intuitive Faddeev-Popov (FP) scheme .
When comparing the scheme which will be developed in this paper with the instanton approach the following basic differences are observed: (i) the approach is formulated in the Minkowski space instead of the Euclidean one, and (ii) the perturbation theory is formulated with physical states instead of the classical vacuum ones.
The paper is organized as follows. In Section 2 we discuss the statement of the problem in detail. Section 3 is devoted to the derivation of the equivalent unconstrained system with the zero-mode in the Yang-Mills theory. In Section 4, the generating functional for Green functions in the form of the Feynman integral is constructed for QCD as the basis for a solution of the problems of the hadronization and confinement.
## 2. The winding number class of functions
We consider the winding number functional
$$X[A]=\frac{1}{8\pi ^2}\underset{V}{}d^3xϵ^{ijk}Tr\left[\widehat{A}_i_j\widehat{A}_k\frac{2}{3}\widehat{A}_i\widehat{A}_j\widehat{A}_k\right]$$
(1)
for the gluon gauge fields $`\widehat{A}_\mu =g\frac{\lambda ^a}{2i}A_\mu ^a`$, in the non-Abelian $`SU_c(3)`$ theory with the action functional
$$W=d^4x\left\{\frac{1}{2}(G_{0i}^{a}{}_{}{}^{2}B_{i}^{a}{}_{}{}^{2})+\overline{\psi }[i\gamma ^\mu (_\mu +\widehat{A}_\mu )m]\psi \right\},$$
(2)
where $`\psi `$ and $`\overline{\psi }`$ are the fermionic quark fields. We use the conventional notations for the non-Abelian electric and magnetic fields
$$G_{0i}^a=_0A_i^aD_i^{ab}(A)A_0^b,B_i^a=ϵ_{ijk}\left(_jA_k^a+\frac{g}{2}f^{abc}A_j^bA_k^c\right),$$
(3)
as well as the covariant derivative $`D_i^{ab}(A):=\delta ^{ab}_i+gf^{acb}A_i^c`$.
The action (2) is invariant with respect to gauge transformations $`u(t,\stackrel{}{x})`$
$$\widehat{A}_i^u:=u(t,\stackrel{}{x})\left(\widehat{A}_i+_i\right)u^1(t,\stackrel{}{x}),\psi ^u:=u(t,\stackrel{}{x})\psi .$$
(4)
It is well-known that the fixation of the gauge in the classical equations of motion leaves the ambiguity in the choice of initial data for the gauge fields due to the invariance of the action with respect to the gauge transformations (4) with stationary matrices $`u(t,\stackrel{}{x})=v(\stackrel{}{x})`$.
The group of the stationary gauge transformations $`v(\stackrel{}{x})`$ in the coordinate space is topologically nontrivial and represents the group of three-dimensional paths lying on the three-dimensional space of the $`SU_c(3)`$-manifold with the homotopy group $`\pi _{(3)}(SU_c(3))=Z`$. The whole group of the stationary gauge transformations is split into topological classes marked by the integer number $`n`$ (the degree of the map) which counts how many times a three-dimensional path turns around the $`SU(3)`$-manifold when the coordinate $`x_i`$ runs over the space where it is defined. The stationary transformations $`v^n(\stackrel{}{x})`$ with $`n=0`$ are called the small ones; and those with $`n0`$
$$\widehat{A}_i^{(n)}:=v^{(n)}(\stackrel{}{x})\widehat{A}_i(\stackrel{}{x})v^{(n)}(\stackrel{}{x})^1+L_i^n,L_i^n=v^{(n)}(\stackrel{}{x})_iv^{(n)}(\stackrel{}{x})^1$$
(5)
the large ones.
In QCD, the winding number functional (1) is not invariant with respect to large gauge transformations (5)
$$X[A^{(n)}]=X[A]+𝒩_1[A,n]+𝒩_2[n],$$
(6)
where
$$𝒩_1[A,n]=\frac{1}{8\pi ^2}d^3xϵ^{ijk}Tr[_i(\widehat{A}_jL_k^n)],$$
(7)
$$𝒩_2[n]=\frac{1}{24\pi ^2}d^3xϵ^{ijk}Tr[L_i^nL_j^nL_k^n]=n.$$
(8)
It determines the degree of a map $`𝒩_2[n]`$ (see Eqs. (3.33), (3.36) in ). The degree of a map $`𝒩_2[n]=n`$ as the condition of the normalization means that the large transformations are given in the class of functions with the spatial asymptotics O$`(1/r)`$. Such a function $`L_i^n`$ (5) is given by
$$v^{(n)}(\stackrel{}{x})=\mathrm{exp}(n\widehat{\mathrm{\Phi }}(\stackrel{}{x})),\widehat{\mathrm{\Phi }}=i\pi \frac{\lambda _A^ax^a}{r}f_0(r),$$
(9)
where the antisymmetric SU(3) matrices are denoted as
$$\lambda _A^1:=\lambda ^2,\lambda _A^2:=\lambda ^5,\lambda _A^3:=\lambda ^7,$$
and $`r=|\stackrel{}{x}|`$. The function $`f_0(r)`$ satisfies the boundary conditions
$$f_0(0)=0,f_0(\mathrm{})=1,$$
(10)
so that the functions $`L_i^n`$ disappear at spatial infinity $``$ O$`(1/r)`$ but can have nonvanishing surface integrals in Eq. (6). We call the class of functions (9) the ’winding number’ class of functions.
The present paper is based on the evident fact that the dynamical field $`A_i`$ and its transformations $`L_i^n`$ both belong to the winding number class of functions.
The second fact is that the winding number class of functions includes the topological excitations of the gluon system as a whole in the form of the zero-mode $`(𝒵)`$ of the non-Abelian Gauss law constraint
$$\frac{\delta W}{\delta A_0}=0D_i^{ab}(A)G_{0i}^b=j_0^a,$$
(11)
where $`j_\mu ^a=g\overline{\psi }\frac{\lambda ^a}{2}\gamma _\mu \psi `$ is the quark current. The Gauss law takes the form of an inhomogeneous equation for the time-like component $`A_0`$ of the gauge field
$$(D^2(A))^{ac}A_{0}^{}{}_{}{}^{c}=D_i^{ac}(A)_0A_i^c+j_0^a.$$
(12)
A general solution of this inhomogeneous equation is a sum of the solution $`𝒵^a`$ of the homogeneous equation
$$(D^2(A))^{ab}𝒵^b=0,$$
(13)
i.e., a zero mode of the Gauss law constraint, and a particular solution $`\stackrel{~}{A}_0^a`$ of the inhomogeneous one (12)
$$A_0^a=𝒵^a+\stackrel{~}{A}_0^a.$$
(14)
It is the central point of our paper, that the zero-mode $`𝒵^a`$ at the spatial infinity can be represented in the form of the product of a new topological variable $`\dot{N}(t)`$ and a phase $`\mathrm{\Phi }_0(\stackrel{}{x})`$
$$\widehat{𝒵}(t,\stackrel{}{x})|_{\mathrm{asymptotics}}=\dot{N}(t)\widehat{\mathrm{\Phi }}_0(\stackrel{}{x}),$$
(15)
where the phase $`\mathrm{\Phi }_0(\stackrel{}{x})`$ belongs to the winding number class of functions (9)
$$\widehat{\mathrm{\Phi }}_0=i\pi \frac{\lambda _A^ax^a}{r}f_0(r),f_0(0)=0,f_0(\mathrm{})=1.$$
(16)
It is determined in the asymptotical region by the equation
$$(D^2)^{ab}(\mathrm{\Phi }_i)\mathrm{\Phi }_0^b(\stackrel{}{x})=0,$$
(17)
where $`\mathrm{\Phi }_i(\stackrel{}{x})`$ the asymptotics of the dynamical gluon fields $`A_i`$ given in the same class of functions
$$\widehat{A}_i(t,\stackrel{}{x})|_{\mathrm{asymptotics}}=\widehat{\mathrm{\Phi }}_i(\stackrel{}{x})=i\frac{\lambda _A^a}{2}ϵ_{iak}\frac{x^k}{r^2}f_1(r),$$
(18)
with the boundary conditions
$$f_1(0)=0,f_1(\mathrm{})=1.$$
(19)
In this case, the single one-parametric variable $`N(t)`$ reproduces the topological degeneracy of all field variables, provided the separation of the zero-mode phase factors of the topological degeneracy from the topological invariant variables $`A^I,\psi ^I`$ of perturbation theory (i.e., the variables without topological degeneracy). This separation is fulfilled by the gauge transformations
$$0=U_𝒵(\widehat{𝒵}+_0)U_𝒵^1,$$
(20)
$$\widehat{A}_i=U_𝒵(\widehat{A}^I+_i)U_𝒵^1,\psi =U_𝒵\psi ^I,$$
where the spatial asymptotics of $`U_𝒵`$ is
$$U_𝒵=T\mathrm{exp}[\stackrel{t}{}𝑑t^{}\widehat{𝒵}(t^{},\stackrel{}{x})]|_{\mathrm{asymptotics}}=\mathrm{exp}[N(t)\widehat{\mathrm{\Phi }}_0(\stackrel{}{x})].$$
(21)
A known example of fields $`(\mathrm{\Phi }_i,\mathrm{\Phi }_0)`$ which satisfy Eq. (17) is the Bogomol’ny-Prasad-Sommerfield (BPS) monopole $`(\mathrm{\Phi }_i^{BPS})`$ with the Higgs field $`(\mathrm{\Phi }_0^{BPS})`$
$`\mathrm{\Phi }_i^{BPS}=i{\displaystyle \frac{\lambda _A^a}{2}}ϵ_{iak}{\displaystyle \frac{x^k}{r^2}}f_1^{BPS}(r),f_1^{BPS}(r)`$ $`=`$ $`1{\displaystyle \frac{r}{ϵ\mathrm{sinh}(r/ϵ)}},`$ (22)
$`\widehat{\mathrm{\Phi }}_0^{BPS}=i\pi {\displaystyle \frac{\lambda _A^ax^a}{r}}f_0^{BPS}(r),f_0^{BPS}(r)`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{tanh}(r/ϵ)}}{\displaystyle \frac{ϵ}{r}},`$ (23)
in the limit $`ϵ0`$, where $`ϵ`$ is the size of the monopole. For $`ϵ0`$ the BPS monopole goes over into the Wu-Yang monopole without singularity at the origin (as the solution of classical equations).
In this case, the winding number functional (1) takes the form of the sum of the invariant term and the new topological variable $`N(t)`$
$$X[A]=X[A^I]+𝒩_1[\mathrm{\Phi }_i,N(t)]+𝒩_2[N(t)]=X[A^I]+N(t).$$
(24)
This can be seen from the fact that for noninteger $`n=N(t)`$ the responses of the winding number (1) are equal to
$$𝒩_1[\mathrm{\Phi }_i,N(t)]=\mathrm{sin}[2\pi N(t)]/2\pi ,$$
and
$$𝒩_2[N(t)]=\{N(t)\mathrm{sin}[2\pi N(t)]/2\pi \},$$
so that the sinus terms cancel each other in the sum (24). The function $`N(t)`$ represents a one-dimensional parametrization of the transition between different maps of the homotopy group.
The topological degeneration of initial data means that the points $`N(t)`$ and $`N(t)n`$ are physically equivalent. Thus, the configuration space of the physical variables of QCD has the topology of a cylinder where the role of the degree of freedom which rotates around the cylinder is played by the topological variable $`N(t)`$ as the zero mode of the first class constraint. This zero mode can be revealed by only the explicit resolving the Gauss law constraint . In the contrast to the zero mode of the second class constraint (i.e., gauge) treated as the Gribov ambiguity , the zero mode of the first class constraint is the inexorable consequence of internal dynamics, which can be implicitly lost by the standard gauge-fixing scheme.
The present paper is devoted to the derivation of an equivalent unconstrained system with the topological variable $`N(t)`$ and to the construction of the generating functional for a unitary perturbation theory.
## 3. Yang-Mills theory
### 3.1. Constraining
Recall that for a long time the problem of quantization of non-Abelian constrained systems was considered as the greatest challenge of relativistic quantum field theory . The success in the description of the quantum dynamics of such systems was the unitary perturbation theory in the form of the Faddeev-Popov functional integral used for the proof of the renormalizability of the Standard Model (Nobel Prize of 1999 for ’t Hooft and Veltman). However, the foundation of the intuitive Faddeev-Popov integral for the Yang-Mills (YM) theory was achieved by Faddeev on the way of the construction of an ’equivalent unconstrained system’ compatible with the simplest canonical scheme of quantization in the form of the Feynman integral over independent physical variables
$$Z_{YM}^{}[J]=\underset{t,x}{}\left\{\underset{a=1}{\overset{3}{}}\frac{[d^2A_a^Td^2E_a^T]}{2\pi }\right\}\mathrm{exp}\left\{iW_{YM}^{}+id^4xJ_i^aA_{i}^{a}{}_{}{}^{T}\right\}.$$
(25)
The action of this equivalent unconstrained system
$$W_{YM}^{}=d^4x\left\{E_{k}^{a}{}_{}{}^{T}\dot{A}_k^{aT}\frac{1}{2}\left\{(E_{k}^{a}{}_{}{}^{T})^2+B_k^2(A^T)+(_k\sigma ^a)^2\right\}\right\},$$
(26)
where $`\sigma ^a`$ satisfies the equation
$$D(A)_i^{bc}_i\sigma ^c=gϵ^{abc}A_i^aE_{i}^{c}{}_{}{}^{T},$$
is derived by the constraining
$$W_{YM}^{}=W_{YM}(\mathrm{constraint})$$
(27)
the initial YM action in the first order formalism
$$W_{YM}=d^4x\left\{F_{0i}^aG_{0i}^a\frac{1}{2}[F_{0i}^2+B_i^2]\right\}\left(G_{0i}^a=\dot{A}_i^aD_i^{ab}(A)A_0^b\right).$$
(28)
The constraining (27) means the substitution of the explicit solution of the Gauss law constraint (see in details )
$$D_i^{ab}(A)F_i^b=0$$
(29)
obtained by decomposing the electrical components of the field strength tensor $`F_{0i}`$ into transverse $`F_{0i}^{a}{}_{}{}^{T}=E_{i}^{a}{}_{}{}^{T}`$ and longitudinal $`F_{0i}^{a}{}_{}{}^{L}=_i\sigma ^a`$ parts
$$F_{0i}^a=E_{i}^{a}{}_{}{}^{T}_i\sigma ^a;_kE_{k}^{b}{}_{}{}^{T}=0.$$
(30)
This decomposition is compatible with the perturbative gauge
$$_iA_i^a:=_iA_{i}^{a}{}_{}{}^{T}=0.$$
(31)
### 3.2. Constraining with the zero mode
To formulate an equivalent unconstrained system for the YM theory in the winding number class of functions in the presence the zero mode $`𝒵^b`$ of the Gauss law constraint (29)
$$A_0^a=𝒵^a+\stackrel{~}{A}_0^a;F_{0k}^a=D_k^{ab}(A)𝒵^b+\stackrel{~}{F}_{0k}^a((D^2(A))^{ab}𝒵^b=0),$$
(32)
we shall follow two principles of the Faddeev foundation : i) the constraint-shell action (27)
$$W_{YM}(\mathrm{constraint})=𝒲_{YM}[𝒵]+\stackrel{~}{W}_{YM}[\stackrel{~}{F}],$$
(33)
and ii) the choice of the gauge
$$D_i^{ab}(\mathrm{\Phi }_i)\stackrel{~}{A}_i^b:=0$$
(34)
compatible with the perturbation theory around the SU(2) monopole $`\mathrm{\Phi }_i^a`$
$$\widehat{\stackrel{~}{F}}_{0k}=U_𝒵\widehat{F}_{0k}^IU_𝒵^1,\widehat{A}_i=U_𝒵(\widehat{A}_i^I+_i)U_𝒵^1,\widehat{A}_i^I(t,\stackrel{}{x})=\widehat{\mathrm{\Phi }}_i(\stackrel{}{x})+\widehat{\stackrel{~}{A}}_i(t,\stackrel{}{x}).$$
(35)
To present the action in the form (33), we can use the evident decomposition
$$F^2=(D𝒵+\stackrel{~}{F})^2=(D𝒵)^22\stackrel{~}{F}D𝒵+(\stackrel{~}{F})^2=(𝒵(D𝒵))2(𝒵\stackrel{~}{F})+(\stackrel{~}{F})^2$$
(36)
and the Gauss Eqs. $`D\stackrel{~}{F}=0`$ and $`D^2𝒵=0`$ which show that the zero mode part $`𝒲_{YM}`$ of the constraint-shell action (33) is the sum of two surface integrals
$$𝒲_{YM}[𝒵]=𝑑td^3x[\frac{1}{2}_i(𝒵^aD_i^{ab}(A)𝒵^b)_i(\stackrel{~}{F}_{0i}^a𝒵^a)]=𝒲^0+𝒲^{},$$
(37)
where the first one $`𝒲^0`$ is the kinetic term and the second one $`𝒲^{}`$ describes the coupling of the zero-mode to the local excitations. These surface terms are determined by the asymptotics of the fields $`(𝒵^a,A_i^a)`$ at spatial infinity (15), (18) which we denoted, in the SU(3) case, by $`(\dot{N}(t)\mathrm{\Phi }_0^a(\stackrel{}{x}),\mathrm{\Phi }_i^a(\stackrel{}{x}))`$. The fluctuations $`\stackrel{~}{F}_{0i}^a`$ belongs to the class of multipoles and since the surface integral over monopole-multipole couplings vanishes, the fluctuation part of the second term obviously drops out. The substitution of the solution with the asymptotics (15) into the first surface term Eq. (37) leads to the zero-mode action
$$𝒲^0=\frac{}{2}𝑑t\dot{N}(t)^2,$$
(38)
where
$$=d^3x_i[\mathrm{\Phi }_0^a(\stackrel{}{x})_i\mathrm{\Phi }_0^a(\stackrel{}{x})]=(\frac{2\pi }{g})^24\pi ϵ,$$
(39)
and
$$ϵ=\underset{r_v\mathrm{}}{lim}r_v^2f_0^{}(r_v).$$
(40)
$`r_v`$ is a value of $`r`$ at a boundary of a finite spatial volume.
The action for the equivalent unconstrained system of the local excitations,
$$\stackrel{~}{W}_{YM}[\stackrel{~}{F}]=d^4x\left\{E_k^a\dot{A}_k^{aI}\frac{1}{2}\left\{E_k^2+B_k^2(A^I)+[D_k^{ab}(\mathrm{\Phi })\stackrel{~}{\sigma }^b]^2\right\}\right\},$$
(41)
is obtained by decomposing the electrical components of the field strength tensor $`F_{0i}^I`$ into transverse $`F_{0i}^{a}{}_{}{}^{T}=E_i^a`$ and longitudinal $`F_{0i}^{a}{}_{}{}^{L}=D_i^{ab}(\mathrm{\Phi })\stackrel{~}{\sigma }^b`$ parts, so that
$$F_{0i}^{aI}=E_iD_i^{ab}(\mathrm{\Phi })\stackrel{~}{\sigma }^b;D_k^{ab}(\mathrm{\Phi })E_k^b=0.$$
(42)
Here the function $`\stackrel{~}{\sigma }^b`$ is determined from the Gauss equation (29)
$$\left((D^2(\mathrm{\Phi }))^{ab}+gϵ^{adc}\stackrel{~}{A}_i^dD_i^{cb}(\mathrm{\Phi })\right)\stackrel{~}{\sigma }^b=gϵ^{abc}\stackrel{~}{A}_i^aE_i^c.$$
(43)
### 3.3. Feynman path integral
The Feynman path integral over the independent variables includes the integration over the topological variable $`N(t)`$
$$Z_{YM}^{\mathrm{Feynman}}[J]=\underset{t}{}dN(t)\stackrel{~}{Z}_{YM}[J^U],$$
(44)
where
$$\stackrel{~}{Z}[J^U]=\underset{t,x}{}\left\{\underset{a=1}{\overset{3}{}}\frac{[d^2A_a^Id^2E_a]}{2\pi }\right\}\mathrm{exp}\left\{i𝒲_{YM}+i\stackrel{~}{W}_{YM}+iS[J^U]\right\}.$$
(45)
As we have seen above, the functionals $`\stackrel{~}{W},S`$ are given in terms of the variables which contain the nonperturbative phase factors $`U=U_𝒵`$ (21) of the topological degeneration of initial data. These factors disappear in the action $`\stackrel{~}{W}`$, but not in the source
$$S[J^U]=d^4xJ_i^aA_i^a,\widehat{A}_i=U(\widehat{A}_a^I+_i)U^1$$
(46)
what reflects the fact of the topological degeneration of the physical fields.
Thus, instead of the instanton averaging over ’interpolations between different vacua’, in order to remove the topological degeneracy, all Green functions should be averaged over values of the topological variable and all possible angles of orientation of the monopole unit vector ($`\stackrel{}{n}=\stackrel{}{x}/r`$) (18) in the group space. This averaging leads to complete destructive interference of colored Green functions and corresponding color amplitudes. The complete destructive interference of color amplitudes can be interpreted as confinement in the spirit of the Feynman quark-parton duality .
## 4. Quantum chromodynamics
### 4.1. Feynmann path integral
QCD differs from the YM theory by the group SU(3) and quark fields.
The decomposition of the spatial components of the gluon fields $`A_i^a=\mathrm{\Phi }_i^a+\stackrel{~}{A}_i^a`$ into the static monopole $`\mathrm{\Phi }_i^a`$ and the fluctuations can be used to define a perturbation theory with respect to the dynamical gluon fields $`\stackrel{~}{A}_i^a`$. The adequate gauge for this perturbation theory is the covariant Coulomb gauge defined by
$$D_i^{ac}(\mathrm{\Phi }_i^b)\stackrel{~}{A}_i^c=0.$$
(47)
After the separation of the phase factors of the topological degeneracy, we obtain the similar path integral in terms of perturbative fields ($`E,A^I,\overline{\psi }^I,\psi ^I`$)
$$Z_{QCD}^{\mathrm{Feynman}}[J,\eta ,\overline{\eta }]=\underset{t}{}dN(t)e^{i𝒲_0}\stackrel{~}{Z}[J^U,\eta ^U,\overline{\eta }^U],$$
(48)
where
$$\stackrel{~}{Z}[J^U,\eta ^U,\overline{\eta }^UJ]=\underset{t,x}{}\left\{[d\overline{\psi }^Id\psi ^I]\underset{a=1}{\overset{8}{}}\frac{[d^2A_a^Id^2E_a]}{2\pi }\right\}\mathrm{exp}\left\{i𝒲^{}+i\stackrel{~}{W}+iS[J^U,\eta ^U,\overline{\eta }^U]\right\},$$
(49)
$$\eta ^U=U_𝒵\eta ,$$
(50)
the zero-mode actions $`𝒲_0`$ and $`𝒲^{}`$ are defined by the expressions similar to (37) and (38), in the first order formalism obtained by the Legandre transformation
$$\frac{1}{2}G_{0i}^2=G_{0i}F_{0i}\frac{1}{2}F_{0i}^2.$$
The quark part of $`𝒲^{}`$ will be discussed later. The action for the equivalent unconstrained system of the local excitations
$$\stackrel{~}{W}=d^4x\left\{E_k^a\dot{A}_k^{aI}\frac{1}{2}\left\{E_k^2+B_k^2(A^I)+[D_k^{ab}(\mathrm{\Phi })\stackrel{~}{\sigma }^b]^2\right\}+\overline{\psi }^I[i\gamma _\mu ^\mu m]\psi ^I+j_i^aA_i^{aI}\right\},$$
(51)
is obtained by decomposing the electrical components of the field strength tensor $`F_{0i}^I`$ (similarly (35)) into transverse $`F_{0i}^{a}{}_{}{}^{T}=E_i^a`$ and longitudinal $`F_{0i}^{a}{}_{}{}^{L}=D_i^{ab}(\mathrm{\Phi })\stackrel{~}{\sigma }^b`$ parts, so that
$$F_{0i}^{aI}=E_iD_i^{ab}(\mathrm{\Phi })\stackrel{~}{\sigma }^b;D_k^{ab}(\mathrm{\Phi })E_k^b=0.$$
(52)
The function $`\stackrel{~}{\sigma }^b`$ is determined from the Gauss equation (11)
$$\left((D^2(\mathrm{\Phi }))^{ab}+gf^{adc}\stackrel{~}{A}_i^dD_i^{cb}(\mathrm{\Phi })\right)\stackrel{~}{\sigma }^b=j_{\mathrm{tot},0}^b;j_{\mathrm{tot},0}^b=gf^{abc}\stackrel{~}{A}_i^aE_i^c+j_0^b.$$
(53)
Equation (53) can be formally solved by introducing a Green function $`G^{ab}(\stackrel{}{x},\stackrel{}{y})`$ defined as the solution to
$$(D^2(\mathrm{\Phi }))^{ab}G^{bc}(\stackrel{}{x},\stackrel{}{y})=\delta ^{ac}\delta ^3(\stackrel{}{x}\stackrel{}{y}).$$
(54)
Then $`\stackrel{~}{\sigma }^a`$ can be expanded in a power series in $`g`$
$$\stackrel{~}{\sigma }^b(t,\stackrel{}{x})=d^3yG^{bc}(\stackrel{}{x},\stackrel{}{y})j_{\mathrm{tot},0}^c(t,\stackrel{}{y})d^3yd^3zG^{bc}(\stackrel{}{x},\stackrel{}{y})gf^{cde}\stackrel{~}{A}_k^d(t,\stackrel{}{y})G^{ef}(\stackrel{}{y},\stackrel{}{z})j_{\mathrm{tot},0}^f(t,\stackrel{}{z})\mathrm{}$$
(55)
in analogy to the standard perturbation theory . The calculation of the Green function $`G^{ab}(\stackrel{}{x},\stackrel{}{y})`$ is given in the Appendix.
### 4.2. Dominance of the Wu-Yang monopole
The action $`𝒲^0`$ in the path integral (49) (given by (38)) describes a free ’rotator’ with the momentum spectrum
$$P_0=\dot{N}M_N=(2\pi k+\theta ),k=0,1,2,\mathrm{},$$
(56)
which follows from the constraint on the wave function for physically equivalent points $`N`$ and $`N+1`$
$$\mathrm{\Psi }(N+1)=\mathrm{exp}(i\theta )\mathrm{\Psi }(N).$$
The action for the ’rotator’ (38) compensates the action for the monopole background
$$W[\mathrm{\Phi }_i]=\frac{1}{2}𝑑td^3x[B_i^a(\mathrm{\Phi }_i)]^2$$
in the Minkowski space for finite values of the momentum of the rotator, so that the perturbation theory begins from the zero value of the Minkowski action, in the contrast to the instanton background. This can be interpreted as the dominance of the monopole background.
For the self-dual BPS ansatz (22) the dominant value of the momentum of the rotator is
$$\left|\frac{P_0}{2\pi }\right|=\frac{4\pi }{g^2}.$$
(57)
### 4.3. Hadronization
In the lowest order of this perturbation theory, we can rewrite the instantaneous interaction term in (51) in the form of a current-current interaction
$$𝒲_{\mathrm{int}}=\frac{1}{2}d^4x[D_k^{ab}(\mathrm{\Phi })\stackrel{~}{\sigma }^b]^2=\frac{1}{2}d^4xd^3yj_{\mathrm{tot},0}^a(t,\stackrel{}{x})G^{ab}(\stackrel{}{x},\stackrel{}{y})j_{\mathrm{tot},0}^b(t,\stackrel{}{y}).$$
(58)
Following the QED perturbation theory of ‘radiative corrections’ one can formulate a similar perturbation theory of ’dynamical corrections’ in QCD. This ’dynamical perturbation theory’ is based on the decomposition of the action $`\stackrel{~}{W}`$ for the dynamical variables (47) into the ‘free’ part $`\stackrel{~}{W}_0`$, the instantaneous interaction $`𝒲_{\mathrm{int}}`$ (58) and the residual dynamical interaction $`\stackrel{~}{W}_{\mathrm{int}}`$
$$\stackrel{~}{W}=\stackrel{~}{W}_0+𝒲_{\mathrm{int}}+\stackrel{~}{W}_{\mathrm{int}}.$$
(59)
The lowest order of the ’dynamical perturbation theory’ $`\stackrel{~}{W}_{\mathrm{int}}=0`$ describes gluons and quarks in the monopole field and their instantaneous bound states. Next orders of this perturbation theory contain ’dynamical corrections’ which describe matrix elements of transitions between these states.
Bound states are obtained in a similar way as in QED (i.e., by the Schwinger-Dyson and Bethe-Salpeter equations with the instantaneous interaction $`G^{bc}(\stackrel{}{x},\stackrel{}{y}))`$ given in the Appendix. In the field of the monopole the instantaneous quark-quark potential is the sum of a Coulomb potential and a rising one. It is well-known that the latter one leads to spontaneous chiral symmetry breaking and to mesonic bound states.
### 4.4. Frozen gluon approximation and U$`{}_{A}{}^{}(1)`$\- problem
The ‘frozen gluon’ approximation (FGA) follows from the Feynman path integral if we neglect the dynamical gluon fields: $`\stackrel{~}{A}_i^a=0`$ (47).
The FGA generating functional (45) takes the form
$$\stackrel{~}{Z}_{FGA}[\eta ,\overline{\eta }]=\underset{t,x}{}\left\{[d\overline{\psi }d\psi ]\right\}\mathrm{exp}\left\{i𝒲+i\stackrel{~}{W}_{FGA}+iS[\eta ,\overline{\eta }]\right\},$$
(60)
where
$$\stackrel{~}{W}_{FGA}=\stackrel{~}{W}_0+𝒲_{\mathrm{int}}$$
(61)
depends only on the quark fields.
The bilocal linearization of the four quark interaction $`𝒲_{\mathrm{int}}`$ using the Hubbard-Stratonovich transformation leads to an effective bilocal meson action . This meson action includes Abelian anomalies in the pseudoscalar isosinglet ($`\eta _0`$-meson) channel .
In our case, these anomalies include the time derivative of topological variable $`N(t)`$ due to interaction of the quark fields with the monopole and the zero-mode $`𝒲^{}`$ (see eq. (37)).
Neglecting all mesonic channels except the $`\eta _0`$-meson one, we get in FGA nothing but the well-known gauge-invariant expression
$$\stackrel{~}{W}_{\mathrm{anomaly}}^{\eta _0}[\eta ,\overline{\eta }]=C_\eta 𝑑t\overline{\eta }(t,0)I_c\gamma _5\eta (t,0)\frac{g^2}{16\pi ^2}d^3xG_{\mu \nu }^a{}_{}{}^{}G_{\mu \nu }^{a},$$
(62)
where $`\eta ,\overline{\eta }`$ are fermion sources, $`C_\eta `$ is a constant. In the BPS monopole selfdual field (9), (18), (22)with $`G_{0i}^a=\dot{N}D_i^{ab}(\mathrm{\Phi })\mathrm{\Phi }_0^b`$, and $`2\pi ϵD_i^{ab}(\mathrm{\Phi })\mathrm{\Phi }_0^b=B_i^a(\mathrm{\Phi })`$ we obtain the normalizable zero-mode
$$\frac{g^2}{16\pi ^2}d^3xG_{\mu \nu }^a{}_{}{}^{}G_{\mu \nu }^{a}=\dot{N}(t),$$
(63)
as
$$\frac{g^2}{8\pi ^2}d^3D_i^{ab}(\mathrm{\Phi })\mathrm{\Phi }_0^bB_i^a(\mathrm{\Phi })=1.$$
The physical anomaly term $`\stackrel{~}{W}_{\mathrm{anomaly}}`$ in the $`\eta _0`$-meson channel takes the form
$$\stackrel{~}{W}_{\mathrm{anomaly}}^{\eta _0}[\eta ,\overline{\eta }]=C_\eta 𝑑t\eta _0\dot{N},\eta _0(t)=\overline{\eta }(t)I_c\gamma _5\eta (t).$$
(64)
Following Veneziano , we can identify $`\eta _0`$ with the field of the $`\eta _0`$-meson at rest. The effective action including the anomaly term (62) and the zero-mode one (38) is then
$$W_{\mathrm{eff}}=𝑑t[\frac{\dot{N}^2}{2}+\eta _0C_\eta \dot{N}+\frac{\eta _0^2m_0^2V}{2}],$$
(65)
where $`m_0`$ is the standard current quark mass contribution to the $`\eta _0`$ meson mass. The diagonalization of this effective action leads to an additional mass of $`\eta _0`$-meson
$$\mathrm{\Delta }m^2V=\frac{C_\eta ^2}{}=\frac{C_\eta ^2}{\mathrm{\Delta }m^2V}=\frac{(2\pi )^2ϵ}{g^2}4\pi .$$
(66)
The finite contribution of the zero-mode to the mass of the $`\eta _0`$-meson entails the disappearance of the mass of the collective topological variable and leads to a stable perturbation theory in the infinite volume limit, as in this limit the singularity-free BPS monopole converts into the Wu-Yang one without singularities at the origin.
### 4.5. Relativistic covariance
Recall that Dirac has also obtained the unconstrained form of QED in terms of gauge invariant variables for QED as functionals of the initial gauge fields by explicitly resolving the Gauss law constraint. The resulting unconstrained formulation of QED coincides with the one obtained in the Coulomb gauge with the physical phenomena of electrostatics, ’dressed’ electrons, and two transverse photon degrees of freedom. In QED, in terms of the Dirac variables , the Poincare symmetry is realized which is mixing with the gauge symmetry . This mixing was interpreted in 1930 by Heisenberg and Pauli (with reference to an unpublished note by von Neumann) as the transition from the Coulomb gauge with respect to the time axis in the rest frame ($`l_\mu ^0=(1,0,0,0)`$) to the Coulomb gauge with respect to the time axis in the moving frame $`l_\mu =l_\mu ^0+\delta _Ll_{\mu }^{}{}_{}{}^{0}=(Ll^0)_\mu `$. The Coulomb interaction has the covariant form
$`W_C={\displaystyle 𝑑x𝑑y\frac{1}{2}j_l^T(x)V_C(z^{})j_l^T(y)\delta (lz)},`$ (67)
where $`j_l^T=e\overline{\psi }^T\text{/}l\psi ^T,z_\mu ^{}=z_\mu l_\mu (zl),z_\mu =(xy)_\mu .`$ This transformation law and the relativistic covariance of this formulation of QED has been proven by Zumino on the level of the algebra of generators of the Poincare group.
In QED we can change variables to construct the generating functional of the Green functions in any gauges including the Lorentz invariant ones. The invariance of the corresponding Green functions under a change of variables (which generates the Ward-Taylor-Slavnov identities ) is garanteed by the Dirac factors in source terms, which restore the Coulomb gauge Feynman rules in any Lorentz invariant gauge. So, the Coulomb interaction and electrostatics are consequences of the identification of the physical degrees of freedom which correspond to an explicit solution of the Gauss law, but not primarily to the choice of the gauge. For example, if one would omit the Dirac factors in the source terms in relativistically invariant Lorentz gauge formulations of QED, one would get the Wick-Cutkosky bound states formed by gauge propagators with light-cone singularities with a spectrum different <sup>1</sup><sup>1</sup>1One of the authors (V.P.) thanks W. Kummer who pointed out that in Ref. the difference between the Coulomb atom and the Wick-Cutkosky bound states in QED has been demonstrated. from the observed one which corresponds to the instantaneous Coulomb interaction.
In QCD, the moving Lorentz frame corresponds to the moving Wu-Yang monopole . We should only manage to realize this transformation law on the level of operators.
## 5. Conclusion
In this paper, we have shown that there is a collective topological excitation in the gluon spectrum as the zero-mode of Gauss’ constraint from the winding number class of functions of the large gauge transformations.
This topological excitation leads to the dominance of the Wu-Yang monopole in the Feynman path integral compatible with the equivalent unconstrained system obtained in this work.
In the field of the Wu-Yang monopole the instantaneous quark-quark potential is the sum of a Coulomb-type potential and a rising one. The latter one leads to spontaneous chiral symmetry breaking and to mesonic bound states.
The $`\eta _0`$-meson mixes with the zero-mode so that after diagonalization of this low-energy action a mass shift of the $`\eta _0`$-meson is obtained which resolves the $`U_A(1)`$ problem.
Colored amplitudes contain additional phase factors which depend on the zero-mode. Averaging the colored amplitudes over the zero-mode parameters can lead to the phenomenon of complete destructive interference , so that the color amplitudes disappear. The colorless ones of the type of the expectation values of the electroweak currents do not vanish since the zero-mode phase factors are absent.
According to Heisenberg, Pauli and Zumino , the relativistic covariance is established by a rotation of the timelike axis so that the non-Abelian Coulomb-gauge field moves together with the relativistic bound states. Recently, Faddeev and Niemi constructed a similar relativistically covariant form of the Wu-Yang monopole.
In summary, the present scheme for the introduction of topological global variables is a promising tool for the investigation of the challenging properties of QCD such as the meson spectrum, chiral symmetry breaking, quark and gluon confinement, and the $`U_A(1)`$ anomaly.
## Acknowledgements
We thank Profs. A. V. Efremov, G. A. Gogilidze, N. Ilieva, V. G. Kadyshevsky, A. M. Khvedelidze, E. A. Kuraev, M. Lavelle, D. McMullan and W. Thirring for interesting and critical discussions. One of us (VNP) is grateful for a stipend from the Max-Planck-Gesellschaft for his study visit at the University of Rostock.
## Appendix A The Green function
We can calculate the instantaneous the Green function (68)
$$(D^2(b))^{ab}(\stackrel{}{x})G^{bc}(\stackrel{}{x},\stackrel{}{y})=\delta ^{ac}\delta ^3(xy).$$
(68)
In the presence of the Wu-Yang monopole we have
$$(D^2)^{ab}(\stackrel{}{x})=\delta ^{ab}\mathrm{\Delta }\frac{n^an^b+\delta ^{ab}}{r^2}+2(\frac{n_a}{r}_b\frac{n_b}{r}_a),$$
and $`n_a(x)=x_a/r;r=|\stackrel{}{x}|`$. Let us decompose $`G^{ab}`$ into a complete set of orthogonal vectors in color space
$$G^{ab}(\stackrel{}{x},\stackrel{}{y})=[n^a(x)n^b(y)V_0(z)+\underset{\alpha =1,2}{}e_\alpha ^a(x)e_\alpha ^b(y)V_1(z)];(z=|\stackrel{}{x}\stackrel{}{y}|).$$
Substituting the latter into the first equation, we get
$$\frac{d^2}{dz^2}V_n+\frac{2}{z}\frac{d}{dz}V_n\frac{n}{z^2}V_n=0n=0,1.$$
The general solution for the last equation is
$$V_n(z)=d_nz^{l_1^n}+c_nz^{l_2^n},$$
(69)
where $`d_n`$, $`c_n`$ are constants, and $`l_1^n,l_2^n`$ can be found as roots of the equation $`(l^n)^2+l^n=n`$ , i.e.
$$l_1^n=\frac{1+\sqrt{1+4n}}{2};l_2^n=\frac{1+\sqrt{1+4n}}{2}.$$
(70)
It is easy to see that for $`n=0`$ we get the Coulomb-type potential $`d_0=1/4\pi `$, and for $`n=1`$ the ’golden section’ potential with
$$l_1^1=\frac{1+\sqrt{5}}{2}1.618;l_2^1=\frac{1+\sqrt{5}}{2}0.618.$$
(71)
The last potential (in the contrast with the Coulomb-type one) leads to spontateous chiral symmetry breaking and can be used as the potential for the ’hadronization’ of quarks and gluons in QCD .
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# Thresholds for Linear Optics Quantum Computation
## I Introduction
In knill:qc2000b we proposed that linear optics quantum computation (LOQC) is a viable option for physically realizing quantum computers. The proposal depends on a series of optical protocols that require single photon state preparation and measurement whose outcome can be used to control other optical elements. The basic idea is to kick back the hidden non-linearities in photo-detectors to qubits encoded in pairs of optical modes (photonic qubits). The protocols were shown to implement the necessary quantum gates with arbitrarily high probability of success. In their simplest form, the resources required for implementing the protocols grow rapidly as the desired success probability is increased. We noted that due to the general accuracy threshold theorem aharonov:qc1996a ; kitaev:qc1997a ; knill:qc1998a ; preskill:qc1998a , the model scales efficiently asymptotically, with constant overhead for implementing the standard fault tolerant model of quantum computation. We also suggested that with the use of erasure codes grassl:1996a , this overhead could be significantly reduced.
Here we apply the techniques of of quantum error-detection and correction to greatly increase the efficieny of implementing quantum information processing by LOQC. We first show that for ideal LOQC (where errors in devices are ignored), it is possible to use a two qubit error-detecting code to rapidly eliminate the probability of failed implementations of the two qubit gates. This follows from a threshold analysis demonstrating a threshold of $`.5`$ for an error model where the errors are $`\sigma _z`$ measurements at known locations. One consequence of this result is that the non-deterministic gates of knill:qc2000b need only be implemented using the $`2`$\- or $`3`$-photon prepared states. The next task is to demonstrate that the main sources of device errors can be efficiently eliminated. We accomplish this by adding phase error-correction, and more importantly erasure coding methods. The latter serves the purpose of removing errors caused by particle loss and detector inefficiency, exploiting the built-in loss detection capabilities of basic LOQC. In the process we give a conservative bound on the threshold for the erasure error model.
On the basis of this paper and knill:qc2000b we can propose a roadmap for experimental and theoretical work toward implementing LOQC. It is based on viewing basic LOQC as a fundamental model of computation that can be used to implement standard quantum computation by using a set of layered techniques. For benchmarking purposes, the relevant resources of LOQC can be taken to be the number of particles independently generated, the probability of (detected) failure of the implemented computation (without post-selection), and the measured error in the output conditional on success (i.e. with post-selection). The roadmap for LOQC based quantum computation can be outlined as follows:
* Basic LOQC.
+ Non-deterministic non-linear sign changes.
+ Preparation of the state $`|t_n`$ for teleportation (see knill:qc2000b for the definition of $`|t_n`$).
+ Teleportation using $`|t_n`$ for increasing $`n`$ with probability of success close to $`1/(n+1)`$.
+ Controlled sign changes with good probability of success.
* QC based on LOQC.
+ Boosting the success probability by the use of encoding.
+ Decrease phase errors by applying error correction.
+ Decrease loss errors by encoding with erasure codes.
+ Concatenation (or larger codes) for achieving high accuracy.
The experimental challenges include the ability to use multiple independently generated single photons and to control the photonic qubits using feedback from detectors. The first demonstrations are likely not to involve feedback, but rather to use post-selection and high repetition rates to demonstrate success. Although feedback can be delayed in our schemes, this comes at the cost of high failure rates in the state preparation protocols, particularly for those states that depend on previously prepared states. The ultimate goal is to build the state preparation protocols into state factories with the ability to attempt failure-prone state preparations at high rates (perhaps in parallel), exploiting the ease with which photons can be generated.
Here is the outline of this paper. We begin by recalling the needed properties of the LOQC protocols given in knill:qc2000b . We briefly describe the concatenation method for establishing thresholds, give a code suitable for increasing the success probability in ideal LOQC, and show how to implement encoded operations. The gain in success probability is estimated so as to determine a threshold. The resources that determine the general error propagation behavior are bounded for the goal of approaching the quantum communication threshold. As a result we conservatively estimate that a gate with the necessary reliability depends on at most a few hundred LOQC controlled sign flips implemented non-determinstically with a pair of states each involving three photons. More are used up in unsuccessful state preparations, but at least in principle, this overhead is not much larger. The quantum coding methods are then enhanced for the purpose of dealing with phase and detected-loss errors, including particle detector inefficiencies, and finally for dealing with any residual general errors.
We assume familiarity with quantum computation aharonov:qc1998a ; divincenzo:qc2000a and quantum error-correction via stabilizer codes gottesman:qc1996a ; gottesman:qc1997a . For the basic ideas of LOQC see knill:qc2000b . Most of this paper is written in the language of qubits using products of Pauli operators.
## II Features of LOQC
LOQC is a model of quantum computation where some of the gates are non-deterministic, with detectable failures. The probability of failure of the gates depends on the resources used. Specifically, the model is characterized by enabling the following operations on standard qubits (which are encoded in the physical system as photonic qubits):
* Preparation of $`|0`$: $`P_0`$.
* Measurement in the basis $`|0,|1`$: $`M_0`$.
* Every one qubit rotation.
* Controlled sign flip: c-$`\sigma _z`$, with a probability of failure, which is detected.
In practice errors other than detected failures occur. For the moment, we assume that such errors are significantly less likely than detected failures, so that initially, they can be ignored. Specifically, this leads to designing implementations by dealing with errors in order of their importance. We call LOQC with only the errors due to detected failure of the conditional sign flip “ideal LOQC” (iLOQC). In iLOQC, single photon state preparation, particle number detection and passive linear optical elements are all perfect. As explained in our previous paper knill:qc2000b , when applying $`\text{c-}\sigma _z^{(12)}`$, with probability $`f`$, qubit $`1`$ is measured in the $`\sigma _z`$ basis. If this event does not occur then with probability $`f`$, qubit $`2`$ is measured in the $`\sigma _z`$ basis. Thus the prior probability of the second event is $`f(1f)`$. The measurement outcome is known in either case and the qubit not measured is preserved. We take that to be the error model of iLOQC. Note that the failure behavior is asymmetric, so that the ordering of the labels is significant. We call the first qubit in this operation the “source”, and the second the “target”. The resource overhead for implementing this gate depends on $`f`$. If the methods of knill:qc2000b are used for implementing the c-$`\sigma _z`$ operation, a state with approximately $`(1/f)1`$ photons needs to be adjoined. The preparation of the state needs to be tried several times, using ancillary photons. With the naive method, the total number of photons used in the preparation attempts grows exponentially in $`1/f`$, so it is desirable to show that we can scale with as high a probability of failure as possible.
For the present purposes, it is convenient to use a gate set based on the product operator formalism sorensen:qc1983a . Thus, gates are exponentials of products of Pauli operators. To simplify the notation, define $`X:=\sigma _x,Y:=\sigma _y,Z:=\sigma _z`$. For a product operator $`U`$, write $`U_\theta =e^{iU\pi \theta /360}`$ and note that since $`U^2=I`$, $`U_\theta =\mathrm{cos}(\pi \theta /360)i\mathrm{sin}(\pi \theta /360)U`$. For example, $`X_{180}`$ is a bit flip up to a global phase. When there is no possibility for confusion, we abbreviate $`UVU^{(a)}V^{(b)}`$, where the parenthesized superscripts are system labels.
One reason for why product operators are convenient is because it is straightforward to follow their evolution under $`90^{}`$ rotations by using appropriate triples of axes. For example, a $`Y_{90}`$ rotation takes the $`Z`$ axis to the $`X`$ axis, so that $`(Y_{90})Z(Y_{90})=X`$. This implies that if in a quantum network, a $`Z`$ measurement or a $`Z`$ rotation precedes a $`Y_{90}`$ rotation, then this is equivalent to an $`X`$ measurement or rotation (by the same angle) after the $`Y_{90}`$. A complete set of gates that generates the determinant $`1`$ unitary matrices is given by
* $`X`$ rotations: $`X_\varphi `$ for any angle $`\varphi `$.
* $`Z`$ rotations: $`Z_{90}`$.
* $`ZZ`$ rotations: $`(Z^{(1)}Z^{(2)})_{90}`$.
(With the use of ancillas, $`Z_{90}`$ may be eliminated from the set without loss of completeness.) The third gate is related to c-$`\sigma _z`$ by
$$(Z^{(1)}Z^{(2)})_{90}=e^{i\pi /4}Z_{90}^{}{}_{}{}^{(1)}Z_{90}^{}{}_{}{}^{(2)}\text{c-}\sigma _z^{(12)}.$$
(1)
As a result, the $`ZZ`$ rotation is readily implemented in LOQC up to global phases. Furthermore, we can modify the implementation to achieve the following failure behavior for $`(Z^{(1)}Z^{(2)})_{90}`$: With probability $`f`$, qubit $`1`$ is measured in $`Z`$ and qubit $`2`$ is untouched. With probability $`(1f)f`$, qubit $`2`$ is measured in $`Z`$. If the outcome is $`0`$, qubit $`1`$ experienced a $`Z_{90}`$ and if it is $`1`$, a $`Z_{90}`$. Note that due to the availability of $`X_{90}`$ rotations, we can also use $`(U^{(1)}V^{(2)})_{\pm 90}`$ rotations for $`U`$ and $`V`$ either $`Z`$ or $`Y`$ with similar error behavior, where the measurements commute with the rotation. Similarly, it is straightforward to implement the $`Y_{\pm 90}`$ rotations, $`Y`$ measurement and $`Y`$ eigenstate preparation. We will find that the encoded $`Z`$ rotation also has a probability of failure, where a $`Z`$ measurement occurs if it fails. In this case, $`X`$ eigenstate preparation (using $`|0`$ preparation and an $`X_{90}`$ followed by a $`Z_{\pm 90}`$) may fail. However, since the preparation is successful if failure has not been detected, it is possible to retry it until success is achieved. This makes it possible to get a perfect $`X`$ eigenstate preparation.
## III Establishing accuracy thresholds
For the very special error model introduced in the previous section, accuracy threshold analyses are much simpler. The basic principle is to use a quantum code that permits implementing the basic operations on the encoded information in such a way that the new error model is consistent with the original one, and so that the new error rate is substantially less. The encoded qubits then behave just like the more fundamental ones used to encode them, so that the same coding method can be used with these new types of qubits. This recursive coding method is known as concatenation and has the property that the error-rates decrease super-exponentially with the number of levels of concatenation. The basic components of an accuracy threshold result by concatenation are the following:
* A quantum code.
* Means of implementing each of the basic operations on the encoded qubit.
* Means of recovering from error, which may be part of 2.
* Establishing the error model that applies to the encoded qubits and calculating a bound on the new error rate.
The plan is to show that a two qubit code suffices for establishing a threshold of $`f=.5`$ for iLOQC. After dealing with the errors of iLOQC, it is necessary to consider the contributions of other sources of noise, particularly photon loss and phase error. It turns out that the same family of codes can be used with phase error correction. We recall that LOQC comes with an effective leakage detection scheme and observe that the basic teleportation protocol used in LOQC can be enhanced to allow detection of loss from particle detector inefficiency. As a result, it possible to use erasure code to eliminate errors from photon loss. We point out that due to the special nature of the erasure error model, good accuracy thresholds apply and error rates can be readily improved with relatively simple codes and few levels of encoding.
## IV A two qubit quantum code
Let $`|+:=(|0+|1)/\sqrt{2}`$ and $`|:=(|0|1)/\sqrt{2}`$ be the eigenstates of $`X`$. Up to overall scale factors, the associated projection operators are $`I\pm X`$. We continue to omit these scale factors in identities. To improve the failure probabilities, we use a two qubit quantum code with the encoding
$`|0`$ $``$ $`|0_L`$ (2)
$`=`$ $`(|+++|)/\sqrt{2}`$
$`=`$ $`(|00+|11)/\sqrt{2}`$
$`|1`$ $``$ $`|1_L`$ (3)
$`=`$ $`(|++|)/\sqrt{2}`$
$`=`$ $`(|01+|10)/\sqrt{2},`$
and show how to implement the necessary operations on the encoded states. The encoded states define the state space of the “logical qubit”. In the language of stabilizer codes, the code is defined by the stabilizer $`X^{(1)}X^{(2)}`$, and accordingly, the projection onto the code space is given by $`I+X^{(1)}X^{(2)}`$. We use $`L`$ as the label for logical (encoded) operators and states. With this encoding, logical operators are given by
$`X^{(L)}`$ $`=`$ $`X^{(1)}=_LX^{(2)}`$ (4)
$`Z^{(L)}`$ $`=`$ $`Z^{(1)}Z^{(2)}=_LY^{(1)}Y^{(1)}`$ (5)
$`Y^{(L)}`$ $`=`$ $`Y^{(1)}Z^{(2)}=_LZ^{(1)}Y^{(2)},`$ (6)
where we introduced the notation $`=_L`$ to denote identity when restricted to the code.
For the purposes of establishing a threshold, we use the following set of basic operations and assumptions:
* $`X`$, $`Y`$ and $`Z`$ eigenstate (eigenvalue $`1`$) preparation: $`P_X,P_Y,P_Z`$.
* $`Y`$ and $`Z`$ measurements: $`M_Y`$ and $`M_Z`$. $`M_Z(s)`$ means that the result of the $`Z`$ measurement was eigenvalue $`s`$.
* $`X_{180}`$, $`Y_{180}`$ and $`Z_{180}`$ rotations.
* $`X_\varphi `$ rotations.
* $`Z_{90}`$ rotation, with failure probability $`f`$ after the first encoding.
* $`(Z^{(1)}Z^{(2)})_{90}`$ rotation, with failure probability $`f`$ in our error model.
Operations 1. to 4. are error-free in iLOQC. The $`Z_{90}`$ rotation is error free in iLOQC, but as we will see, it may fail with probability $`f`$ with a $`Z`$ measurement after the first level of encoding. Note that $`Y_{90}`$, $`(YZ)_{90}`$ and $`(YY)_{90}`$ rotations can be implemented by conjugation with error-free $`X`$ rotations.
### IV.1 State preparation
To encode an arbitrary state $`|\psi __1`$ to $`|\psi __L`$, one can use the sequence of operations given by
$$E^{(12)}=Y_{}^{(1)}{}_{90}{}^{}(Z^{(1)}Y^{(2)})_{90}Y_{}^{(1)}{}_{90}{}^{}P_{Z}^{}{}_{}{}^{(2)}.$$
(7)
To see that this works, follow the effect of the unitary gates on the initial operators $`X^{(1)}`$, $`Z^{(1)}`$ (basic operators associated with the state to be encoded) and $`I+Z^{(2)}`$ (the projection onto the prepared state). It can be seen that $`X^{(1)}X^{(1)}=_LX^{(L)}`$, $`Z^{(1)}Y^{(1)}Y^{(2)}=_LZ^{(L)}`$ and $`I+Z^{(2)}I+X^{(1)}X^{(2)}`$ (the projection onto the code).
The sequence $`E^{(12)}`$ can be used to prepare encoded eigenstates of $`X`$, $`Y`$ or $`Z`$ by preparing the eigenstate on qubit $`1`$ first. The process fails with probability $`f+(1f)f=f(2f)`$, and can be repeated an expected $`1/(1f(2f))`$ times to successfully prepare it. Later, it will be the case that the $`Y`$ rotation in the sequence can fail, which changes the failure probability to $`f(1+(1f)+(1f)^2+(1f)^3+(1f)^4)`$. The resource usage can be improved by reusing qubits not affected by the measurement in the failure.
### IV.2 Measurement
To measure the logical qubit in the logical basis, measure both of the supporting qubits. A $`Z`$ measurement on both yields a logical $`Z`$ measurement via the total parity of the two outcomes. Similarly, a $`Y`$ measurement on the first qubit and a $`Z`$ measurement on the second gives a logical $`Y`$ measurement. A logical $`X`$ measurement is accomplished by measuring $`X`$ on the first qubit. All of these measurements are without error.
### IV.3 Logical qubit rotations
The logical $`180^{}`$ rotations are implemented by applying basic ones to each qubit. Thus
$`X_{}^{(L)}{}_{180}{}^{}`$ $`=`$ $`X_{}^{(1)}{}_{180}{}^{}`$ (8)
$`Z_{}^{(L)}{}_{180}{}^{}`$ $`=`$ $`Z_{}^{(1)}{}_{180}{}^{}Z_{}^{(2)}{}_{180}{}^{}`$ (9)
$`Y_{}^{(L)}{}_{180}{}^{}`$ $`=`$ $`Y_{}^{(1)}{}_{180}{}^{}Z_{}^{(2)}{}_{180}{}^{}.`$ (10)
These are error-free. The ability of implementing $`180`$’s in this way is generic for stabilizer codes gottesman:qc1997a . The rotations $`X_{}^{(L)}{}_{\varphi }{}^{}`$ are obtained by applying $`X_{}^{(1)}{}_{\varphi }{}^{}`$ and are error-free. To implement $`Z_{}^{(L)}{}_{90}{}^{}`$, apply $`(Z^{(1)}Z^{(2)})_{90}`$. This is not error-free, and we show later how to use recovery from $`Z`$-measurement to get a smaller probability of failure.
### IV.4 $`ZZ`$ rotation
If the logical qubits are encoded in qubits $`1`$,$`2`$ and $`3`$,$`4`$, respectively, the logical $`ZZ`$ $`90^{}`$ rotation can be obtained by applying $`(Z^{(1)}Z^{(2)}Z^{(3)}Z^{(4)})_{90}`$. This can be done by the sequence
$$(Z^{(L_1)}Z^{(L_2)})_{90}=\begin{array}{c}(Y^{(1)}Z^{(2)})_{90}(Y^{(1)}Z^{(4)})_{90}\hfill \\ (Z^{(1)}Z^{(3)})_{90}\hfill \\ (Y^{(1)}Z^{(4)})_{90}(Y^{(1)}Z^{(2)})_{90},\hfill \end{array}$$
(11)
where we use the convention that the order of application is right to left within a line and top to bottom for multiple lines. Unfortunately, this does not readily yield a logical gate with significantly less error. To do that requires using the teleportation techniques of gottesman:qc1997a ; gottesman:qc1999a .
## V Robust teleportation
### V.1 Basic teleportation
The basic teleportation protocol transfers an arbitrary state from qubit $`1`$ to qubit $`3`$ by first preparing a state on qubits $`2`$ and $`3`$, then making a measurement on qubits $`1`$ and $`2`$, and finally correcting qubit $`3`$ by applying one of the $`180^{}`$ rotations. Here is a sequence $`T^{(123)}`$ for a variant of the usual protocol that has better error behavior for our purposes.
$$T^{(123)}=\begin{array}{c}P_{Z}^{}{}_{}{}^{(2)}P_{Y}^{}{}_{}{}^{(3)}\hfill \\ (Y^{(2)}Z^{(3)})_{90}\hfill \\ (Z^{(1)}Y^{(2)})_{90}\hfill \\ M_{Z}^{}{}_{}{}^{(2)}(s_1)M_{Y}^{}{}_{}{}^{(1)}(s_2)\hfill \\ (U(s_1,s_2)^{(3)})_{180}\hfill \end{array}$$
(12)
The source in the second coupling evolution is chosen to be qubit $`2`$. The necessary correction $`U(s_1,s_2)`$ can be derived by determining the effect of the process on an input operator. Using the projection operators $`I\pm Z`$ and $`I\pm Y`$ for the prepared states and for the effects of the measurement the transformation for an initial $`Z^{(1)}`$ operator is
$`Z^{(1)}(I+Z^{(2)})(I+Y^{(3)})`$ $``$ $`(I+s_2Y^{(1)})(I+s_1Z^{(2)})`$ (13)
$`\times Z^{(1)}(IZ^{(1)}Z^{(2)}Z^{(3)})(IY^{(2)}X^{(3)})`$
$`\times (I+s_1Z^{(2)})(I+s_2Y^{(1)}).`$
Using the rules
$`(I+Z^{(a)}U^{(b)})(I+sZ^{(a)})=(I+sZ^{(a)})(I+sU^{(b)})`$ (14)
$`(I+sZ^{(a)})(I+Y^{(a)}U^{(b)})(I+sZ^{(a)})=(I+sZ^{(a)})`$ (15)
and their variations (continuing to omit constants in the identities), this evaluates to $`s_1Z^{(3)}(I+s_2Y^{(1)})(I+s_1Z^{(2)})`$. Similarly,
$`X^{(1)}(I+Z^{(2)})(I+Y^{(3)})`$ $``$ $`(I+s_2Y^{(1)})(I+s_1Z^{(2)})`$ (16)
$`\times Y^{(1)}Y^{(2)}(IZ^{(1)}Z^{(2)}Z^{(3)})(IY^{(2)}X^{(3)})`$
$`\times (I+s_1Z^{(2)})(I+s_2Y^{(1)}),`$
which evaluates to to $`s_2X^{(3)}(I+s_2Y^{(1)})(I+s_1Z^{(2)})`$. This implies that
$$\begin{array}{ccc}\hfill U(1,1)& =& Y\hfill \\ \hfill U(1,1)& =& X\hfill \\ \hfill U(1,1)& =& Z\hfill \\ \hfill U(1,1)& =& I\hfill \end{array}.$$
(17)
For a nice group theoretic treatment of teleportation, see braunstein:qc2000a .
The state obtained on qubits $`2`$ and $`3`$ before the last rotation and measurements is a prepared entanglement denoted by $`|t_e_{_{23}}`$. The idea is to use the built in error-detection and multiple tries to obtain the state without error.
In order to analyze the propagation of errors, we need to understand the effect of an unintended $`Z^{(1)}`$ or $`Y^{(2)}`$ measurement before the protocol’s end. Both of these commute with the two rotations in the protocol. A $`Y^{(2)}`$ measurement implies that the net effect of the applied rotations on the third qubit is a $`Z_{}^{(3)}{}_{\pm 90}{}^{}`$ (the sign depends on the measurement outcome). Nothing happens to the first qubit. If a $`Z^{(1)}`$ measurement happens, this directly applies to the first qubit. The second qubit experiences a $`Y_{}^{(2)}{}_{\pm 90}{}^{}`$ rotation. To simplify matters we intentionally perform the $`Y`$ measurement on this qubit to return to the first case as far as qubit $`3`$ is concerned.
### V.2 Logical $`ZZ`$ rotations by teleportation
One method for implementing the logical $`ZZ`$ rotation on qubits encoded in qubits $`1,2`$ and $`3,4`$ respectively is to first teleport the four qubits, then apply the rotation $`R_Z=(Z^{(1)}Z^{(2)}Z^{(3)}Z^{(4)})_{90}`$. Since the final step of the teleportation protocol involves a number of $`180^{}`$ rotations, and $`R_Z`$ is in the normalizer of the Pauli group, we can instead apply $`R_Z`$ to the destination qubits in the four copies of $`|t_e`$ and apply appropriately modified $`180^{}`$ rotations after the teleportation measurement. Actually, it is better to apply $`R_Z`$ first, use the original $`180^{}`$ corrections and note that the overall effect is equivalent to $`R_Z`$ or $`R_Z^{}`$ after teleportation. Which one actually occurred can be determined from the measurement outcomes. To go from $`R_Z^{}`$ to $`R_Z`$ it is sufficient to apply $`Z_{180}`$’s to each qubit. Because of the ability to retry the state preparation until it succeeds, this reduces the problem of reliably implementing $`R_Z`$ to the problem of reliable teleportation.
### V.3 Error recovery by teleportation
Our methods are designed so that the only error from which it is necessary to recover is a $`Z`$ measurement with known outcome on one of the qubits. It is desirable to implement the recovery so that at worst, a $`Z`$ measurement occurs on the logical qubit. By symmetry, it is sufficient to consider a $`Z`$ measurement on qubit $`1`$. The effect of the measurement is to apply $`(I+sZ^{(1)})`$, where the sign $`s`$ is known. One way of restoring the encoded qubit is to notice that
$$(I+sZ^{(1)})(I+X^{(1)}X^{(2)})=(I+siY^{(1)}X^{(2)})(I+X^{(1)}X^{(2)}).$$
(18)
Since the first operator is a $`90^{}`$ rotation around $`Y^{(1)}X^{(2)}`$, it can be undone by applying its inverse.
A method more easily made reliable is based on the syndrome measurement technique of error correction. From this perspective, the unintended $`Z`$ measurement creates a superposition of states with two syndromes, that is eigenvalues of $`X^{(1)}X^{(2)}`$. The encoded qubit can be recovered by measuring $`S^{(12)}X^{(1)}X^{(2)}`$ and if the eigenvalue is $`1`$, applying $`Z^{(1)}`$. To measure $`S^{(12)}`$ we again use teleportation, reducing the measurement problem to a state preparation and teleportation problem (see steane:qc1999a for similar ideas used to solve the more difficult problem of correcting unknown errors). The idea is to measure $`S^{(12)}`$ on the destination qubits of two copies of $`|t_e`$ before completing the protocol, and then infer the eigenvalue from the combination of all measurement outcomes. The correction operations $`U`$ of the protocol are unchanged. To see how this works, implement the protocol by teleporting qubit $`1`$ with $`|t_e_{_{36}}`$ and qubit $`2`$ with $`|t_e_{_{47}}`$, so that the teleported state ends up in qubits $`6`$ and $`7`$. The measurement on qubits $`6`$ and $`7`$ can be implemented by the sequence
$$(Y^{(6)}X^{(7)})_{90}M_{Z}^{}{}_{}{}^{(6)}(s)(Y^{(6)}X^{(7)})_{90}.$$
(19)
Note that the two $`X`$ operators on qubit $`7`$ occurring in the two-qubit rotations can be obtained by conjugating $`Z`$ operators by a $`Y_{90}`$ rotation. If the measurement of $`XX`$ results in $`s=1`$, we correct the state by applying $`Y_{}^{(6)}{}_{180}{}^{}Z_{}^{(3)}{}_{180}{}^{}`$. (To see that this correction works, examine the quantum network, moving the correction operator back to the beginning by appropriately changing orientations of rotations by anti-commuting operators, and then absorbing it at the commuting state preparation steps.) Again, we can retry this state preparation until it succeeds. The prepared state is now given by $`(I+S^{(67)})|t_e_{_{36}}|t_e_{_{47}}`$, which is already a logical qubit state on bits $`6`$ and $`7`$. Suppose teleportation concludes successfully, with correction operators $`U(s_1,s_2)^{(6)}`$ and $`U(s_3,s_4)^{(7)}`$. The resulting state is the same as if we had applied $`(I+s^{}S^{(67)})`$ after the protocol, where $`s^{}`$ depends on the $`s_i`$ ($`180^{}`$ rotations only change the sign of Pauli operators when changing the order of events in a sequence). If the state to be teleported is in the code, then $`s^{}=1`$ is impossible (as the projection is otherwise orthogonal to the state, resulting in a zero probability event), while if it is in the $`1`$ eigenvalue space of $`S^{(12)}`$, then $`s^{}=1`$ is impossible. Because of the relationship of $`s^{}`$ to the $`s_i`$, which event occurred can be determined from the combination of $`s_i`$’s that resulted from the teleportation measurement, and the necessary correction can be applied.
## VI Error analysis
We begin by considering errors that occur in the encoded operations and how one should respond to such errors.
### VI.1 Errors in recovery
The recovery procedure is applied when an unintended $`Z`$ measurement occurs. Suppose this occurs at qubit $`1`$. Since the state of qubit $`1`$ is now known, we can take advantage of this to prepare a state with the first teleportation step (that is the one involving qubit $`1`$) already completed. This can be done error-free, given a number of attempts. Specifically, the first teleportation protocol can be replaced by preparing the destination qubit in the appropriate eigenstate of $`Z`$, and then applying the $`XX`$ measurement to this and the target of the second teleportation before completing the latter. As before, one can use either outcome of the $`XX`$ measurement, in this case by applying only a $`Z_{180}`$ correction, if necessary.
If the teleportation of the second qubit fails at the source of the relevant rotation, the second qubit is untouched, and we try again. If it fails at the target, then the second qubit is measured in $`Z`$, which implies that the logical qubit is measured in $`Z`$. Let $`F_r`$ be the probability of failure. By following the different possible outcomes in the attempts, we get
$`F_r`$ $`=`$ $`fF_r+(1f)f`$ (20)
$`F_r`$ $`=`$ $`f.`$ (21)
### VI.2 Errors in the implementation of $`Z_{}^{(L)}{}_{90}{}^{}`$
When applying $`(Z^{(1)}Z^{(2)})_{90}`$ to the qubits, the following can happen: 1. The first qubit is measured in the $`Z`$ basis. In this case, apply the recovery procedure and if it succeeds, attempt the operation again. If it fails, the logical qubit is measured in $`Z`$. 2. The second qubit is measured in the $`Z`$ basis with outcome $`s`$ and the first qubit experiences a $`Z_{}^{(1)}{}_{s90}{}^{}`$. The effect is the same as if a $`Z_{}^{(L)}{}_{90}{}^{}`$ had been applied before the measurement, so if the subsequent recovery procedure succeeds, then the desired operation has been applied. Otherwise, the logical qubit has been measured in $`Z`$. The probability of failing once case 2 is entered is $`f`$. The probability of entering case 1 is $`f`$. By following the re-attempts going through case 1, we obtain the equation for the probability of failure
$`F_Z`$ $`=`$ $`(1f)f^2+f^2+f(1f)F_Z`$ (22)
$`F_Z`$ $`=`$ $`f^2(2f)/(1f(1f)).`$ (23)
### VI.3 Errors in the implementation of $`(Z^{(L_1)}Z^{(L_2)})_{90}`$
To avoid having to retry the operation we modify the protocol for implementing the logical $`ZZ`$ rotation slightly. Instead of preparing $`4`$ copies of $`|t_e`$ with the $`ZZZZ`$ rotation already applied, prepare $`k4`$ such copies. In the end, the destination qubits of the unused copies are measured in $`Z`$, so that the effective applied rotation becomes $`(Z^{(L_1)}Z^{(L_2)})_{\pm 90}`$, where the sign depends on the measurement outcomes. Compensate for a minus sign by applying $`Z_{180}`$’s to each qubit. Let the qubits be encoded in qubits $`1,2`$ and $`3,4`$ respectively. Perform the teleportation protocols for the qubits in this order. When a protocol fails at the source of the critical rotation, try again with the next available pair of qubits in the prepared state. When the protocol fails at the target, the procedure depends on whether it is the first or second member of a pair of encoding qubits. If it is the first, attempt recovery using the second qubit as usual. If it is the second, do the same, but using the already teleported qubit as the first member. In this case, we do not need to re-attempt teleportation for implementing the rotation, as in the case of the logical $`Z`$ rotation. Computing the probability of failure that the procedure fails by the first logical qubit being measured gives
$`F_{ZZ}`$ $`=`$ $`f^2+(1f)f^2+f(1f)F_{ZZ}`$ (24)
$`F_{ZZ}`$ $`=`$ $`F_Z.`$ (25)
The probabilities for the second logical qubit being measured given successful completion of the first two steps is the same, as required by the assumptions of the model.
### VI.4 The threshold
The threshold $`T_d`$ for obtaining an improvement in the failure parameter can be determined by solving $`F_{ZZ}=f`$, which gives
$$T_d=.5.$$
(26)
## VII Resource analysis
Resource analysis can be used to estimate the effect of residual errors (not fitting the error model) in the basic operations and to determine the total overhead of implementing an accurate quantum gate for standard quantum computation or communication. The total overhead includes the expected number of attempts required to prepare the requisite states. In an actual system, the state preparation attempts can be arbitrarily parallelized and implemented in independent high-throughput state factories. In the system as proposed here, the states are relatively simple in terms of the lower level implementation, and success probabilities are reasonable. An explicit total resource analysis is left as a problem for future work. For now, our primary concern is how general errors can propagate from the physical implementation to the encoded qubits. This depends only on the operations that directly contribute toward the state used in the final gate via their errors conditional on success. This property can be exploited to largely eliminate the problem of inefficient detectors, see Sect. VIII.2.
The analysis that follows is intended as an example for how this can be done and is completed with an explicit example. We begin by counting resources in terms of the operations of iLOQC, counting separately first the error-free one qubit rotations, state preparations and measurements in the $`Z`$ or $`Y`$ basis ($`R_0`$), second the $`90^{}`$ $`Z`$ and $`Y`$ rotations ($`R_1`$), and third the $`90^{}`$ two qubit rotations with $`Y`$ or $`Z`$ operators ($`R_2`$). This separation helps with the resource estimate due to the fact that they differ in resource requirements at the first and later levels. For our purposes two or three levels are expected to suffice. Note that we are not counting steps that are required to temporarily store a qubit. This is necessary in principle, as imperfect memories without parallelism imply that truly scalable quantum or classical computing is impossible aharonov:qc1996b . In particular, it is beneficial to parallelize implementations as much as possible.
As we proceed, we will comment on the expected number of tries for state preparations. For this purpose, define $`q`$ as the probability of failure of the one qubit rotations contributing to $`R_1`$ and let $`p=1f(2f)`$ be the total failure probability of the two qubit rotations contributing to $`R_2`$. In a total resource analysis, one can exploit the fact that $`q=0`$ at the first level.
### VII.1 Resources for teleportation
The preparation of $`|t_e`$ requires
$`R_0(t_e)`$ $`=`$ $`2`$ (27)
$`R_1(t_e)`$ $`=`$ $`0`$ (28)
$`R_2(t_e)`$ $`=`$ $`1.`$ (29)
Since the probability of success is $`(1p)`$, the expected number of attempts to assure success is $`1/(1p)`$.
To prepare $`2k`$ copies of $`|t_e`$ with the $`(ZZZZ\mathrm{})_{90}`$ rotation applied to the targets using the method given in Sect. IV.4 requires
$`R_0(Z^{2k})`$ $`=`$ $`2kR_0(t_e)=4k`$ (30)
$`R_1(Z^{2k})`$ $`=`$ $`2kR_1(t_e)=0`$ (31)
$`R_2(Z^{2k})`$ $`=`$ $`2kR_2(t_e)+4k+1=6k+1.`$ (32)
The probability of successful preparation is only $`(1p)^{4k+1}`$, so the preparation method needs to be improved. An efficient (in $`k`$) scheme is based on the idea of using a parity containing ancilla to kick back the desired rotation, and to generate this ancilla in a tree like fashion cmoore:qc1998a . The sequence is defined recursively by: $`S_{1}^{}{}_{}{}^{(ab)}`$ consists of preparing $`|t_e__a`$ and $`|0__b`$, then applying c-$`\sigma _x`$ from the target of $`|t_e__a`$ to $`|0__b`$. The qubit $`b`$ is the “parity” qubit. $`S_{l+1}^{}{}_{}{}^{(abd)}`$ applies $`\text{c-}\sigma _x^{(bd)}`$ to the outputs of $`S_{l}^{}{}_{}{}^{(ab)}`$ and $`S_{l}^{}{}_{}{}^{(cd)}`$, then measures parity qubit $`b`$ in the $`X`$ basis by applying $`Y_{90}`$ and then measuring $`Z`$. If the outcome is $`1`$, apply a $`Z_{180}`$ to each of the target qubits of the $`|t_e`$ that make up $`a`$. The new parity qubit is $`d`$. The controlled-not operation $`\text{c-}\sigma _x^{(ab)}`$ is applied with
$$Z_{}^{(a)}{}_{90}{}^{}X_{}^{(b)}{}_{90}{}^{}Y_{}^{(b)}{}_{90}{}^{}(Z^{(a)}Z^{(b)})_{90}Y_{}^{(b)}{}_{90}{}^{}.$$
(33)
To kick back the desired rotation to $`2^l`$ copies of $`|t_e`$ after $`S_l`$ has been successfully completed, apply $`Z_{90}`$ to the parity qubit and measure it in the $`X`$ basis, applying a $`Z_{180}`$ correction as before, if necessary. The last step can be deferred until after the teleportation when using this for implementing the logical $`ZZ`$ operation. The failure response of the algorithm can be optimized by recovering states as much as possible. For simplicity, we assume that the state associated with a failed $`S_l`$ is discarded.
The resources required for implementing $`S_l`$ can be determined recursively.
$`R_0(S_1)`$ $`=`$ $`R_0(t_e)+2=4`$ (34)
$`R_1(S_1)`$ $`=`$ $`R_1(t_e)+3=3`$ (35)
$`R_2(S_1)`$ $`=`$ $`R_2(t_e)+1=2`$ (36)
$`R_0(S_{l+1})`$ $`=`$ $`2R_0(S_l)+2+2^{l1}`$ (37)
$`R_1(S_{l+1})`$ $`=`$ $`2R_1(S_l)+4`$ (38)
$`R_2(S_{l+1})`$ $`=`$ $`2R_2(S_l)+1,`$ (39)
which one can solve to obtain
$`R_0(S_l)`$ $`=`$ $`8\times 2^{l1}4`$ (40)
$`R_1(S_l)`$ $`=`$ $`7\times 2^{l1}4`$ (41)
$`R_2(S_l)`$ $`=`$ $`3\times 2^{l1}1.`$ (42)
The probability of success of $`S_{l+1}`$ using two independent outputs of $`S_l`$ is $`(1q)^4(1p)`$, which can be shown to imply polynomial total resource use. We now obtain new expressions for the $`Z^k`$ preparation resources (assuming that the last correcting series of $`Z_{180}`$’s is deferred):
$`R_0(Z^{2^l})`$ $`=`$ $`R_0(S_l)+1=7\times 2^{l1}2`$ (43)
$`R_1(Z^{2^l})`$ $`=`$ $`R_1(S_l)+2=7\times 2^{l1}2`$ (44)
$`R_2(Z^{2^l})`$ $`=`$ $`R_2(S_l)=3\times 2^{l1}1.`$ (45)
The success probability given the output of $`S_l`$ is $`(1q)^2`$.
To prepare the state needed for recovery after one of the qubits has been measured, follow the part of the protocol that does not involve the remaining qubit. There are two measurements that need to be made, and we note that the state can be used for completing the protocol regardless of the outcome. As explained in Sect. VI.1, the preparation can be decomposed into making a copy of $`|t_e`$ and of an eigenstate of $`Z`$ and then measuring $`XX`$, correcting the outcome with a $`Z_{180}`$ if ncessary. Using the implementation of the $`XX`$ measurement above and counting the two $`Y`$ rotations needed to obtain the $`X`$ operators from $`Z`$’s in the couplings, we get
$`R_0(r_e)`$ $`=`$ $`R_0(t_e)+3=5`$ (46)
$`R_1(r_e)`$ $`=`$ $`R_1(t_e)+2=2`$ (47)
$`R_2(r_e)`$ $`=`$ $`R_2(t_e)+2=3.`$ (48)
The success probability is $`(1q)^2(1p)^2`$.
### VII.2 Resources for operations
To follow the resource usage through several levels of concatenation, it is necessary to determine the maximum resource usage for each category of operations. First are the one qubit $`180^{}`$ rotations, state preparations and measurements in the $`Z`$ or $`Y`$ basis. Of these, state preparation has the highest resource requirements for $`R_1`$ and $`R_2`$. $`R_0`$ is highest for the $`180^{}`$ rotations. This gives the following estimates:
$`R_0(0)`$ $`=`$ $`2`$ (49)
$`R_1(0)`$ $`=`$ $`3`$ (50)
$`R_2(0)`$ $`=`$ $`1.`$ (51)
The next category consists of the $`90^{}`$ $`Z`$ or $`Y`$ rotations. For the $`Y`$ rotation, we conjugate a $`Z`$ rotation by logical $`X_{90}`$’s. The resource analysis is complicated by the need for using recovery operations on partial failures. The expected number of $`ZZ`$ rotations that need to be retried after failure and successful recovery of the first qubit is $`1/(1f(1f))`$, as $`f(1f)`$ is the probability of failing and then successfully recovering. To get a better bound on the number of directly contributing operations, note that the $`ZZ`$ couplings are implemented in such a way that if the source fails, the target is not touched. Thus teleportation steps that fail at the source of the coupling that precedes the measurements need not be counted except when estimating total resources. Thus, the expected number of contributing recovery operations can be bounded by $`1/(1f(1f))1=f(1f)/(1f(1f))`$ for failures at the first qubit and $`f(1f)(1f^2/(1f(1f)))=f(1f)^2/(1f(1f))`$ for failures at the second qubit. Thus
$`R_0(1)`$ $`=`$ $`2+f(1f)(2f)/(1f(1f))R_0(r_e)`$ (52)
$``$ $`2+5f(1f)(2f)/(1f(1f))`$
$`R_1(1)`$ $`=`$ $`f(1f)(2f)/(1f(1f)))R_1(r_e)`$ (53)
$``$ $`2f(1f)(2f)/(1f(1f))`$
$`R_2(1)`$ $`=`$ $`1/(1f(1f))`$ (54)
$`+f(1f)(2f)/(1f(1f)))R_2(r_e)`$
$``$ $`1/(1f(1f))`$
$`+3f(1f)(2f)/(1f(1f)).`$
The final category has the $`90^{}`$ couplings. Except for at most $`4`$ $`X_{90}`$ rotations needed to get $`Y`$ operators, it suffices to determine the requirements for the logical $`ZZ`$ operation. Let $`2^l`$ be the number of copies of $`|t_e`$ used in the prepared state. The calculation is similar to that for $`R_x(1)`$. Noting that the probability of a teleportation failing at the target and the subsequent recovery succeeding is $`f(1f)^2`$, one can bound the expected number of recovery attempts by $`4f(1f)^2/(1f(1f)^2)`$. Some of the teleportations are attempted multiple times (just like the rotation is attempted multiple times in the previous case) and we need to account for the correction steps of the teleportation.
$`R_0(2)`$ $`=`$ $`8+2/(1f(1f))+R_0(S_l)`$
$`+(4f(1f)^2/(1f(1f)^2))R_0(r_e)`$
$``$ $`8\times 2/(1f(1f))+2^{l1}`$
$`+20f(1f)^2/(1f(1f)^2)+4`$
$`R_1(2)`$ $`=`$ $`R_1(S_l)`$
$`+(4f(1f)^2/(1f(1f)^2))R_1(r_e)`$
$``$ $`2^{l1}+(8f(1f)^2/(1f(1f)^2))2`$
$`R_2(2)`$ $`=`$ $`1/(1f(1f))+R_2(S_l)`$
$`+(4f(1f)^2/(1f(1f)^2))R_2(r_e)`$
$``$ $`1/(1f(1f))+3\times 2^{l1}`$
$`+12f(1f)^2/(1f(1f)^2)1.`$
### VII.3 An explicit example
Suppose the goal is to have a $`5\%`$ failure probability per qubit, which is not far from the current best estimates for the communication threshold briegel:qc1998a . If we use pairs of three-photon entanglements to generate the controlled sign flip at the first level in LOQC, the initial failure probability per qubit is $`f=1/4`$. Using the expression for $`F_Z`$ in Eq. 23 gives $`f=0.135`$ after one level and $`f=0.038`$ after two. For simplicity, we choose $`l=3`$ in both levels (hopefully a safe choice, though in principle this effects the probabilities a bit). Other useful values at the first and second levels are $`f/(1f)=0.333,0.156`$, $`1/(1f(1f))=1.23,1.132`$ and $`f/(1f(1f))=.308,.152`$. We evaluate the values for the coupling evolution at the second level. Using parenthesized superscripts to denote the levels gives, for example,
$`R_{0}^{}{}_{}{}^{(2)}(2)`$ $`=`$ $`(7\times 2^2+20\times 0.156+9)R_{0}^{}{}_{}{}^{(1)}(0)`$ (58)
$`+(2^2+8\times 0.1562)R_{0}^{}{}_{}{}^{(1)}(1)`$
$`+(3\times 2^2+12\times 0.156+.5)R_{0}^{}{}_{}{}^{(1)}(2)`$
$``$ $`656`$
Similarly,
$`R_{1}^{}{}_{}{}^{(2)}(2)`$ $``$ $`169`$ (59)
$`R_{2}^{}{}_{}{}^{(2)}(2)`$ $``$ $`239`$ (60)
Much of the inefficiency comes from the lack of optimization for implementing the logical $`ZZ`$ rotation.
The expected total resource usage including those needed for the failed state preparation attempts can be determined from the probabilities of successful preparation and requires recalculating the expressions for the various resources. For example, at the first level, $`p=p_1=0.44`$ and $`q=q_1=0`$ and at the second, $`p=p_2=.25`$ and $`q=q_2=.135`$. Thus the expected number of attempts required to make the state needed for error recovery is approximately $`3.16`$ at the first level and $`2.38`$ at the second.
The resource values imply that a controlled sign flip at the top level depends on less than $`250`$ controlled sign flips implemented with pairs of three photon entanglements in iLOQC. This can be used to bound the effect of errors that cannot otherwise be controlled. The resource bounds improve if four or five photon states $`|t_n`$ can be reliably generated, thus giving better initial values of $`f`$ and substantial gains in efficiency at the higher levels.
## VIII Compensating for other errors
So far we have shown how to use a simple code with careful implementation of the basic operation to rapidly boost the probability of successful completion of gates. The next step is to consider the contribution of other errors and how to correct for them. Two important types of errors are phase errors and loss of particles.
### VIII.1 Phase errors
The occurrence of a phase error can be detected for the code used above by following the teleported $`XX`$ measurement procedure in the absence of a failure. A $`1`$ eigenvalue indicates an error. At this point, it is necessary to return the state to the code space by applying an appropriate $`180^{}`$ rotation. Although it is not possible to correct for the error (we don’t know which qubit is faulty), its detection can be used as information for correction in a higher level erasure code. An alternative is to use the generalization of the code to three qubits, which does have the capability of correcting for phase errors. In fact, the $`k`$-fold concatenation of the two qubit code with itself is actually a $`2^k`$ qubit code that corrects for up to $`2^{k1}1`$ phase errors. That property can be exploited by introducing an appropriate error detection/correction procedure periodically after the first few levels, without changing the overall concatenation scheme or the basic methods for implementing operations.
### VIII.2 Loss of particles and detector inefficiency
The methods of LOQC include one that can, with some probability of failure, detect loss of a photon used to define a photonic qubit. The failure mode is again one involving a $`Z`$ measurement, so the scheme can be used with the two qubit code. This will have an effect on the overall error behavior. We leave the calculations as an open problem.
One observation not made in knill:qc2000b is that if after any of the teleportation steps used for basic LOQC, we measure the modes not containing the teleported qubits, and the total number of photons detected is not equal to the number initially prepared in the entanglement, then a photon was lost, possibly due to detector inefficiency. Such an event can be declared as a qubit loss. Doing so turns the problem of detector inefficiency into one of having to handle detected qubit loss at some rate. It is a useful task to determine the relationship between detector inefficiency and the probability of detected loss.
If total loss is detected, the two qubit code is insufficient for restoring the encoded information. Instead, it is necessary (perhaps after a few levels of two qubit encodings) to use erasure codes. These are codes with the property that one (or a few) lost qubits can be restored without loss of information. The accuracy threshold for the error model where each operation satisfies that the target qubits are lost with independent probability $`s`$ appears to be very good also, perhaps below $`95\%`$. A brief explanation based on conservative calculations giving a value close to $`99\%`$ is below. The erasure model is one of few for which it is possible to establish exact quantum communication quantities, such as the quantum channel capacity bennett:qc1997b . We therefore suggest that calculating the threshold for the erasure error-model is an excellent open problem to solve.
### VIII.3 An erasure code
A useful erasure code encoding one qubit into four is defined by the stabilizer group generated by
$$S_E=\{X^{(1)}X^{(2)}X^{(3)},Y^{(2)}Y^{(3)}Y^{(4)},Z^{(1)}Z^{(3)}Z^{(4)}\}$$
(61)
One can define encoded operators by
$`Z^{(L)}`$ $`=`$ $`Z^{(2)}Z^{(3)}`$ (62)
$`X^{(L)}`$ $`=`$ $`X^{(3)}X^{(4)}.`$ (63)
We chose the erasure code and the logical operators for their small support (number of non-identity Pauli operators in the products). If it is desirable to encode two qubits into four, the erasure code with stabilizer generated by $`XXXX`$ and $`ZZZZ`$ can be used.
It turns out that it is easier to analyze thresholds for erasure errors than for $`Z`$ measurements. Implementations of operations can be based on teleportation the same way we did before, again relying on the ability to guarantee prepared states by retrying after failure. For the present discussion, we implement the teleportations required for an operation all in one step, and then follow up with error recovery if necessary. Of the operations required for quantum computing, state preparation can be implemented error-free at all levels by retrying the process until success. To measure $`Z^{(L)}`$, first measure $`Z^{(2)}`$, and if this fails, restore the logical qubit. After successfully measuring $`Z^{(2)}`$, measure $`Z^{(3)}`$, and if that fails, $`Z^{(1)}`$ and $`Z^{(4)}`$. The value of the $`Z^{(3)}`$ measurement can be inferred by using the parity constraint associated with the third generator of $`S_E`$. The other logical Pauli operators can be measured similarly. The probability that the measurement fails is bounded by $`sr+3s^2`$, where $`r`$ is the probability that the error recovery procedure fails.
Suppose that qubit $`1`$ is lost. The error recovery procedure requires measuring the first and third operators of $`S_E`$, which, by reintroducing a fixed state for qubit $`1`$, can be done with two full teleportation steps each. We ignore the fact that some failures in these steps can be recovered and estimate that each teleportation step fails with probability at most $`3s`$ (due to the coupling rotation and two measurements, not including the correction for now). The recovery step concludes with a correcting operation on qubit $`1`$, so that we can estimate $`r13s`$ (ignoring the possibility of retrying the recovery step if one of the last corrections fail). The probability of a failed measurement is therefore bounded by $`16s^2`$.
Since we are unable to implement arbitrary rotations directly, it is necessary to implement (say) a $`45^{}`$ rotation by teleportation. (The compensation step is now a $`90^{}`$ coupling rotation, which can be implemented either by teleportation or directly, exploiting the independence assumption on loss of qubits.) Since the most complex operation involves coupling two logical qubits with a $`90^{}`$ rotation, we finish by giving a rough estimate of this failure probability. The Hamiltonians to be evolved have weight four (two on each pair of encoding qubits), so it suffices to teleport four qubits after a suitable state preparation. If one of the teleportations fail, the recovery procedure is followed. We may need to compensate for a negative rotation induced by the teleportations’ last correction step, but that can be absorbed into the procedure. Thus the error probability for the first logical qubit is bounded by $`(8s)(13s)=104s^2`$, whence an accuracy threshold better than about $`99\%`$.
### VIII.4 General errors
The techniques discussed so far compensate for all of the primary errors that are expected to occur in an LOQC system. Additional errors can be attributed to improper settings of beam splitters, undetected loss, stray photons, etc. The first of these depends on accurate calibration of classical control parameters. Happily, this type of error affects the probabilities quadratically, so that it can be minimized well by engineering. To deal with other errors eventually requires more powerful quantum error correction, and the goal of any implementation is to minimize the need for these techniques.
## IX Conclusion
This work is a first attempt at reducing the overhead and need for efficient devices for implementing LOQC. It shows that even without much optimization of the operations or tight estimates of errors and resources, there are techniques that can be used to obtain useful quantum operations in LOQC with overheads within two orders of magnitude of those required for other reliable implementations of quantum computers. The advantages of LOQC include the ability to compensate readily for the primary errors in optics while preparing large numbers of the basic quantum systems, which are photons in a superposition of two modes. It is necessary to further optimize the methods and to better analyze the error behavior, particularly the effects of detector inefficiencies, single photon state preparation errors and timing or overlap problems for photons. A fruitful area of further investigation is to determine whether larger codes and the method of encoding multiple qubits at once can be used to improve efficiency and error behavior.
Our proposal consists of a multi-level system with various types of error-control gradually introduced and supported by high-output state preparation factories that exploit the ease with which many photons can be produced. Although in the very long term, solid state or molecular computing methods are the preferred implementation for large scale quantum computing, LOQC is now a viable alternative to achieving the capability of non-trivial quantum information processing. One area where LOQC has a long term future is in communication. Photon based systems are currently the only reasonable proposals for long distance quantum communication. Since the necessary accuracies for successfully exchanging entanglement over arbitrary distances are well below $`99\%`$, this may also be the first application of LOQC to quantum information processing to be experimentally implemented.
Acknowledgments: We thank the Aspen Center for Physics for its hospitality. E.K. and R.L. received support from the NSA and from the DOE (contract W-7405-ENG-36).
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# 1 Introduction
## 1 Introduction
In recent years, a relation between superstring theory and a deformation quantization has been explored. D-branes, boundaries of open strings, are non-perturbative object of superstring theory. Matrix Models \[BFSS, IKKT\] were proposed as a D-brane action a few years ago. It is shown that a stable solution of this action is non-commutative manifold \[CDS, AIIKKT\]. This non-commutativity, however, comes from the Moyal quantization. So, this solution implies a flat D-brane We also derive a noncommutative gauge theory on a fuzzy sphere from the matrix model\[IKTW\] . If we regard some D-brane as space-time, we should study the deformation quantization of curved spaces in order to realize the quantum gravity.
In order to proceed further it will be useful to clarify mathematical background of the deformation quantization. The star product was first introduced by Groenewold\[Gr\], which is now known as the Moyal product\[Mo\]. They associate an operator product to a noncommutative product of functions. Here, the operators are mapped into the functions by taking account the Weyl ordering. The Weyl ordering means a skew-symmetric definition as we will see later. Also Berezin tried to quantize curved phase spaces about 25 years ago and succeeded to quantize some Kähler manifold e.g. sphere\[Be\]. Recently the Berezin quantization has been generalized to arbitrary Kähler manifold\[RT\]. However the Berezin quantization is defined without skew symmetry, hence is not a generalization of Moyal one.
The correlation between these methods of quantization is, however, not clear at all. In this paper, we attempt to skew-symmetrize the Berezin quantization by means of the multiple star product method. The multiple star products reduce to the path-integral form in large $`N`$ limit. (See \[Sh, Al\] for the original ideas.) As a result, our formulation turns similar to the path-integral form of the Kontsevich quantization which is defined perturbatively on Poisson manifold\[Ko\] but also can be described by a bosonic string path-integral\[CF\]. Especially in the flat case, our star product coincides with the Kontsevich star product.
This paper consists of the following sections. In section 2, we review the deformation quantization e.g. Moyal \[Mo\], Berezin \[Be, MM\] and Kontsevich \[Ko, CF\] quantization. In section 3, we first construct the multiple star product method and explain the symmetrized Berezin (or Wick type) star product\[SW, Ma\]. We also study its associativity in detail. In section 4, we derive the path integral form of the Kontsevich star product on the flat plane from the multiple star product method. Section 5 is devoted to discussions. Appendix includes some examples that the multiple star product method is available.
## 2 Deformation Quantization
This section includes the definition and properties of the deformation quantization. We also review Moyal, Berezin and Kontsevich quantization briefly as examples.
### 2.1 General Definition and Property
The deformation quantization\[BFFLS, St\] is provided by a star product, which is defined by
$$fg=\underset{m=0}{\overset{\mathrm{}}{}}B_m(f,g)\lambda ^m,$$
(1)
where
* $`\lambda `$ is a deformation parameter,
* $`f,gA=C^{\mathrm{}}(M)[[\lambda ]]`$
$`C^{\mathrm{}}(M)[[\lambda ]]`$ means that the coefficients of $`\lambda `$ power series are $`C^{\mathrm{}}`$ functions on $`M`$ ,
* $`B_m`$ are bi-differential operators ($`B_m:A\times AA`$),
The deformation quantization has the following properties:
1. associativity
$$f(gh)=(fg)h.$$
(2)
2. $`m=0`$
$$B_0(f,g)=fg.$$
(3)
3. $`m=1`$
$$B_1(f,g)B_1(g,f)=\{f,g\}=2\underset{i,j}{}\alpha ^{ij}_if_jg,$$
(4)
where $`i,j=1,2,\mathrm{},d=dim(M)`$ and $`\{,\}`$ is a Poisson bracket which satisfies
$$\{f,\{g,h\}\}+\{g,\{h,f\}\}+\{h\{f,g\}\}=0,$$
(5)
so the skew-symmetric bivector field $`\alpha `$ satisfies
$$\alpha ^{il}_l\alpha ^{jk}+\alpha ^{jl}_l\alpha ^{ki}+\alpha ^{kl}_l\alpha ^{ij}=0.$$
(6)
It can be shown that one or more star product determined by (2)(3)(4) exist. The deformation quantization has the following equivalence called a gauge equivalence. $``$ and $`^{}`$ are identified if
$$f^{}^{}g^{}=D(fg),$$
(7)
where $`f^{}=D(f),g^{}=D(g)`$ and $`D`$ is a differential operator ($`D:AA`$). However, we can take two simple gauges. One is the skew-symmetric gauge
$$B_1(f,g)=\frac{1}{2}\{f,g\}=\underset{i,j}{}\alpha ^{ij}_if_jg.$$
(8)
The other gauge is
$$B_1(f,g)=\underset{i,j}{}\beta ^{ij}_if_jg,$$
(9)
where $`\beta `$ is the upper triangle matrix of $`\alpha `$ which satisfies $`\alpha ^{ij}=\beta ^{ij}\beta ^{ji}`$ . So we call the star product determined by (8) and (9) the skew-symmetric and asymmetric product respectively. We have three concrete examples of deformation quantization, which are put together in the following table.
| Manifold | flat plane | Kähler | Poisson |
| --- | --- | --- | --- |
| Quantization | Moyal | Berezin | Kontsevich |
| Symbol | $``$ | | $``$ |
| $`m=1`$ | skew-symmetric | asymmetric | skew-symmetric |
We consider real two dimensional manifolds for the sake of simplicity from now on.
### 2.2 Moyal Quantization
The Poisson bracket on the flat plane is defined by
$$\{f,g\}=\underset{i,j}{}\epsilon ^{ij}_if_jg=_xf_pg_pf_xg.$$
(10)
Thus $`\alpha ^{ij}=\epsilon ^{ij}/2`$ by (4). So we obtain the associative star product on the flat plane i.e. the Moyal star product as
$`fg(x,p)`$ $`=`$ $`f(x,p)e^{\frac{\lambda }{2}(\stackrel{}{_x}\stackrel{}{_p}\stackrel{}{_p}\stackrel{}{_x})}g(x,p)`$ (11)
$`=`$ $`fg+\lambda {\displaystyle \frac{1}{2}}\{f,g\}+O(\lambda ^2).`$
This star product agrees with eq.(3) and (8), and satisfies the associativity from the following results
$$e_1(e_2e_3)=e^{\frac{\lambda }{2}(m_1n_2+m_2n_3+n_3m_1n_1m_2n_2m_3m_3n_1)}e_1e_2e_3=(e_1e_2)e_3,$$
(12)
where $`e_i`$’s are the Fourier series
$$e_i=e^{i(m_ix+n_ip)}.$$
Here we require a usual canonical commutation relation
$$[x,p]_{}=xppx=i\mathrm{},$$
(13)
so that we obtain $`\lambda =i\mathrm{}`$ . Also this star product can be written by the integral form\[Ba\], because we have the following relations
$$e^{i(mx+np)}e^{i(m^{}x+n^{}p)}=e^{\frac{i\mathrm{}}{2}(mn^{}m^{}n)}e^{i(mx+np)}e^{i(m^{}x+n^{}p)}$$
(14)
and
$$\frac{dwd\eta }{\pi \mathrm{}}\frac{dw^{}d\eta ^{}}{\pi \mathrm{}}e^{\frac{2i}{\mathrm{}}S}e^{i(mw+n\eta )}e^{i(m^{}w^{}+n^{}\eta ^{})}=e^{\frac{i\mathrm{}}{2}(mn^{}m^{}n)}e^{i(mx+np)}e^{i(m^{}x+n^{}p)},$$
(15)
where
$$S=\left|\begin{array}{ccc}1& 1& 1\\ x& w& w^{}\\ p& \eta & \eta ^{}\end{array}\right|.$$
The left hand sides of eq.(14) and (15) are equivalent, so that we obtain the integral form of the Moyal star product after multiplying arbitrary Fourier coefficients and integrating over $`m,n`$ as
$$fg(x,p)=\frac{dwd\eta }{\pi \mathrm{}}\frac{dw^{}d\eta ^{}}{\pi \mathrm{}}e^{\frac{2i}{\mathrm{}}S}f(w,\eta )g(w^{},\eta ^{}).$$
(16)
### 2.3 Berezin Quantization
The Poisson bracket on the Kähler manifold is given as
$$\{f,g\}=\frac{2}{i}h^{z\overline{z}}(_zf_{\overline{z}}g_{\overline{z}}f_zg),$$
(17)
where $`h^{z\overline{z}}`$ is the inverse of a Kähler metric
$$h_{z\overline{z}}=_z_{\overline{z}}K(z,\overline{z}),$$
and $`K(z,\overline{z})`$ is a Kähler potential. The factor $`2/i`$ in eq.(17) is necessary in order that the Poisson bracket becomes eq.(10) in the flat case. The original Berezin quantization covers only Ricci flat Kähler manifold. The Berezin star product is defined by
$$f\text{}g(z,\overline{z})=𝑑\mu _\nu (v,\overline{v})e^{\frac{1}{\nu }\mathrm{\Phi }(z,\overline{z},v,\overline{v})}f(z,\overline{v})g(v,\overline{z}),$$
(18)
where $`\mathrm{\Phi }(z,\overline{z},v,\overline{v})`$ is called the Calabi function and defined by the Kähler potential as
$$\mathrm{\Phi }(z,\overline{z},v,\overline{v})=K(z,\overline{v})+K(v,\overline{z})K(z,\overline{z})K(v,\overline{v}),$$
(19)
and the measure $`d\mu _\nu `$ is determined by the metric as
$$d\mu _\nu (z,\overline{z})=h_{z\overline{z}}\frac{idzd\overline{z}}{2\pi \nu }.$$
(20)
This star product can be expanded around $`\lambda =\nu /2i=0`$ as follows
$$f\text{}g(z,\overline{z})=fg+\lambda \left(B_1^+(f,g)+B_1^{}(f,g)\right)+O(\lambda ^2),$$
(21)
where $`B_1^+`$ and $`B_1^{}`$ are a symmetric part and a skew-symmetric part respectively as
$$B_1^+=2iAfg\frac{1}{2}\{f,g\}_+,B_1^{}=\frac{1}{2}\{f,g\}$$
(22)
and
$$A=\frac{1}{2}h^{z\overline{z}}_z_{\overline{z}}\mathrm{log}h_{z\overline{z}},\{f,g\}_+=\frac{2}{i}h^{z\overline{z}}(_zf_{\overline{z}}g+_{\overline{z}}f_zg).$$
If the manifold $`M`$ is the Kähler manifold, $`A=0`$ \[RT\]. Thus (3) and (4) are satisfied. Also the associativity is shown as the following.
$$\left((f\text{}g)\text{}h\right)(z,\overline{z})=𝑑\mu _\nu (v,\overline{v})𝑑\mu _\nu (u,\overline{u})f(z,\overline{v})g(v,\overline{u})h(u,\overline{z})e^{\frac{1}{\nu }\left(\mathrm{\Phi }(z,\overline{u},v,\overline{v})+\mathrm{\Phi }(z,\overline{z},u,\overline{u})\right)}.$$
$$\left(f\text{}(g\text{}h)\right)(z,\overline{z})=𝑑\mu _\nu (v,\overline{v})𝑑\mu _\nu (u,\overline{u})f(z,\overline{v})g(v,\overline{u})h(u,\overline{z})e^{\frac{1}{\nu }\left(\mathrm{\Phi }(v,\overline{z},u,\overline{u})+\mathrm{\Phi }(z,\overline{z},v,\overline{v})\right)}.$$
The Calabi function clearly satisfies
$$\mathrm{\Phi }(z,\overline{u},v,\overline{v})+\mathrm{\Phi }(z,\overline{z},u,\overline{u})=\mathrm{\Phi }(v,\overline{z},u,\overline{u})+\mathrm{\Phi }(z,\overline{z},v,\overline{v}),$$
so the associativity is shown:
$$\left((f\text{}g)\text{}h\right)(z,\overline{z})=\left(f\text{}(g\text{}h)\right)(z,\overline{z}).$$
(23)
### 2.4 Kontsevich Quantization
The Kontsevich quantization covers the Poisson manifold ($`M`$) which is a general manifold with the Poisson structure. He perturbatively solved $`B_m(f,g)`$’s under the conditions (2),(3) and (8) as
$$B_m(f,g)=\underset{\mathrm{\Gamma }G_m}{}w_\mathrm{\Gamma }B_\mathrm{\Gamma }(f,g),$$
(24)
where $`G_m`$ is a set of diagrams related to the number $`m`$ , $`B_\mathrm{\Gamma }(f,g)`$ is a bi-differential operator determined by the Feynman diagram and $`\omega _\mathrm{\Gamma }`$ is a weight \[Ko\]. Thus Kontsevich defines the star product on the Poisson manifold by a formal power series of $`\lambda `$ as
$$fg=\underset{m=0}{\overset{\mathrm{}}{}}\lambda ^m\underset{\mathrm{\Gamma }G_m}{}w_\mathrm{\Gamma }B_\mathrm{\Gamma }(f,g).$$
(25)
Also Cattaneo and Felder have shown that the Kontsevich star product (25) coincides with the path integral form of a topological bosonic string (non-linear sigma model):
$$fg(x)=_{X(\mathrm{})=x}f(X(1)g(X(0))e^{\frac{i}{\mathrm{}}S[X,\eta ]}𝒟X𝒟\eta ,$$
(26)
where the action is defined on a disk $`D`$ as
$$S[X,\eta ]=_D\eta _i(u)dX^i(u)+\frac{1}{2}\alpha ^{ij}(X(u))\eta _i(u)\eta _j(u),$$
and
* $`D=\{uR^2,|u|1\}`$,
* $`X`$ and $`\eta `$ are real bosonic fields,
* $`X:DM`$,
* $`\eta `$ is a differential 1-form on $`D`$ : $`X^{}(T^{}M)T^{}D`$.
In the symplectic case, the action can be integrated over $`\eta `$ and becomes a boundary integration by the Stokes’s theorem as
$$f_{\mathrm{symp}}g(x)=_{\gamma (\pm \mathrm{})=x}f(\gamma (1))g(\gamma (0))e^{\frac{i}{\mathrm{}}_\gamma d^1\omega }𝑑\gamma ,$$
(27)
where $`\gamma `$ is a loop trajectory from $`x`$ to $`x`$ .
## 3 Symmetrized Berezin star product
In this section, we first explain the multiple star product method. Next we define the S-star product to clarify the relationship between Moyal and Berezin quantization. However the S-star product is not associative in the curved space. So using the multiple star product method, we derive an associative star product i.e. the O-star product.
### 3.1 Multiple star product method
Generally the integral forms of the star products can be written as
$$fg(\alpha )=𝑑\mu _\lambda (\beta ,\gamma )e^{𝒦_\lambda (\alpha ,\beta ,\gamma )}f(\beta )g(\gamma ),$$
(28)
where $`\alpha ,\beta ,\gamma M`$, $`𝒦_\lambda =𝒦/\lambda `$ is an integral kernel and $`d\mu _\lambda =d\mu /\lambda ^2`$ is a measure which relates two points on $`M`$. We assume that this star product $``$ satisfies the followings:
$$fg=fg+\lambda \frac{\{f,g\}}{2}+O(\lambda ^2),$$
(29)
$$f1=1f=f,$$
(30)
$$d\mu _\lambda (\beta ,\gamma )=d\mu _\lambda (\gamma ,\beta ).$$
(31)
We also add a assumption $`𝒦_\lambda (\alpha ,\beta ,\beta )=0`$ in particular. Note that we don’t require this star product $``$ is associative. Then we call this product $``$ the non-associative star product.
Next we define the multiple star product of $``$ as
$$A^N(f)=f_{N/N}f_{N1/N}\mathrm{}f_{2/N}f_{1/N}.$$
(32)
An equivalence of the forward product $`\stackrel{}{A}^N`$ and the backward product $`\stackrel{}{A}^N`$ is necessary at least in order that $`A^N`$ is well-defined where
$`\stackrel{}{A}^N(f)`$ $`:=`$ $`(f_{N/N}(f_{N1/N}(\mathrm{}(f_{1/N}1)\mathrm{})))`$ (33)
$`=`$ $`{\displaystyle \left(\underset{i=1}{\overset{N}{}}d\mu _{\lambda N}(\beta _{i/N},\gamma _{i/N})f_{i/N}(\beta _{i/N})\right)\mathrm{exp}\underset{i=1}{\overset{N}{}}\frac{1}{N}𝒦_\lambda (\gamma _{i+1/N},\beta _{i/N},\gamma _{i/N})},`$
$`\stackrel{}{A}^N(f)`$ $`:=`$ $`(((\mathrm{}(1f_{N/N})\mathrm{})f_{2/N})f_{1/N})`$ (34)
$`=`$ $`{\displaystyle \left(\underset{i=1}{\overset{N}{}}d\mu _{\lambda N}(\beta _{i/N},\gamma _{i/N})f_{i/N}(\beta _{i/N})\right)\mathrm{exp}\underset{i=1}{\overset{N}{}}\frac{1}{N}𝒦_\lambda (\gamma _{i1/N},\gamma _{i/N},\beta _{i/N})},`$
$$\alpha =\beta _0=\beta _{N+1/N}=\gamma _0=\gamma _{N+1/N}.$$
(35)
Note that we change the deformation parameter $`\lambda `$ to $`\lambda N`$. ¿From this equivalence, we obtain a condition
$$\frac{1}{N}\underset{i=1}{\overset{N}{}}\left(𝒦_\lambda (\gamma _{i+1/N},\beta _{i/N},\gamma _{i/N})𝒦_\lambda (\gamma _{i1/N},\gamma _{i/N},\beta _{i/N})\right)=0.$$
(36)
Using the boundary condition (35) and the additional condition $`𝒦(\alpha ,\beta ,\beta )=0`$, this condition (36) is also deformed as
$$\frac{1}{N}\underset{i=0}{\overset{N}{}}\left(𝒦_\lambda (\gamma _{i+1/N},\beta _{i/N},\gamma _{i/N})𝒦_\lambda (\gamma _{i/N},\gamma _{i+1/N},\beta _{i+1/N})\right)=0.$$
(37)
This condition corresponds to the associativity condition in the case of $`f_{i/N}=1`$ except for three $`f_{i/N}`$’s. We denote $`A^N(f):=\stackrel{}{A}^N(f)=\stackrel{}{A}^N(f)`$ when $`𝒦_\lambda `$ satisfies eq.(37).
### 3.2 Relationship between Moyal and Berezin Star Product
The Berezin star product in the flat case, coincides with the Moyal’s one except for skew-symmetry or asymmetry. This difference is explained as follows. First, we write the Moyal star product in complex variables to make clear the correspondence to Berezin star product,
$$fg(z,\overline{z})=f(z,\overline{z})e^{\mathrm{}(\stackrel{}{_z}\stackrel{}{_{\overline{z}}}\stackrel{}{_{\overline{z}}}\stackrel{}{_z})}g(z,\overline{z}),$$
(38)
where $`z=x+ip`$. This star product is gauge equivalent(7) to $`_{st}`$ and $`_{ar}`$ where
$$_{st}=e^{2\mathrm{}\stackrel{}{_z}\stackrel{}{_{\overline{z}}}}\mathrm{and}_{ar}=e^{2\stackrel{}{_{\overline{z}}}\stackrel{}{_z}},$$
(39)
because the gauge equivalent condition is satisfied in the case of
$$D=e^{\mathrm{}\stackrel{}{_z}\stackrel{}{_{\overline{z}}}}\mathrm{and}D=e^{\mathrm{}\stackrel{}{_{\overline{z}}}\stackrel{}{_z}},$$
(40)
respectively\[Vo, Be, APS\]. Thus we obtain a star product relation,
$$=(_{st}_{ar})^{\frac{1}{2}}.$$
(41)
Here $`_{st}^{\frac{1}{2}}`$ and $`_{ar}^{\frac{1}{2}}`$ can be written in the integral forms\[APS\] as
$`f_{st}^{\frac{1}{2}}g(z,\overline{z})`$ $`=`$ $`{\displaystyle \frac{idwd\overline{w}}{2\pi \theta }e^{\frac{1}{\theta }|wz|^2}f(w,\overline{z})g(z,\overline{w})},`$
$`f_{ar}^{\frac{1}{2}}g(z,\overline{z})`$ $`=`$ $`{\displaystyle \frac{idvd\overline{v}}{2\pi (\theta )}e^{\frac{1}{\theta }|vz|^2}f(z,\overline{v})g(v,\overline{z})},`$ (42)
where $`\theta =\mathrm{}`$. Thus we obtain
$$fg(z,\overline{z})=\frac{idvd\overline{v}}{2\pi \theta }\frac{idwd\overline{w}}{2\pi \theta }e^{\frac{1}{\theta }(|vz|^2+|wz|^2)}f(w,\overline{v})g(v,\overline{w}).$$
(43)
The star products (42) are two types of the Berezin star product on the flat plane i.e. the term $`|vz|^2`$ is the Calabi function on the flat plane. Taking this result into account, we generalize the Moyal star product $``$ to the S-star product on the Ricci flat Kähler manifold <sup>§</sup><sup>§</sup>§In \[Ma\], it is discussed that the S-star product may be available to general Kähler manifold., which is defined by
$$f\text{}g(z,\overline{z}):=𝑑\mu _\theta (v,\overline{v})𝑑\mu _\theta (w,\overline{w})\mathrm{exp}\frac{1}{\theta }\left(\mathrm{\Phi }(z,\overline{z};v,\overline{v})\mathrm{\Phi }(z,\overline{z};w,\overline{w})\right)f(w,\overline{v})g(v,\overline{w}),$$
(44)
where $`d\mu _\theta (z,\overline{z})=h_{z\overline{z}}idzd\overline{z}/2\pi \theta `$ similarly to the definition (20). However, this star product is not associative unless flat. This complication is overcome by using the multiple star product method.
### 3.3 Associativity of Symmetrized Berezin Star Product
In this section, we attempt to recover the associativity of the S-star product. First, we show that the S-star product is non-associative star product. In the case of the S-star product, we know the following correspondence:
$$\alpha =(z,\overline{z}),\beta =(w,\overline{v}),\gamma =(v,\overline{w}),$$
(45)
$$\lambda =\theta ,d\mu _\lambda (\beta ,\gamma )=d\mu _\theta (v,\overline{v})d\mu _\theta (w,\overline{w}),$$
(46)
$$𝒦_\lambda (\alpha ,\beta ,\gamma )=\frac{1}{\theta }\left(\mathrm{\Phi }(z,\overline{z};v,\overline{v})\mathrm{\Phi }(z,\overline{z};w,\overline{w})\right).$$
(47)
Thus the S-star product clearly satisfies the conditions (31) and $`𝒦_\lambda (\alpha ,\beta ,\beta )=0`$ . Also it is derived in \[Ma\] that this product satisfies the conditions (29) and (30).
Next, we survey whether the S-star product satisfies the condition (37) or not. In the S-star product, $`lhs`$ of (37) is written in terms of the Kähler potential $`K`$ as
$`lhs`$ $`=`$ $`{\displaystyle \frac{1}{\theta }}{\displaystyle \frac{1}{N}}{\displaystyle \underset{i=0}{\overset{N}{}}}\left(K(v_{i+1/N},\overline{v}_{i/N})+K(v_{i/N},\overline{v}_{i+1/N})2K(v_{i/N},\overline{v}_{i/N})\right)`$ (48)
$`\left(K(w_{i+1/N},\overline{w}_{i/N})+K(w_{i/N},\overline{w}_{i+1/N})2K(w_{i/N},\overline{w}_{i/N})\right).`$
This result is non-zero but becomes zero in the large $`N`$ limit as
$$lhs\frac{1}{\theta }_0^1𝑑\tau d\left(K(v,\overline{v})K(w,\overline{w})\right)=0,$$
(49)
where $`v,w=v(\tau ),w(\tau )`$ and the boundary conditions (35) become
$$v(0)=v(1)=w(0)=w(1)=z,\overline{v}(0)=\overline{v}(1)=\overline{w}(0)=\overline{w}(1)=\overline{z}.$$
(50)
Thus $`A^N(f)`$ is ill-defined but $`A(f)=lim_N\mathrm{}A^N(f)`$ is well-defined. Then we call this star product pseudo-associative product. By using $`A(f)`$ and
$$f_{i/N}=\{\begin{array}{ccc}g\hfill & (i/N=i_1/N\tau _1)\hfill & \\ f\hfill & (i/N=i_2/N\tau _2)\hfill & \\ 1\hfill & (i/N\tau \tau _1,\tau _2)\hfill & \end{array},$$
(51)
we construct an associative star product in terms of the path integral form as
$`f\text{}g(z,\overline{z}):=A(f)`$ (52)
$`=`$ $`\mathrm{}1\text{}1\text{}f\text{}1\text{}1\mathrm{}1\text{}1\text{}g\text{}1\text{}1\mathrm{}`$
$`=`$ $`\underset{N\mathrm{}}{lim}{\displaystyle \underset{i=1}{\overset{N}{}}d\mu _{\theta N}(v_{i/N},\overline{v}_{i/N})d\mu _{\theta N}(w_{i/N},\overline{w}_{i/N})\mathrm{exp}\left[\frac{i}{\theta }S_i\right]f(v_{i_2/N},\overline{w}_{i_2/N})g(v_{i_1/N},\overline{w}_{i_1/N})}`$
$`=`$ $`{\displaystyle 𝒟\mu _\theta (v,\overline{v})𝒟\mu _\theta (w,\overline{w})\mathrm{exp}\left[\frac{i}{\theta }S\right]f(v(\tau _2),\overline{w}(\tau _2))g(v(\tau _1),\overline{w}(\tau _1))},`$
where actions are written as follows
$$S_i=\frac{1}{i}\frac{1}{N}\underset{i=1}{\overset{N}{}}\left(\mathrm{\Phi }(v_{i+1/N},\overline{w}_{i+1/N};v_{i/N},\overline{v}_{i/N})\mathrm{\Phi }(v_{i+1/N},\overline{w}_{i+1/N};w_{i/N},\overline{w}_{i/N})\right),$$
(53)
$$S=\frac{1}{i}_{\mathrm{}}^{\mathrm{}}\left[\left(\psi (v,\overline{v})\psi (v,\overline{w})\right)\frac{v}{\tau }\left(\overline{\psi }(w,\overline{w})\overline{\psi }(v,\overline{w})\right)\frac{\overline{v}}{\tau }\right]𝑑\tau ,$$
(54)
and the path integral measure is defined by
$$𝒟\mu _\theta (v,\overline{v}):=\underset{N\mathrm{}}{lim}\underset{i=1}{\overset{N}{}}d\mu _{\theta N}(v_{i/N},\overline{v}_{i/N}).$$
(55)
Note that $`\psi (z,\overline{z})`$ is a canonical conjugation of $`z`$, which is defined by
$$\psi (z,\overline{z}):=\frac{K(z,\overline{z})}{z}.$$
(56)
As above, associative symmetrized Berezin star product is defined as O-star product correctly. The associativity is satisfied as illustrated in Figure 1.
## 4 Construction of Kontsevich Star Product from Multiple Star Product Method
In this section, we show that the multiple Moyal star product corresponds to the path integral form of the Kontsevich star product on the flat plane. First preparatory to this derivation, we write the multiple Moyal star product in large $`N`$ limit as
$`A_{}(f)`$ $`:=`$ $`\underset{N\mathrm{}}{lim}f_{N/N}(x,p)f_{N1/N}(x,p)\mathrm{}f_{2/N}(x,p)f_{1/N}(x,p)`$ (60)
$`=`$ $`\underset{N\mathrm{}}{lim}{\displaystyle \underset{i=1}{\overset{N}{}}\frac{d\xi _{i/N}d\eta _{i/N}}{\pi \mathrm{}N}\frac{d\xi _{i/N}^{}d\eta _{i/N}^{}}{\pi \mathrm{}N}f_{i/N}(\xi _{i/N},\eta _{i/N})\mathrm{exp}\left[\frac{2i}{\mathrm{}N}\underset{i=1}{\overset{N}{}}\left|\begin{array}{ccc}1& 1& 1\\ \xi _{i+1/N}^{}& \xi _{i/N}& \xi _{i/N}^{}\\ \eta _{i+1/N}^{}& \eta _{i/N}& \eta _{i/N}^{}\end{array}\right|\right]}`$
$`=`$ $`{\displaystyle 𝒟\xi 𝒟\eta 𝒟\xi ^{}𝒟\eta ^{}\underset{\tau =0}{\overset{1}{}}f(\tau ;\xi ,\eta )\mathrm{exp}\frac{2i}{\mathrm{}}_0^1𝑑\tau \left[\frac{d\xi ^{}}{d\tau }(\eta \eta ^{})(\xi \xi ^{})\frac{d\eta ^{}}{d\tau }\right]},`$ (61)
where real fields $`\xi ,\eta ,\xi ^{},\eta ^{}`$ have boundary conditions
$$x=\xi (0)=\xi (1)=\xi ^{}(0)=\xi ^{}(1),p=\eta (0)=\eta (1)=\eta ^{}(0)=\eta ^{}(1),$$
(62)
and functional measures are defined as follows
$$𝒟\xi :=\underset{N\mathrm{}}{lim}\underset{i=1}{\overset{N}{}}\frac{d\xi _{i/N}}{\pi \mathrm{}N},\mathrm{}etc.$$
(63)
In eq.(61), we can integrate out $`\xi ^{},\eta ^{}`$ by using partial integration and obtain simplified form
$$A_{}(f)=𝒟\xi 𝒟\eta \underset{\tau =0}{\overset{1}{}}f_\tau (\xi (\tau ),\eta (\tau ))\mathrm{exp}\frac{i}{\mathrm{}}_0^1\eta \frac{\xi }{\tau }𝑑\tau .$$
(64)
Here we can change the integration area of $`\tau `$ $`(0,1)`$ to $`(\mathrm{},\mathrm{})`$ by a reparametrization of $`\tau `$ . Thus eq.(64) and the boundary conditions (62) are changed as
$$A_{}(f)=𝒟\xi 𝒟\eta \underset{\tau =\mathrm{}}{\overset{\mathrm{}}{}}f_\tau (\xi (\tau ),\eta (\tau ))\mathrm{exp}\frac{i}{\mathrm{}}_{\mathrm{}}^{\mathrm{}}\eta \frac{\xi }{\tau }𝑑\tau .$$
(65)
$$x=\xi (\pm \mathrm{})=\xi ^{}(\pm \mathrm{}),p=\eta (\pm \mathrm{})=\eta ^{}(\pm \mathrm{}).$$
(66)
Next by using eq.(65), we show that the multiple Moyal star product is included in the path integral form of the Kontsevich star product(27). If we put
$$f_\tau (x,p)=\{\begin{array}{cc}f(x,p)\hfill & (\tau =1)\hfill \\ g(x,p)\hfill & (\tau =0)\hfill \\ 1\hfill & (\tau 0,1)\hfill \end{array},$$
(67)
then
$`fg(x,p)`$ $`=`$ $`A_{}(f)`$ (68)
$`=`$ $`\underset{N\mathrm{}}{lim}\mathrm{}1f(x,p)1\mathrm{}1g(x,p)1\mathrm{}`$
$`=`$ $`{\displaystyle 𝒟\xi 𝒟\eta f(\xi (1),\eta (1))g(\xi (0),\eta (0))\mathrm{exp}\frac{i}{\mathrm{}}_\gamma d^1\omega _0},`$
where
$$d^1\omega _0:=(d\xi )\eta =\eta \frac{d\xi }{d\tau }d\tau ,$$
(69)
and
$$\omega _0=d(d^1\omega _0)=d\xi d\eta .$$
(70)
Eq.(68) corresponds to eq.(27) on the flat plane.
## 5 Discussions
In this paper, we have proposed the multiple star product method. It is useful for constructing associative star products from pseudo-associative star products. We have shown that a pseudo-associative S-star product becomes an associative O-star product by using the multiple star product method. It is explained as follows. Although the pseudo-associative products break associativity condition a little, the multiple star product method, which is a set of infinite pseudo-associative product, smoothes and overcomes this risk and the associativity condition (37) is satisfied. In consequence, the pseudo-associative product turns to an associative product within the frame work of the path integral formalism.
The multiple star product method also has been available to well-known associative products e.g. the Moyal star product. The multiple Moyal star product coincides with the path integral form of the Kontsevich star product on the flat plane. This result implies a justice of the multiple star product method.
## Appendix A Other Examples
By using eq.(65), we can obtain the transition amplitude in quantum dynamics and the bosonic string generating function. If we put in eq.(65)
$$f_\tau (x,p)=\{\begin{array}{cc}\psi _I(x)\hfill & (\tau =t_I)\hfill \\ \overline{\psi _F}(x)\hfill & (\tau =t_F)\hfill \\ e^{\frac{i}{\mathrm{}}H(x,p)}\hfill & (t_I<\tau <t_F)\hfill \\ 1\hfill & (\tau <t_I,t_F<\tau )\hfill \end{array},$$
(71)
we obtain
$`A_{}(f)`$ $`=`$ $`\underset{N\mathrm{}}{lim}\mathrm{}1\overline{\psi }_F(x)e^{\frac{i\epsilon }{\mathrm{}}H(x,p)}\mathrm{}e^{\frac{i\epsilon }{\mathrm{}}H(x,p)}\psi _I(x)1\mathrm{}`$ (72)
$`=`$ $`{\displaystyle 𝒟\xi 𝒟\eta \overline{\psi }_F(\xi (t_F))\psi _I(\xi (t_I))e^{\frac{i}{\mathrm{}}_{t_I}^{t_F}\left(\eta \frac{\xi }{\tau }H(\xi ,\eta )\right)𝑑\tau }}`$
$`=`$ $`\psi _F,t_F|\psi _I,t_I,`$
where integration of $`\tau <t_I,t_F<\tau `$ vanishes because of $`f(\tau ;x,p)=1`$ .
Next, a bosonic string generating function can be derived from infinite dimensional multiple Moyal star products. In (65), we put
$$f_\tau (x,p)=e^{i\left(k(\tau )x+m(\tau )p\right)},$$
(73)
and obtain
$`A_{}(f)`$ $`=`$ $`\underset{N\mathrm{}}{lim}e^{i\left(k(\mathrm{})x+m(\mathrm{})p\right)}\mathrm{}e^{i\left(k(0)x+m(0)p\right)}\mathrm{}e^{i\left(k(\mathrm{})x+m(\mathrm{})p\right)}`$ (74)
$`=`$ $`{\displaystyle 𝒟\xi 𝒟\eta e^{\frac{i}{\mathrm{}}_{\mathrm{}}^{\mathrm{}}\left(\eta \frac{\xi }{\tau }+\mathrm{}(k\xi +m\eta )\right)𝑑\tau }}.`$
Here, we generalize $`xx_{n,\mu },pp_n^\mu `$ where $`n`$ runs from $`0`$ to $`\mathrm{}`$ and $`\mu `$ runs from $`1`$ to dimension $`d`$ . $`n`$ can be changed to continuous parameter $`\sigma `$ by the following definitions and relations
$$x_\mu (\sigma ):=\underset{n=0}{\overset{\mathrm{}}{}}x_{n,\mu }\mathrm{cos}n\sigma ,p^\mu (\sigma ):=\underset{n=0}{\overset{\mathrm{}}{}}\frac{p_n^\mu }{n}\mathrm{sin}n\sigma ,$$
$$\underset{n=0}{\overset{\mathrm{}}{}}x_{n,\mu }p_n^\mu =2_0^{2\pi }x_\mu \frac{p^\mu }{\sigma }𝑑\sigma .$$
Substituting above relation into $`A_{mn}`$ , we obtain
$`A_{}(f)={\displaystyle 𝒟X\mathrm{exp}\left[\frac{2i}{\mathrm{}}_{\mathrm{}}^{\mathrm{}}𝑑\tau _0^{2\pi }𝑑\sigma \left(\frac{X^\mu }{\sigma }\frac{X_\mu }{\tau }+J^\mu X_\mu \right)\right]},`$ (75)
$$X_\mu =\frac{\xi _\mu +\eta _\mu }{\sqrt{2}},J^\mu =\mathrm{}\left(\frac{n^\mu }{\sigma }\frac{m^\mu }{\sigma }\right),𝒟X=𝒟\xi 𝒟\eta .$$
$`A_{}(f)`$ becomes the bosonic string generating function. More details are in \[SW\].
Acknowledgments
I’m indebted to Professors Y.Maeda and S.Saito for useful discussions.
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# 1 Introduction
## 1 Introduction
Nowadays there are sufficient indications to believe that neutrinos are massive particles mixing with each other . These indications come from both experimental and theoretical sides. The solar neutrinos deficit, the atmospheric neutrino anomaly and the results of the LSND neutrino oscillation experiment, all can be explained in terms of neutrino oscillations implying non-zero neutrino masses and mixings. On the theoretical side almost all phenomenologically viable models of the physics beyond the standard model(SM) predict non-zero masses for neutrinos which can be either Majorana or Dirac particles.
Majorana masses violate the conservation of total lepton number by two units $`\mathrm{\Delta }L=2`$. Therefore lepton number violating($`L/`$) processes represent a most appropriate tool to address the question of whether neutrinos are Majorana or Dirac particles. Various $`L/`$processes have been studied in the literature in this respect. Among them there are the neutrinoless nuclear double beta ($`0\nu \beta \beta `$) decay , the decay $`\mathrm{K}^+\mu ^+\mu ^+\pi ^{}`$ , nuclear muon to positron or to antimuon conversion, trimuon production in neutrino-nucleon scattering , the process $`e^+p\overline{\nu }l_1^+l_2^+X`$, relevant for HERA , as well as direct production of heavy Majorana neutrinos at various colliders . The analysis made in the literature leads to the conclusion that if these processes are mediated by the Majorana neutrino exchange then, except $`0\nu \beta \beta `$-decay, they can hardly be observed experimentally. This analysis relies on the current neutrino oscillation data, and on certain assumptions about the neutrino mass matrix. In the present paper we concentrate on the neutrinoless double muon decay of kaon $`\mathrm{K}^+\mu ^+\mu ^+\pi ^{}`$. We will show that despite the above conclusion being true for contributions of the neutrino states much lighter or much heavier than the typical energy of the $`\mathrm{K}^+\mu ^+\mu ^+\pi ^{}`$ decay, there is still a special window in the neutrino sector which can be efficiently probed by searching for this process. This window is in the neutrino mass range $`245\text{ MeV}m_{\nu _j}389`$ MeV, where the s-channel neutrino contribution to the $`\mathrm{K}^+\mu ^+\mu ^+\pi ^{}`$ decay is resonantly enhanced, therefore making this decay very sensitive to the neutrinos in this mass domain. If neutrinos with masses in this region exist, then from present experimental data we can extract stringent limits on their mixing with $`\nu _\mu `$. We derive these limits from the upper bound on the branching ratio of $`\mathrm{K}^+\mu ^+\mu ^+\pi ^{}`$ recently obtained by E865 experiment at BNL .
## 2 $`\mathrm{K}^+\mu ^+\mu ^+\pi ^{}`$ decay in Standard Model with Majorana neutrinos
In the SM extension with Majorana neutrinos there are two lowest order diagrams, shown in Fig.1, which contribute to $`\mathrm{K}^+\mu ^+\mu ^+\pi ^{}`$ decay. These diagrams were first considered long ago in Refs. . Here we are studying previously overlooked aspects of this decay. We concentrate on the s-channel neutrino exchange diagram in Fig. 1(a) which plays a central role in our analysis. The t-channel diagram in Fig. 1(b) requires in general a detailed hadronic structure calculation. In Ref. this diagram was evaluated in the Bethe-Salpeter approach and shown to be an order of magnitude smaller than the diagram in Fig. 1(a), for light and intermediate mass neutrinos. As we will see, in the neutrino mass domain of our main interest, the diagram in Fig. 1(a) absolutely dominates over the t-channel diagram in Fig. 1(b), independently of hadronic structure.
The contribution from the factorizable s-channel diagram in Fig. 1(a) can be calculated in a straightforward way, without referring to any hadronic structure model. A final result for the $`\mathrm{K}^+\mu ^+\mu ^+\pi ^{}`$ decay rate is given by
$`\mathrm{\Gamma }(\mathrm{K}^+\mu ^+\mu ^+\pi ^{})=c{\displaystyle \underset{s_1^{}}{\overset{s_1^+}{}}}𝑑s_1\left|{\displaystyle \underset{k}{}}{\displaystyle \frac{U_{\mu k}^2m_{\nu k}}{s_1m_{\nu k}^2}}\right|^2G({\displaystyle \frac{s_1}{m__K^2}})+`$ (1)
$`2{\displaystyle \frac{c}{m__K^2}}\mathrm{Re}{\displaystyle \underset{k,n}{}}[{\displaystyle \underset{s_1^{}}{\overset{s_1^+}{}}}𝑑s_1{\displaystyle \frac{U_{\mu k}^2m_{\nu k}}{s_1m_{\nu k}^2}}{\displaystyle \underset{s_2^{}}{\overset{s_2^+}{}}}𝑑s_2\left({\displaystyle \frac{U_{\mu n}^2m_{\nu n}}{s_2m_{\nu n}^2}}\right)^{}H({\displaystyle \frac{s_1}{m__K^2}},{\displaystyle \frac{s_2}{m__K^2}})].`$
The unitary mixing matrix $`U_{ij}`$ relates $`\nu _i^{}=U_{ij}\nu _j`$ weak $`\nu ^{}`$ and mass $`\nu `$ neutrino eigenstates. The numerical constant in Eq. (1) is $`c=(G_F^4/32)(\pi )^3f_\pi ^2f__K^2m__K^5|V_{ud}|^2|V_{us}|^2`$, where $`f__K=1.28f__\pi `$, $`f__\pi =0.668m_\pi `$ and $`m_K=494`$ MeV is the K-meson mass. The functions $`G(z)`$ and $`H(z_1,z_2)`$ in Eq. (1) after the phase space integration can be written in an explicit algebraic form
$`G(z)={\displaystyle \frac{\varphi (z)}{z^2}}\left[h_+(z)h_{}(z)x_\pi ^2h_+(z)\right]\left[x_\mu ^2+z(x_\mu ^2z)^2\right]`$ (2)
$`H(z_1,z_2)=h_{}(z_1)h_{}(z_2)+x_\pi ^2[r_+(z_1z_2)x_\mu ^2t(z_1,z_2,1)]r_{}(z_1z_2)t(z_1,z_2,x_\mu ).`$
Here we defined $`x_i=m_i/m__K`$ and $`h_{\pm \pm }(z)=z\pm x_\pi ^2\pm x_\mu ^2`$, $`r_\pm (z_1z_2)=z_1z_2x_\pi ^2\pm x_\mu ^4`$, $`t(z_1,z_2,z_3)=z_1+z_22z_3^2`$, $`\varphi (z)=\lambda ^{1/2}(1,x_\mu ^2,z)\lambda ^{1/2}(z,x_\mu ^2,x_\pi ^2)`$ with $`\lambda (x,y,z)=x^2+y^2+z^22xy2yz2xz`$. The integration limits in Eq. (1) are
$`s_1^{}=m__K^2(x_\pi +x_\mu )^2,s_1^+=m__K^2(1x_\mu )^2,`$ (3)
$`s_2^\pm ={\displaystyle \frac{m__K^2}{2y}}\left[2y(1+x_\mu ^2)(1+yx_\mu ^2)h_+(y)\pm \varphi (y)\right]`$
with $`y=s_1/m__K^2`$.
First we assume that neutrinos can be separated into light $`\nu _k`$ and heavy $`N_k`$ states, with masses $`m_{\nu i}<<\sqrt{s_1^{}}`$ and $`\sqrt{s_1^+}<<M_{Nk}`$. Then the Eq. (1) can be approximately written as
$`\mathrm{\Gamma }(\mathrm{K}^+\mu ^+\mu ^+\pi ^{})=\left|m_\nu _{\mu \mu }\right|^2m__K^1𝒜_\nu +\left|M_N^1_{\mu \mu }\right|^2m__K^3𝒜_N`$ (4)
$`2\text{Re}\left[m_\nu _{\mu \mu }M_N^1_{\mu \mu }\right]m__K𝒜_{\nu N},`$
where the average neutrino masses are determined in the standard way
$`m_\nu _{\mu \mu }={\displaystyle \underset{k=light}{}}U_{\mu k}^2m_{\nu k},M_N^1_{\mu \mu }={\displaystyle \underset{k=heavy}{}}U_{\mu k}^2M_{Nk}^1.`$ (5)
The approximate formula (4) can be used for extracting limits on the average neutrino masses from the experimental data. This leads to the limits: $`|m_\nu |Exp(\nu )`$, $`|M_N^1|Exp(N)`$. We point out that in this case the upper bounds must satisfy the consistency conditions $`Exp(\nu )<<\sqrt{s_1^{}}m__K,Exp(N)^1>>\sqrt{s_1^+}m__K`$. Otherwise the so derived limits are not applicable to $`m_\nu `$, $`M_N^1`$ as it is in Refs. . If consistency conditions are not satisfied one has to use the initial formula (1). Similar consistency conditions take place for the other $`L/`$processes.
The dimensionless coefficients in Eq. (4) are $`𝒜_\nu =4.0\times 10^{31},𝒜_N=7.0\times 10^{32},𝒜_{\nu N}=1.7\times 10^{31}`$. With these numbers we can estimate the current upper bound on the $`\mathrm{K}^+\mu ^+\mu ^+\pi ^{}`$ decay rate from the experimental data on other processes.
Atmospheric and solar neutrino oscillation data, combined with the tritium beta decay endpoint, allow one to set upper bounds on the masses of the known three neutrinos $`m_{e,\mu ,\tau }3`$ eV. Thus in the three neutrino scenario one gets $`|m_\nu _{\mu \mu }|9\text{ eV}`$. In this case we derive from Eq. (4) the following branching ratio
$`_{\mu \mu }={\displaystyle \frac{\mathrm{\Gamma }(\mathrm{K}^+\mu ^+\mu ^+\pi ^{})}{\mathrm{\Gamma }(\mathrm{K}^+all)}}3.0\times 10^{30}(3\mathrm{light}\mathrm{neutrino}\mathrm{scenario}).`$ (6)
Assuming the existence of heavy (i.e. order $``$ GeV) mass neutrinos N, we may obtain an upper bound on $`\mathrm{\Gamma }(\mathrm{K}^+\mu ^+\mu ^+\pi ^{})`$ using the current LEP limit on heavy stable neutral leptons $`M_N39.5`$ GeV , which leads to $`|M_N^1_{\mu \mu }|n\left(39.5\text{ GeV}\right)^1`$, where $`n`$ is the number of heavy neutrinos. This limit being substituted in Eq. (4) results in the upper bound
$`_{\mu \mu }2.0\times 10^{19}(3\mathrm{light}+1\mathrm{heavy}\mathrm{neutrino}\mathrm{scenario}).`$ (7)
Direct searches for $`\mathrm{K}^+\mu ^+\mu ^+\pi ^{}`$ decay by E865 experiment at BNL give
$`_{\mu \mu }3.0\times 10^9(\text{ 90\%CL},\text{Ref. }\text{[15]})`$ (8)
Comparison of this experimental bound with the theoretical predictions in Eqs. (6), (7) clearly shows that both cases are far from being ever detected. On the other hand experimental observation of $`\mathrm{K}^+\mu ^+\mu ^+\pi ^{}`$ decay at larger rates would indicate some new physics beyond the SM, or, as we will see, the presence of an extra neutrino state $`\nu _j`$ with the mass in the hundred MeV domain. Let us note that such neutrinos are not excluded by the LEP neutrino counting experiments measuring the Z-boson invisible width. These experiments set limits not on the number of light massive neutrinos but on the number of active neutrino flavors $`N_\nu =3`$. Thus extra massive neutrino states $`\nu _j`$ can appear as a result of mixing of the three active neutrinos with certain number of sterile neutrinos. These massive neutrinos are searched for in many experiments . The $`\nu _j`$ states would manifest themselves as peaks in differential rates of various processes, and can give rise to significant enhancement of the total rate if their masses lie in an appropriate region.
## 3 Hundred-MeV neutrinos in $`\mathrm{K}^+\mu ^+\mu ^+\pi ^{}`$-decay. Resonant case.
Assume there exists a massive Majorana neutrino $`\nu _j`$ with the mass $`m_j`$ in the range
$`\sqrt{s_1^{}}245\text{ MeV}m_j\sqrt{s_1^+}389\text{ MeV}.`$ (9)
Then the s-channel diagram in Fig. 1(a) blows up because the integrand of the first term in Eq. (1) has a non-integrable singularity at $`s=m_j^2`$. Therefore, in this resonant domain the total $`\nu _j`$-neutrino decay width $`\mathrm{\Gamma }_{\nu j}`$ has to be taken into account. This can be done by the substitution $`m_jm_j(i/2)\mathrm{\Gamma }_{\nu j}`$. The total decay width $`\mathrm{\Gamma }_{\nu j}`$ of the Majorana neutrino $`\nu _j`$ with the mass in the resonant domain (9) receives contributions from the following decay modes:
$`\nu _j\{\begin{array}{c}e^+\pi ^{},e^{}\pi ^+,\mu ^+\pi ^{},\mu ^{}\pi ^+,\hfill \\ e^+e^{}\nu _e^c,e^+\mu ^{}\nu _\mu ^c,\mu ^+e^{}\nu _e^c,\mu ^+\mu ^{}\nu _\mu ^c\hfill \\ e^{}e^+\nu _e,e^{}\mu ^+\nu _\mu ,\mu ^{}e^+\nu _e,\mu ^{}\mu ^+\nu _\mu .\hfill \end{array},`$ (13)
Since $`\nu _j\nu _j^c`$ it can decay in both $`\nu _jl^{}X(\mathrm{\Delta }L=0)`$ and $`\nu _jl^+X^c(\mathrm{\Delta }L=2)`$ channels. Calculating partial decay rates we obtain
$`\mathrm{\Gamma }(\nu _jl\pi )=|U_{lj}|^2{\displaystyle \frac{G_F^2}{4\pi }}f_\pi ^2m_j^3F(y_l,y_\pi )|U_{lj}|^2\mathrm{\Gamma }_2^{(l)},`$ (14)
$`\mathrm{\Gamma }(\nu _jl_1l_2\nu )=|U_{l_1j}|^2{\displaystyle \frac{G_F^2}{192\pi ^3}}m_j^5H(y_{l1},y_{l2})|U_{l_1j}|^2\mathrm{\Gamma }_3^{l_1l_2},`$ (15)
where $`y_i=m_i/m_j`$ and
$`F(x,y)=\lambda ^{1/2}(1,x^2,y^2)[(1+x^2)(1+x^2y^2)4x^2],`$ (16)
$`H(x,y)=12{\displaystyle \underset{z_1}{\overset{z_2}{}}}{\displaystyle \frac{dz}{z}}(zy^2)(1+x^2z)\lambda ^{1/2}(1,z,x^2)\lambda ^{1/2}(0,y^2,z).`$ (17)
The integration limits are $`z_1=y_{l_2}^2,z_2=(1y_{l_1})^2`$ and $`F(0,0)=H(0,0)=1`$. Summing up all the decay modes in (13) one gets for the total $`\nu _j`$ width
$`\mathrm{\Gamma }_{\nu j}=2|U_{\mu j}|^2(\mathrm{\Gamma }_2^{(\mu )}+\mathrm{\Gamma }_3^{(\mu e)}+\mathrm{\Gamma }_3^{(\mu \mu )})+2|U_{ej}|^2(\mathrm{\Gamma }_2^{(e)}+\mathrm{\Gamma }_3^{(ee)}+\mathrm{\Gamma }_3^{(e\mu )})`$
$`|U_{\mu j}|^2\mathrm{\Gamma }_\nu ^{(\mu )}+|U_{ej}|^2\mathrm{\Gamma }_\nu ^{(e)}.`$ (18)
In the resonant domain (9) $`\mathrm{\Gamma }_{\nu j}`$ reaches its maximum value at $`m_j=\sqrt{s_1^+}`$. Assuming for the moment $`|U_{\mu j}|=|U_{ej}|=1`$, we estimate this maximum value to be $`\mathrm{\Gamma }_{\nu j}4.7\times 10^{10}`$ MeV. Since $`\mathrm{\Gamma }_{\nu j}`$ is so small in the resonant domain (9) the neutrino propagator in the first term of Eq. (1) has a very sharp maximum at $`s=m_j^2`$. The second term, being finite in the limit $`\mathrm{\Gamma }_{\nu j}=0`$, can be neglected in the considered case. Thus, with a good precision we obtain from Eq. (1)
$`\mathrm{\Gamma }^{res}(\mathrm{K}^+\mu ^+\mu ^+\pi ^{})c\pi G(z_0){\displaystyle \frac{m_j|U_{\mu j}|^4}{|U_{\mu j}|^2\mathrm{\Gamma }_\nu ^{(\mu )}+|U_{ej}|^2\mathrm{\Gamma }_\nu ^{(e)}}}`$ (19)
with $`z_0=(m_j/m_K)^2`$. This equation allows one to derive, from the experimental bound of Eq. (8), the constraints on $`\nu _j`$ neutrino mass $`m_j`$ and the mixing matrix elements $`U_{\mu j},U_{ej}`$ in a form of a 3-dimensional exclusion plot. However one may reasonably assume that $`|U_{\mu j}||U_{ej}|`$. Then from the experimental bound (8) we derive a 2-dimensional $`m_j|U_{\mu j}|^2`$ exclusion plot given in Fig. 2. For comparison we also present in Fig. 2 the existing bounds taken from . As shown in the figure, the experimental data on the $`\mathrm{K}^+\mu ^+\mu ^+\pi ^{}`$ decay exclude a region unrestricted by the other experiments. The constraints can be summarized as
$`|U_{\mu j}|^2(5.6\pm 1)\times 10^9\text{for}245\text{ Mev}m_j385\text{ MeV},`$ (20)
The best limit $`|U_{\mu j}|^24.6\times 10^9`$ is achieved at $`m_j300`$ MeV. Note that these limits are compatible with our assumption that $`|U_{\mu j}||U_{ej}|`$ since in this mass domain, typically $`|U_{ej}|^210^9`$ .
The constraints from $`\mathrm{K}^+\mu ^+\mu ^+\pi ^{}`$ in Fig. 2 and Eq. (20) can be significantly improved in the near future by the experiments E949 at BNL and E950 at FNAL . It is important to notice that in the resonant domain we have $`\mathrm{\Gamma }^{res}(\mathrm{K}^+\mu ^+\mu ^+\pi ^{})|U|^2`$, while outside $`\mathrm{\Gamma }(\mathrm{K}^+\mu ^+\mu ^+\pi ^{})|U|^4`$. Thus in the resonant mass domain the $`\mathrm{K}^+\mu ^+\mu ^+\pi ^{}`$ decay has a significantly better sensitivity to the neutrino mixing matrix element. In forthcoming experiments the upper bound on the ratio in Eq. (8) can be improved by two orders of magnitude or even more. Then this experimental bound could be translated to the limit $`|U_{\mu j}|^2<10^{11}`$ and stronger.
Finally, it is important to mention the possible cosmological and astrophysical consequences of the hundred-MeV neutrinos considered in the present paper. Massive neutrinos contribute to the mass density of the universe, participate in cosmic structure formation, big-bang nucleosynthesis, supernova explosions, imprint themselves in the cosmic microwave background etc. (see for a review). This implies certain constrains on the neutrino masses and mixings. Currently, for massive neutrinos in the mass region of Eq. (9), the only available cosmological constraints arise from the mass density of the universe, $`\tau _{\nu _j}<(10^{14})`$ sec, and cosmic structure formation, $`\tau _{\nu _j}<(10^7)`$ sec, (taken from ref. ) where $`\tau _{\nu _j}`$ is the lifetime of massive neutrinos. On the other hand, on the basis of the formula (3), assuming $`|U_{\mu j}|^2|U_{ej}|^24.6\times 10^9`$ as in Eq. (20), we find conservatively $`10^2`$ sec $`<\tau _{\nu _j}`$. Thus massive neutrinos with masses in the interval (9) do not contradict the known cosmological constraints and there remains a wide open interval of allowed mixing matrix elements: $`(10^{18})<|U_{\mu j}|^2,|U_{ej}|^2<(10^9)`$. Big-bang nucleosynthesis and the SN 1987A neutrino signal may lead to much more restrictive constraints. Unfortunately, as yet the analysis of these constraints does not involve the mass region (9). It may happen that these constraints, in combination with our constrains in Eq. (20), close the window for neutrinos with masses in the interval (9). Then the only physics left to be studied using the $`\mathrm{K}^+\mu ^+\mu ^+\pi ^{}`$ searches would be physics beyond the SM other than neutrino issues. Nevertheless, significant model dependence of all the cosmological constraints should be carefully considered before such a determining conclusion is finally drawn.
## 4 Conclusion
We studied the potential of the K-meson neutrinoless double muon decay as a probe of the Majorana neutrino masses and mixings. We found that this process is very sensitive to the handred-MeV neutrinos $`\nu _j`$ in the resonant mass range (9). We analyzed the contribution of these neutrinos to the $`\mathrm{K}^+\mu ^+\mu ^+\pi ^{}`$ decay rate and derived stringent upper limits on the Majorana neutrino mixing matrix element $`|U_{\mu j}|^2`$ from the current experimental data. In Fig. 2 we presented these limits in the form of a 2-dimensional exclusion plot, and compared them with existing limits. The $`\mathrm{K}^+\mu ^+\mu ^+\pi ^{}`$ decay excludes a domain previously unrestricted experimentally. We point out that the current and near future experimental searches for this decay are not able to provide any information on the light eV-mass ($`m_\nu `$ eV) or heavy GeV-mass ($`m_\nu `$ GeV) neutrinos, since in those cases the required experimental sensitivities are by many orders of magnitude far from the realistic ones. Finally, we notice that the decay $`\mathrm{K}^+\mu ^+\mu ^+\pi ^{}`$ can in principle probe lepton number violating interactions beyond the standard model. A well known example is given by R-parity violating interactions in supersymmetric models. However these aspects of the $`\mathrm{K}^+\mu ^+\mu ^+\pi ^{}`$ decay are yet to be studied.
Acknowledgments
We thank Ya. Burdanov, O. Espinosa and R. Rosenfeld for helpful discussions. This work was supported in part by Fondecyt (Chile) under grants 1990806, 1000717, 1980150 and 8000017, and by a Cátedra Presidencial (Chile).
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# Tidal Tails Around 20 Galactic Globular Clusters: Based on observations made at the European Southern Observatory, La Silla, Chile,Plate scanning done with the MAMA (Machine Automatique à Mesurer pour l’Astronomie), a facility developed and operated by the INSU (Institut National des Sciences de l’Univers) at the Observatoire de Paris, France
## 1 Introduction
Recent works have emphasized that globular clusters (GCs) are stellar systems which exhibit strong dynamical evolution in the potential well of their host galaxy (Gnedin & Ostriker 1997, hereafter GO97, Murali & Weinberg 1997). Since their birth, GCs suffer internal evolution on their own: at the very beginning GCs evolve rapidly because of the fast stellar evolution of their massive stars (Vesperini & Heggie 1997, Portegies Zwart et al. 1997). The finite size of the GCs produces two-body relaxation between the member stars; equipartition of energy heats the lighter stars which diffuse outwards into the halo while the heavier stars sink slowly towards the contracting core (Spitzer & Hart 1971). Typically, the relaxation time $`t_{rh}`$ (hereafter we will refer to the relaxation time at the half mass radius; see, e.g., Binney & Tremaine 1987) in a GC is of the order of $`t_{rh}=10^9`$ yr (GO97), which is significantly less than the age of all galactic globular clusters.
Owing to the negative specific heat in self-gravitating stellar systems (Antonov 1962, Lynden-Bell & Wood 1968), related to the virial equilibrium, the core is contracting monotonically during a process called the “gravothermal catastrophe”, leading to core collapse, characterized by very high stellar densities (up to $`10^6\text{M}_{}\text{pc}^3`$). It has been eventually recognized that this process is not so catastrophic after all, since the cluster core does not collapse for ever but bounces back towards lower stellar density phases (Hénon 1975). The key role of binaries (either primordial or formed via encounters during the high stellar density phase) has been emphasized in the internal dynamics of the GCs: these binaries act as a heating source and slow down the core collapse (Goodman & Hut 1989) and even reverse it. Many authors have studied numerically the post core-collapse phase, using conducting-gas-sphere, Fokker-Planck, and N-body codes, in simulations computed well into core collapse and beyond, leading to the discovery of possible post-collapse gravothermal oscillations (see Meylan & Heggie 1997 for a review). Globular clusters evolve dynamically, even when considering only relaxation, which causes stars to escape, consequently cluster cores to contract and envelopes to expand.
Any galaxy through its gravitational potential well influences the dynamical evolution of its globular clusters, accelerating their destruction. The stars in the globular cluster halo are stripped by the tidal field of the galaxy: the outward diffusion of the stars towards the sub-thermalised halo is speeded up and the core contracts even more (Spitzer & Chevalier 1973). Moreover the gravitational shocks heat up the outer parts of the globular cluster, increasing its loss of stars (Aguilar et al. 1988, Weinberg 1994, Kundic & Ostriker 1995, Leon et al. 2000). Shocks are caused by the tidal field of the galaxy: interactions with the disk, the bulge and, at a lower level, with the giant molecular clouds (GMCs, see Spitzer 1958), heat up the outer regions of each star clusters. The disk-shocking occurs when the GC crosses the thin disk where it is compressed by the varying z-component of the galactic plane potential; this has been found to dominate the heating of GCs (Chernoff et al. 1986). A GC globally gains energy during the crossing and exhibits peculiar transient deformation (Leon et al. 2000). The shocks with the bulge and the GMCs are similar in their physical processes: the GCs suffer an elongation aligned parallel to the density gradient in the bulge or the GMC. These processes combined with the internal dynamical evolution have probably destroyed an important fraction of the primordial GCs and are still currently at play. GO97 estimate that “half of the present clusters are to be destroyed within the next Hubble time”.
All GCs are expected to have already lost an important fraction of their mass, deposited in the form of individual stars in the halo of the Galaxy. The mass-loss rate is a function of the total mass of the cluster, its structural parameters like the concentration $`c`$ (where $`c`$ = log ($`r_t/r_c`$) with $`r_t`$ and $`r_c`$ are the tidal and core radii), and its orbital motion around the galactic center. Till recently, the only way to investigate the orbital history of a globular cluster, apart from proper motions, was to derive its tidal radius $`r_t`$, which gives an indication of its perigalactic distance. Unfortunately, there are difficulties in defining the tidal radius, both theoretically and observationally; and, as expected, there are discrepancies between the theoretical and observational values of the tidal radius (see, e.g., Odenkirchen et al. 1997, Scholz et al. 1998).
N-body simulations of globular clusters embedded in a realistic galactic potential (Oh & Lin 1992; Johnston et al. 1997; Combes, Leon & Meylan 1999, hereafter CLM99) have been performed in order to study the amount of mass loss for different kinds of orbits and different kinds of clusters, along with the dynamics and the mass function in tidal tails. The 2-D structure of such tidal tails appears to be a good tracer of gravitational shocks and should be a tracer of the potential well. Moreover the detection of unbound stars released by the clusters is the only way to measure directly the mass loss rate of the cluster.
Grillmair et al. (1995) in an analysis of star-count in the outer parts of a few galactic globular clusters found extra-cluster overdensities that they associated partly with stars stripped into the Galaxy field. Similar tidal interaction remnants around globular clusters have also been found in external galaxies: Grillmair et al. (1996) observed three clusters in M31 which exhibit departures from a King profiles. Leon et al. (1999) find tidal extensions in the outskirts of interacting binary clusters and isolated clusters in the Large Magellanic Cloud (LMC). Not surprisingly, galactic tidal forces play also an essential role in the evolution of smaller-scale galactic clusters, the open clusters: they exhibit tidal tails in their neighborhood (Odenkirchen 1998; Bergond et al. 2000).
In this work we study the 2-D structure of the tidal tails associated with 20 galactic globular clusters by using the wavelet transform to detect weak structures at large scale and filter the strong background noise for the low galactic latitude clusters. We analyze with great care the observational bias which can be very important. In Section 2 we present the observations, in Section 3 we describe the data reduction, in Section 4 we present the results with, for each globular cluster, detailed comments about related observational bias. Section 5 presents the general discussion of all results.
## 2 Observations
We chose to observe, for the present study, a number of clusters sharing different properties or locations in the Galaxy, with various masses and structural parameters. In order to detect the tidal extensions around globular clusters, it is necessary to have very wide field observations. During the years 1996 and 1997, optical observations were performed at the ESO Schmidt telescope with photographic films. The field of view is of $`5.5{}_{}{}^{}\times 5.5^{}`$ with a scale of 67.5″/mm. In combination with the UG1, BG12 and RG630 filters, the Kodak Tech Pan emulsion 4415 provides passbands close to $`U`$, $`B`$, and $`R`$, respectively. Eventually, we use only the $`R`$ and $`B`$ filters for this analysis, since the $`U`$ band images are significantly shallower than the images in the other two filters. The sample of clusters is enlarged by using some ESO and SERC survey plates available at the Centre d’Analyse des Images (CAI) of the Paris Observatory, along with a few POSS I plates kindly provided by Leiden Observatory. All the observed globular clusters studied in this paper are listed in Table 1 with corresponding plates/films identification numbers, filters, and exposure times. All these photographic plates and films are digitized using the MAMA (Machine Automatique à Mesurer pour l’Astronomie) scanning machine at the CAI, which provides a pixel size of 10 $`\mu m`$ . The astrometric performances of the machine are described in Berger et al. (1991).
Table 2 displays for each cluster, its J2000 equatorial coordinates ($`\alpha `$,$`\delta `$), its galactic coordinates (l,b), its equatorial rectangular coordinates (X,Y,Z), and, when available, its velocity components (U,V,W) in the equatorial rectangular system.
## 3 Data reduction
### 3.1 Source extraction
Once the plates/films are digitized, the next step consists of identifying all point sources in these frames. The source extraction is performed on each frame using SExtractor (Bertin & Arnouts 1996), a software dedicated to the automatic analysis of astronomical images using a multi-threshold algorithm allowing good object deblending. The detection of the stars is done at a 3-$`\sigma `$ level above the background. This software, which can deal with huge amount of data (up to 60,000 $`\times `$ 60,000 pixels) is not suited for very crowded field like the centers of the globular clusters. Since the radial surface density is so much unreliable towards the very crowded parts of these globular clusters, we just ignore it in all crowded inner areas. From the catalogues in $`B`$ and $`R`$ filters, produced by SExtractor for each cluster, we construct a color $`B`$ and color index $`BV`$ catalogue with the instrumental magnitudes. We do not calibrate our data, except in the case of NGC 5139, since we need only relative magnitudes and colors for the purpose of establishing cluster membership. The magnitude error from the photographic plate is found to be up to 0.2 mag for the faintest stars. Typically, we get, for each field, a total number of stars from $`7\times 10^4`$ up to $`2\times 10^6`$ for the richest fields. We do not apply any crowding correction to our stellar counts, first, because crowding is nearly constant and weak in the outer areas (the only ones we consider) surrounding of the clusters (see also Grillmair et al. 1995), and, second, because crowding is completely dominated by the observational biases for the overdensities.
### 3.2 Star/Galaxy separation
In case of low fore- and background densities towards a considered globular cluster, i.e., for GCs located at high Galactic latitudes, we perform a star/galaxy separation by using the method of star/galaxy magnitude vs. log(star/galaxy area) - which is shown to work well down to the 18 instrumental magnitude, as displayed in Fig. 1. The galaxies have a lower surface brightness than the stars and in the magnitude vs. log(area) plane, the two classes of objects follow different loci as clearly shown on Fig. 1. The star sequence is fitted by a fifth or sixth order polynomial and the objects more than 3 $`\sigma `$ fainter than the fit are considered as galaxies, through an iterative process which stops as soon as no new galaxies are detected (typically after about 10 iterations).
We detect some background clusters of galaxies from the overdensities at high spatial frequencies, using the Wavelet Transform (WT, cf. Slezak et al. 1994 and hereafter). As a check, the galaxy catalogue obtained with this method is shown to correlate very well, with the Abell cluster catalog (Abell et al. 1989). We present in Fig. 2 the detection of clusters of galaxies in a wide field observed with the plate SRC678: the use of the WT allows a clear detection of about 50 clusters and substructures of clusters of galaxies at the typical scale of these structures. We point out that we refer to Abell cluster without any distinction between north and south Abell clusters. We emphasize that the star catalogues produced from such galaxy/star separation is polluted by galaxies at faint magnitudes, as visible in Fig. 2 because of confusion at faint magnitude for the detection. Our problem being the genuine detection of star overdensities, we look at the correlation between the galaxy clusters detected and the star overdensities, in order to disentangle the confusion, assuming that the faint magnitude galaxies follow the distribution of the bright galaxies detected. This assumption is justified by comparing some fields heavily polluted by galaxies with the so-called stellar overdensities (see, e.g., Fig. 25 hereafter). In order to remove the bad classification on the globular cluster region because of the crowding, we set the surface density of this area ($`r<r_t`$) to zero.
### 3.3 Star selection
Following the method of Grillmair et al. (1995) we perform, for each considered globular cluster, a star selection from the color-magnitude diagram (CMD) in which cluster stars and field stars exhibit different colors. In this way we can differentiate present and past cluster members from the fore- and background field stars by identifying in the CMD the area occupied by cluster stars. The envelope of this area is empirically chosen so as to optimize the ratio of cluster stars to field stars in the relatively sparsely populated outer regions of each cluster.
This search for a mask is done by subdividing the color-magnitude plane into a 50 $`\times `$ 50 array in which individual sub-areas are about 0.1 mag wide in color and 0.15 mag high in $`B`$ mag depending slightly on the cluster. Assuming that the color-magnitude distribution of the field stars is constant across the plate, a color-magnitude sequence for each cluster can be estimated from:
$$f_{cl}(i,j)=n_{cl}(i,j)gn_f(i,j)$$
(1)
in the notation of Grillmair et al. (1995) where $`n_{cl}`$ and $`n_f`$ refer to the number of stars with color index i and the instrumental magnitude index j counted within the central region of the cluster and in an annulus outside the cluster, respectively. The factor g is the ratio of the area of the cluster annulus to the field annulus. We compute the “signal-to-noise” ratio for each color-magnitude sub-area:
$$s(i,j)\frac{f_{cl}(i,j)}{\sqrt{n_{cl}(i,j)+g^2n_f(i,j).}}$$
(2)
Given the magnitude resolution of about 0.2 mag, for both plates and films, we perform a wavelet transform (WT, cf. Slezak et al. 1994) of the s(i,j) function on 5 planes, with $`W(l,i,j)=WT(s(i,j))`$. We remove the planes 0 and 1 to obtain the S/N function
$$\stackrel{~}{s}(i,j)=\underset{l=2}{\overset{5}{}}W(l,i,j)$$
(3)
with a magnitude resolution of at most 0.2 mag, see below for more details on the WT.
Grillmair et al. (1995) estimated a CMD mask for the selection of stars by computing a cumulative function from the S/N function for each cluster. This cumulative function reaches a maximum for a sub-area of the color-magnitude plane. Then by selecting all sub-areas with S/N values higher than this maximum, it is possible to construct the mask. In the present case, we depart slightly from Grillmair et al. (1995) procedure by selecting, for some fields, a subset of the mask, with a higher S/N value relative to the background stars (see Fig. 3 and Fig. 4). It must be a compromise between the S/N ratio and the number of stars selected, in order to get a sufficient spatial resolution which is lowered by Poissonian noise for small star counts. Given the S/N chosen, we have been able to eliminate from 50 % up to 99 % of the field stars. We show in Fig. 4 the CMD selected for the less contaminated fields.
On the CMD-selected star-count map $`M(x,y)`$, we fit a background map $`Z(x,y)`$, following Grillmair et al. (1995), by masking the GC (1 to 2 $`r_t`$) and using a blanking value inside, equal to the mean between 1.5 and 2.5 $`r_t`$ to get a smooth background. We fit, on a $`128\times 128`$ binned grid, a low-order bivariate polynomial surface $`Z(x,y)`$, mainly first- or second-order surface, to avoid to erase some local variation:
$`Z(x,y)`$ $`=`$ $`{\displaystyle \underset{i,j}{}}a_{ij}x^iy^j\text{with}0i,j2.`$ (4)
We subtract this background from the CMD-selected map to get a surface-density map $`T_r(x,y)`$ of the overdensities that we can attribute to the tidal extension of the GC:
$`T_r(x,y)`$ $`=`$ $`M(x,y)Z(x,y)`$ (5)
after having analyzed the potential observational biases that could create the fluctuations in the star-counting analysis.
### 3.4 Wavelet analysis
The wavelet transform is a powerful signal processing technique which provides a decomposition of the signal into elementary local contribution labeled by a scale parameter (Grossman & Morlet 1985). They are the scalar products with a family of shifted and dilated functions of constant shape called wavelets. The data are unfolded in a space-scale representation which is invariant with respect to dilation of the signal. Such an analysis is particularly suited to study signals which exhibit space-scale discontinuities and/or hierarchical features. Its ability to detect structures at particular scales has already been used in several astrophysical problems (Gill & Henriksen 1990, Slezak et al. 1994, Cambrézy, L. 1999, Chereul et al. 1999)
#### 3.4.1 “A trous” algorithm
We perform on the raw tidal map $`T_r(x,y)`$ a wavelet analysis using the “à trous” algorithm (see Bijaoui 1991). It allows to get a discrete wavelet decomposition within a reasonable CPU time. The kernel function $`B_s(x,y)`$ for the convolution is a $`B_3`$ spline function. The wavelet decomposition $`W(i,x,y)`$ is obtained from the following steps:
$`c_o(x,y)`$ $`=`$ $`image(x,y),`$ (6)
$`c_i(x,y)`$ $`=`$ $`c_{i1}B_s({\displaystyle \frac{x}{2^i}},{\displaystyle \frac{y}{2^i}}),`$ (7)
$`W(i,x,y)`$ $`=`$ $`c_i(x,y)c_{i1}(x,y).`$ (8)
The last plane, called Last Smoothed Plane (LSP), is the residuals of the last convolution and not a wavelet plane, but afterwards we will abusively speak of wavelet plane for all these planes. Each plane $`W(i,x,y)`$ represents the details of the image at the scale i. We divide each image in $`128\times 128`$ bins, a process which changes the spatial resolution of each cluster according to the different sizes of the fields: typically we get star-count maps of 3′ to 16′ resolution. The spatial resolution $`\sigma _\theta `$ for each wavelet-rebuilt cluster field can be computed, in arcmin, from the following relation: $`\sigma _\theta =0.0538\theta _{field}`$, with $`\theta _{field}`$ being the total field size in arcmin. It is so far possible to do a filtering of each plane to get only the relevant wavelet component. One problem is to find the noise level for each plane. We know that the raw tidal map is blurred by the Poissonian noise of the background objects; this is especially true at low galactic latitudes. We could perform an Anscombe transformation (Murtagh et al. 1995) to transform the Poissonian noise into a Gaussian noise on each scale. Actually we choose to perform Monte-Carlo simulations, because of the varying Poissonian noise through the field, in order to follow easily the spatial variation of the rms noise at each scale. The contours for the surface density are computed to be above 3 $`\sigma `$ level, with $`\sigma `$ being the rms fluctuation of the selected wavelet coefficients computed in an area avoiding the central cluster.
#### 3.4.2 Filtering of high varying density background noise
In case of strong gradient density of the galactic background, the noise is varying with the location in the field. To filter properly this noise, we perform N Poissonian simulations from the fitted background star counts $`Z(x,y)`$ and we take the WT of the N realizations and perform statistics on each pixel for the whole set of wavelet planes,
$`W_n(i,x,y)`$ $`=`$ $`WT(Poisson(Z(x,y)))\text{with}n=1,N.`$
Practically we have taken N=100. Then we fit by a low-order bivariate polynomial surface a rms noise map ($`\sigma _{bck}(i,x,y)`$), for each wavelet plane, to obtain an estimate of the rms fluctuation on the N realizations. This allows a good estimate of the rms noise without the need of performing a great number of simulations, which are CPU time-consuming because of the WT:
$$\sigma _{bck}(i,x,y)=\underset{k,l}{}a_{k,l}^{(i)}x^ky^l\text{with}0k,l2$$
We filter each component of the raw map $`W(i,x,y)`$, above a given threshold $`\alpha `$, using the rms noise $`\sigma _{bck}(i,x,y)`$ map on each wavelet scale $`i`$ to get the filtered wavelet planes $`W_f(i,x,y)`$ of our image:
$`W_f(i,x,y)`$ $`=`$ $`W(i,x,y)\text{ If }\left|W(i,x,y)\right|>\alpha \sigma _{bck}(i,x,y)`$
$`=`$ $`0\text{ Otherwise}`$
In this study the coefficients are filtered at the 3-sigma level. We show the case of the globular cluster NGC 5139, located at low galactic latitude, for which we present the raw surface density map and the filtered map at different resolution (see Fig. 5 and Fig. 6). We remind that a wavelet plane $`i`$ has a typical resolution of $`0.86\times 2^i`$ pixels, which is the Gaussian-equivalent resolution of the wavelet function at the scale $`i`$. We have to point out that our filtered solution is not the optimum solution (cf. Starck et al. 1997, e.g. for 1-D optimum solution) since it is only a selection of significant coefficients from the raw map. The “à trous” algorithm implemented permits a WT transform with a rate of about 2 kPixel/s/plane on a DecAlpha500 workstation <sup>1</sup><sup>1</sup>1The IDL and C procedures for the “A trous” algorithm are available at ftp://smart.asiaa.sinica.edu.tw/pub/sleon/wavelet.tar.gz or under request to sleon@asiaa.sinica.edu.tw
The final tidal map $`T_f(x,y)`$ is built as follows:
$$T_f(x,y)=\underset{i}{}W_f(i,x,y)\text{with}i_1ii_2$$
(9)
where $`W_f`$ is the filtered WT in case of strong background noise. The lower and upper indexes $`i_1`$ and $`i_2`$ constrain the resolution of the rebuilt map by filtering the higher and/or the lower space scale wavelet planes. For the adopted binning, we find that the map with the planes 3 to 7 (LSP) gives the best compromise between the spatial resolution and the Poissonian noise of the star-counting after the filtering of the background noise. It provides, in most cases, a higher spatial resolution than in Grillmair et al. (1995), because the wavelet decomposition extract the energy only at the useful scales. We have to point out nevertheless that the wavelet basis used here is not orthogonal, mixing up slightly the scale energy on different planes, but this does not affect our rebuilt map.
## 4 Results
In this section we present all results related to the observations of stars surroundings of each GC. We discuss individually each cluster for the particular observational biases which could affect its results. Grillmair et al. (1995) found that the clusters in their sample with obvious tidal extensions showed a break in their surface density profiles, becoming pure power law at large radii. We try to link in a systematic way the shape of each observed tidal tail to the orbital phase of the corresponding cluster. For this we define $`Q_a^b`$ as the slope of the radial surface density between $`a\times r_t`$ and $`b\times r_t`$. The slope of the radial surface density is computed when the tidal tails are not dominated by the noise which would lead to a flat slope. In practice, we choose to compute the three slopes $`Q_1^3`$, $`Q_3^6`$, and $`Q_1^6`$, and give in Table 3 their values only for clusters where the signal/noise ratio for these azimuth averaged parameters is sufficiently high. Practically, this means that we remove all clusters with $`Q_1^3`$ shallower than -0.5. Whenever possible, our surface density profiles <sup>2</sup><sup>2</sup>2The 1-d and 2-d star counts are available by request to the authors. are extended inwards with the surface-brightness profiles from Trager et al. (1995), assuming a linear relation between light emission and stellar surface density through the globular cluster. Crowding and saturation problems in our plates/films make the inner parts of our density profiles highly unreliable. Consequently, the adjustment between Trager et al.’s profiles and ours is done, in the short radius range where both profiles are reliable, by adjusting a constant $`K`$ in the following way:
$$\mathrm{log}(\text{surface density})=\mu /2.5+K$$
(10)
where $`\mu `$ is the fitted surface brightness at r (see Trager et al., 1995). For Trager et al. (1995) profiles, only data outside the radius r=1′ are shown. We point out that differences between the two profiles in the very outer parts can partly be explained by mass segregation in the cluster, unveiled by different limiting magnitudes.
It is worth mentioning that the measured slope will be flattened at small radii since the closer to the cluster the larger the crowding and since a azimuthal averaged value is more sensitive to noise at large radii. For a power law dependence, with a slope $`\alpha `$, of the tidal tail surface density, the tail/noise surface density ratio scales as $`r^\alpha /\sigma _{bck}`$, where $`\sigma _{bck}`$ is the background surface density. We choose the quantity $`Q_1^3`$ — the only one we are able to determine in a large enough number of GCs — as a quantitative estimator for comparing the outer structures of these clusters. Table 4 gives, for all GCs in our sample, the dynamical and structural parameters (from GO97 and references therein) which are representative of the internal and external dynamical evolution of these globular clusters; $`t_{rh}`$ is the relaxation time at the half-mass radius; $`\nu _{evap}`$ is the destruction rate due to evaporation; the ratio $`\nu _{tot}/\nu _{evap}`$ of the total destruction rate to the destruction rate due to evaporation illustrates the importance of the galaxy-driven evolution suffered by the clusters; $`c=\mathrm{log}(r_t/r_c)`$ is the concentration of the cluster, where $`r_c`$ and $`r_t`$ are the core and tidal radii, respectively; M is the cluster mass and $`V_{HB}`$ is the $`V`$ magnitude of the horizontal branch stars. Given the low surface density and low S/N of some cluster tidal tails, the radial surface density profile is not shown for all the clusters.
### 4.1 NGC 104 $``$ 47 Tucanae
NGC 104 is at a distance of 4.1 kpc from the sun, with its horizontal branch (HB) at $`V`$ = 14.06 mag. It has a tidal radius of about 55 pc with a rather high concentration $`c`$ = log ($`r_t/r_c`$) = 2.04 (see Table 4). It is one of the most massive and nearby GCs (Meylan & Mayor 1986, Meylan 1989). In their study, Odenkirchen et al. (1997) estimate that NGC 104 has experienced a disk crossing about 60 Myr ago and suffers frequent disk-crossing events with a period of about 160 Myr between each passages. The rotation of this cluster (Meylan & Mayor 1986) should enhance its mass-loss rate by a factor of about 1.5 (Longaretti & Lagoute 1996). Unfortunately the detection of extra-tidal material is made difficult by the strong pollution along the line of sight due to the Small Magellanic Cloud (SMC) background stars: in the CMD of NGC 104, the cluster sequences are superposed, with a vertical translation, with the sequences drawn by the stellar populations in the SMC. In order to discriminate efficiently between cluster and SMC stars, we compute the tail/background S/N function relative to the SMC by taking the background field for the s(i,j) function on the SMC (see Fig. 4. We fit a 3$`\times `$3 bivariate surface in order to reproduce the SMC overdensity. A strong residual of the SMC is still present (see Fig. 7), fortunately located in the S-E corner. It represents a high surface density of stars of the SMC, well correlated with the dust emission seen in the IRAS 100-$`\mu m`$ map. About the globular cluster itself, tidal extensions towards the N-W and S-W directions are marginally present around NGC 104. The IRAS 100-$`\mu m`$ map does not exhibit any anticorrelation with these tidal extensions. No fit of the surface density profile of these tidal tails has been performed because of their poor statistical significance. A steep mass function must be present in the tidal tails because of the mass segregation observed in this cluster (Anderson & King, 1996)
### 4.2 NGC 288
NGC 288 is at a distance of 8.1 kpc from the sun, with its horizontal branch (HB) at $`V`$ = 15.38 mag. It has a tidal radius of about 32 pc (see Table 4). Its concentration is low, with $`c`$ = log ($`r_t/r_c`$) = 0.96. It is located close to the South Galactic Pole, at 8 kpc from the Sun (Harris 1996), with a retrograde orbit (Dinescu et al. 1997). From GO97, NGC 288 is a cluster with a dynamical evolution strongly driven by the galactic tidal field (see Table 4). NGC 288 was already observed by Grillmair et al. (1995), who found tidal extensions on a field $`200\mathrm{}\times 200\mathrm{}`$ smaller than ours, but with the same spatial resolution (16′). In Fig. 8, the wavelet decomposition clearly reveals some wide structures missed by Grillmair et al. (1995), especially towards the south. (The arrows indicating the direction of the Galactic center in Grillmair et al.(1995) for NGC 288, NGC 362, and NGC 1904 are in error \- Grillmair, private communication). No dust emission from the IRAS 100-$`\mu m`$ survey is detected. NGC 288 is nearly free of observational biases, apart from some galaxy clusters. A few such clusters of galaxies are clearly detected (see Fig. 8). We suggest that the tidal radius determination could be overestimated because of the presence of the clusters Abell 118 and 122, as already pointed out by Scholz et al. (1998) for Pal 5.
Tidal tails are well separated in two directions: first, towards the galactic center (dashed arrow), second, aligned with the orbit of the GC (dotted arrow). NGC 288 has recently undergone a gravitational shock (Odenkirchen 1998). It is very likely that the tidal tails visible in Fig. 8a are the results, in projection on the plane sky, of both the disk shocking and the relics of the bulge shocking from the last passage close to the bulge (Dauphole et al. 1996). NGC 288 exhibits very important tidal tails, extending up to 350 pc from the cluster: this has to be related to its strong interaction with the Galaxy, as found by GO97 (see Table 4). We count about 1200 stars outside the tidal radius of the cluster but we did not attempt an estimate of the mass in the outer parts of the cluster because of the poor photometry. Further CCD observations with deep and precise photometry should provide very accurate mass loss rates. From its orbital motion, this cluster appears to be a very good candidate for tracing the local galactic potential (disk scale height and surface density).
### 4.3 NGC 1261
NGC 1261 is a remote cluster at a distance of 15.1 kpc from the sun, with its horizontal branch (HB) at $`V`$ = 16.70 mag. It has a tidal radius of about 34 pc and a concentration $`c`$ = log ($`r_t/r_c`$) = 1.27. Its evolution is probably driven by its internal dynamics (Zoccalli et al. 1998, GO97). Its field is not polluted by strong dust extinction (E($`BV`$) = 0.02), but the main bias is coming from the extra-galactic object overdensities. Although no Abell cluster is present in the field, we detect the presence of galaxy clusters which are strongly correlated with some stellar extensions as visible in Fig. 9. We can conclude here that the N-E extension of the extra-tidal material, which is aligned with the direction of the galactic center (dashed arrow), is a real tidal feature of the GC, because there is no strong galaxy cluster at this location. The slope $`Q_3^6`$ (see Table 3) is probably highly contaminated by back- and foreground stars and not useful. Zoccali et al. (1998) find evidence for mass segregation in the cluster ($`x_{obs}=1.7\pm 0.5,r>4.4r_c`$), segregation which should affect the tidal tails as discussed in Section 5.
### 4.4 NGC 1851
NGC 1851 is a remote cluster at a distance of 11.7 kpc from the sun, with its horizontal branch (HB) at $`V`$ = 16.15 mag. It has a tidal radius of about 49 pc and a very high concentration $`c`$ = log ($`r_t/r_c`$) = 2.24. The western part and the S-W part of NGC 1851 extension are contaminated by galaxy clusters (Abell 514 and anonymous) and by a bright star also observed by the IRAS 100-$`\mu m`$ map (see Fig. 10). The extinction is not important towards NGC 1851, with E($`BV`$) = 0.02. Stars unbound from the cluster are likely tracing the orbital path, here these tails seem to have a preferential direction towards the galactic center (dashed arrow and S-E extension). The cluster position indicates that it is not suffering a strong shock, as confirmed by the ratio $`\nu _{tot}/\nu _{evap}=1.0`$ from GO97, which indicates that the evolution of this cluster is mainly internally driven. Consequently, the surface density profile in the outer parts of the cluster is mainly shaped by evaporation and tidal stripping at its location in the Galaxy. Saviane et al. (1998) found a slight mass segregation in this cluster which affects the tidal tail detection by lowering the mean mass of the unbound stars (Section 5).
### 4.5 NGC 1904 $``$ M79
NGC 1904 is a remote cluster located at a distance of 12.2 kpc from the sun, with its horizontal branch (HB) at $`V`$ = 16.15 mag. It has a tidal radius of about 32 pc and a concentration $`c`$ = log ($`r_t/r_c`$) = 1.72. NGC 1904 is surrounded by a halo of unbound stars (see Fig. 11), as previously seen by Grillmair et al. (1995), on a wider field but with a lower spatial resolution (them with 16′, us with 6.5′) which blurred all the small structures we observe around the cluster. We do not find evidence for a large southern tidal extension as observed by Grillmair et al. (1995). The difference here could be accounted to the lower resolution used by them, one part of this large tail could be due to the southern galaxy clusters not well separated. We point out that in their and our work we select stars below the completeness limit ($`19`$ mag) , completeness fluctuation are another possibility to explain some differences, but not on such a large scale. The tail is oriented in the direction of the galactic center (dashed arrow). As in the case of NGC 288, the tidal radius determination may be overestimated because of the presence of galaxy clusters close to NGC 1904. Nevertheless, the tidal tails of this cluster do not appear to be correlated with the distribution of the extra-galactic objects. The dust extinction is low towards this cluster (E($`BV`$) = 0.01) and the fluctuations of the dust emission are low as traced by the IRAS 100-$`\mu m`$ map. Because of the short relaxation time of NGC 1904 ($`t_{rh}=8.8\times 10^8`$ yr), the mass segregation should affect as well the stellar populations in the tidal tails. Since, following GO97, $`\nu _{tot}`$ is about 30 % higher than $`\nu _{evap}`$, this may indicate a slight influence of the galaxy on this cluster.
### 4.6 NGC 2298
NGC 2298 is a remote cluster located at a distance of 10.4 kpc from the sun, with its horizontal branch (HB) at $`V`$ = 16.11 mag. It has a tidal radius of about 19 pc and a concentration $`c`$ = log ($`r_t/r_c`$) = 1.40. There are background fluctuations owing to the dust along the line of sight, as clearly traced by the IRAS 100-$`\mu m`$ map (see Fig. 12). We perform a quite high tail/background S/N CMD selection because of the high background density (see Fig. 4, but there is still a bias because of the dust extinction, as seen in Fig. 12. There is a southern extension towards the galactic center (dashed arrow) which is interrupted by dust absorption. Some parts of the Eastern extension located at (x = –50, y = –10) of the tidal tails may be questionable, because of the stronger dust presence, nevertheless the lower absorption can hardly explain all these overdensities, since their distribution does not follow the minimum IRAS 100-$`\mu m`$ emission map. Clearly, the overdensity at (x = 60, y = 60) is associated with a low IRAS emission area. Given the position and the distance of NGC 2298 from the galactic center (15.1 kpc), the southern extension is likely tracing its orbital path and not the result of gravitational shock, as indicated by the ratio $`\nu _{tot}/\nu _{evap}=1.0`$ from GO97. The low value $`Q_3^6`$ = –0.25, likely due to the small extensions in the outer parts, is questionable.
### 4.7 NGC 4372
NGC 4372 is a nearby globular cluster located at a distance of 4.6 kpc from the sun, with its horizontal branch (HB) at $`V`$ = 15.30 mag. It has a tidal radius of about 52 pc and a concentration $`c`$ = log ($`r_t/r_c`$) = 1.30. The presentation of the detection of the overdensities around this cluster illustrates the dramatic influence of varying dust extinction (see Fig. 13). Strangely enough, the very elongated dust filament observed in the IRAS 100-$`\mu m`$ map ends very close to the cluster: this may suggest an interaction of the cluster with the interstellar medium currently at play. Following GO97 (see Table 4), this cluster has an evolution strongly driven by the galaxy ($`\nu _{tot}/\nu _{evap}=3.8`$).
### 4.8 NGC 5139, $`\omega `$ Cen
NGC 5139, the most massive galactic globular cluster (Meylan 1987, Meylan et al. 1995, Merritt et al. 1997), currently crossing the disk plane, is a nearby globular cluster located at a distance of 5.0 kpc from the sun, with its horizontal branch (HB) at $`V`$ = 14.53 mag. It has a tidal radius of about 65 pc and a concentration $`c`$ = log ($`r_t/r_c`$) = 1.24. Its relative proximity allows to reach the main sequence for the star count selection.
Given the very good tail/background S/N ratio, we perform an absolute calibration of the photometry using the data from Cannon & Stobie (1973) and Alcaino & Liller (1987) with an error which is still about $`\sigma `$ = 0.2 mag. Although obvious biases by dust absorption affect the star counts, as seen, e.g., at the positions (x = 50, y = –25) and (x = –50, y = –70) on the IRAS 100-$`\mu m`$ map (Fig. 14), there are two large and significant tidal tails: NGC 5139 is releasing currently some large amounts of stars. The tidal tail extensions are perpendicular to the galactic plane (see Fig. 14), which is a clear sign of disk-shocking, as observed in our numerical simulations (CLM99). By considering star-counts with magnitude $`B<19`$ (the completeness limit), we found about 7000 $`\pm 600`$ stars outside one tidal radius, in the $`4^{}\times 4^{}`$ field. This magnitude corresponds to a 0.63 M star at a distance of 5 kpc. Assuming the same mass function in the cluster and in the tidal extensions, because of its large relaxation time, we estimate a total of 1.9 $`10^4`$ M for the escaped stars, with the assumption of a Salpeter law ($`\alpha `$ = –2.35) mass function for the stars down to 0.1 M. Thus the tidal tails represent about 0.6 % of the cluster mass for total cluster mass of about 5.1 10<sup>6</sup> M. This is consistent with the numerical simulations (CLM99, Johnston et al. 1998) given the high uncertainties on the mass function, the photometric calibration, the mass-luminosity relation used (see e.g., Saviane et al. 1998), and the possible steeper mass function in the tidal tails, as discussed in Section 5 for a slope $`\alpha `$= –2.8. We point out that a steeper mass function has been observed in the halo of NGC 5139 (Anderson, 1998)
The $`Q_1^2`$ parameter value and the position of the cluster in the galaxy indicate that NGC 5139 is presently experiencing a disk shocking, with an important mass loss of stars, whose presence is clearly observed in the immediate neighborhood of the cluster. The observed proper motion of NGC 5139 indicates that this cluster is in the early phases of its disk crossing. This confirms the high value of the ratio $`\nu _{tot}/\nu _{evap}=9.4`$ estimated by GO97. In the case of NGC 5139, the disk-shocking consequences are combined with the bulge-shocking ones, since the cluster orbit goes as close as 1.8 kpc from the galactic center (Dauphole et al. 1996).
We choose to present here the same wavelet planes that those for the other clusters, but given the high density — significance – of NGC 5139 tidal tails, we illustrate, in Fig. 6, the different spatial resolutions for $`\omega `$ Cen after filtering of the background noise. It is worth mentioning that, because of the internal rotation of this cluster (Meylan & Mayor 1986, Merritt, Meylan & Mayor 1997), the global mass loss rate is enhanced by a factor of 2 with respect to the N-body simulations (CLM99) and Fokker-Planck estimates (Longaretti & Lagoute 1996). In the discussion we consider the effect of the mass segregation on the mass loss derivation.
### 4.9 NGC 5272 $``$ M3
NGC 5272 is a globular cluster located at a distance of 9.7 kpc from the sun, with its horizontal branch (HB) at $`V`$ = 15.65 mag. It has a tidal radius of about 105 pc and a concentration $`c`$ = log ($`r_t/r_c`$) = 1.85. The cluster is near the edge of the plate, preventing the study of its Eastern side (see Fig. 15). The field is polluted only by 2 small galaxy clusters, viz. Abell 1781 and Abell 1769, the former being detected only at 2.5-$`\sigma `$ level. Unfortunately, a defect on the plate E131 (POSS) have blurred the extra-galactic object detection (peak at x = 120, y = –20). We emphasize that point-source detection with SExtractor is less affected by this defect. There is no anticorrelation at all between the tidal tails and the dust emission, as we checked with the IRAS 100-$`\mu m`$ map, which is at a low level (E($`BV`$) = 0.01). The extension at (x = –30, y = –50), towards the galactic center (dashed arrow), is the more reliable. Thus from the low value of the slope $`Q_1^3`$ = –0.35, we can infer that the field pollution bias must be quite strong, providing a rather constant radial surface density. The comparison with the data from Trager et al. (1995), which obtained star-count values smaller than our data near the tidal radius, confirms this point. The long relaxation time of NGC 5272, viz. $`t_{rh}=7.3\times 10^9`$ yr, implies that the mass segregation should not affect strongly the mass function of the unbound stars. Gunn & Griffin (1979) found some weak rotation in this globular cluster which should slightly enhance the mass loss rate by a factor 1.1-1.2 (Longaretti & Lagoute 1996). There is no apparent correlation between the tidal tail direction and the proper motion of the cluster (dotted arrow).
### 4.10 NGC 5694
NGC 5694 is a very remote globular cluster located at a distance of 33 kpc from the sun, with its horizontal branch (HB) at $`V`$ = 18.50 mag, which is a strong limitation for star counts. Given its large distance from the galactic center, namely 27.5 kpc, this cluster is not expected to suffer strong gravitational shocks ($`\nu _{tot}/\nu _{evap}`$ = 1.0, GO97). It has a tidal radius of about 41 pc and a concentration $`c`$ = log ($`r_t/r_c`$) = 1.84. We select the stars on the giant branch only, with a higher tail/background S/N ratio in order to avoid as much as possible the galaxies which are the strongest bias in this field (see Fig. 16). The lower dust extinction, mapped through IRAS 100-$`\mu m`$ emission, could induce an artificial extension in the S-W part of the cluster, at the position (x = 20, y = –15). But the huge extension in the S-E part can be attributed to extra-tidal material, with high confidence since it is correlated with higher dust extinction and there is only one small galaxy cluster at the position (x = –20, y = –3). It must be stars tidally stripped from the cluster, material which is now trailing/leading along the orbit of the cluster. As in the other clusters, it is aligned towards the galactic center direction (dashed arrow), but it might also be a projection effect of its orbital plane with the galactic center direction. The size of this extension is about 300 pc in the sky and is probably even much larger because of the shallow photometry available on this distant cluster.
### 4.11 NGC 5824
NGC 5824 is a very remote globular cluster located at a distance of 32.2 kpc from the sun, with its horizontal branch (HB) at $`V`$ = 18.60 mag, which is a strong limitation for star counts. At a large distance from the galactic center, namely 26 kpc, this cluster is not expected to suffer strong gravitational shocks ($`\nu _{tot}/\nu _{evap}`$ = 1.6, GO97). It has a tidal radius of about 147 pc and a concentration $`c`$ = log ($`r_t/r_c`$) = 2.45. Because of a low tail/background S/N ratio, a consequence of the faint $`V_{HB}`$ magnitude, the overdensity map around NGC 5824 appears to be very noisy (Fig. 17). Grillmair et al. (1995) find around this cluster more extended structures, aligned with the N-S direction, than we do in the same field: this may be partly due to our rather shallow photographic films. There are some strong biases due to dust extinction as it can be seen on Fig. 17 at the position (x = –20, y = –20) and due also to some galaxies spread mainly over the Southern part. GO97 indicate that NGC 5824 should experience important interactions with the tidal galactic field, a prediction we are not able to confirm because of the tangled observational biases. But it appears that a preferential direction of the cluster extension could be perpendicular to the galactic plane (solid arrow), either due to a disk shocking or tracing the orbital motion of the cluster. Nevertheless a bulge shocking effect cannot be ruled out in the case of a very eccentric orbit.
### 4.12 NGC 5904 $``$ M5
NGC 5904 is a globular cluster located at a distance of 7 kpc from the sun, with its horizontal branch (HB) at $`V`$ = 15.06 mag. It has a tidal radius of about 66 pc and a concentration $`c`$ = log ($`r_t/r_c`$) = 1.87. We present only the S-E part of the tidal extensions (Fig. 18) because of its position on the plate. No bias due to dust is reported towards this field (E($`BV`$) = 0.03). The presence on a galaxy cluster close to the tidal radius of the cluster enhances artificially and locally the tidal tail pointing towards the galactic center (dashed arrow) and the direction perpendicular to the galactic plane (solid arrow). Nevertheless, it is obvious that an extension is present towards this direction, since the galaxy cluster size is significantly smaller than the size of the globular cluster extension, as shown on Fig. 18b with the same resolution used for the star and galaxy surface densities. Lehmann & Scholz (1996) found already indication of tidal tail around this cluster from its surface brightness profile which departs from a King profile; this may be explained as well by the galaxy cluster near the tidal radius. From Odenkirchen et al. (1997), NGC 5904 is just beginning its crossing through the disk and towards the galactic center. Consequently, we could observe the first effect of the gravitational shocking on this cluster, with the tail aligned towards the tidal directions (see CLM99) after being compressed during the crossing. Indeed the momentum transfer to the cluster stars is in the Z direction during the disk shocking. From GO97, NGC 5904 suffers strong interactions with the galaxy, with $`\nu _{tot}/\nu _{evap}`$ = 26, a high value due to the use of the Bahcall et al. (1983) galactic model which enhances the gravitational shocks because of its nuclear component and the form of the disk potential which does not vanishes at the center as it is the case for the model from Ostriker & Caldwell (1983).
### 4.13 NGC 6205 $``$ M13
NGC 6205 is a globular cluster located at a distance of 6.8 kpc from the sun, with its horizontal branch (HB) at $`V`$ = 14.90 mag. It has a tidal radius of about 56 pc and a concentration $`c`$ = log ($`r_t/r_c`$) = 1.49. The bias towards NGC 6205 are not strong as shown by the weak IRAS 100-$`\mu m`$ flux and the relatively high tail/background S/N ratio in the CMD. Given the position of the cluster on the survey plates, we extract a field of 90′ in size. There is no strong bulk of tidal stars (see Fig. 19) due to any shock and the field is too small to detect any large scale structure corresponding to the orbital path. An extension can be seen towards the galactic center (dashed arrow) at the position (x = –10, y = –25), although located inside the tidal radius, which highlights the limitation of an azimuthally averaged radial surface density. This extension is not correlated with the proper motion. We note that an extended default on the plate center worsens the cluster/background star separation.
### 4.14 NGC 6254 $``$ M10
NGC 6254 is a nearby globular cluster located at a distance of 4.1 kpc from the sun, with its horizontal branch (HB) at $`V`$ = 14.65 mag. It has a tidal radius of about 26 pc and a concentration $`c`$ = log ($`r_t/r_c`$) = 1.40. This cluster is a striking case because of a strong gradient in the dust extinction as seen on Fig. 20 with the IRAS 100-$`\mu m`$ map. The southern extension anticorrelates quite well with the dust emission which is the sign of a possible bias. An obvious decrease of the stellar surface density is correlated with the dust emission at the position (x = –40, y = –30). Nevertheless the inner NE-SW and the northern extensions are not anticorrelated with the dust emission. The second break at $`\mathrm{log}(r)1.6`$, apart from the one around the tidal radius, in the radial density profile (see Fig. 20c) must correspond to a very recent disk-shocking, with the diffusing stars (cf. CLM99) still close to the cluster. This conclusion is strengthenedd by the proper motion of the cluster whose direction (dotted arrow) is opposite to the disk direction (solid arrow): Odenkirchen et al. (1998 ) indicate that NGC 6254 suffered its last disk crossing about 20 Myr ago. Considering the northern extension as a genuine tidal tail made of stars from NGC 6254, we can give a lower limit for the diffusion velocity, which is a projected expansion velocity of the tidal material in the cluster reference frame: at a distance of 4.1 kpc, for a projected distance of 150 pc, we obtain about 7 km s<sup>-1</sup> as a lower limit of the diffusion velocity. We note that the velocity dispersion of stars in NGC 6254 is similar, with $`\sigma _0`$ = 6.6 km s<sup>-1</sup> (Pryor & Meylan 1993). This velocity diffusion probes the differential diffusion of stars released in the Galaxy along with the global dynamical friction of the cluster which is not felt by the unbound stars. Actually this diffusion velocity is surprisingly high compared to the dispersion velocity, where we would expect low velocity dispersion for the unbound stars: a misclassification of these clumps as genuine cluster stars or an underestimation of the last crossing time cannot be ruled out. Given the quite short relaxation time $`t_{rh}=7.6\times 10^8`$ yr, the mass segregation must be present in this cluster, even though Hurley et al. (1989) did not find any evidence. Such a mass segregation must lead to a steep mass function in the tidal tails.
### 4.15 NGC 6397
NGC 6397 is a very nearby globular cluster located at a distance of 2.2 kpc from the sun, with its horizontal branch (HB) at $`V`$ = 12.87 mag. It has a tidal radius of about 66 pc and a concentration $`c`$ = log ($`r_t/r_c`$) = 2.50. It is the only post core-collapsed cluster in our sample, although NGC 1851 and NGC 5824 have also rather large concentrations. This is the second example, with NGC 4372, of overdensities strongly biased by dust extinction fluctuation as it can be seen in Fig. 21. All the overdensities found in the northern and eastern parts cannot be disentangled from dust extinction. Only the S-E extension, at the position (x = –100, y = –100), could be a genuine tidal tail, but with a somewhat low confidence in spite of the fact that these star counts are more than 3 $`\sigma `$ above the background because the dust extinction fluctuations in this field are quite high ($`\sigma ^2(S_{100})\overline{S_{100}}20`$ MJy/sr for the IRAS-100$`\mu m`$ flux). Nevertheless, we emphasize that this extension is perpendicular to the plane (solid arrow) as expected for disk shocking (CLM99), thanks to the momentum transfer in the Z direction. During the disk crossing the gained acceleration for the cluster stars is directed towards the cluster equatorial plane parallel to the galactic plane. Then the energy gained is released in this direction, perpendicular to the galactic plane. The mass segregation found by Mould et al. (1996) in this cluster will affect the mass function of the tidal tails. A weak rotation of NGC 6397 has been found (Meylan & Mayor 1991) which should enhance the mass loss rate by about 20 %, using Fig. 7 of Longaretti & Lagoute (1996).
### 4.16 NGC 6535
NGC 6535 is a globular cluster located at a distance of 6.6 kpc from the sun, with its horizontal branch (HB) at $`V`$ = 15.73 mag. It has a tidal radius of about 17 pc and a concentration $`c`$ = log ($`r_t/r_c`$) = 1.30. In spite of the high tail/background S/N color selection on the CMD in order to avoid the high background density, the cluster density remains lower than the background (see Fig. 22). The dust extinction induces a strong bias towards this field (E($`BV`$) = 0.32) as seen in the IRAS 100-$`\mu m`$ map on Fig. 22. Part of the northern extension could be artificially enhanced by the local lower dust extinction. It is likely that the Southern extension is lowered by local higher dust extinction at the position (x = 0, y = –20). As indicated by GO97, the evolution of this cluster is influenced by the galactic potential ($`\nu _{tot}/\nu _{evap}`$ = 1.4). Currently, NGC 6535 is experiencing a strong bulge shocking and disk shocking as indicated by the correlation of the tail with the disk/bulge direction (solid and dashed arrows, respectively) and confirmed by its location in the galaxy, viz. 1.2 kpc above the plane and 4 kpc from the galactic center (Harris 1996).
### 4.17 NGC 6809 $``$ M55
NGC 6809 is a globular cluster located at a distance of 5.1 kpc from the sun, with its horizontal branch (HB) at $`V`$ = 14.40 mag. It has a tidal radius of about 23 pc and a concentration $`c`$ = log ($`r_t/r_c`$) = 0.76. The overdensities (see Fig. 23) are strongly anticorrelated with the dust emission traced by the IRAS 100-$`\mu m`$ map, e.g. at the position (x = –50, y = –50), where the dust extinction is probably disturbing the tail surface density. It may also be possible that the Western extension (x = 90, y = 0) towards the galactic center (dashed arrow) could be associated with the cluster, because it is not anticorrelated with the dust emission; the same remark applies for the extension at (x = –30, y = –30). The study of such a cluster should greatly benefit from the better transparency to the dust absorption offered in $`J`$ and $`K`$ bands. In a previous study in V and I bands, Zaggia et al. (1997) found already evidence for cluster stars in the halo of this object.
### 4.18 NGC 7492
NGC 7492 is a remote globular cluster located at a distance of 24.3 kpc from the sun, with its horizontal branch (HB) at $`V`$ = 17.63 mag. It has a tidal radius of about 62 pc and a concentration $`c`$ = log ($`r_t/r_c`$) = 1.0. There is no dust emission towards this field and the background galaxy clusters are located far from the cluster, as indicated in Fig. 24. Obviously, the overdensity at the position (x = 18, y = 25) is associated with the galaxy cluster Abell 2533. Because of the low mass of this cluster, GO97 find a fast evolution in the Galaxy field, with $`\nu _{tot}/\nu _{evap}=77.8`$, compared to its intrinsic evolution. Clearly, a tiny extension is visible, pointing towards the galactic center (dashed arrow). This lack of tidal extension is not in contradiction with the conclusion drawn by GO97, given its current location far from the center of the Galaxy (23.5 kpc). A higher tail/background S/N ratio selection, using high-quality CCD data, may allow the detection of very low surface density extension related to tidal tails extending away from NGC 7492.
### 4.19 Palomar 5
Palomar 5 is a remote globular cluster located at a distance of 21.8 kpc from the sun, with its horizontal branch (HB) at $`V`$ = 17.63 mag. It has a tidal radius of about 107 pc and a concentration $`c`$ = log ($`r_t/r_c`$) = 0.74. It is one of the most remote cluster with measured proper motions (Schweitzer et al. 1993, Scholz et al. 1998). The tidal radius could be lower than previously measured, down to 7′, in agreement with its orbit (Scholz et al. 1998). In Fig. 25 we present the overdensities, which are strongly biased by the background galaxy clusters present in the field. Because of the unreliable star/galaxy separation above $`R18`$, the confusion is quite strong for this remote and faint cluster. As pointed out already by Scholz et al. (1998), the galaxy cluster Abell 2050 could be responsible for the previous (commonly adopted) overestimate of the tidal radius. The dust extinction is very weak in this field (E($`BV`$) = 0.03), and do not exhibit any anticorrelation with the overdensities, as checked on the IRAS 100-$`\mu m`$ map. The background galaxy distribution and the large distance to this cluster make difficult any conclusion on the genuine location, if any, of stars stripped from the cluster.
### 4.20 Palomar 12
Palomar 12 is a remote globular cluster located at a distance of 17.8 kpc from the sun, with its horizontal branch (HB) at $`V`$ = 17.13 mag. It has a tidal radius of about 49 pc and a concentration $`c`$ = log ($`r_t/r_c`$) = 0.90. The dust extinction is very low (E($`BV`$) = 0.02), but the contamination by background galaxy clusters is very important (see Fig. 26), although only two Abell galaxy clusters are reported in this field. The N-S oriented very long tail is contaminated by some galaxies as shown at position (x = 15, y = 30) in Fig. 26. Nevertheless this tail is a genuine feature made of stars tidally stripped, as shown by the distribution of the galaxies as the same resolution. The western and eastern overdensities are related mainly to galaxies. A higher tail/background S/N CMD selection confirmed the Pal 12 membership of the top and bottom clumps at positions (x = 0, y = $`\pm `$60). To get an estimate of the time of the last, if any, gravitational shock on this cluster, we assume that these two latter star clumps are remains of the last shock. Adopting a diffusion velocity for the tidal stars equal to 1 km s<sup>-1</sup>, similar to the velocity dispersion (Djorgovski & Meylan 1994) of such a low mass cluster ($`2\times 10^4\text{ M}\text{}`$) and assuming the distance in projection between the clumps and the cluster to be 350 pc, we estimate 350 Myr as the time since the last shock. This is a lower limit because of the projection effect and the limited field. Contrary to most other tidal tail directions, the extension is perpendicular to the galactic center direction (dashed arrow) and is in a plane parallel to the galactic disk.
## 5 Discussion
The detection of stars tidally stripped from globular clusters emphasizes strongly the importance of the interactions of these stellar systems with the Galaxy. In the light of the possible biases present in the above observed fields, it is possible to give an estimate of the physical status of the clusters relative to the gravitational shocks they suffer in the Galaxy (see Table 5). In the case of a cluster experiencing disk-shocking only, it will be first compressed in the direction perpendicular to the galactic plane, during the short time of the crossing; then the tidally released stars form tails perpendicular to the galactic plane (see, e.g., NGC 5139). In the case of a cluster experiencing bulge-shocking only, i.e., not too far from the Galaxy center, the tails are elongated mainly along the galactic density gradient (spherical symmetry) and one can expect a correlation between the tidal tail direction and the galactic center. This is true also for the more general case of galaxy-shocking, when bulge- and disk-shocking are both at play, i.e., when the cluster is close to the galactic center. If the cluster has not experienced for a long time a gravitational shock, its tidal tails are on a large scale oriented along its orbit. However Grillmair (1992) showed using N-body simulations, without any disk potential, that strong “bars” orthogonal to the orbital path will develop naturally near the apogalactica of the cluster’s orbit. In Table 5 we, tentatively, give the processes at play for creating the recent mass loss in the clusters: it is based on tidal tails shapes, but, as well, on their positions in the Galaxy and their orbit and proper motions, when they are available. It explains the discrepancy between some tidal tail orientation and the type of physical process. We point out that the projection effect must be important in some case (e.g NGC 288). It has to be noted that a combination of disk- and bulge-shocking are expected to confuse the above simplified scenario (e.g. NGC 6535).
The case of NGC 5904 is interesting since its proper motion is known: the small tidal extension observed is perfectly aligned towards the galactic center and the direction perpendicular to the galactic plane and not with its motion along its orbit. Given its position in the Galaxy, this cluster is probably suffering a weak disk and bulge shocking.
In Fig. 27 we show the slope values $`Q_1^3`$ (between $`r_t`$ and 3$`r_t`$) versus $`Q_3^6`$ (between 3$`r_t`$ and 6$`r_t`$) for the few clusters where it is possible to measure these two parameters. We emphasize that these slope values are probably overestimated, especially for $`Q_1^3`$, because of central crowding. As found in dynamical simulations by Johnston et al. (1998) and in other observations by Grillmair et al. (1998) on different radius ranges, the mean slope value for $`Q_1^6`$ is $`0.91\pm 0.24`$. The coefficient $`Q_3^6`$ presents a strong scatter (1.69) around its mean value equal to $`1.0`$. Its determination is difficult because the stars no more bound to the cluster have a very low density in the outer parts where the noise dominates (see, e.g., NGC 2298). The quantity $`Q_1^3`$ must be a reliable indicator of the recent mass loss from the cluster, with a steep slope for the cluster suffering shocks. Then the diffusion of the heated stars will flatten the surface density profile.
In Fig. 28, we present, from our N-body simulations (CLM99), the variation with time of the surface density slope fitted on a power law for two different ranges of radii. It is remarkable to note the strong variation of the slopes during the crossing of the galactic plane. Moreover there is a delay between the variation of the $`\alpha (30\text{}40\text{pc})`$ slope and $`\alpha (40\text{–50}\text{pc})`$ slope. Here the dumping frequency of the simulations is too low to allow any estimate of the diffusion velocity of the bulk of stars stripped during the crossing. It appears nevertheless to be lower than the velocity dispersion of the simulated globular cluster ($``$ 8 km s<sup>-1</sup>).
We may link the case of NGC 6254 (see Fig. 20) to the surface density profile computed from our N-body simulations (CLM99) before and after the crossing and displayed in Fig. 29). Clearly the second break, at a radius $`r>r_t`$, in the observed cluster density profile ($`\mathrm{log}(r)1.6`$) and simulated cluster density profile ($`\mathrm{log}(r)1.9`$) indicates that disk shocking is currently at play on NGC 6254 and the halo of unbound stars has not yet diffused outwards. Even if other mechanisms could produce such break (e.g. “bars” at the apogalactica radius) we note that this NGC 6254 is currently just 1.6 kpc above the galactic plane.
The variations of the $`Q_1^3`$ coefficient between clusters with a strong galaxy-driven evolution are expected to be important as observed in the simulations. Nevertheless this coefficient is, as well, dependent on the orbital phase as shown by Grillmair (1992).
From our N-body simulations (CLM99), using multi-mass King-Michie models, we show that the tidal tails are populated mainly by the lighter stars of the pruned globular cluster, because of its mass segregation. In Fig. 30, we present the evolution of the mass function slope (assumed to be a power-law) for a simulation with a globular cluster on polar orbit (CLM99). The duration of this simulation is too short to observe strong changes in the mass spectrum through the cluster itself, nevertheless it can be seen that the radius of constant mass function slope is slightly expanding during the 800 Myr of the simulation. This is especially true in the inner parts of the cluster (see e.g. the isocontour $`\alpha =2`$ on Fig. 30).
Let us compare the amount of tidally stripped stars obtained in our simulations and observations. Because of the magnitude limitation of the plates and films, it is likely that we underestimate the observed tidal tails. The mass of the tidal tails in the case of NGC 5139 has been computed for a Salpeter law: if we assume a steep slope $`\alpha `$= –2.8 for the mass function in the tidal tails, we get a tail mass equal to about 1 % of the total mass of the cluster, equal to 5.1 10<sup>6</sup> M. In spite of the great uncertainty on the star counts of the tidal tails, such a large mass is an upper limit for the tidal tail mass from the simulations (CLM99). It confirms both that NGC 5139 has a genuinely large total mass (Meylan et al. 1995) and that the spectrum mass is likely less steep than $`\alpha =2.8`$ in the outer part. All these considerations point toward a mass for the tidal tails between 0.6 and 1.0 % of the total mass of NGC 5139. From N-body simulation performed by Moore (1996) we can note that the presence of tidal tails is the indication of low dark matter content in globular clusters.
NGC 7492 and Pal 12 provide an interesting comparison because of their similar characteristics: low masses (6 and 2 $`\times `$ $`10^4`$ M, respectively), low concentrations (1.0 and 0.9, respectively), and large distances from the galactic center (23.5 and 14.7 kpc, respectively). They are also both strongly influenced by the Galaxy: GO97 compute $`\nu _{tot}/\nu _{evap}`$ = 77.8 and 17.9, respectively. Nevertheless, their tidal tails appear strongly different, with a very extended structure for Pal 12 and very tiny one for NGC 7492. The last gravitational shock suffered by NGC 7492, if any, has to have occured much before the last one for Pal 12. The long time since the last gravitational shock suffered NGC 7492 has allowed the surface density of the unbound stars to fade along the cluster orbit. We estimate the last tidal shock suffered by Pal 12 to be about 350 Myr.
## 6 Conclusions
We have observed 20 galactic globular clusters with multi-color Schmidt plates and films on wide fields. Field and cluster stars are sorted in the color-magnitude plane. A star-count analysis is performed on the color selected stars to study the overdensities that can be attributed to the stars stripped from the globular clusters by tidal shocks (disk/bulge) as well as from internal dynamical evolution. We use the wavelet transform in order to enhance the weak tidal structures at large scales and in order to filter the high background noise at low galactic latitude. After highlighting the observational biases resulting from dust extinction and background galaxy clustering at faint magnitudes, we reach the following conclusions:
* All the clusters observed, which do not suffer from strong observational biases, present tidal tails, tracing their dynamical evolution in the Galaxy (evaporation, tidal shocking, tidal torquing, and bulge shocking).
* The clusters in the following sub-sample (viz. NGC 104, NGC 288, NGC 2298, NGC 5139, NGC 5904, NGC 6535, and NGC 6809) exhibit tidal extensions resulting from a shock, i.e. tails aligned with the tidal field gradient.
* The clusters in another sub-sample (viz. NGC 1261, NGC 1851, NGC 1904, NGC 5694, NGC 5824, NGC 6205, NGC 7492, Pal 5, and Pal 12) present extensions which are likely tracing the orbital path of the cluster with various degrees of mass loss.
* NGC 7492 is a striking case because of its very small extension and its high destruction rate driven by the galaxy as computed by GO97. Its dynamical “twin” for such an evolution, namely Pal 12, exhibits, on the contrary, a large extension tracing its orbital path, with a possible shock which happened more than 350 Myr.
* The velocity diffusion of the stripped stars is tentatively estimated, in one case (viz. NGC 6254), to be similar to the cluster velocity dispersion.
* Thanks to the relatively small distance of NGC 5139 and its high release of unbound stars during its current disk shocking, we estimate the mass loss to be between 0.6 and 1 % of the cluster total mass, taking into account a possible mass segregation in the cluster halo. This mass loss rate is consistent with our estimates from N-body simulations (CLM99).
* The second break in the surface density slope, apart from the break at the tidal radius (cf. the case of NGC 6254) could be an indicator of some recent gravitational shocks, with the $`Q_1^3`$ indicator displaying a range of values between -0.9 and -2. The latter is likely overestimated because of the uncorrected crowding towards the clusters.
The use of better quality data, e.g. wide-field CCD observations, combined with the present star-count method will allow in a near future to get rid easily of the observational biases and to obtain a better color selection thanks to a more accurate photometry. These observations will provide more precise observational estimates of the mass loss rates for different regimes of galaxy-driven cluster evolution. With the help of numerical simulations and accurate proper motions, it will be possible to constrain efficiently the parameters describing the galactic potential (disk scale-height, surface density, bulge size). In case of a flat dark matter halo (Pfenniger et al. 1994), the tidal shock on the globular clusters would be enhanced, depending on the surface density and the scale-height of this dark matter halo. Pal 2 is a good candidate to probe such dark matter halo flattening, because of its small distance to the galactic plane (Z = -2.2 kpc) and its relatively large distance to the galactic center (21.6 kpc), placing this cluster in a region where the tidal shocking by the disk only is expected to be low.
###### Acknowledgements.
We acknowledge warmly the ESO Schmidt Telescope operators O. and G. Pizarro, and the efficiency of the Schmidt astronomers B. Reipurth and J. Brewer for the conduction of the observations. We would like to think the whole MAMA team for its efficient support. We are greatful to R. Le Poole (Sterrenwacht Leiden) for facilitating the use of POSS I plates. We thank E. Bertin, L. Cambresy, J. Guibert, C. Loir and M. Odenkirchen for very helpful discussions. We acknowledge the W.E. Harris’s catalog of parameters for the galactic globular clusters. O. Gnedin gave us kindly an electronic version of his galactic globular cluster destruction rates. And we are very grateful to the referee, C. Grillmair, for his enlightening suggestions and comments which helped to improve this paper.
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# References
International Journal of Modern Physics B,
c World Scientific Publishing Company
A SURVEY OF NONADIABATIC SUPERCONDUCTIVITY
IN CUPRATES AND FULLERIDES
EMMANUELE CAPPELLUTI
Dipartimento di Fisica, Università “La Sapienza”, P.le Aldo Moro 2
Roma, 00185, Italy
and INFM, Unità Roma1
CLAUDIO GRIMALDI
École Polytechnique Fédérale, Départment de microtechnique IPM
Lausanne, CH-1015, Switzerland
LUCIANO PIETRONERO
Dipartimento di Fisica, Università “La Sapienza”, P.le Aldo Moro 2
Roma, 00185, Italy
and INFM, Unità Roma1
and
SIGFRID STRÄSSLER
École Polytechnique Fédérale, Départment de microtechnique IPM
Lausanne, CH-1015, Switzerland
High-$`T_c`$ superconductors are characterized by very low carrier densities. This feature leads to two fundamental consequences: on one hand the Fermi energies are correspondingly small and they can be of the same order of phonon frequencies. In such a situation nonadiabatic corrections arising from the breakdown of Migdal’s theorem can not be longer neglected. In addition, small carrier densities imply poor screening and correlation effects have to be taken into account. We present a comprehensive overview of the theory of superconductivity generalized into the nonadiabatic regime which is qualitatively different from the conventional one. In this framework some of the observed properties of the cuprates and the fullerene compounds can be naturally accounted for, and a number of theoretical predictions are proposed that can be experimentally tested.
After almost $`15`$ years and in spite of the huge amount of work spent in the field, no definitive description of the high temperature superconductivity phenomenon has been still achieved. In the meanwhile, a lot of “exotic” features have been discovered to compose the extremely rich phase diagram of these compounds. Most recently, theoretical and experimental research has addressed the issues of the pseudogap onset and of a possible stripe ordering. Many models have been proposed in order to account at least for some of the several anomalies in the high temperature superconductivity compounds (HTSC). Among them, the role of the electron-phonon (el-ph) coupling has received alternate fortune. In this contribution we review the
| conventional materials | HTSC |
| --- | --- |
| | |
| $`T_c^{\mathrm{max}}20`$ K | $`T_c40÷100`$ K |
| | |
| $`\alpha _{T_c}=0.5`$ | $`\alpha _{T_c}0.1÷0.8`$ |
| | |
| $`\alpha _m^{}=0`$ | $`\alpha _m^{}0.6÷0.8`$ |
| | |
| $`\rho (T)T^5`$ | $`\rho (T)T`$ |
| | |
| $`s`$ wave | $`d`$ wave |
| | |
| large bands | narrow bands |
| (high density of | (low density of |
| charge carriers) | charge carriers) |
| | |
| phononic pairing | pairing (?) |
| | |
main experimental evidences of a relevant role of the el-ph interaction on the superconductive pairing. We discuss these evidences in the context of the nonadiabatic theory of superconductivity and of the normal state.<sup>?,?,?</sup> The anomalies of the el-ph phenomenology, which were interpreted as hints of negligible conventional el-ph interaction, acquire now a natural explanation in the nonadiabatic regime. We also point briefly out analogies and differences between the cuprate family and the fullerides.
From the point of view of the conventional theory of superconductivity, described by Migdal-Eliashberg (ME) equations, the high value of the critical temperature $`T_c`$ in the HTSC is by itself a puzzle. The well-known McMillan formula relates $`T_c`$ essentially to two microscopic parameters, $`\mathrm{\Omega }_b`$ and $`\lambda _b`$, which represent respectively the energy scale of the intermediate boson and its coupling strength with electrons. In order to achieve in conventional theory $`T_c`$’s as high as $`100`$ K one should assume an anomalous large $`\lambda _b`$, physically prevented by structural distortions, or alternatively high energy bosons such as electronic excitations. This latter idea was initially supported by the discovery of a negligible isotope effect on $`T_c`$, $`\alpha _{T_c}`$, at optimal doping. However, later works found a drastic increase of $`\alpha _{T_c}`$, up to $`\alpha _{T_c}0.8`$, as the materials were underdoped.<sup>?</sup> Moreover, recent studies reported a finite, large and negative isotope effect on the electronic mass $`\alpha _m^{}0.6÷0.8`$, whereas $`\alpha _m^{}=0`$ is expected in conventional ME theory.<sup>?</sup> These experimental results point out the relevance of el-ph scattering in determining superconducting and normal state properties. The small value of $`\alpha _{T_c}`$ as well as the linear dependence on temperature of the resistivity, both peculiarities of the optimal doping, could be understood better as signature of some anomalous feature of optimal doping on the top of a phononic pairing scenario rather than characteristic of the whole phase diagram. Is $`d`$-wave symmetry of the superconductive gap compatible with this outlined picture? Conventional el-ph pairing does not usually show any momentum structure, yielding an isotropic $`s`$-wave superconductivity. However, it has been shown by different techniques that strong electronic correlation induces a momentum structure with a predominance of forward scattering (small $`𝐪`$’s).<sup>?</sup> In this situation $`d`$-wave symmetry can be favoured with respect to $`s`$-wave even within a phonon driven superconductivity.
As it is clear from the above brief overview of experimental results, el-ph coupling seems to play a major role in the normal and superconductive phenomenology, but at the same time it can not be properly understood within the framework of conventional ME theory. The new perspective we propose is the nonadiabatic theory of Fermi liquid.<sup>?,?,?</sup>
Electronic structure in high-$`T_c`$ materials is characterized by narrow bands crossing the Fermi level. This situation has two fundamental consequences with respect to both phononic and electronic scattering: ($`i`$) on one hand, electronic bands in these compounds can be so narrow to be of the same order of the phonon energy scale; ($`ii`$) on the other hand, strong electronic correlation effects are also definitively important in narrow band systems. Theoretical works has focused mainly on the second point. In our opinion, the first one can be equally and even more important.
The adiabatic parameter, defined as ratio between phonon frequencies and Fermi energy $`\mathrm{\Omega }_{ph}/E_F`$, can be as large as $`\mathrm{\Omega }_{ph}/E_F0.3`$ in cuprates. In this regime Migdal’s theorem, on which conventional el-ph ME theory rests, breaks down.<sup>?</sup> In order to investigate the nonadiabatic regime, in the past years we have developed a new theory using a different perturbation approach based on $`\lambda \mathrm{\Omega }_{ph}/E_F`$ instead of $`\mathrm{\Omega }_{ph}/E_F`$.<sup>?,?,?</sup> A sketch of the Cooper channel interaction in the nonadiabatic regime is depicted in Fig. 1.
A point which we would like to stress is that nonadiabaticity defines a qualitatively new theory wherein el-ph interaction plays a different role than in conventional ME theory. In this context, strong correlation acts as a positive factor with respect to el-ph coupling. As we have shown, the forward scattering predominance induced by the strong electronic correlation selects the phase space where nonadiabatic effects enhance the effective el-ph interaction.
In a schematized picture, we can say nonadiabaticity provides a enlarged scenario where corrections to the conventional ME theory, arising from the breakdown of Migdal’s theorem, can favour or depress superconducting pairing depending on the degree of electronic correlation in the system. This latter can be parametrized by the quantity $`q`$<sub>c</sub> representing the momentum selection induced by the strong correlation.<sup>?</sup> More correlated the system smaller $`q`$<sub>c</sub>.<sup>?</sup>
The issue of a phonon or non-phonon driven superconductivity is now re-analyzed in the light of the nonadiabatic theory of superconductivity. The evidences previously discussed of an “anomalous” el-ph pairing can now be understood in a natural way in the context of phonon mediated nonadiabatic superconductivity sustained by strong electronic correlation.
* $`T_c`$: critical temperatures as large as the experimental ones (up to 100 K) are accompanied in HTSC by a strong coupling phenomenology (large ratio $`2\mathrm{\Delta }/T_c`$, anomalous dip in tunnelling, etc …). Conventional phonon based ME theory requires an unphysically large el-ph coupling ($`\lambda 4`$) to account for these features.<sup>?</sup> In several works we have shown as nonadiabatic corrections modify the el-ph interaction and the structure itself of the Eliashberg equations generalized in nonadiabatic regime. For small $`q`$<sub>c</sub>’s (strong correlation) nonadiabatic corrections enhance the effective el-ph pairing reproducing high critical temperatures and strong coupling phenomenology with reasonable values of $`\lambda `$ ($`\lambda 1`$).<sup>?,?</sup>
* $`\alpha _{T_c}`$: By simple scaling analyses, the isotope effect $`\alpha _{T_c}`$ in ME theory (with no Coulomb repulsion: $`\mu =0`$) is shown to be $`\alpha _{T_c}=0.5`$ for any $`\lambda `$. However, as above discussed, nonadiabatic theory does not yield just an “effective” enhanced el-ph coupling $`\lambda `$, but defines a qualitatively new theory where a strong coupling phenomenology arises from normal value of $`\lambda `$. In similar way, evaluation of the isotope effect is also deeply different in nonadiabatic superconductivity. In particular, $`\alpha _{T_c}`$ shows strong fluctuations ($`\alpha _{T_c}0.20.8`$ for $`\mu =0`$) for small $`q`$<sub>c</sub>’s, precisely where the enhancement of $`T_c`$ is the largest, in agreement with measurement data.<sup>?,?</sup> An additional role can be played by the Van Hove singularity experimentally observed.<sup>?</sup>
* $`\alpha _m^{}`$: the different structure of the nonadiabatic equations of superconductivity is reflected also in a finite and negative isotope effect on the effective electronic mass $`m^{}`$,<sup>?</sup> as experimentally observed.<sup>?</sup> This can be therefore considered as a trademark of nonadiabaticity. Other possible explanations, as polaron band narrowing or closeness of a Van Hove singularity, do not seem satisfactory, although Van Hove singularity is certainly present.
* $`\rho (T)T`$: one of the most striking features of HTSC is the linear dependence of $`\rho (T)`$ in a wide range of temperature (up to 1100 K in La214).<sup>?</sup> At so high temperature resistivity is expected to be dominated by phonon scattering. A linear resistivity with no change of slope suggests therefore a common, phonon based, scattering mechanism for high and low temperature. As we have shown in Ref. 11, the linear behaviour can be related to the presence of a Van Hove singularity located at optimal doping in the context of nonadiabatic phonon scattering.<sup>?</sup>
* $`d`$-wave: it is often argued in literature that $`d`$-wave superconductivity is not compatible with phonon pairing while it points out a spin fluctuation interaction. Reason for this belief is the structureless el-ph interaction of conventional ME materials which leads to $`s`$-wave symmetry. HTSC compounds, however, are characterized by strong correlation inducing an important momentum selection with predominance of forward scattering. In such a situation $`d`$-wave superconductivity is favoured with respect to $`s`$-wave. We find a crossover from $`s`$\- to $`d`$-wave symmetry by lowering the momentum selection $`q`$<sub>c</sub>, or, in other words, by increasing the rate of electronic correlation. Nonadiabatic effects are shown to increase both the kind of ordering.<sup>?</sup>
* low charge carrier density: Another puzzling feature of HTSC materials, as shown by Uemura’s plot,<sup>?</sup> is the extreme low density of charge carriers, at least one order of magnitude lower than in low temperature superconductors. It is therefore quite surprising that the best superconductors are those with the poorest number of carriers whereas in conventional ME theory $`T_c`$ increases with the number of carriers through the density of states. This inconsistency is solved in natural way in the nonadiabatic theory of superconductivity. Small number of carriers leads to small Fermi energies driving the system into nonadiabatic regime where correction beyond Migdal’s theorem become important. Moreover, low density of charges implies poor screening of long-range el-ph interactions (small $`q`$<sub>c</sub>’s) and predominance of small exchanged momenta.
From the above analysis of different “exotic” features of HTSC materials the fundamental role of el-ph interaction appears evident. However, these fature can not be properly explained in the conventional framework of ME theory. The nonadiabatic theory of superconductivity and normal state provides a coherent picture where the above discussed anomalies arise as natural hallmarks of nonadiabatic effects. The primary actors are two: on one hand small Fermi energies determine the nonadiabaticity rate of the system and turn on nonadiabatic corrections due to the breakdown of Migdal’s theorem. Within this enriched phonon based scenario an additional but fundamental role is played by the strong electronic correlation that, through the induced momentum selection, amplifies the nonadiabatic effects and yields an effective enhanced el-ph coupling. This modified el-ph interaction, generalized in nonadiabatic regime, defines a new theory qualitatively different both from the conventional ME one and from the polaronic picture. Secondary effects can be also related to supporting actors, as the Van Hove singularity, that need in any case to be considered to account for peculiar details on the phase diagram.
All through this contribution, we have mainly focused on the physics of cuprates, the most studied HTSC compounds by theoretical and experimental means. However, a look at Uemura’s plot, completed by later materials like the fullerides, suggests a common origin of superconductivity for the different families of “HTSC” materials, as cuprates, bismuthates, fullerides and heavy fermion systems.<sup>?</sup> All these compounds are characterized by small density of charge carriers in contrast to conventional ME materials, and nonadiabaticity stands out as the natural candidate to explain superconductivity in these systems. The relevant parameter will be the ratio between intermediate boson frequencies (spin fluctuations in heavy fermions or phonons in fullerides and bismuthates) and the Fermi energies. In particular, alkali-doped C<sub>60</sub> compounds appear, for their relatively simple phase diagram, as the best materials where to check the nonadiabatic Fermi liquid picture. A detailed study shows how the conventional ME theory can not explain the experimental data available for these materials (high values of $`T_c`$, of $`2\mathrm{\Delta }/T_c`$, small value of $`\alpha _{T_c}`$) while the nonadiabatic theory is able to reproduce the experimental scenario with quite realistic values of the microscopical parameters.<sup>?</sup> In addition, our proposed description is liable to be tested in experimental different ways. For instance, we predict a finite isotope effect on the electronic specific heat<sup>?</sup> and on the spin susceptibility<sup>?</sup> in fullerides and more generally in any nonadiabatic superconductor. An anomalous reduction of $`T_c`$ by paramagnetic impurity scattering is also expected.<sup>?</sup> Experimental accuracy for these kind of measurements in nowadays already available and any experimental research along this line is welcome.
Acknowledgements
E. C. acknowledges the support of the INFM PRA-HTSC project.
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# Spectral Dependence of Polarized Radiation due to Spatial Correlations
## Abstract
Abstract: We study the polarization of light emitted by spatially correlated sources. We show that in general polarization acquires nontrivial spectral dependence due to spatial correlations. The spectral dependence is found to be absent only for a special class of sources where the correlation length scales as the wavelength of light. We further study the cross correlations between two spatially distinct points that are generated due to propagation. It is found that such cross correlation leads to sufficiently strong spectral dependence of polarization which can be measured experimentally.
PACS: 42.25.Ja, 42.25.Kb,42.25.Hz
In a series of interesting papers Wolf showed that in general the spectrum of electromagnetic radiation does not remain invariant under propagation, even through vacuum. The effect arises if the source has spatial correlations. The phenomenon was later confirmed experimentally and has been a subject of considerable interest . Further investigations of the source correlation effects have been done in the time domain theoretically and experimentally . Several applications of the effect have also been proposed . In a related development it has been pointed out that spectral changes also arise due to static scattering and dynamic scattering .
In the current paper we study the spectral dependence of polarized radiation that can arise due to spatial correlations. We are interested in sources where the emission from different points in the source are correlated i.e. the phase and amplitude shows systematic dependence on the position at the source. This subject has attracted considerable attention recently . In the present paper we are interested in obtaining the form of the correlation matrix for which the polarization in the far zone will not show any spectral dependence. One physical example we have in mind is radiation from a plasma in the presence of background magnetic field. The motion of charged particles at different spatial locations is in general correlated and will lead to spatially correlated radiation. In a separate paper we have also studied the angular dependence of the polarization in the far zone due to spatial correlations.
Consider a spatially extended 3-D source of polarized radiation, as shown in figure 1, characterized by the charge density $`\rho ^{(r)}(𝐫,t)`$ and current density $`𝐉^{(r)}(𝐫,t)`$. The electric and magnetic field vectors, $`𝐄^{(r)}`$ and $`𝐁^{(r)}`$, generated by this source, can be written as ,
$`𝐄^{(r)}(𝐑,t)`$ $`=`$ $`{\displaystyle \frac{1}{4\pi ϵ_0}}{\displaystyle }d^3r({\displaystyle \frac{𝐒}{S^3}}\left[\rho ^{(r)}(𝐫,t^{})\right]_{\mathrm{ret}}+{\displaystyle \frac{𝐒}{cS^2}}`$ (2)
$`\left[{\displaystyle \frac{\rho ^{(r)}(𝐫,t^{})}{t^{}}}\right]_{\mathrm{ret}}{\displaystyle \frac{1}{c^2S}}\left[{\displaystyle \frac{𝐉^{(r)}(𝐫,t^{})}{t^{}}}\right]_{\mathrm{ret}})`$
$`𝐁^{(r)}(𝐑,t)`$ $`=`$ $`{\displaystyle \frac{\mu _0}{4\pi }}{\displaystyle }d^3r(\left[𝐉^{(r)}(𝐫,t^{})\right]_{\mathrm{ret}}\times {\displaystyle \frac{𝐒}{S^3}}`$ (3)
$`+`$ $`\left[{\displaystyle \frac{𝐉^{(r)}(𝐫,t^{})}{t^{}}}\right]_{\mathrm{ret}}\times {\displaystyle \frac{𝐒}{cS^2}})`$ (4)
where $`𝐒=𝐑𝐫`$, $`S=|𝐒|`$, $`[f(𝐫,t^{})]_{\mathrm{ret}}=f(𝐫,tS/c)`$, $`c=`$ speed of light and $`𝐑`$ and $`𝐫`$ are the position vectors of the observation point $`Q`$ (Fig 1) and a particular point $`P`$ on the source respectively. Let $`𝐄(𝐑,\omega )`$, $`𝐁(𝐑,\omega )`$, $`\rho (𝐫,\omega )`$ and $`𝐉(𝐫,\omega )`$ be the analytic signals associated with the Fourier transforms of the electric field, magnetic field, charge and current densities respectively. In the radiation zone, $`𝐄(𝐑,\omega )`$ and $`𝐁(𝐑,\omega )`$ are given by,
$`𝐄(𝐑,\omega )`$ $`=`$ $`{\displaystyle \frac{i\omega }{4\pi ϵ_0c^2}}{\displaystyle d^3r\frac{e^{i\omega S/c}}{R}\left[𝐉(𝐫,\omega )𝐉(𝐫,\omega )\widehat{𝐑}\widehat{𝐑}\right]}`$ (5)
$`𝐁(𝐑,\omega )`$ $`=`$ $`{\displaystyle \frac{i\omega \mu _0}{4\pi c}}{\displaystyle d^3r\frac{e^{i\omega S/c}}{R}𝐉(𝐫,\omega )\times \widehat{𝐑}}`$ (6)
where $`\widehat{𝐑}=𝐑/R`$ and $`R=|𝐑|`$. In obtaining this result we have used the current conservation equation and kept only the leading order terms in $`1/R`$. The dimensions of the source (Fig. 1) are assumed to be much smaller than its distance to the observation point $`Q`$ and hence dropping higher orders in $`1/R`$ is reasonable. We see that the electric and magnetic fields are orthogonal to one another as well as to the propagation vector $`𝐤=\widehat{𝐑}\omega /c`$. Therefore at the far away point $`Q`$ only the $`\theta ,\varphi `$ components of the electric and magnetic field vector are nonzero.
We are interested in evaluating the coherency matrix $`J_{ij}(𝐑,\omega )`$ at the point Q in the radiation zone. We can relate this in terms of the source correlation function
$$W_{ij}^S(𝐫,𝐫^{};\omega )\delta (\omega \omega ^{})=<J_i^{}(𝐫,\omega )J_j(𝐫^{},\omega ^{})>,$$
(7)
where the angular brackets represent ensemble averages. Here we have made the standard assumption that the source fluctuations are represented by a stationary statistical ensemble, atleast in the wide sense . The correlation function of the electric field $`W_{ij}`$ is then also given by,
$$W_{ij}(𝐑,𝐑^{};\omega )\delta (\omega \omega ^{})=<E_i^{}(𝐑,\omega )E_j(𝐑^{},\omega ^{})>$$
(8)
The coherency matrix $`J_{ij}(𝐑,\omega )`$ is given by,
$$J_{ij}(𝐑,\omega )=W_{ij}(𝐑,𝐑,\omega ).$$
(9)
A straightforward calculation using Eqs. 5 and 6 gives,
$`J_{ij}(𝐑,\omega )`$ $`=`$ $`Z(\omega ){\displaystyle d^3rd^3r^{}\frac{e^{ik(SS^{})}}{SS^{}}}`$ (11)
$`\xi _{il}W_{lm}^S(𝐫,𝐫^{},\omega )\xi _{mj}`$
where,
$$\xi _{ij}=\widehat{R}_i\widehat{R}_j+\delta _{ij},$$
(12)
$`k=\omega /c`$, $`S^{}=|𝐑𝐫^{}|`$ and $`Z(\omega )`$ is an overall normalization factor which will not play any role in our analysis. In Eq. 11 as well in the rest of the paper summation over repeated indices is understood unless otherwise stated. We point out that since only the transverse $`(\theta ,\varphi )`$ components of the electric field vector are nonzero, the subscripts on $`J`$ as well as $`W^S`$ refer only to these components. The $`(SS^{})`$ term in the exponent can be replaced by $`\widehat{𝐑}𝚫`$, where $`𝚫=𝐫^{}𝐫`$.
The matrix $`W_{ij}^S(𝐫,𝐫^{},\omega )`$ measures the correlations between two spatially distinct points. We are interested in studying the polarization at a point $`Q`$ in the far zone when the cross correlation matrix has nontrivial dependence on $`𝐫`$ and $`𝐫^{}`$. It is convenient to express the matrix $`W_{ij}^S`$ in terms of the variables $`𝐫_𝐚=(𝐫+𝐫^{})/2`$ and $`𝚫`$ instead of $`𝐫`$ and $`𝐫^{}`$. We rewrite Eq. 11 in terms of $`𝐫_𝐚`$ and $`𝚫`$,
$`J_{ij}(𝐑,\omega )`$ $`=`$ $`Z(\omega ){\displaystyle d^3r_ad^3\mathrm{\Delta }\frac{e^{ik\widehat{𝐑}𝚫}}{R^2}}`$ (14)
$`\xi _{il}W_{lm}^S(𝐫_𝐚,𝚫,\omega )\xi _{mj}.`$
The limits of integration over the variables $`𝐫_𝐚`$ and $`𝚫`$ are assumed to extend over all space. In obtaining Eq. 14 we have ignored higher order terms in $`|𝚫|/R`$ in the exponent which is reasonable if the observation point is far enough away.
In this paper we are primarily interested in studying the $`𝚫`$ dependence of the correlation matrix and hence consider only sources for which the $`𝐫_𝐚`$ dependence factorizes, that is,
$$W_{ij}^S(𝐫_𝐚,𝚫,\omega )=j_{il}(𝐫_𝐚)G_{lj}(𝚫,\omega ).$$
(15)
This equation defines the two matrices, $`j_{ij}(𝐫_𝐚)`$ and $`G_{ij}(𝚫,\omega )`$. Physically this factorization means that the spatial correlations as well as the spectral response of different regions of the source are identical to one another. However the polarization of light emitted by different points on the source need not be identical. Furthermore in order to focus on the spectral dependence of polarization which arises due to spatial correlations we assume that,
$$G_{ij}(𝚫=0,\omega )=N𝒮(\omega )\delta _{ij}$$
(16)
where $`N`$ is a normalization factor, $`𝒮(\omega )`$ is the spectrum of light emitted by any point on the source and $`\delta _{ij}`$ is the standard Kronecker delta. Physically Eq. 16 means that if we ignore spatial correlations, the polarization of light emitted by the source has no spectral dependence. If the spatially distinct points are independent then,
$$G_{ij}(𝚫,\omega )=\delta ^3(𝚫)𝒮(\omega )\delta _{ij}.$$
(17)
In this case we find that the resultant matrix $`J`$ is given by
$$J_{ij}(𝐑,\omega )=\frac{𝒮(\omega )}{R^2}d^3r_a\xi _{il}j_{lm}(𝐫_𝐚)\xi _{mj}$$
(18)
i.e. an incoherent integral over the entire source. In this case we find, as expected, that the resulting polarization at $`Q`$ has no spectral dependence.
In general, however, the polarization observed at the far away point $`Q`$ does acquire spectral dependence purely due to spatial correlations. We next obtain conditions on the correlation matrix $`W^S`$ under which this spectral polarization dependence is absent. We consider only those sources which satisfy the conditions given in Eq. 15 and 16 since we are interested in isolating the spectral dependence that arises only due to spatial correlations.
If the coherency matrix factorizes into a function of $`\omega `$ times a matrix independent of $`\omega `$, that is,
$$J_{ij}(𝐑,\omega )=M_{ij}(𝐑)h(\omega ),$$
(19)
where $`M_{ij}`$ is a matrix independent of $`\omega `$ and $`h(\omega )`$ is a function of $`\omega `$, then the polarization observed at $`Q`$ will have no spectral dependence. In order to obtain the functional form of $`W_{lm}^S(𝐫_𝐚,𝚫,\omega )`$ which can lead to coherency matrix of the form given in Eq. 19 we substitute Eq. 15 into Eq. 14. We can then express $`J_{ij}`$ as
$`J_{ij}(𝐑,\omega )`$ $`=`$ $`Z(\omega )\xi _{il}{\displaystyle \frac{J_{ln}^0\stackrel{~}{G}_{nm}(𝐤,\omega )}{R^2}}\xi _{mj}`$ (20)
where
$$J_{ln}^0=d^3r_aj_{ln}(𝐫_𝐚)$$
(21)
and
$$\stackrel{~}{G}_{nm}(𝐤,\omega )=d^3\mathrm{\Delta }G_{nm}(𝚫,\omega )e^{i𝐤\mathrm{\Delta }},$$
(22)
with $`𝐤=k\widehat{𝐑}`$. Hence we see that $`\stackrel{~}{G}_{nm}(𝐤,\omega )`$ is the Fourier transform of $`G_{nm}(𝚫,\omega )`$. We can therefore also write
$$G_{nm}(𝚫,\omega )=\frac{d^3k}{(2\pi )^3}\stackrel{~}{G}_{nm}(𝐤,\omega )e^{i𝐤𝚫}$$
(23)
We point out that the integration in the above equation is performed treating $`\omega `$ to be independent of $`𝐤`$. In order that the polarization at the point $`Q`$ has no spectral dependence, $`\stackrel{~}{G}_{nm}(𝐤,\omega )`$ has to be of the form
$$\stackrel{~}{G}_{nm}(𝐤,\omega )=\stackrel{~}{A}_{nm}(𝐤/\omega )h(\omega )$$
(24)
or
$$\stackrel{~}{G}_{nm}(𝐤,\omega )=\delta _{nm}\stackrel{~}{H}(𝐤,\omega ),$$
(25)
where the matrix $`\stackrel{~}{A}_{nm}`$ depends on $`𝐤`$ and $`\omega `$ only through the combination $`𝐤/\omega `$ and $`H(𝐤,\omega )`$ is some function of $`𝐤,\omega `$. Substituting Eq. 24 into Eq. 23 we find that,
$$G_{lm}(𝚫,\omega )=h(\omega )G_{lm}(\omega 𝚫),$$
(26)
which is analogous to the scaling law obtained by Wolf in his analysis of spectral shifts from spatially correlated sources. Alternatively substituting the factorized form Eq. 25 into Eq. 23 we find that,
$$G_{lm}(𝚫,\omega )=\delta _{lm}H(𝚫,\omega ),$$
(27)
where $`H(𝚫,\omega )`$ is the Fourier transform of $`\stackrel{~}{H}(𝐤,\omega )`$. We therefore find that in order that the polarization in the radiation zone does not acquire spectral dependence, the correlation matrix $`G_{lm}(𝚫,\omega )`$ has to be of the form given in Eq. 26 or Eq. 27.
We next consider a specific example and calculate the spectral dependence arising due to correlations. We consider a planar circular source of radius $`a`$, which is spatially uncorrelated. The source emits polarized radiation such that its coherency matrix is given by,
$$J(\rho ,\varphi )=A\left(\begin{array}{cc}\mathrm{sin}^2\varphi & \mathrm{sin}\varphi \mathrm{cos}\varphi \\ \mathrm{sin}\varphi \mathrm{cos}\varphi & \mathrm{cos}^2\varphi \end{array}\right)$$
(28)
where $`A`$ is a constant and $`\rho ,\varphi `$ are the polar coordinates of any point at the position vector $`𝐫`$ on the source. The source luminosity is independent of position and the polarization vectors point along $`\widehat{\varphi }`$. We point out that the source has been constructed such that at any point close to the axis of symmetry of the source the integrated polarization is zero.
As is well known, although the source is uncorrelated, the cross correlation between any two points $`P_1`$ and $`P_2`$ need not be zero due to the Van Cittert-Zernike theorem . We consider the experimental arrangement shown in Fig. 2, where the light after being reflected from $`P_1`$ and $`P_2`$ is observed at the point $`Q`$. We are interested in the spectral dependence of the polarization observed at $`Q`$. The cross correlation matrix between any two points $`P_1(𝐑_\mathrm{𝟏})`$ and $`P_2(𝐑_\mathrm{𝟐})`$ in the far zone close to symmetry axis of the source, is given by,
$$W_{ij}(𝐑_\mathrm{𝟏},𝐑_\mathrm{𝟐},\omega )=\left(\frac{k}{2\pi }\right)^2d^2rJ_{ij}(𝐫)\frac{e^{ik(\rho /R)L\mathrm{cos}(\varphi \psi )}}{R^2}$$
(29)
where $`L\mathrm{cos}\psi =x_2x_1`$, $`L\mathrm{sin}\psi =y_2y_1`$, $`(x_1,y_1)`$ and $`(x_2,y_2)`$ are the cartesian coordinates of the projections of $`𝐑_\mathrm{𝟏}`$ and $`𝐑_\mathrm{𝟐}`$ respectively on the plane of the source and $`R=R_1=R_2`$ is the distance of the point $`O`$ on the source from points $`P_1`$ and $`P_2`$ which have been assumed to be placed symmetrically for simplicity. In obtaining Eq. 29 we have followed the treatment given in Ref. for the calculation of cross correlation between two points $`P_1`$ and $`P_2`$ close to the symmetry axis of a spatially uncorrelated source. Lack of spatial correlation implies that the cross correlation matrix at the source $`W_{ij}^S(𝐫_\mathrm{𝟏},𝐫_\mathrm{𝟐},\omega )=J_{ij}(𝐫_\mathrm{𝟏},\omega )\delta ^2(𝐫_\mathrm{𝟐}𝐫_\mathrm{𝟏})`$. We have further assumed that $`J_{ij}(𝐫_\mathrm{𝟏},\omega )`$ has no spectral dependence. The integral in Eq. 29 is over the source and we are using polar coordinates $`x=\rho \mathrm{cos}\varphi `$ and $`y=\rho \mathrm{sin}\varphi `$.
Following the treatment given in we can calculate the cross correlation matrix and hence the Stokes parameters at the point of observation. For the source under consideration the result can be obtained analytically. We find, upto an overall common factor $`Aa^2k^2/2\pi R^2`$
$$s_0=1+2J_1(v)/v,$$
$$s_1=\frac{2}{v^2}\frac{(y_2y_1)^2(x_2x_1)^2}{L^2}[vJ_1(v)2(1J_0(v)],$$
$$s_2=\frac{4}{v^2}\frac{(x_2x_1)(y_2y_1)}{L^2}[vJ_1(v)2(1J_0(v)],$$
$$s_3=0$$
where $`v=kLa/R`$. We therefore find that the wave is linearly polarized at the point of observation $`Q`$ with the degree of polarization given by,
$$P=2\frac{|vJ_1(v)+22J_0(v)|}{v^2+2vJ_1(v)}$$
(30)
which has nontrivial spectral dependence. The orientation of the linear polarization vector, given by,
$$\mathrm{tan}(2\psi )=\frac{2(x_2x_1)(y_2y_1)}{(x_2x_1)^2(y_2y_1)^2},$$
(31)
no spectral dependence. The calculated degree of polarization for this example is plotted in Fig. 3. We clearly see that it is a very significant effect and can be observed experimentally. The orientation of the linearly polarized component depends on the positions of $`P_1`$ and $`P_2`$ and contains information about the polarization profile of the source.
We next study a somewhat more complicated source for which the coherency matrix at any point $`𝐫=(\rho ,\varphi )`$ is same as that given in Eq. 28, with $`\varphi `$ replaced by $`\varphi +\alpha \rho `$, where $`\alpha `$ is a parameter. In this case we numerically calculate the Stokes parameter. The spectral dependence of the degree of polarization and the orientation of the linear polarization is shown in Fig. 4 for some representative choices of the parameter $`\alpha `$. The plot uses $`x_2x_1=1.0`$ and $`y_2y_1=0.2`$ in arbitrary units. The state of polarization in the far zone only depends on the dimensionless quantities $`v=kLa/R`$, $`(x_2x_1)/L`$, $`(y_2y_1)/L`$ and $`\alpha a`$ where $`L^2=(x_2x_1)^2+(y_2y_1)^2`$. We find that in this case both the orientation angle of the linear polarization vector and the degree of polarization show a dramatic spectral variation. The Stokes parameter $`s_3`$ vanishes in this case also showing that there is no circularly polarized component at the point $`Q`$.
The effect discussed in this paper is observable experimentally by using a primary source which has spatial correlations or by generating the correlations due to propagation as shown in the above example. We assume an experimental arrangement shown in Fig. 2. We first measure the spectral dependence of polarization at the point $`Q`$ due to the waves emerging from the secondary sources $`P_1`$ and $`P_2`$ individually. The effect of spatial correlations can then be determined by measuring the polarization at $`Q`$ due to interference of the waves emerging from $`P_1`$ and $`P_2`$.
We conclude that spatially correlated sources of polarized radiation generically display nontrivial spectral dependence of the state of polarization in the far zone. This dependence goes away if the correlation matrix displays a scaling law or factorizes into a constant matrix and a function of the relative coordinate and frequency.
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# Comments on Noncommutative Open String Theory: V-duality and Holography
## I Introduction
A noncommutative space or spacetime is the one with noncommuting coordinates, satisfying
$$[x^\mu ,x^\nu ]=i\theta ^{\mu \nu }\mu ,\nu =0,1,2,\mathrm{},$$
(1)
where $`\theta ^{\mu \nu }`$ are antisymmetric and real parameters of dimension length squared. A field theory on such a space can be formulated using a representation, in which the coordinates $`x^\mu `$ are the same as usual, but the product of any two fields of $`x^\mu `$ is deformed to the Moyal star-product:
$$fg(x)=e^{(i/2)\theta ^{\mu \nu }_\mu ^x_\nu ^y}f(x)g(y)|_{y=x},$$
(2)
while the commutator in Eq. (1) is understood as the Moyal bracket with respect to the star product:
$$[x^\mu ,x^\nu ]x^\mu x^\nu x^\nu x^\mu .$$
(3)
Recently it has been shown that Yang-Mills theory (or open string theory) on such noncommutative space (or spacetime), which we will abbreviate as NCYM (or NCOS), arises naturally in string or M(atrix) theory on coincident D-brane world-volume in anti-symmetric tensor backgrounds in certain scaling limits (decoupling or DLCQ limits) . (In order to obtain a nontrivial theory defined only on the brane world-volume, these scaling limits require that in addition to the usual $`\alpha ^{^{}}0`$ limit, certain components of the closed-string metric and/or those of the background parallel to the brane world-volume should also be scaled in appropriate way. For details, see refs. .) This strongly suggests that space-space or even space-time noncommutativity could be a general feature of the unified theory of quantum gravity at a generic point inside the moduli space of string/M theory. Though perhaps not every noncommutative field (or string) theory is a consistent quantum theory on its own, there is a belief that noncommutative field (or string) theories that can arise as effective limits in fundamental string theory should be consistent quantum theory on their own. Up to now, only NCYM with space-space noncommutativity and NCOS with space-time noncommutativity have been obtained by taking certain decoupling limits in string theory. It is important to clarify whether there exist decoupling limits in string theory with backgrounds that lead to either NCYM or NCOS with both space-space and space-time noncommutativity. Namely one wants to know how big the moduli space is for NCYM and NCOS that can arise from string theory.
Constant bulk B-background in string theory in topologically trivial spacetime can be gauged away, while inducing constant gauge field background on the D-brane world-volume. One might wonder whether the nature, electric or magnetic, of the gauge background would affect the scaling limit that decouples the theory on the D-branes from closed strings in the surrounding bulk. From the Born-Infeld action it is known that on the D-brane no electric field can be stronger than a critical electric field, while the same is not true for magnetic fields. Indeed recent careful reexamination of the decoupling limits shows that though the decoupling limit in a magnetic background always results in NCYM with only space-space noncommutativity , in an electric background the decoupling limit becomes different and leads to NCOS with only space-time noncommutativity . Moreover, theory with space-time noncommutativity is expected to behave very differently from one with only space-space noncommutativity. Recently whether a field theory with space-time noncommutativity is unitary has been questioned in the literature . All these inspire the following questions: What would happen if there are both electric and magnetic backgrounds? Could an NCYM with space-time noncommutativity, or an NCOS with both space-space and space-time noncommutativity, arises in favorable situations?
The present paper will address the problem of the interplay of constant electric and magnetic backgrounds in determining the decoupling limit towards a noncommutative theory on the D-brane world-volume. To simplify, we will restrict ourselves to the case of the D3-brane(s). Generalizing to other D$`p`$-branes should be straightforward. We will consider two special cases, in which the electric and magnetic backgrounds are either parallel or perpendicular to each other. It is known that the endpoints of an open string behave like (opposite) charges on the D3-brane, and the motion of a charge in the above two background configurations is very different. So we expect that there should be important differences between the decoupling limits in the above two cases. As we will show in Sec. 2, in either case no decoupling limit can be found to lead to an NCYM with space-time noncommutativity. On the other hand, in Sec. 3 we will show that in favorable situations appropriate decoupling limit may result in NCOS with both space-space and space-time noncommutativity.
In electrodynamics it is known that Lorentz boosts act on constant electromagnetic backgrounds. Through the decoupling limit the latter, in turn, affects the noncommutativity parameters that define the resulting NCOS. Thus, the NCOS that result from Lorentz-boost related backgrounds should be equivalent to each other, describing the same decoupled D-brane system. We will call this exact equivalence among NCOS as V-duality, which can be viewed as the fingerprint of the antecedent Lorentz-boost action surviving the decoupling limit. In Sec. 4 we will identify some orbits of V-duality in the moduli space of NCOS (with both space-space and space-time noncommutativity).
Previously Li and one of us have shown that there is a running holographic correspondence between NCYM and its gravity dual. Namely, the radial dependence of the profile of NSNS fields in the gravity dual of an NCYM can be derived from the Seiberg-Witten relations between close string moduli and open string moduli, provided that the string tension is running with a simply prescribed dependence on the energy scale, which is identified with the radial coordinate by the well-known UV-IR relation . We will show in Sec. 5 that the Li-Wu holography argument can be generalized to NCOS, though with a different prescription for the running string tension.
## II Decoupling limit for NCYM
In this section, we concentrate on the decoupling limit for NCYM, when the gauge background $`B_{\mu \nu }`$ on a flat D$`p`$-brane world volume (with a constant metric $`g_{\mu \nu }`$) has both electric and magnetic components. For definiteness, we consider the case with $`p=3`$. To be specific, we restrict ourselves to the special cases with the electric and magnetic fields are either perpendicular or parallel to each other. The generalization to the most general configuration should be straightforward.
A constant $`B`$-background on the D-brane does not affect the equations of motion for open strings, while it changes the open string boundary conditions to
$$g_{\mu \nu }_nX^\nu +2\pi \alpha ^{}B_{\mu \nu }_sX^\nu |_\mathrm{\Sigma }=0,$$
(4)
where the operators $`_n`$ and $`_s`$ are the derivatives normal and tangential to the worldsheet boundaries $`\mathrm{\Sigma }`$. For the disc topology, the propagator along the boundary is known to be
$$<x^\mu (\tau )x^\nu (0)>=\alpha ^{}G^{\mu \nu }\mathrm{ln}(\tau ^2)+i\frac{\theta ^{\mu \nu }}{2}\epsilon (\tau ).$$
(5)
As emphasized by Seiberg and Witten in Ref. , the physics behind these equations is that the moduli ($`G_{\mu \nu }`$, $`\theta ^{\mu \nu }`$, $`G_s`$) seen by open string ends on the D-brane are very different from those ($`g_{\mu \nu }`$, $`B_{\mu \nu }`$, $`g_s`$) seen by close strings; they are related by the following elegant relations
$$G_{\mu \nu }=g_{\mu \nu }(2\pi \alpha ^{})^2(Bg^1B)_{\mu \nu },$$
(6)
$$G^{\mu \nu }=\left(\frac{1}{g+2\pi \alpha ^{}B}\right)_S^{\mu \nu },$$
(7)
$$\theta ^{\mu \nu }=2\pi \alpha ^{}\left(\frac{1}{g+2\pi \alpha ^{}B}\right)_A^{\mu \nu },$$
(8)
$$G_s=g_s\left(\frac{\text{det}G_{\mu \nu }}{\text{det}(g_{\mu \nu }+2\pi \alpha ^{}B_{\mu \nu })}\right)^{\frac{1}{2}},$$
(9)
where $`()_S`$ and $`()_A`$ denote, respectively, the symmetric and anti-symmetric parts, and $`G_s`$, $`g_s`$ the open string and closed string coupling.
For the purely magnetic case (with $`B_{0i}=0`$), the scaling limit that decouples the theory on the D-brane from closed strings in the bulk has been analyzed in Ref. , and may be summarized as a limit subject to the following conditions: (1) $`\alpha ^{}0`$; (2) $`G_{\mu \nu }`$ is finite; (3) $`\theta ^{\mu \nu }`$ is finite. The decoupling limit results in an NCYM on the D-brane world volume. For the purely electric case (with $`B_{ij}=0`$), the above decoupling limit has been shown not to exist, because of the existence of a critical electric field-strength. In the following, we will carry out an analysis for the case with both electric and magnetic components present in the antisymmetric tensor $`B_{\mu \nu }`$. The interplay between the electric background ($`𝐄`$) and the magnetic background ($`𝐁`$) is worthwhile to explore, since in the presence of both an electric field and a magnetic field the dynamical behavior of a point charge, representing an endpoint of the open string, is known to be very different from the case in either a purely magnetic or a purely electric field. For simplicity, we assume that either $`𝐄𝐁`$ or $`𝐄𝐁`$. In this section, we discuss whether a decoupling limit leading to NCYM exists in these two cases.
First, let us consider the case with $`B_{01}=E`$ and $`B_{12}=B`$, all other components being zero; namely, the tensor $`B_{\mu \nu }`$ takes the form (for $`\mu ,\nu =0,1,2`$)
$$B_{\mu \nu }=\left(\begin{array}{ccc}0& E& 0\\ E& 0& B\\ 0& B& 0\end{array}\right).$$
(10)
The closed string metric $`g_{\mu \nu }`$ is taken to be of the diagonal form
$$g_{\mu \nu }=\left(\begin{array}{ccc}g_0& 0& 0\\ 0& g_1& 0\\ 0& 0& g_2\end{array}\right).$$
(11)
For convenience, we follow Ref. and to introduce the critical value, $`E_c`$, of the electric field
$$E_c=\frac{\sqrt{g_0g_1}}{2\pi \alpha ^{}}.$$
(12)
Substituting Eq. (10) and Eq. (11) into the Seiberg-Witten relations Eqs. (6)-(9), we get
$$G_{\mu \nu }=\left(\begin{array}{ccc}g_0(1e^2)& 0& g_0eb\\ 0& g_1(1e^2)+\frac{g_0g_1}{g_2}b^2& 0\\ g_0eb& 0& g_2+g_0b^2\end{array}\right)$$
(13)
$$\theta ^{\mu \nu }=\frac{1}{[g_2(1e^2)+g_0b^2]E_c}\left(\begin{array}{ccc}0& g_2e& 0\\ g_2e& 0& g_0b\\ 0& g_0b& 0\end{array}\right),$$
(14)
$$G_s=g_s\sqrt{1+\frac{g_0}{g_2}b^2e^2}$$
(15)
where the dimensionless electric and magnetic field strength are given by
$$e=\frac{E}{E_c},b=\frac{B}{E_c}.$$
(16)
To get an NCYM, we need to take $`\alpha ^{}0`$ to decouple massive open string excitations, while keeping the open string moduli $`G_{\mu \nu }`$, $`\theta ^{\mu \nu }`$ and $`G_s`$ finite. Inspection of Eqs. (13) and (14) shows that the following conditions provide the only possible solution for the NCYM limit:
1. $`|e|<1`$;
2. $`B=bE_c=1/\theta `$ finite;
3. $`g_0=1`$, $`g_1=g_2=g(\alpha ^{})^2`$, so that formally $`E_c`$ is a finite parameter; for later convenience, to normalize open string metric to $`G_{11}=G_{22}=1`$, one may take $`g=(2\pi \alpha ^{}B)^2`$;
4. $`g_s\alpha ^{}`$ to keep $`G_s`$ finite.
This solution is unique up to finite separate rescaling for $`g_0`$, $`g_1`$ and $`g_2`$. It is easy to verify that in this limit
$$\theta ^{0i}=0,\theta ^{12}=\theta .$$
(17)
Therefore the resulting field theory has only space-space noncommutativity. Though $`E`$ or $`e`$ does not affect the noncommutativity parameters $`\theta ^{\mu \nu }`$, it does make the open string metric $`G_{\mu \nu }`$ non-diagonal, i.e. it makes the $`x_0`$ and $`x_2`$ axes oblique with respect to open string metric. The appearance of the off-diagonal $`G_{02}`$ is not surprising: the open string endpoint, behaving like a charge, acquires a drift velocity in the $`x_2`$-direction in the present cross-field background with $`E_1=E`$ and $`B_3=B`$.
In this way, we see that the scaling limit of NCYM is incompatible with space-time noncommutativity. This is just right, since field theory with space-time noncommutativity is potentially non-unitary . A similar analysis can be done for the case with $`𝐄𝐁`$, again resulting in an NCYM with vanishing $`\theta _{0i}`$.
## III Decoupling limit of NCOS
In this section, we present a new decoupling limit of NCOS to demonstrate the interplay between the electric and magnetic components of the background.
### A The E$``$B case
To achieve this goal, we take in the closed string metric Eq. (11) $`g_0=g_1=g`$, this leads to corresponding open string moduli by using Eqs. (12)-(15).
$$G_{\mu \nu }=\left(\begin{array}{ccc}g(1e^2)& 0& geb\\ 0& g(1e^2)+\frac{g^2b^2}{g_2}& 0\\ geb& 0& g_2+gb^2\end{array}\right),$$
(18)
$$\theta ^{\mu \nu }=\frac{2\pi \alpha ^{}}{g_2g(1e^2)+g^2b^2}\left(\begin{array}{ccc}0& g_2e& 0\\ g_2e& 0& gb\\ 0& gb& 0\end{array}\right).$$
(19)
$$G_s=g_s\sqrt{1e^2+\frac{gb^2}{g_2}}$$
(20)
In taking the decoupling limit for NCOS, $`\alpha ^{}`$ is kept fixed, while $`G_{\mu \nu }`$ and $`\theta ^{\mu \nu }`$ have to have a finite limit. To achieve this goal, we introduce the following scaling limit
1. $`e1`$, with $`g(1e^2)=\frac{2\pi \alpha ^{}}{\theta _0}`$ finite;
2. $`g_2`$ is finite; for convenience, we take $`g_2=1`$;
3. $`b0`$, with $`gb=\frac{2\pi \alpha ^{}}{\theta _1}`$ finite.
With this scaling limit, we get the moduli of the resulting NCOS as follows: the metric
$$G_{\mu \nu }=\left(\begin{array}{ccc}\frac{2\pi \alpha ^{}}{\theta _0}& 0& \frac{2\pi \alpha ^{}}{\theta _1}\\ 0& \frac{2\pi \alpha ^{}}{\theta _0}+\left(\frac{2\pi \alpha ^{}}{\theta _1}\right)^2& 0\\ \frac{2\pi \alpha ^{}}{\theta _1}& 0& 1\end{array}\right),$$
(21)
and the noncommutativity matrix
$$\theta ^{\mu \nu }=\frac{2\pi \alpha ^{}}{\frac{2\pi \alpha ^{}}{\theta _0}+\left(\frac{2\pi \alpha ^{}}{\theta _1}\right)^2}\left(\begin{array}{ccc}0& 1& 0\\ 1& 0& \frac{2\pi \alpha ^{}}{\theta _1}\\ 0& \frac{2\pi \alpha ^{}}{\theta _1}& 0\end{array}\right).$$
(22)
This scaling limit is striking in that it results in NCOS with both space-time and space-space noncommutativity. Note that this is different from the NCYM limit, where space-time noncommutativity can not result from the decoupling limit. Again, the appearance of nonzero off-diagonal elements $`G_{02}=G_{20}`$ in the NCOS metric is very natural: The end points of the open string behave like opposite charges which, in a cross-field with $`E_1`$ and $`B_3`$, acquire a drift velocity in the $`x_2`$-direction independent of the sign of charges. (This is nothing but the classical picture of the Hall effect in condensed matter physics.)
In the above scaling limit, the open string coupling $`G_s`$ vanishes. To have an interacting theory, we can follow Ref. to consider $`N`$ coincident $`D`$-branes, so that the effective open string coupling is
$$G_{eff}=NG_s=Ng_s\sqrt{1e^2}.$$
(23)
In the large $`N`$ limit, if we scale $`N`$ as
$$N\frac{1}{\sqrt{1e^2}},$$
(24)
we can keep the effective open string coupling $`G_{eff}`$ finite.
In passing, we emphasize that the decoupling conditions $`(1)`$ and $`(3)`$ imply that the ratio between the electric and magnetic field strength is greater than $`1`$. In other words, in our decoupling scheme, the magnetic field is held to a finite value. (In fact, the parameter $`\theta _1`$ is just $`1/B`$). One may wonder what will be the NCOS scaling limit if one assumes $`|B|>|E|`$. The answer is that in this case, we do not have a consistent NCOS limit; rather we should take the NCYM limit, just as we have discussed in last section.
### B The E$``$B case
In this subsection, we study the other special case where the electric field is parallel to the magnetic field. The motivation is to show once more that the magnetic effects can survive the scaling limit for NCOS, resulting in space-space noncommutativity. To do so, we choose the closed string metric $`g_{\mu \nu }`$ and anti-symmetric tensor field $`B_{\mu \nu }`$ as
$$g_{\mu \nu }=\left(\begin{array}{cccc}g& 0& 0& 0\\ 0& g& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 1\end{array}\right),$$
(25)
$$B_{\mu \nu }=\left(\begin{array}{cccc}0& E& 0& 0\\ E& 0& 0& 0\\ 0& 0& 0& B\\ 0& 0& B& 0\end{array}\right).$$
(26)
Again, by using Seiberg-Witten relations Eqs. (6)-(9), we get the open string moduli
$$G_{\mu \nu }=\left(\begin{array}{cccc}g(1e^2)& 0& 0& 0\\ 0& g(1e^2)& 0& 0\\ 0& 0& 1+g^2b^2& 0\\ 0& 0& 0& 1+g^2b^2\end{array}\right),$$
(27)
$$G^{\mu \nu }=\left(\begin{array}{cccc}\frac{1}{g(1e^2)}& 0& 0& 0\\ 0& \frac{1}{g(1e^2)}& 0& 0\\ 0& 0& \frac{1}{1+g^2b^2}& 0\\ 0& 0& 0& \frac{1}{1+g^2b^2}\end{array}\right),$$
(28)
$$\theta ^{\mu \nu }=2\pi \alpha ^{}\left(\begin{array}{cccc}0& \frac{e}{g(1e^2)}& 0& 0\\ \frac{e}{g(1e^2)}& 0& 0& 0\\ 0& 0& 0& \frac{gb}{1+g^2b^2}\\ 0& 0& \frac{gb}{1+g^2b^2}& 0\end{array}\right),$$
(29)
$$G_s=g_s\sqrt{1e^2}\sqrt{1+g^2b^2}.$$
(30)
Here we adopted the same conventions for $`E_c`$, $`e`$, and $`b`$ as in the previous section. From the above open string moduli, we see that the same decoupling limit as that in $`𝐄𝐁`$ case can be applied. We also get the NCOS with both space-time and space-space noncommutativity. In contrast to the $`𝐄𝐁`$ case, the effects of the magnetic field is to increase the effective open string coupling constant $`G_s`$ by a factor $`\sqrt{1+g^2b^2}`$ after we take the large $`N`$ limit, without inducing a drift motion in other directions.
The most general configuration of $`B_{\mu \nu }`$ can be considered as a superposition of the two cases we have discussed, with $`𝐄𝐁`$ and $`𝐄𝐁`$ respectively. So we conclude that in general, by decoupling procedure, we can obtain NCYM with only space-space noncommutativity, or NCOS with both space-time and space-space noncommutativity.
## IV V-duality of NCOS
In the previous section, in the case with $`𝐄𝐁`$, we have managed to get a decoupling limit that leads to NCOS with both space-space and space-time noncommutativity, provided that $`|𝐄|`$ is greater than $`|𝐁|`$. In electrodynamics it is known that in this case by a Lorentz boost one can go to a favorable inertial frame in which the electromagnetic background becomes purely electric. If we start with this frame, the decoupling limit will give us an NCOS with only space-time noncommutativity. Before the decoupling limit, our string theory is known to have Lorentz symmetry, which allows us to transform the gauge field background on the D-brane world volume without changing the physics. So the above argument implies that the NCOS theory with both space-space and space-time noncommutativity that we obtained in the previous section for the case with $`𝐄𝐁`$ and $`|𝐄|>|𝐁|`$ should be equivalent to an NCOS with only space-time noncommutativity. More generally, this argument suggests that NCOS theories resulting from electromagnetic backgrounds on the D-brane that are related by Lorentz boosts should be equivalent to each other. This is a duality among NCOS with different open string moduli, and it is related to Lorentz boosts depending on the relative Velocity of the inertial frames. We call it V-duality, so that alphabetically it follows the S, T, U dualities we have had already.
An immediate question is: how V-duality acts on the open string moduli of NCOS? Now let us try to determine the orbit of the V-duality action in the moduli space of NCOS that we obtained in the last section. Let us start with two inertial frames $`K`$ and $`K^{}`$ on the world volume of $`D3`$-branes, with $`K^{}`$ moving relative to $`K`$ in $`x_2`$-direction with velocity $`v`$. Suppose the anti-symmetric tensor field $`B_{\mu \nu }`$ in $`K`$ is purely electric:
$$B_{\mu \nu }=\left(\begin{array}{cccc}0& E& 0& 0\\ E& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 0& 0\end{array}\right).$$
(31)
Then the corresponding $`B_{\mu \nu }^{}`$ in $`K^{}`$ has the form
$$B_{\mu \nu }^{}=\left(\begin{array}{cccc}0& E^{}& 0& 0\\ E^{}& 0& B^{}& 0\\ 0& B^{}& 0& 0\\ 0& 0& 0& 0\end{array}\right).$$
(32)
To relate $`E^{}`$ and $`B^{}`$ with $`E`$, we need to know the transformation between $`K`$ and $`K^{}`$. Note that we have taken the metric in both $`K`$ and $`K^{}`$ to be
$$ds^2=g(dx_0^2dx_1^2)+(dx_2^2+dx_3^2).$$
(33)
To make this metric invariant, the transformation should be the following ”adapted” Lorentz one:
$`x_2^{}`$ $`=`$ $`\gamma (x_2v\sqrt{g}x_0),`$ (34)
$`x_0^{}`$ $`=`$ $`\gamma (x_0{\displaystyle \frac{v}{\sqrt{g}}}x_2),`$ (35)
where $`\gamma =1/\sqrt{1v^2}`$. It is easy to check that the transformed metric is
$$g_{00}^{}=g_{11}^{}=g.$$
(36)
Using the invariance of the 2-form $`F=\frac{1}{2}B_{\mu \nu }dx^\mu dx^\nu `$, we get the transformed $`B_{\mu \nu }^{}`$ in $`K^{}`$ as
$`E^{}=\gamma E,B^{}={\displaystyle \frac{\gamma v}{\sqrt{g}}}E.`$ (37)
Note that in both $`K`$ and $`K^{}`$, the definition of $`E_c=g/2\pi \alpha ^{}`$ is the same. Thus we have the transformation law for dimensionless electric and magnetic fields:
$`e^{}{\displaystyle \frac{E^{}}{E_c}}=\gamma e,b^{}{\displaystyle \frac{B^{}}{E_c}}={\displaystyle \frac{\gamma v}{\sqrt{g}}}e.`$ (38)
To study the $`V`$-duality of NCOS, now we need to take the decoupling limit for NCOS in both frame $`K`$ and $`K^{}`$. In frame $`K`$, the decoupling limit dictates
$$g(1e^2)=\frac{2\pi \alpha ^{}}{\theta _0},$$
(39)
Correspondingly, in frame $`K^{}`$ we have
$$g(1e^2)=\frac{2\pi \alpha ^{}}{\theta _0^{}},$$
(40)
$$gb^{}=\frac{2\pi \alpha ^{}}{\theta _1^{}}.$$
(41)
Thus we can establish the following relation by using Eq. (41)
$$\gamma v\sqrt{g\frac{2\pi \alpha ^{}}{\theta _0}}=\frac{2\pi \alpha ^{}}{\theta _1^{}}.$$
(42)
The decoupling limit is the one in which
$`e1,g\mathrm{},`$ (43)
So taking the decoupling limit reduces Eq. (42) to
$$v\sqrt{g}=\frac{2\pi \alpha ^{}}{\theta _1^{}},$$
(44)
Therefore, we conclude that the boost velocity $`v0`$.
On the other hand, the boost transformation Eq. (40) leads to
$`{\displaystyle \frac{2\pi \alpha ^{}}{\theta _0^{}}}`$ $`=`$ $`g(1e^2)`$ (45)
$`=`$ $`g(1e^2)+ge^2(1\gamma ^2)`$ (46)
$`=`$ $`g(1e^2){\displaystyle \frac{gv^2e^2}{1v^2}}`$ (47)
$``$ $`{\displaystyle \frac{2\pi \alpha ^{}}{\theta _0}}\left({\displaystyle \frac{2\pi \alpha ^{}}{\theta _1^{}}}\right)^2,`$ (48)
where the arrow $`{}_{}{}^{\prime \prime }_{}^{\prime \prime }`$ means taking the decoupling limit. Thus, we have proved the $`V`$-duality action for NCOS
$$\frac{2\pi \alpha ^{}}{\theta _0}=\frac{2\pi \alpha ^{}}{\theta _0^{}}+\left(\frac{2\pi \alpha ^{}}{\theta _1^{}}\right)^2.$$
(49)
More generally, with other noncommutativity parameters vanishing, the following gives us an invariant under $`V`$-duality action:
$$\frac{2\pi \alpha ^{}}{\theta _0}+\left(\frac{2\pi \alpha ^{}}{\theta _1}\right)^2=\mathrm{invariant}.$$
(50)
This invariance gives us some orbits for $`V`$-duality. The displacements on an orbit are determined by the action of group elements. This invariance can be viewed as a descendant of the Lorentz invariance with a boost parameter $`v0`$, and this is the signal of the Galilean group. Therefore, we suggest that the $`V`$-duality should be characterized by a Galilean group or its deformation. The invariant (50) is for one of its abelian subgroup.
## V Holography in NCOS
Previously in Ref. a holographic correspondence between NCYM and its supergravity dual was suggested. Namely the radial profile of the on-shell close string moduli (string-frame metric, NSNS $`B`$-tensor and dilaton) in the supergravity dual of an NCYM can be easily derived through the Seiberg-Witten relations between close string moduli and open string moduli, provided a simple ansatz for the running string tension as the function of the energy scale is assumed. In this section, we generalize this link between holography and noncommutativity to NCOS.
For convenience of making a contrast between NCYM and NCOS, we first briefly recall the case of NCYM. Suppose only $`B_{23}0`$ on a stack of D3-branes. The central suggestion made in Ref. is that in the supergravity dual the UV limit (from the NCYM perspective) $`u\mathrm{}`$ is identified with the NCYM ”scaling limit” or ”decoupling limit” in Ref. . In this limit, $`\alpha ^{}`$ should approach zero, as in the AdS/CFT correspondence . To implement this, the overall factor $`R^2u^2`$, appeared in the 4d geometry along D3-branes, is interpreted as a running string tension
$$\alpha _{run}^{}=\frac{1}{R^2u^2},$$
(51)
which obviously runs to zero in the UV limit. Note that the manner it approaches zero compared to $`g_{22}`$ and $`g_{33}`$ agrees with the NCYM scaling limit taken in Ref. . The holographic correspondence suggested in Ref. is that the radial profiles of the on-shell NSNS fields in the gravity dual should satisfy the Seiberg-Witten relations Eqs. (6), (8) and (9), with $`\alpha ^{}`$ being replaced by the running $`\alpha _{run}^{}`$ given by Eq. (51) and with constant (unrenormalized) open string moduli.
In Ref. , the same holographic correspondence was shown to hold for all cases in which decoupling leads to an NCYM with space-space noncommutativity and with gravity dual known. These include high dimensional D$`p`$-branes in a magnetic background <sup>*</sup><sup>*</sup>*When $`p3`$, the open string (or NCYM) coupling constant is no longer $`u`$-independent: $`G_s^2=g^2u^{(7p)(p3)/2}`$. But this just means that the open string coupling runs in the same way as in the case when there is no B field, in agreement with the result of Ref. in the large-N limit. and Euclidean D3-branes in a self-dual $`B`$-background. In the following, we would like to examine whether a similar holographic correspondence holds as well between NCOS (with space-time noncommutativity) and its gravity dual, despite that the NCOS limit is very different from the NCYM limit.
Let us consider the case with only $`B_{01}0`$, In this case, because of the existence of a critical electric field on the D3-branes, to decouple the closed strings, one can no longer take $`\alpha ^{}0`$. Instead, $`\alpha ^{}`$ is fixed, leading to an NCOS. Certainly, the above ansatz Eq. (51) for the running string tension $`\alpha _{run}^{}`$ should no longer hold. We will see that indeed an appropriate modification of the ansatz exists, so that the above holographic correspondence remains to hold for NCOS.
The supergravity dual (with Lorentz signature) with only $`B_{01}`$ nonvanishing was given in Ref. :
$`ds_{str}^2`$ $`=`$ $`H(u)^{1/2}\left[{\displaystyle \frac{u^4}{R^4}}H(u)(dt^2+dx_1^2)+(dx_2^2+dx_3^2)+H(u)(du^2+u^2d\mathrm{\Omega }_5^2)\right],`$ (52)
$`B_{01}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle \frac{u^4}{R^4}},`$ (53)
$`e^{2\varphi }`$ $`=`$ $`g^2{\displaystyle \frac{u^4}{R^4}}H(u),`$ (54)
where we have $`\alpha ^{}=1`$, and $`R=4\pi gN`$, $`H(u)1+R^4/u^4`$. (Again, we omit the RR fields.) Recall that in the previous case with $`B_{23}0`$, the close string metric $`g_{ij}`$ (with $`i,j=2,3`$) shrinks to zero in the UV limit $`u\mathrm{}`$. In contrast, in the present case the close string metric $`g_{\mu \nu }`$ (with $`\mu ,\nu =0,1`$) goes to infinity in the UV limit, being consistent with the NCOS limit . So in the spirit of Ref. , we again identify the UV limit $`u\mathrm{}`$ with the NCOS “scaling limit” or “decoupling limit”, assuming the running string tension of the form
$$\alpha _{run}^{}=H(u)^{1/2}\left(1+\frac{R^4}{u^4}\right)^{1/2},$$
(55)
which is nothing but the inverse of the overall factor in front of the bracketInside the bracket the transverse metric $`g_{ij}`$ for $`i,j=2,3`$ is taken to be $`\delta _{ij}`$. in the close string metric in the gravity dual (52), in accordance with the same prescription as before for Eq. (51).
Now we want to show that the close string moduli in Eq. (52) can be derived from the Seiberg-Witten relations (6) and (9), in which $`\alpha ^{}`$ is replaced by a running one given by Eq. (55). Introduce the ansatz
$$g_{\mu \nu }=f(u)\eta _{\mu \nu },2\pi B_{\mu \nu }=h(u)ϵ_{\mu \nu },$$
(56)
for $`\mu ,\nu =0,1`$, due to the boost symmetry in the $`(x_0,x_1)`$ plane, and assuming constant open string moduli:
$$G_{\mu \nu }=\eta _{\mu \nu },\theta _{\mu \nu }=2\pi ϵ_{\mu \nu }.$$
(57)
Then two equations in Eq. (6) yield
$`1`$ $`=`$ $`f{\displaystyle \frac{h^2H}{f}},`$ (58)
$`1`$ $`=`$ $`{\displaystyle \frac{hH}{f^2h^2H}};`$ (59)
Namely,
$$f=f^2h^2H,f=hH.$$
(60)
Solving these two equations, one obtains
$`f(u)`$ $`=`$ $`\left(1H^1\right)^1={\displaystyle \frac{u^4}{R^4}}\left(1+{\displaystyle \frac{R^4}{u^4}}\right),`$ (61)
$`h(u)`$ $`=`$ $`{\displaystyle \frac{u^4}{R^4}}.`$ (62)
Similarly, substituting the above solution into the relation (9) with the identification $`G_s=g`$, one obtains the $`u`$-dependent closed string coupling:
$$g_s(u)=g\left(det(g+\alpha _{run}^{}B)\right)^{1/2},$$
(63)
or
$$e^{2\varphi }=g^2\left(1+\frac{u^4}{R^4}\right).$$
(64)
The results (61) and (64) are precisely what appeared in the gravity dual (52), which was previously obtained as a solution to classical equations of motion in IIB supergravity. Note that the closed string coupling approaches unity in the UV limit, in agreement with the decoupling of closed strings for NCOS.
Certainly this derivation adds more evidence to the universality of the link between holography and noncommutativity observed in Ref. . Thus, we have seen that the relations among the closed string moduli and the open string moduli contain much more than we could have imagined. With appropriate ansatz for the input $`\alpha _{eff}`$, they determine the closed string dual of both NCYM and NCOS! This demonstrates a simple and direct connection between holography and noncommutativity, either of which is believed to play a role in the ultimate theoretical structure for quantum gravity.
###### Acknowledgements.
One of us, GHC, thanks the Institute for Theoretical Physics, University of California at Santa Barbara, for an ITP Graduate Fellowship, and for the warm hospitality he receives during his stay. GHC also acknowledges stimulating discussions with Ian Low and Miao Li, while YSW thanks Feng-Li Lin for discussion. This research was supported in part by the National Science Foundation under Grants No. PHY94-07194(GHC) and PHY-9970701(YSW).
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# Spin Waves and Electronic Interactions in La2CuO4
## Abstract
The magnetic excitations of the square-lattice spin-1/2 antiferromagnet and high-Tc parent La<sub>2</sub>CuO<sub>4</sub> are determined using high-resolution inelastic neutron scattering. Sharp spin waves with absolute intensities in agreement with theory including quantum corrections are found throughout the Brillouin zone. The observed dispersion relation shows evidence for substantial interactions beyond the nearest-neighbor Heisenberg term, which can be understood in terms of a cyclic or ring exchange due to the strong hybridization path around the Cu<sub>4</sub>O<sub>4</sub> square plaquettes.
While there is consensus about the basic phenomenology - electron pairs with non-zero angular momentum, unconventional metallic behavior in the normal state, tendencies towards inhomogeneous charge and spin density order - of the high temperature copper oxide superconductors, there is no agreement about the microscopic mechanism. After over a decade of intense activity, there is not even consensus as to the simplest “effective Hamiltonian”, which is a short-hand description of the motions and interactions of the valence electrons, needed to account for cuprate superconductivity. Because much speculation is centered on magnetic mechanisms for the superconductivity, it is important to identify the interactions among the spins derived from the unfilled Cu<sup>2+</sup> $`d`$-shells. The present experiments show that there are significant (on the scale of the pairing energies for high-Tc superconductivity) interactions coupling spins at distances beyond the 3.8 Å separation of nearest-neighbor Cu<sup>2+</sup> ions. Cyclic or ring exchange due to a strong hybridization path around the Cu<sub>4</sub>O<sub>4</sub> squares (see Fig. 1A), from which the cuprates are built, provides a natural explanation for the measured dispersion relation. CuO<sub>2</sub> planes are thus the second example of an important Fermi system (<sup>3</sup>He is the other ) where significant cyclic exchange terms have been deduced.
Magnetic interactions are revealed through the wavevector dependence or dispersion of the magnetic excitations. In magnetically ordered materials, the dominant excitations are spin waves which are coherent (from site to site as well as in time) precessions of the spins about their mean values. The lower frame of Fig. 1B shows the dispersion relation calculated using conventional linear spin-wave theory in the classical large-$`S`$ limit, where the only magnetic interaction is a strong nearest-neighbor superexchange coupling $`J`$ . We identify wavevectors by their coordinates ($`h,k`$) in the two-dimensional (2D) reciprocal space of the square lattice. Spin waves emerge from the wavevector (1/2,1/2) characterizing the simple antiferromagnetic (AF) unit cell doubling in La<sub>2</sub>CuO<sub>4</sub> , and disperse to reach a maximum energy $`2J`$ that is a constant along the AF zone boundary marked by dashed squares in Fig. 1B. Longer-range interactions manifest themselves most simply at the zone boundary. The upper frame of Fig. 1B shows the dispersion calculated with modest interactions between next nearest-neighbors. Virtually the only visible effect of the additional interactions is the dispersion of the spin waves along the zone edge. Thus, experiments to test for such interactions must measure the spin waves along the zone boundary. Only inelastic neutron scattering with high energy and wavevector resolution can accomplish this, although photon spectroscopy has led to suspicions of such interactions.
For La<sub>2</sub>CuO<sub>4</sub>, a requirement that complicates meeting the resolution goals is the need to use neutrons with energies in the epithermal, 0.1-1.0 eV, range rather than in the more conventional cold and thermal , 2-50 meV, regimes. An early high energy neutron scattering experiment revealed well-defined spin-wave excitations throughout the Brillouin zone which could be modeled using a nearest-neighbor Heisenberg exchange $`J`$=136 meV. The directions of the scattered neutrons were specified only to within the solid angle determined by the large detector dimensions. Thus, the measured spectra represented averages over large portions of the reciprocal space, so that dispersion along the zone boundary was unresolvable and only an upper bound could be placed on further neighbor couplings. The advance enabling the present investigation is the use of position-sensitive detectors for the scattered neutrons, which increases the wavevector resolution by an order of magnitude. The new detector bank is installed in the direct-geometry High-Energy Transfer (HET) time-of-flight spectrometer at the ISIS proton-driven pulsed neutron spallation source.
Fig. 2A shows data in the form of constant energy scans for wavevectors around the antiferromagnetic zone center. As $`E`$ increases, counter-propagating modes become apparent.
As the zone boundary is approached and there is less dispersion, inspection of Fig. 1B reveals that it should be easier to locate the spin waves via energy scans performed at fixed wavevector. Fig. 2B shows a series of such scans collected at various points along the zone boundary. The spin waves have a clearly noticeable dispersion, from a minimum of 292$`\pm `$7 meV near $`𝑸`$=(3/4,1/4) to a maximum of 314$`\pm `$7 meV near (1/2,0). This is in obvious contrast to the dispersion-less behavior of linear spin-wave theory for the nearest-neighbor Heisenberg model. We have collected data throughout the Brillouin zone and Fig. 3A shows the resulting dispersion along major symmetry directions obtained from cuts of the type shown in Fig. 2. Fig. 3B displays the corresponding spin-wave intensities, in absolute units calibrated using acoustic phonon scattering from the sample.
To understand our results, we consider a Heisenberg Hamiltonian including higher order couplings
$``$ $`=`$ $`J{\displaystyle \underset{i,j}{}}𝑺_i𝑺_j+J^{}{\displaystyle \underset{i,i^{}}{}}𝑺_i𝑺_i^{}+J^{\prime \prime }{\displaystyle \underset{i,i^{\prime \prime }}{}}𝑺_i𝑺_{i^{\prime \prime }}`$ (3)
$`+J_c{\displaystyle \underset{i,j,k,l}{}}\{(𝑺_i𝑺_j)(𝑺_k𝑺_l)+(𝑺_i𝑺_l)(𝑺_k𝑺_j)`$
$`(𝑺_i𝑺_k)(𝑺_j𝑺_l)\},`$
where $`J`$, $`J^{}`$ and $`J^{\prime \prime }`$ are the first-, second- and third-nearest-neighbor magnetic exchanges where the paths are illustrated in Fig. 1A. $`J_c`$ is the ring exchange interaction coupling four spins (labelled clockwise) at the corners of a square plaquette. Each spin coupling is counted once in Eq. (3). Using classical (large-$`S`$) linear spin-wave theory the dispersion relation is $`\omega _𝑸`$=$`2Z_c(𝑸)\sqrt{A_𝑸^2B_𝑸^2}`$, $`A_𝑸`$=$`JJ_c/2(J^{}J_c/4)(1\nu _h\nu _k)J^{\prime \prime }\left[1(\nu _{2h}+\nu _{2k})/2\right]`$, $`B_𝑸`$=$`(JJ_c/2)(\nu _h+\nu _k)/2`$, $`\nu _x=\mathrm{cos}(2\pi x)`$ and $`Z_c(𝑸)`$ is a renormalization factor that includes the effect of quantum fluctuations. Within linear spin-wave theory all three higher-order spin couplings ($`J^{},J^{\prime \prime }`$ and $`J_c`$) have similar effects on the dispersion relation and intensity dependence, therefore they cannot be determined independently from the data without additional constraints. We first assume that only $`J`$ and $`J^{}`$ are significant as in , i.e. $`J^{\prime \prime }`$=$`J_c`$=0. The solid lines in Fig. 2 are fits to a one-magnon cross-section and Fig. 3 shows fits to the extracted dispersion relation and spin-wave intensity. As can be seen in the figures, the model provides an excellent description of both the spin-wave energies and intensities. The extracted nearest-neighbor exchange $`J`$=111.8$`\pm `$4 meV is antiferromagnetic, while the next-nearest-neighbor exchange $`J^{}`$=-11.4$`\pm `$3 meV across the diagonal is ferromagnetic. A wavevector-independent quantum renormalization factor $`Z_c`$=1.18 was used in converting spin-wave energies into exchange couplings. The zone-boundary dispersion becomes more pronounced upon cooling as shown in Fig. 3A and the dispersion at $`T`$=10 K can be described by the couplings $`J`$=104.1$`\pm `$4 meV and $`J^{}`$=-18$`\pm `$3 meV.
A ferromagnetic $`J^{}`$ contradicts theoretical predictions , which give an antiferromagnetic superexchange $`J^{}`$. Wavevector-dependent quantum corrections to the spin-wave energies can also lead to a dispersion along the zone boundary even if $`J^{}=0`$, but with sign opposite to our result. Another problem with a ferromagnetic $`J^{}`$ comes from measurements on Sr<sub>2</sub>Cu<sub>3</sub>O<sub>4</sub>Cl<sub>2</sub> . This material contains a similar exchange path between Cu<sup>2+</sup> ions to that corresponding to $`J^{}`$ in La<sub>2</sub>CuO<sub>4</sub> and analysis of the measured spin-wave dispersion leads to an antiferromagnetic exchange coupling for this path .
While we can cannot definitively rule out a ferromagnetic $`J^{}`$ we can obtain a natural description of the data in terms of a one-band Hubbard model , an expansion of which yields the spin Hamiltonian in Eq. (3) where the higher-order exchange terms arise from the coherent motion of electrons beyond nearest-neighbor sites . The Hubbard Hamiltonian has been widely used as a starting point for theories of the cuprates and is given by
$$=t\underset{i,j,\sigma =,}{}\left(c_{i\sigma }^{}c_{j\sigma }+\text{H.c.}\right)+U\underset{i}{}n_in_i,$$
(4)
where $`i,j`$ stands for pairs of nearest-neighbors counted once. Eq. (4) has two contributions: the first is the kinetic term characterized by a hopping energy $`t`$ between nearest-neighbor Cu sites and the second the potential energy term with $`U`$ being the penalty for double occupancy on a given site. At half filling, the case for La<sub>2</sub>CuO<sub>4</sub>, there is one electron per site and for $`t/U0`$, charge fluctuations are entirely suppressed in the ground state. The remaining degrees of freedom are the spins of the electrons localized at each site. For small but non-zero $`t/U`$, the spins interact via a series of exchange terms, as in Eq.(3), due to coherent electron motion touching progressively larger numbers of sites. If the perturbation series is expanded to order $`t^4`$ (i.e. 4 hops) one regains the Hamiltonian (3) with the exchange constants $`J=4t^2/U24t^4/U^3`$, $`J_c=80t^4/U^3`$ and $`J^{}=J^{\prime \prime }=4t^4/U^3`$ . We again fitted the dispersion and intensities of the spin-wave excitations using these expressions for the exchange constants and linear spin wave theory. The fits are indistinguishable from those for variable $`J`$ and $`J^{}`$. Again assuming $`Z_c`$=1.18, we obtained $`t`$=0.33$`\pm `$0.02 eV and $`U`$=2.9$`\pm `$0.4 eV ($`T`$=295 K), in agreement with $`t`$ and $`U`$ determined from photoemission and optical spectroscopy . The corresponding exchange values are $`J`$=138.3$`\pm `$4 meV, $`J_c`$=38$`\pm `$8 meV and $`J^{}`$=$`J^{\prime \prime }`$=$`J_c/20`$=2$`\pm `$0.5 meV (the parameters at $`T`$=10 K are $`t`$=0.30$`\pm `$0.02 eV, $`U`$=2.2$`\pm `$0.4 eV, $`J`$=146.3$`\pm `$4 meV and $`J_c`$=61$`\pm `$8 meV). Using these values, the higher-order interactions amount to $``$11% ($`T`$=295 K) of the total magnetic energy 2($`J`$-$`J_c`$/4-$`J^{}`$-$`J^{\prime \prime }`$) required to reverse one spin on a fully-aligned Nèel phase.
Many results on oxides of copper fall into place when cyclic exchange of the size extracted from our experiments is taken into account. First, the relative magnitude of the cyclic exchange $`J_c/J`$=0.27$`\pm 0.06`$ at $`T`$=295 K (0.41$`\pm `$0.07 at $`T`$=10 K) is similar to the ratio of 0.30 estimated from numerical simulations on finite clusters taken from the Cu-O square lattice. Second, magnetic Raman scattering and infrared absorption experiments show an unusual broadening towards higher energies that cannot be accounted for by a simple (quadratic) Heisenberg Hamiltonian, but can be attributed to a cyclic term. Finally, in the related compound Sr<sub>14</sub>Cu<sub>24</sub>O<sub>41</sub>, which has square plaquettes stacked to form a ladder, the exchange constants corresponding to the nearly-equal-length rungs and legs of the ladder are 130 meV and 72 meV respectively when no cyclic exchange is included. The inclusion of a ring exchange term $`J_c`$=34 meV allows the rungs and legs of the ladder to have similar exchange constants of 121 meV.
We have used a new high-wavevector-resolution epithermal-neutron scattering technique to discover that interactions beyond those coupling nearest-neighbor Cu<sup>2+</sup> ions are needed to account for the magnetism of La<sub>2</sub>CuO<sub>4</sub>. The observed further neighbor couplings may be explained by a four-spin cyclic interaction, which arises because the large orbital hybridization in the CuO<sub>2</sub> planes provides an exchange path to include all four spins at the corners of elementary Cu<sub>4</sub>O<sub>4</sub> plaquettes. Thus, La<sub>2</sub>CuO<sub>4</sub> joins the nuclear magnet <sup>3</sup>He as a system where there is good evidence for substantial ring exchange. Ring exchange occurs even for very simple models, such as the single band Hubbard model, which contains only hopping ($`t`$) and on-site Coulomb ($`U`$) terms. We have determined $`t`$ and $`U`$ using only knowledge - in the form of our spin wave dispersion relation - of charge-neutral excitations and find values in excellent agreement with those obtained via charge-sensitive spectroscopies as well as numerical work on finite clusters. Thus, our results demonstrate that a one-band Hubbard model is an excellent starting point for describing the magnetic interactions in the cuprates, and that even when considering relatively low energy spin excitations, charge fluctuations involving double occupancy must be taken into account. In addition, the scale of the cyclic exchange interactions, which are comparable to pairing energies in the high-Tc materials, implies that they themselves or related electronic currents might be important for superconductivity in the doped cuprates.
We are grateful to Risø National Laboratory for help in preparing this experiment and to J.F. Annett, H.-B. Braun, A.V. Chubukov, B. Roessli, H.M. Rønnow, G. Sawatzky, Z-X. Shen and R.R.P. Singh for very helpful discussions. ORNL is managed for the US DOE by UT-Battelle, LLC, under contract DE-AC05-00OR22725.
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# Discovery of a Classic FR-II Broad Absorption Line Quasar from the FIRST Survey
## 1 Introduction
It was long believed that a quasar could not be radio-loud and simultaneously have broad absorption lines in its optical spectrum (Stocke et al. 1992; Hamann, Korista, & Morris 1993; others). Over the past several years, the existence of quasars exhibiting both properties has been firmly established (Becker et al. 1997; Brotherton et al. 1998; White et al. 2000; Becker et al. 2000). We report here the most striking example to date of a radio-loud BAL quasar at z = 2.455, FIRST J101614.3+520916, which is not only extremely radio-loud but has a classic Fanaroff-Riley Type II (FR-II; Fanaroff & Riley, 1974) morphology with very bright radio lobes.
About 10% of the known quasars exhibit BAL outflows in their spectra (though the actual fraction depends heavily on how samples are selected, see Becker et al. 2000). This has been generally interpreted as BAL quasars being ordinary quasars viewed along a line of sight grazing the optically thick torus which surrounds the massive black hole engine (Weymann et al. 1991). The broad, highly blueshifted absorption is postulated to arise from clouds which are evaporated from the torus and accelerated outward by radiation pressure. In this picture, however, the apparent dichotomy of strong radio emission and the BAL phenomenon is extremely puzzling. There is no obvious reason in the orientation picture for radio emission to be suppressed, especially that from extended radio lobes. Additionally, typical BAL quasars have emission line and continuum spectral properties indistinguishable from ordinary quasars, in both the optical and radio (Stocke et al. 1992; Weymann et al. 1991; Barvainis & Lonsdale 1997). A number of possible explanations for the lack of radio-loud BALs – free-free absorption, frustrated jets, small scale structure in the accretion region – have been suggested (Stocke et al. 1992; Begelman, de Kool, & Sikora 1991; Boroson, Persson, & Oke 1985), but all are problematic. So perhaps it is not surprising that radio-loud BALs have finally been found.
This shift in understanding is being driven by studies which are probing new parts of parameter space: the FIRST Bright Quasar survey (FBQS; Gregg et al. 1996; White et al. 2000) relies on unprecedentedly faint radio fluxes of 1 mJy for candidate selection. The only other significant sample of radio-loud BALs, that of Brotherton et al. (1998) which found 5, also used a faint radio flux-limited sample (down to 2.5 mJy) selected from the NRAO VLA Sky Survey (NVSS, Condon et al. 1998) to select quasar candidates. These studies are the first which are radio-sensitive enough to probe the transition region between radio-quiet and radio-loud quasars with good statistics. The radio properties of the BAL quasars, both radio-quiet and loud, present serious challenges to the current understanding of the quasar phenomenon.
FIRST J101614.3+520916 (hereafter J1016+5209) is a BAL quasar which we have found in an ongoing project to extend the FBQS from a limiting magnitude of E=17.8 to E=19. With the discovery of this object, any remaining uncertainty over whether a BAL quasar can be radio-loud is completely removed. Further, because of its FR-II morphology, some constraints can be placed on the viewing angle and orientation of the system. The properties of J1016+5209, taken together with the general results of the BAL quasar study of Becker et al. (2000), may call for a fundamental alteration of the commonly accepted unification-by-orientation scheme, or at least how the BAL phenomenon fits into the general quasar model.
The only other object known which may be similar to J1016+5209 in combining the characteristics of BALs with FR-II radio morphology is the z=0.240 quasar PKS 1004+13 (Wills, Brandt, & Laor 1999). PKS 1004+13 is not as extreme in its radio-loudness nor in the strength of its broad absorption (about which there is some remaining doubt due to the low S/N of the defining IUE spectrum), but if real, there are now two members of this hybrid class.
## 2 Observations
A POSS-I stellar counterpart with E and O magnitudes of 18.6 and 20.2, respectively, lies within 0$`\stackrel{}{\mathrm{.}}`$25 of the radio source J1016+5209.<sup>1</sup><sup>1</sup>1These magnitudes are on the “corrected” APM system of White et al. (2000); photographic O and E are comparable to the more familiar B and R bands. Galactic reddening in this direction is insignificant, A$`{}_{V}{}^{}=0.005`$. A 9Å resolution optical spectrum was obtained using the Low Resolution Imaging Spectrograph (LRIS, Oke et al. 1995) at the Keck Observatory in December, 1998. A longer exposure but lower resolution ($``$ 20Å) LRIS spectrum was taken at Keck in June, 1999 (Figure 1). Both spectra clearly show prominent broad absorption features. The emission line redshift is 2.454 based on the fitted peak of C III\] 1909 which is unaffected by BAL features, though it may be contaminated by Fe III, Si III\], and Al III emission.
A second, higher S/N, 9Å resolution spectrum was obtained with LRIS in November, 1999. A close look at these data plotted in velocity space (Figure 2), shows an overall similarity between the C IV and Si IV BAL systems, though they differ in some details. Both C IV and Si IV exhibit a very broad system at $`15,000`$ $`\mathrm{km}\mathrm{s}^1`$and deeper but more complex absorption system at velocities from 0 to -8000 $`\mathrm{km}\mathrm{s}^1`$. The broad C IV feature extends from $`8500`$ $`\mathrm{km}\mathrm{s}^1`$to velocities of at least -17,200 $`\mathrm{km}\mathrm{s}^1`$from the line center, possibly to -20,000 $`\mathrm{km}\mathrm{s}^1`$, at a depth $`10\%`$ of the continuum. Both Si IV and C IV also have significant absorption to the red of the line centers, by $``$ 1000 $`\mathrm{km}\mathrm{s}^1`$. The C IV ($`\lambda \lambda `$ 1548.2, 1550.8) doublet is separated by only $`500`$ $`\mathrm{km}\mathrm{s}^1`$while the Si IV ($`\lambda \lambda `$ 1393.8, 1402.8) doublet is separated by a much greater 1920 $`\mathrm{km}\mathrm{s}^1`$; three obvious velocity systems seen in both species are marked in Figure 2. The LRIS spectral resolution is comparable to the redshifted separation of the C IV doublet (8.5Å), but even the relatively narrow absorption trough at -6500 $`\mathrm{km}\mathrm{s}^1`$ has a FWHM of 23Å, about twice as broad as the instrumental resolution convolved with a single C IV doublet, indicative of broad or multiple components. The Si IV doublet is easily split and the widths of the narrow individual lines are consistent with the instrumental resolution. In addition to the narrow features, the general continuum level in the velocity interval 0 to -10,000 $`\mathrm{km}\mathrm{s}^1`$is depressed by very broad absorption comparable in depth to the higher velocity BAL.
Using the C IV region, we have calculated the “BALnicity” index of J1016+5209 to be 2401 $`\mathrm{km}\mathrm{s}^1`$ (see Appendix A of Weymann et al. 1991). The BALnicity index (BI), though a quantitative measurement, is sensitive to subjective considerations, particularly continuum placement. For J1016+5209, the regions around -9000 and -5000 $`\mathrm{km}\mathrm{s}^1`$(Figure 2) are particularly important; if counted as continuum, then the BI drops to $`1600`$ $`\mathrm{km}\mathrm{s}^1`$, based solely on the very broad absorption extending from -17,200 to -8500 $`\mathrm{km}\mathrm{s}^1`$. Even this lower value is well within the ranks of other BAL quasars (Weymann et al. 1991).
As an ultimate test of its BAL nature, we obtained a high resolution spectrum of J1016+5209 in 2000 April, using the Echelle Spectrograph and Imager (ESI; Epps & Miller 1998) on the Keck II telescope. With a 1″ slit, the instrument delivers a dispersion of 0.15Å to 0.3Å per pixel over a wavelength range of 3900 to 10900Å, highly oversampling the $`1.5`$Å resolution spectrum. In Figure 3, we show the 4800 to 5500Å (restframe 1390 to 1590Å) region from the 1800s integration ESI spectrum, smoothed with a 9-pixel box, appropriate for the high oversampling, to improve the S/N. Overplotted is the LRIS spectrum from Figure 2. The S/N of the ESI spectrum is somewhat lower, and there is some additional structure in the depths of the BAL features, as might be expected from higher resolution data. But it is apparent that over the velocity range -17200 to -1500 $`\mathrm{km}\mathrm{s}^1`$ ($``$1460-1540Å rest wavelengths), none of the broad absorption breaks up into discrete narrow-line systems, confirming the BAL nature of J1016+5209. An example of what might be expected if the BAL regions in J1016+5209 did break up into discrete narrow-line clouds can be seen in the right panel inset of Figure 3 which shows the striking increase in resolution provided by the ESI spectrum for the intervening Mg II system, a truly narrow-line absorber. The only place where the BAL troughs do resolve into discrete components is at 5330Å (observed), but this is within $`1400`$ $`\mathrm{km}\mathrm{s}^1`$ of the C IV peak and hence does not contribute to the “BALnicity” index.
We have also obtained spectropolarimetry with Keck and LRIS in 2000 January, as part of a broader program to obtain spectropolarimetry for all of the FBQS BAL quasars. J1016+5209 is polarized at the 2.5% level, rising gradually from less than 2% at 8000Å (2350Å rest) to about 3% at 4200Å (1250Å rest). The polarization position angle varies from 85° to 75° over the same wavelength interval. These polarization characteristics are typical of BAL quasars (Hines & Wills 1995; Goodrich & Miller 1995; Cohen et al. 1995). We will discuss the polarization properties of J1016+5209 in more detail in a future paper.
A contour plot of the FIRST survey 1400 MHz radio image of J1016+5209 is displayed in Figure 3a. This field has the typical FIRST <sup>2</sup><sup>2</sup>2The FIRST Survey World Wide Web homepage is http://sundog.stsci.edu survey image characteristics: 5″ resolution and 0.15 mJy RMS noise. The core radio component associated with the quasar is marginally resolved, with a deconvolved size of 4″, and has a flux density of 6.5 mJy. It is bracketed by two bright radio sources of 131 and 39 mJy, both slightly resolved in the FIRST data. We interpret these as edge-brightened radio lobes, making J1016+5209 a classic triple radio source with FR-II morphology and a total flux density of 177 mJy. The contours in Figure 3a suggest that there is a physical connection between the core and the brighter lobe. The total angular distance between lobe centers is 45″ at a position angle of 146°. The radio source is bright enough to appear in several other radio surveys. The NVSS (Condon et al. 1998) lists a 20cm flux of 174 mJy for this object, indicating that FIRST adequately detects all of the flux and also that the source is probably not highly variable on timescales of a few years. The 92 cm WENSS survey (Rengelink et al. 1997) measured a total flux density of 850 mJy. The source is also detected in the 6 cm Greenbank (Becker, White, & Edwards 1991) survey with a total flux density of 44 mJy. The WENSS and Greenbank data yield a global spectral index of $`\alpha =1.1`$ (where $`S_\nu \nu ^\alpha `$), dominated by the bright lobes.
In 1999 July, a 0$`\stackrel{}{\mathrm{.}}`$25 resolution Very Large Array<sup>3</sup><sup>3</sup>3The Very Large Array is a facility of the National Radio Astronomy Observatory, operated by Associated Universities, Inc., under cooperative agreement with the National Science Foundation. (VLA) image of J1016+5209 was obtained in the A-configuration at 3.6 cm wavelength. The core is just marginally resolved with a fitted flux density of 1.84 mJy in this new image (RMS of 0.078 mJy); however, the northwest lobe shows an extended hotspot (Figure 3b) with a flux of 13.8 mJy and deconvolved size of $`0\stackrel{}{\mathrm{.}}59\times 0\stackrel{}{\mathrm{.}}15`$. These are lower limits as flux on scales greater than a few arcseconds will be resolved out. The position angle of the major axis of the resolved hotspot is $`140\stackrel{}{\mathrm{.}}3\pm 0\stackrel{}{\mathrm{.}}6`$; the position angle of the quasar from the hotspot location is a nearly identical 140$`\stackrel{}{\mathrm{.}}`$8, indicating that the radio lobe emanates from the quasar and is not a separate source. The southeast lobe is not reliably detected in the A-array data, probably because it is too diffuse.
We observed J1016+5209 yet again with the VLA in 1999 November using the B-configuration at a wavelength of 3.6 cm, this time obtaining polarization information as well (Figure 4). Because data were taken at only one frequency, we are unable to correct for Faraday rotation; however, since our measurements were at high frequency, this is expected to be small since the angle of rotation $`\theta \lambda ^2`$. In fact, the orientation of the magnetic field lines is as expected for a double-lobed FR-II source, parallel to the jet axis until reaching the hotspots, where it becomes perpendicular to the jet, indicating a shock-compressed field at the ends of the source.
Flux measurements at 8.46 GHz are: North hotspot/lobe = 18.6 mJy, Core = 2.1 mJy ($`4\%`$ polarized with a position angle of $`54\mathrm{°}\pm 9\mathrm{°}`$), South hotspot/lobe = 3.8 mJy. The lobe fluxes are lower limits as there will be some flux missing from the map on scales $`>10\mathrm{}`$. The core is strongly polarized, 4% at 8.4 GHz, perhaps greater if there is any depolarization. If the rotation measure is high, which could be the case if J1016+5209 is embedded in a thick shroud (see §3.1), depolarization of the core radio source could be significant. Comparing these flux density estimates to the FIRST survey numbers, the spectral indices are -1.10, -0.63, and -1.31 for the North lobe, Core, and South lobe, respectively.
## 3 Analysis and Discussion
The observed and derived properties of J1016+5209 are summarized in Table 1; we adopt H$`{}_{}{}^{}=50`$ $`\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ and q$`{}_{}{}^{}=0.5`$. The radio properties are extreme in several respects: J1016+5209 has a total radio luminosity of $`10^{34.4}`$ ergs cm<sup>-1</sup>s<sup>-1</sup>Hz<sup>-1</sup>, the highest of any known BAL quasar. Using the definition of Stocke et al. (1992), we compute log(R), the ratio of radio-to-optical brightness, as 3.4, using the global radio spectral index of -1.1 and optical spectral index of -1. This value of log(R) is also an extreme for known BAL quasars, and is at the high end of the distribution even for radio-selected non-BAL quasars (White et al. 2000). In fact, even if we consider only the core radio flux of J1016+5209, it is still radio-loud with log(R) = 2.0.
It is possible to estimate the angle between our line of sight and the jet axis in J1016+5209 using the “core-dominance” measure defined by Wills & Brotherton (1995) as $`\mathrm{log}()=\mathrm{log}(L_R)+0.4M_V13.69`$, where $`L_R`$ is the 5 GHz radio luminosity (of the core alone, Table 1) in units of W Hz<sup>-1</sup>. The $`M_B=26.2`$ (Table 1) (we have assumed the $`B`$ and $`O`$ passbands to be equivalent); we adopt $`BV=1`$ as a reasonable estimate of its color based on the $`OE`$ of 1.6. This yields $`\mathrm{log}()1.3`$, placing it at the extreme of orientations in the Wills & Brotherton sample where viewing angles are large, $`>40\mathrm{°}`$, but not well constrained. Even so, it is additional evidence that J1016+5209 is viewed well away from the jet axis.
### 3.1 The orientation model cannot explain BAL quasars
A popular explanation for the presence of BALs in a minority ($`10\%`$) of the quasar population is that quasars must be viewed at a preferred orientation to exhibit BALs in their spectra, along a line of sight roughly in the plane of the quasar accretion disk. To test this hypothesis, it is necessary to establish the orientation of BAL quasars; this can be done through radio observations. The FBQS is finding a surprisingly large number of BAL quasars for a radio-selected sample (Becker et al. 2000) with a frequency at least as great as that of the optically selected Large Bright Quasar Survey (Hewett, Foltz, & Chafee 1995). A large sample of radio-emitting BAL quasars offers the chance to investigate BAL viewing orientations.
The present sample of $`25`$ FBQS BAL quasars consists predominantly of unresolved compact radio sources, and, even the few that are slightly resolved still show no outright structure (such as lobes) on the angular scale of the FIRST survey (a few arcseconds). This differs markedly from the radio morphologies of a non-BAL subsample from the FBQS, matched in redshift and radio flux, where 30% of the non-BAL quasars show extended radio structure (Becker et al. 2000). Even without direct evidence from the radio morphology, some information on orientation can be obtained from radio spectra. The radio spectral indices of the 29 BALs in the FBQS vary widely, from -1.2 to +0.7, with 8 flatter than -0.5, and several more with indices of -0.5 (Becker et al. 2000). The scatter in spectral index is consistent with the findings of Barvainis & Lonsdale (1997) for the radio spectral indices of a smaller sample of 15 BAL quasars. If the bulk of the radio emission is from a relativistic jet, the objects with flat spectra are naturally interpreted as viewed close to the radio jet axis, while the steep-spectrum objects are those seen at larger viewing angles to the jet axis. This analysis is at odds with the model in which quasars must be viewed at a particular orientation to see BALs (e.g., Weymann, et al. 1991; Hines & Wills 1995; Goodrich & Miller 1995; Cohen et al. 1995).
Our results for J1016+5209 support the conclusion that a preferred viewing angle is not necessary to produce a BAL quasar. This object has unambiguous properties of both a BAL and a radio-loud FR-II. The radio spectral indices of the core and lobes exhibit the expected behavior, with the core flatter than the lobes, and the orientation is clearly well away from the jet axis, quite the opposite from the flat-spectrum BAL quasars. We conclude that a preferred line of sight is not necessary to observe BALs in quasars, and suggest that the alternative view that the BAL phenomenon is a stage, probably early, in the development of the quasar is more consistent with our data. Becker et al. (2000) speculated that the compact radio size of the BAL quasars implied a relatively young age for the BAL phenomenon: any radio jet is still in the process of emerging from a thick cocoon of material, and extended radio lobes on the 100 kpc scale have not had time to develop, as has been suggested by Briggs, Turnshek, & Wolfe (1984), Voit et al. (1993), Egami et al. (1996), and others. This suggests a possible link to compact steep spectrum (CSS) or gigahertz-peaked Spectrum (GPS) radio sources (e.g., O’Dea 1998; also see below, §3.2).
Hamann et al. (1993) argue from detailed modeling of BAL spectra that the covering factor $`q`$, the fraction of the sky covered by BAL regions as seen from the central source, is $`0.1`$. Coupled with the statistical result that $`10\%`$ of all quasars have BALs, their modeling result strongly implies that the conditions which give rise to the BAL phenomenon are present in every quasar and we simply do not see it 90% of the time because our line of sight does not pass through the BAL region. One possible way to reconcile this with the conclusion drawn above – that BAL quasars are not seen at any particular orientation – is to relax the usual assumption that BAL clouds are spatially concentrated near the plane of the obscuring torus surrounding the central engine. Arranging this could be problematic, however, if the BAL clouds originate in the torus region and are accelerated radially outward, which would naturally work to confine them to the plane of the torus.
Perhaps a more likely explanation of the radio results for the FBQS sample, which imply covering factors of approximately unity, is that the critical assumption of the Hamann et al. models that photons are conserved is not applicable. Voit et al. (1993) argue that BAL quasars which have very weak or absent C IV emission cannot be plausibly explained by small covering factors. Rather, the C IV photons are destroyed by repeated scatterings during their passage through a spherical shell of gas and dust. Such a shell does not even need to be very optically thick to reduce the C IV and other resonance emission lines such as Mg II to a negligible amount, as long as the dust and C IV ions are co-spatial. J1016+5209 certainly has weak resonance emission lines and its continuum is quite red for a quasar. Adopting the “starburst” reddening law of Calzetti et al. (1994), we estimate that J1016+5209 has A$`{}_{\mathrm{V}}{}^{}0.75`$, by comparing the shape of the rest frame continuum between 1600Å and 2200Å with that of a composite quasar spectrum (Brotherton et al. 2000) from the FBQS. Figure 5 displays the observed and dereddened spectrum of J1016+5209 and the quasar composite. The derived A<sub>V</sub> implies an optical depth from dust at 1550Å of $`1.6`$. In the simple scattering/absorption model that Voit et al. propose, this is sufficient to destroy the large majority of resonance line photons and provides a covering factor $`1`$.
The significant reddening from dust in J1016+5209 has further implications. First, high reddening is generally associated with BAL quasars which show absorption from low ionization species such as Mg II 2800, the “LoBAL” quasars (Sprayberry & Foltz 1992). Such objects generally also have strong Fe II emission (Weymann et al. 1991). Inspection of Figures 1 and 5 reveals that the Mg II emission is certainly much weaker in J1016+5209 than in the composite quasar, perhaps because of weak BALs. There is also a noticeable enhancement of Fe II emission to the blue of Mg II, which may in fact be partly responsible for filling in any possible Mg II BALs. In the higher resolution Keck spectrum of J1016+5209, there are weak absorption lines which correspond to Al III $`\lambda \lambda 1854.7,1862.8`$. The C III\] emission feature in J1016+5209 is broader and not as peaked as that of the composite quasar; this can be attributed to emission from Fe III $`\lambda \lambda 1895,1926`$. We conclude that J1016+5209 has some properties in common with LoBAL quasars.
All of these considerations again lead us to prefer a picture where BAL quasars are emerging from a dusty cocoon of material, probably at an early phase in their history. The statistic that BAL quasars make up 10% of the quasar population suggests that this phase lasts about 10% of the total quasar lifetime. As LoBALs are generally more highly reddened, they are an earlier period in the emergence of a quasar in this model than HiBALs.
Correcting for the dust extinction makes J1016+5209 brighter at B by 1.1 magnitude, and so reduces its radio-loudness from log(R) = 3.4 to 3.0. This still leaves it as the most radio-loud BAL known.
### 3.2 J1016+5209 as a Transition or Hybrid Object
J1016+5209 is the only FR-II quasar among the $`50`$ BAL quasars which have been discovered in follow-up to the FIRST survey ($`25`$ from Becker et al. 2000, plus an additional $`25`$ in subsequent follow-up, unpublished). In the FBQS, $`12\%`$ of z$`>0.5`$ quasars exhibit double-lobe morphology (Becker et al. 2000), and an additional $`10\%`$ show at least some radio structure. Why are BAL quasars with large radio lobes so rare? One possibility is that J1016+5209, and its potential low-z counterpart PKS 1004+13, are transition objects, on the way to becoming normal (non-BAL) FR-II quasars, caught in a relatively brief period during which the two phases co-exist. Another possibility is that J1016+5209 is a hybrid object, perhaps an old FR-II source which has recently been rejuvenated as a CSS/BAL source in its core.
Even though the high resolution ESI spectrum shows that the BAL features of J1016+5209 in general do not break up into myriad cloudlets, the lowest velocity trough does exhibit more structure in its depths than does the trough at -15000 $`\mathrm{km}\mathrm{s}^1`$, and at one location, v=-1400 $`\mathrm{km}\mathrm{s}^1`$, a narrow inter-cloud continuum is nearly resolved (Figure 3). The absorption within 9000 $`\mathrm{km}\mathrm{s}^1`$ of the emission line redshift is reminiscent of the more extreme examples of the class of “associated absorber” (AA) quasars (Foltz 1987). Were it not for its prominent BAL trough at -15000 $`\mathrm{km}\mathrm{s}^1`$, J1016+5209 might fall more naturally into the AA class, though it would be by far the most extreme example. A somewhat similar AA is PKS 1157+014 (Wright et al. 1979), a z=1.9, radio-loud quasar. It has two moderately broad absorption troughs at -6500 and 0 $`\mathrm{km}\mathrm{s}^1`$, not unlike the corresponding but more extreme spectral regions of J1016+5209, even though BI=0 for PKS 1157+014. Whether such AA features, which are not broad enough to gain distinction as true BALs in the quantitative BALnicity definition of Weymann et al. (1991), are intrinsic to the quasar or generated in an intervening object has been debated for some time (Morris et al. 1986; Foltz et al. 1986). Recently, Aldcroft, Bechtold, & Foltz (1998) presented evidence for variability of the higher velocity outflow in PKS1157+014, perhaps resolving the argument in favor of the intrinsic case, at least for this well-studied example. This is consistent with the growing evidence that many systems which are currently thought to be intervening, especially in radio-loud quasars, are really intrinsic (Richards et al. 1999).
It may be that J1016+5209 is in transition from a BAL to a more normal radio-loud quasar and PKS1157+014 is representative of the next evolutionary phase of a radio-loud object such as J1016+5209. If the highest velocity absorber in J1016+5209, already not as deep, were to fade away first, leaving behind the lower velocity troughs, the result would be similar to PKS1157+014. Perhaps we have just happened to catch J1016+5209 in a relatively rare, short-lived state in which it exhibits BAL features while having already developed strong radio emission. This could occur in a brief period at the end of the evolution of a BAL in which the radio emission finally manages to erupt from confinement but the dense cocoon has not completely dissipated. During such dissipation, the BALs may eventually evolve into distinct cloudlets as hinted at here, driven by the ensuing outflows of ionized plasma accompanying the radio emission. That the central region of J1016+5209 is completely surrounded by turbulent absorbing material is supported by the presence of absorption occurring at velocities to the red of the rest frame Si IV and C IV by $`1000`$ $`\mathrm{km}\mathrm{s}^1`$.
This picture is consistent with the radio core of J1016+5209 having an unusually steep spectrum, $`\alpha =0.63`$, and being unresolved at the 0$`\stackrel{}{\mathrm{.}}`$25 ($`2`$ kpc) scale. These properties are reminiscent of CSS or GPS sources: the leading interpretation of CSS and GPS sources is that they are young radio objects, confined to a small region by dense gas but which evolve with time into extended radio sources with lobes as they escape confinement (O’Dea 1998), much like the picture of BAL quasars emerging from cocoons (Voit et al. 1993). This coincidence of attributes in J1016+5209 supports the notion that BAL quasars are an early evolutionary phase in the life cycle of a quasar. The polarization of CSS sources is typically higher than that of GPS sources, $`5\%vs\mathrm{.\; 0.2}\%`$ at 6cm (O’Dea 1998), so the $`4\%`$ polarization at 3.6cm for the core of J1016+5209 suggests that it is a CSS object, but low frequency data for the core alone are needed to confirm this. If it is a CSS, then J1016+5209 may exhibit the so-called “alignment effect” between its radio and optical structure (McCarthy 1993 and references therein); any possible connection with its BAL nature would be interesting in this context. The misalignment of the optical polarization and large-scale radio jet axes could be explained if on subarcsecond scales J1016-5209 has been reborn with a different jet axis.
Steep-spectrum cores, however, are not uncommon in high redshift, lobe-dominated quasars (Lonsdale, Barthel, & Miley 1993). An alternative possibility is that J1016+5209 is a normal FR-II quasar in a very low density environment which allows rapid expansion of the radio lobes. If the radio lobes were expanding unimpeded at relativistic speeds, then the brighter, jet-side lobe should be significantly farther from the core, whereas just the opposite is seen. The arm-length ratio for J1016+5209, however, is $`Q=0.64`$, at the extreme low end of the distribution found by Scheuer (1995) for a sample of radio-luminous, double-lobed quasars. The asymmetry of J1016+5209 then is probably due to environmental rather than relativistic effects, implying that the lobes are not expanding freely and rapidly, and hence are not particularly young. If the radio source is expanding at speeds typical of FR-II radio sources, $`0.1c`$ (Arshakian & Longair 2000), the large extent of the lobes of J1016+5209, $`350`$ kpc, suggests a fairly advanced age (for a radio source) of $`10^7`$ yr; PKS 1004+413 is also a large source, $`475`$kpc in size, and so of comparable age.
If the core of J1016+5209 is a CSS or GPS object, the presence of larger scale, presumably older, very bright radio lobes (Figures 3 and 4) at a large distance from the central engine supports the hypothesis that quasars can be “reborn” and that perhaps both the BAL and CSS/GPS properties can occur repeatedly in a given object, but always early in any “on” cycle of AGN activity. In support of the rejuvenation picture, about 10% of GPS/CSS sources have extended emission (O’Dea 1998 and references therein), possibly from an earlier period of activity, now dissipated. It may be that J1016+5209 is in an early phase of rejuvenation, having particularly compact inner lobes which will grow, becoming a double-double radio source; a number of such objects are known (Schoenmakers et al. 2000). Higher resolution mapping of the core of the J1016+5209 is needed to test its GPS/CSS nature.
The interpretation of J1016+5209 as a rejuvenated quasar suggests that it may not be correct to compute its radio-loudness using the entire radio flux, at least not in the context of evaluating the “BAL-related” radio-loudness in its present incarnation. With the reddening correction and counting only the core flux, log(R$`{}_{}{}^{})1.6`$, still formally radio-loud, but not as exceptional as $`3`$ for the total radio flux.
## 4 Summary
The properties of the quasar FIRST J1016+5209 stand out in several respects. It exhibits bona fide BALs in its optical spectrum while also having a classic FR-II radio-loud morphology. J1016+5209 is the most radio-loud and radio-luminous BAL quasar known. The only other object which may be of a similar nature is the less extreme, low redshift quasar PKS 1004+13 (Wills et al. 1999). The presence of distinct bright radio lobes and its low “core dominance” parameter implies that J1016+5209 is viewed well away from the jet axis, at an angle of $`40^{}`$. Based on the large scatter in radio spectral indices, Becker et al. (2000) argue that BAL quasars are not viewed at any particular orientation, contrary to the popular orientation model. The relatively steep spectrum ($`\alpha 0.6`$) and compact size ($`<0\mathrm{}3`$ at 3.6cm) of the radio core of J1016+5209 suggest that it is a CSS source, suggesting that J1016+5209 is young. This supports the alternate model of BAL quasars in which they are an early phase in the evolution of quasars. The 20°–30° misalignment of the optical and radio polarization axes is further evidence that J1016+5209 does not easily fit the orientation model for BAL quasars.
The large scale (350 kpc) FR-II radio lobes of J1016+5209 do not easily fit the picture of it being young, so we postulate that it is a rejuvenated quasar, possibly through a merger or interaction. If there is a newly created – perhaps even episodic – CSS source at the core of J1016+5209, higher resolution imaging at various wavelengths should reveal interesting connections among the various attributes (BALs, CSS, radio-loudness, FR-II morphology) that have come together in this one object.
We gratefully acknowledge D. Stern and H. Spinrad for obtaining the low resolution Keck spectrum of FIRST J1016+5209 and Willem De Vries for helpful comments. Some of the data presented here were obtained at the W.M. Keck Observatory, which is operated as a scientific partnership among the California Institute of Technology, the University of California and the National Aeronautics and Space Administration. The Keck Observatory was made possible by the generous financial support of the W.M. Keck Foundation. The FIRST Survey is supported by grants from the National Science Foundation (grant AST-98-02791), NATO, the National Geographic Society, Sun Microsystems, and Columbia University. Part of the work reported here was done at the Institute of Geophysics and Planetary Physics, under the auspices of the U.S. Department of Energy by Lawrence Livermore National Laboratory under contract No. W-7405-Eng-48.
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# Correlation property of length sequences based on global structure of complete genome
## I Introduction
Recently, there has been considerable interest in the finding of long-range correlation (LRC) in DNA sequences $`[116]`$. Li et al$`^{\text{[1]}}`$ found that the spectral density of a DNA sequence containing mostly introns shows $`1/f^\beta `$ behaviour, which indicates the presence of LRC. The correlation properties of coding and noncoding DNA sequences were first studied by Peng et al$`^{\text{[2]}}`$ in their fractal landscape or DNA walk model. The DNA walk defined in is that the walker steps “up” if a pyrimidine ($`C`$ or $`T`$) occurs at position $`i`$ along the DNA chain, while the walker steps “down” if a purine ($`A`$ or $`G`$) occurs at position $`i`$. Peng et al$`^{\text{[2]}}`$ discovered that there exists LRC in noncoding DNA sequences while the coding sequences correspond to a regular random walk. By doing a more detailed analysis, Chatzidimitriou, Dreismann and Larhammar$`^{\text{[5]}}`$ concluded that both coding and noncoding sequences exhibit LRC. A subsequent work by Prabhu and Claverie$`^{\text{[6]}}`$ also substantially corroborates these results. If one considers more details by distinguishing $`C`$ from $`T`$ in pyrimidine, and $`A`$ from $`G`$ in purine (such as two or three-dimensional DNA walk model$`^{\text{[8]}}`$ and maps given in ), then the presence of base correlation has been found even in coding sequences. In view of the controversy about the presence of correlation in all DNA or only in noncoding DNA, Buldyrev et al$`^{\text{[14]}}`$ showed the LRC appears mainly in noncoding DNA using all the DNA sequences available. Alternatively, Voss$`^{\text{[10]}}`$, based on equal-symbol correlation, showed a power law behaviour for the sequences studied regardless of the percent of intron contents. Investigations based on different models seem to suggest different results, as they all look into only a certain aspect of the entire DNA sequence. It is therefore important to investigate the degree of correlations in a model-independent way.
Since the first complete genome of the free-living bacterium Mycoplasma genitalium was sequenced in 1995$`^{\text{[17]}}`$, an ever-growing number of complete genomes has been deposited in public databases. The availability of complete genomes induces the possibility to ask some global questions on these sequences. The avoided and under-represented strings in some bacterial complete genomes have been discussed in . A time series model of CDS in complete genome has also been proposed in . Maria de Sousa Vieira$`^{\text{[22]}}`$ have done the low-frequency analysis of complete DNA of 13 microbial genomes and shown that its fractal behaviour not always prevails through the entire chain, the autocorrelation function have a rich variety of behaviours including the presence of anti-correlations.
For the importance of the numbers, sizes and ordering of genes along the chromosome, one can refer to Part 5 of the famous book of Lewin (Ref.). Hence one may ignore the composition of the four kinds of bases in coding and noncoding segments and only considers the rough structure of the complete genome or long DNA sequences. Provata and Almirantis $`^{\text{[24]}}`$ proposed a fractal Cantor pattern of DNA. They map coding segments to filled regions and noncoding segments to empty regions of random Cantor set and then calculate the fractal dimension of the random fractal set. They found that the coding/noncoding partition in DNA sequences of lower organisms is homogeneous-like, while in the higher eucariotes the partition is fractal. This result seems too rough to distinguish bacteria because the fractal dimensions of bacteria they gave out are all the same. The classification and evolution relationship of bacteria is one of the most important problem in DNA research. Yu and Anh$`^{\text{[25]}}`$ proposed a time series model based on the global structure of the complete genome and considered three kinds of length sequences. After calculating the correlation dimensions and Hurst exponents, it was found that one can get more information from this model than that of fractal Cantor pattern. Some results on the classification and evolution relationship of bacteria were found in . Naturally it is desirable to know if there exists LRC in these length sequences. The quantification of these correlations could give an insight of the role of the ordering of genes on the chromosome, which is far to be irrelevant for gene function. We attempt to answer this question in this paper.
Viewing from the level of structure, the complete genome of an organism is made up of coding and noncoding segments. Here the length of a coding/noncoding segment means the number of its bases. Based on the lengths of coding/noncoding segments in the complete genome, we can get three kinds of integer sequences by the following ways.
i) First we order all lengths of coding and noncoding segments according to the order of coding and noncoding segments in the complete genome, then replace the lengths of noncoding segments by their negative numbers. This allows to distinguish lengths of coding and noncoding segments. This integer sequence is named whole length sequence.
ii) We order all lengths of coding segments according to the order of coding segments in the complete genome. We name this integer sequence coding length sequence.
iii) We order all lengths of noncoding segments according to the order of noncoding segments in the complete genome. This integer sequence is named noncoding length sequence.
We can now view these three kinds of integer sequences as time series. In the following, we will investigate the correlation property through Detrended Fluctuation Analysis (DFA)$`^{\text{[26]}}`$ and spectral analysis.
## II Detrended fluctuation analysis and spectral analysis
We denote a time series as $`X(t),t=1,\mathrm{},N`$. First the time series is integrated as $`y(k)=_{t=1}^k[X(t)X_{ave}]`$, where $`X_{ave}`$ is the average over the whole time period. Next, the integrated time series is divided into boxes of equal length, $`n`$. In each box of length $`n`$, a least-squares line is fit to the data, representing the trend in that box. The $`y`$ coordinate of the straight line segments is denoted by $`y_n(k)`$. We then detrend the integrated time series, $`y(k)`$, by subtracting the local trend, $`y_n(k)`$, in each box. The root-mean-square fluctuation of this integrated and detrended time series is calculated as
$$F(n)=\sqrt{\frac{1}{N}\underset{k=1}{\overset{N}{}}[y(k)y_n(k)]^2}$$
(1)
Typically, $`F(n)`$ will increase with box size $`n`$. A linear relationship on a double log graph indicates the presence of scaling
$$F(n)n^\alpha .$$
(2)
Under such conditions, the fluctuations can be characterised by the scaling exponent $`\alpha `$, the slope of the line relating $`\mathrm{ln}F(n)`$ to $`\mathrm{ln}n`$. For uncorrelated data, the integrated value, $`y(k)`$ corresponds to a random walk, and therefore, $`\alpha =0.5`$. An value of $`0.5<\alpha <1.0`$ indicates the presence of LRC so that a large interval is more likely to be followed by a large interval and vice versa. In contrast, $`0<\alpha <0.5`$ indicates a different type of power-law correlations such that large and small values of time series are more likely to alternate. For examples, we give the DFA of the coding length sequence of A. aeolicus in FIG. 1.
Now we analyse the time series using the quantity $`M(L)`$, the mean distance a walker spanned within time $`L`$. Dunki and Ambuhl$`^{\text{[27, 28]}}`$ used this quantity to discuss the scaling property in temporal patterns of schizophrenia. Denote
$$W(j):=\underset{t=1}{\overset{j}{}}[X(t)X_{ave}],$$
(3)
from which we get the walks
$$M(L):=<|W(j)W(j+L)|>_j,$$
(4)
where $`<>_j`$ denotes the average over $`j`$, and $`j=1,\mathrm{},NL`$. The time shift $`L`$ typically varies from $`1,\mathrm{},N/2`$. From a physics viewpoint, $`M(L)`$ might be thought of as the variance evolution of a random walker’s total displacement mapped from the time series $`X(t)`$. $`M(L)`$ may be assessed for LRC $`^{\text{[29]}}`$ (e.g. $`M(L)L^\alpha ^{}`$, $`\alpha ^{}=1/2`$ corresponding to the random case). We give some examples to estimate the scale parameter $`\alpha ^{}`$ in FIG. 2.
Dunki et al$`^{\text{[28]}}`$ proposed the following scale which seems to perform better than the scale $`\alpha ^{}`$. The definition
$$W^{}(j):=\underset{t=1}{\overset{j}{}}|X(t)X_{ave}|$$
(5)
leads to
$$M^{}(L):=<|W^{}(j)W^{}(j+L)|>_j.$$
(6)
Analyses of test time series showed that (6) are more robust against distortion or discretization of the corresponding amplitudes $`X(t)`$ than (4). From the $`\mathrm{ln}(L)`$ v.s. $`\mathrm{ln}(M^{}(L))`$ plane, we find the relation
$$M^{}(L)L^\gamma .$$
(7)
The exponent $`\gamma `$ measures only the presence of nonlinear correlations and remains equal to unity for all sequences with only linear correlations. We carried out this kind of analysis on coding length sequences of A. aeolicus, B. burgdorferi and T. maritima. The results are reported in the left figure of FIG. 3.
We also consider the discrete Fourier transform$`^{\text{[30]}}`$ of the time series $`X(t),t=1,\mathrm{},N`$ defined by
$$\widehat{X}(f)=N^{\frac{1}{2}}\underset{t=0}{\overset{N1}{}}X(t+1)e^{2\pi ift},$$
(8)
then
$$S(f)=|\widehat{X}(f)|^2$$
(9)
is the power spectrum of $`X\left(t\right)`$. In recent studies, it has been found $`^{\text{[31]}}`$ that many natural phenomena lead to the power spectrum of the form $`1/f^\beta `$. This kind of dependence was named $`1/f`$ noise, in contrast to white noise $`S(f)=const`$, i.e. $`\beta =0`$. Let the frequency $`f`$ take $`k`$ values $`f_k=k/N,k=1,\mathrm{},N`$. From the $`\mathrm{ln}(f)`$ v.s. $`\mathrm{ln}(S(f))`$ graph, the existence of $`1/f^\beta `$doesn’t seem apparent. For example, we give the figure of the coding length sequence of A. aeolicus on the right of FIG. 3.
When we use the least squares line to fit data, we need to consider the errors. If the data are $`\{(x_i,y_i)\}_{i=1}^n`$, we can define the coefficient of linear correlation as$`^{\text{[32]}}`$
$$r=\frac{_{i=1}^n(x_i\overline{x})(y_i\overline{y})}{\sqrt{(}_{i=1}^n(x_i\overline{x})^2_{i=1}^n(y_i\overline{y})^2)},$$
(10)
where $`\overline{x}`$ and $`\overline{y}`$ are the average of the values $`\{x_i\}_{i=1}^n`$ and $`\{y_i\}_{i=1}^n`$ respectively. If $`r=\pm 1`$, then the points lie exactly on a straight line; that is, there is a perfect linear relationship between $`x`$ and $`y`$. If $`r=0`$, there is no linear relationship. The quantity $`r`$ measures the strength of linear relationships between $`x`$ and $`y`$. The values of $`r`$ in figures of obtaining exponents $`\alpha `$, $`\alpha ^{}`$ and $`\beta `$ are 0.987685, 0.9949939 and 3.7918E-03 respectively. Hence we can see the informations given by exponents $`\alpha `$ and $`\alpha ^{}`$ are more convincing than that given by exponent $`\beta `$.
## III Data and results.
More than 21 bacterial complete genomes are now available in public databases. There are five Archaebacteria: Archaeoglobus fulgidus (aful), Pyrococcus abyssi (pabyssi), Methanococcus jannaschii (mjan), Aeropyrum pernix (aero) and Methanobacterium thermoautotrophicum (mthe); four Gram-positive Eubacteria: Mycobacterium tuberculosis (mtub), Mycoplasma pneumoniae (mpneu), Mycoplasma genitalium (mgen), and Bacillus subtilis (bsub). The others are Gram-negative Eubacteria. These consist of two Hyperthermophilic bacteria: Aquifex aeolicus (aquae) and Thermotoga maritima (tmar); six Proteobacteria: Rhizobium sp. NGR234 (pNGR234), Escherichia coli (ecoli), Haemophilus influenzae (hinf), Helicobacter pylori J99 (hpyl99), Helicobacter pylori 26695 (hpyl) and Rickettsia prowazekii (rpxx); two Chlamydia: Chlamydia trachomatis (ctra) and Chlamydia pneumoniae (cpneu), and two Spirochaete: Borrelia burgdorferi (bbur) and Treponema pallidum (tpal).
We calculate scales $`\alpha `$, $`\alpha ^{}`$, $`\beta `$ of low frequencies ($`f<1`$) and $`\gamma `$ of three kinds of length sequences of the above 21 bacteria. The estimated results are given in Table I ( where we denote by $`\alpha _{whole}`$, $`\alpha _{cod}`$ and $`\alpha _{noncod}`$ the scales of DFA of the whole, coding and noncoding length sequences, from top to bottom, in the increasing order of the value of $`\alpha _{whole}`$ ), Table II ( where we denote by $`\alpha _{whole}^{}`$, $`\alpha _{cod}^{}`$ and $`\alpha _{noncod}^{}`$ the scales of $`M(L)`$ of the whole, coding and noncoding length sequences, from top to bottom, in the increasing order of the value of $`\alpha _{whole}^{}`$ ) and Table III ( where we denote by $`\beta _{whole}`$, $`\beta _{cod}`$ and $`\beta _{noncod}`$ the scales of spectral analysis of the whole, coding and noncoding length sequences, from top to bottom, in the decreasing order of the value of $`\beta _{whole}`$; we denote by $`\gamma _{whole}`$, $`\gamma _{cod}`$ and $`\gamma _{noncod}`$ the scales of $`\gamma `$ of the whole, coding and noncoding length sequences).
From the right figure of FIG. 3 it is seen that $`S(f)`$ does not display clear power-law $`1/f`$ dependence on the frequencies when $`f<1`$. Although the meaning of region $`f>1`$ of the power spectrum is not clear, whether $`S(f)`$ displays perfect power-law $`1/f`$ in this region is important. When one considers the electrical characteristics of polysilicon emitter bipolar transistors, for high frequency analog applications the transistor $`1/f`$ noise is also an important parameter since it can degrade the spectral purity of the circuit$`^{\text{[33]}}`$. There is also some evidence that $`1/f`$ noise spectral density in the low and in the high current region have a different physical origion (the reader can refer Ref. and reference therein). We want to know if there is another region of frequencies in which $`S(f)`$ displays perfect power-law $`1/f`$ dependence on the frequencies. We carried out the spectral analysis for $`f>1`$, and found that $`S(f)`$ displays almost a perfect power-law $`1/f`$ dependence on the frequencies in some interval:
$$S(f)\frac{1}{f^\beta }.$$
(11)
We give the results for coding length sequences of M. genitalium, A. fulgidus, A. aeolicus and E. coli (their lengths are 303, 1538, 891 and 3034 respectively) in FIG. 4, where we take $`k`$ values $`f_k=3k(k=1,\mathrm{},1000)`$ of the frequency $`f`$. From FIG. 4, it is seen that the length of the interval of frequency in which $`S(f)`$ displays almost a perfect power-law $`1/f`$ depends on the length of the length sequence. The shorter sequence corresponds to the larger interval.
From FIG. 4, one can see that the power spectrum exhibit some kind of periodicity. But the period seems to depend on the length of the sequence. We also guess that the period also depends on the complexity of the sequence. To support this conjecture, we got a promoter DNA sequence from the gene bank, then replaced $`A`$ by -2, $`C`$ by -1, $`G`$ by 1 and $`T`$ by 2 (this map is given in ); so we obtained a sequence on alphabet $`\{2,1,1,2\}`$. Then a subsequences was obtained with the length the same as the coding length sequences of A. aeolicus, A. fulgidus and M. genitalium (their lengths are 891, 1538 and 303 respectively). A comparison is given in FIG. 5, but the results are not clear-cut.
## IV Discussion and conclusions
Although the existence of the archaebacterial urkingdom has been accepted by many biologists, the classification of bacteria is still a matter of controversy$`^{\text{[34]}}`$. The evolutionary relationship of the three primary kingdoms (i.e. archeabacteria, eubacteria and eukaryote) is another crucial problem that remains unresolved$`^{\text{[34]}}`$.
From Table I, we can roughly divide bacteria into two classes, one class with $`\alpha _{whole}`$ less than 0.5, and the other with $`\alpha _{whole}`$ greater than 0.5. All Archaebacteria belong to the same class except Pyrococcus abyssi. All Proteobacteria belong to the same class except E. coli; in particular, the closest Proteobacteria Helicobacter pylori 26695 and Helicobacter pylori J99 group with each other. In the class with $`\alpha _{whole}<0.5`$, we have $`\alpha _{cod}<\alpha _{noncod}`$ except H. pylori J99 and M. genitalium; but in the other class we have $`\alpha _{cod}>\alpha _{noncod}`$.
Using the exponent $`\alpha ^{}`$, we can also divide bacteria into two class as in Table II. In one class, $`\alpha _{cod}^{}<\alpha _{noncod}^{}`$. In another class, we have $`\alpha _{cod}^{}>\alpha _{noncod}^{}`$ except Treponema pallidum and Borrelia burgdorferi. Two Hyperthermophilic bacteria Aquifex aeolicus and Thermotoga maritima group with each other.
From Tables I and II, we can see the similar rules as above if we use the exponents $`\alpha _{cod}`$ and $`\alpha _{cod}^{}`$. This follows the fact that the coding sequences occupy the main part of space of the DNA chain of bacteria. This coincides with the conclusion of Ref..
Although from Table III, we can see the values of all $`\beta `$ are not far from $`0`$. From FIG.s 1, 2 and 3, one can see exponents $`\alpha `$ and $`\alpha ^{}`$ are more convincing than the exponent $`\beta `$ because the error of estimating $`\alpha `$ and $`\alpha ^{}`$ using the least-squares linear fit is much less than that of the exponent $`\beta `$ (The values of $`r`$ in figures of obtaining exponents $`\alpha `$, $`\alpha ^{}`$ and $`\beta `$ are 0.987685, 0.9949939 and 3.7918E-03 respectively). From Tables I and II, we can see most of values $`\alpha `$ and $`\alpha ^{}`$ are not equal to 0.5, hence we can conclude that most of these length sequences exhibit long-range correlations. We can also see the correlation have a rich variety of behaviours including the presence of anti-correlations. Hence the length sequences have a same character as the DNA sequences$`^{\text{[22]}}`$. Further more, from Table III, we get $`\gamma =1.0\pm 0.03`$. Hence we can conclude that the long-range correlations exist in most length sequences are linear.
We find in an interval of frequency ($`f>1`$), $`S(f)`$ displays perfect power-law $`1/f`$ dependence on the frequencies (see FIG. 4)
$`S(f){\displaystyle \frac{1}{f^\beta }}.`$
The length of the interval of frequency in which $`S(f)`$ displays almost a perfect power-law $`1/f`$ depends on the length of the length sequence. The shorter sequence corresponds to the larger interval. The shape of the graph of power spectrum in $`f>1`$ also exhibits some kind of periodicity. The period seems to depend on the length and the complexity of the length sequence.
## ACKNOWLEDGEMENTS
Authors Zu-Guo Yu and Bin Wang would like to express their thanks to Prof. Bai-lin Hao of Institute of Theoretical Physics of Chinese Academy of Science for introducing them into this field and continuous encouragement. They also wants to thank Dr. Guo-Yi Chen of ITP for useful discussions. Research is partially supported by Postdoctoral Research Support Grant No. 9900658 of QUT. The authors also want to thank the referee for telling the property of exponent $`\gamma `$ and many useful suggestions to improve this paper.
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# Sinks of Light Elements in Stars - Part II
## 1. Introduction
The study of lithium in metal-poor stars has implications for stellar structure, galactic chemical evolution, and Big Bang nucleosynthesis (BBN). There has thus been considerable observational and theoretical effort on this question, and we will therefore attempt to draw together the main themes and outstanding questions rather than undertake a detailed analysis of the fine points. This paper will be organized as follows. In section 2, we recall the main trends of the theoretical expectations. In section 3, we discuss the major features of the observational data; in particular we stress the similarities and differences with the population I pattern. Section 4 is devoted to the comparison of different classes of theoretical models with the data, and our conclusions are summarized in section 5.
## 2. Main Trends of the Theoretical Expectations
Let us first briefly recall that the so-called standard case refers to models which exclude any kind of transport processes of the chemicals in the radiative zones, and thus consider only convection as a mixing mechanism. Lithium is destroyed at moderate temperature by stellar interior standards; for typical main sequence (MS) and pre-MS densities, Li<sup>6</sup> burns at a time scale comparable to or shorter than the evolutionary time scale at around 2 million K and Li<sup>7</sup> burns at around 2.6 million K. Both beryllium and boron are less fragile, with characteristic burning temperatures of order 3.5 million K and 5 million K respectively. Because the observational data points to modest lithium depletion in halo stars (and even this conclusion is controversial!) we will restrict our discussion of light element depletion to lithium.
The major predictions of classical (sometimes referred to as standard) stellar evolution models for halo stars are summarized in Deliyannis et al. (1990). They depend mainly on the variations of the depth of the stellar convective envelope (and thus of the temperature at its base) with the stellar mass, metallicity and evolutionary stage. Pre-MS depletion increases with decreased mass, and there will therefore be a strong decrease of lithium with decreased $`T_{eff}`$ for cool stars. Main sequence depletion is predicted to be minimal for all but the coolest stars; the absolute degree of lithium depletion decreases with decreased metal abundance. The net effect predicted is that hot halo dwarfs should exhibit little or no dependence of their surface lithium abundance on effective temperature, and for the lowest metal abundances there should be a minimal dependence on \[Fe/H\], in the sense that lower abundances would be predicted for higher metallicities. This implies a small dispersion in lithium at fixed effective temperature; in the Population I case, it also implies a weak dependence of lithium on age which can be tested in open clusters.
Classical stellar models neglect some physical processes which are known to be important for interpreting the surface lithium abundances of stars. The linked phenomena of gravitational settling and thermal diffusion, which are solidly based in our knowledge of plasma physics, are among the most important. Atomic diffusion is a fundamental process which must occur in the stellar gas unless some macroscopic motions counteract it, causing heavy elements to sink with respect to light ones under the conditions applicable for halo dwarfs. The time scale for this process decreases as the depth of the surface convection zone decreases.
Theoretical models which include pure atomic diffusion<sup>1</sup><sup>1</sup>1By ”pure atomic diffusion” we mean that it is not counteracted by any macroscopic process in the stellar radiative zones (see Michaud et al. 1984 for the first computations for halo stars) therefore predict that lithium sinks below the surface convection zone for the conditions appropriate for subdwarfs, and furthermore that the degree of diffusion increases with increased $`T_{eff}`$. Microscopic diffusion will not generate a dispersion in abundance at fixed $`T_{eff}`$, and the effects at a given surface temperature are not strongly metallicity dependent.
Stellar winds can counteract and prevent atomic diffusion without leading to nuclear destruction (Vauclair & Charbonnel 1995). The corresponding mass loss necessary in the hottest halo stars is in excess of the value inferred from an extrapolation of the solar Mdot to halo stars by a factor of about 10 to 30, but not by a degree that can be ruled out observationally. For even larger mass loss rates the outer layers containing lithium can be removed; as noted by Swenson & Faulkner (1992) the finite depth of the surface convection zone must be accounted for and models with strong (stronger than inferred by Vauclair & Charbonnel 1995) mass loss alone are incompatible with the Population I lithium pattern. For hot halo stars the combined effect of diffusion and mass loss produce both lithium depletion and a small dispersion in abundance.
There are known mechanisms for mild mixing in the radiative envelopes of low mass stars. The two most frequently studied are rotationally-induced mixing and turbulence induced by gravity waves (see Pinsonneault 1997 for a review and DPC). Gravity waves can be produced by turbulence in the surface convection zone; because the convection zone depth is a strong function of mass, this can produce mass-dependent lithium depletion that is a function of time on the main sequence. It would not produce a dispersion in abundance for a sample of uniform age and composition, but abundance differences could be generated by a range of age and composition as seen in field halo stars.
The degree of rotational mixing depends on several major factors. Low mass Population I stars are observed to have a range of surface rotation rates and stars of the same mass, composition, and age could therefore have different initial angular momenta, different rotation rates as a function of time, and different degrees of rotational mixing. A dispersion in surface lithium abundance, even at fixed effective temperature in clusters, is therefore expected. However, the observed distribution of rotation rates in young clusters is nongaussian which strongly affects the expected lithium depletion pattern. We also cannot directly observe the initial conditions for Population II stars. To predict the detailed distribution of abundances the best we can to is to infer the distribution of initial conditions from young open cluster stars (see Pinsonneault et al. 1999, hereafter PWSN). The degree of mixing is also directly linked to the internal transport of angular momentum (e.g., Zahn 1992, Maeder 1995). Rotational mixing can also be inhibited by gradients in mean molecular weight - induced by nuclear burning in the cores of stars and possibly also by gravitational settling of helium in their outer layers (see also Michaud and Vauclair in these proceedings). In contrast with classical models, more modern models including mild mixing below the envelope on the main sequence can simultaneously produce modest depletions of species, such as lithium and beryllium, that burn at very different temperatures. Rotational mixing will also be a function of age, and some classes of models predict trends with \[Fe/H\] and effective temperature that can be tested in halo stars.
Realistic models should include the possible interactions between the above, since mass loss can counteract atomic diffusion and diffusion can interact with mixing (also see Michaud, these proceedings). We note that Vauclair (these proceedings) has proposed a nonlinear interaction between mixing and microscopic diffusion which would permit negligible halo star lithium depletion; note, however, that the details of such an interaction need to be computed and that such a cancellation would still have to be consistent with the globular cluster and Population I star data.
## 3. Observational Pattern
As we have just seen the surface lithium abundances of stars are sensitive to a variety of effects, both those accounted for in classical stellar models and those caused by physically well-motivated but still so-called non-standard effects. Uniqueness is thus a real issue when interpreting the observational data. We will therefore begin with the overall conclusions from studies of Population I stars, and then proceed to the current status of the observational data for Population II stars.
### 3.1. Population I Properties
The properties of Population I stars are summarized in DPC. Here we briefly recall those that are important for the problem of lithium depletion in Population II stars. In progressively older open cluster stars there is clear evidence for increased lithium depletion with age and a dispersion in abundance at fixed $`T_{eff}`$ which is inconsistent with classical models. The predicted dropoff of lithium for cool stars from pre-MS burning is clearly seen. The overall properties favor mild mixing below the envelope on the MS; there is also evidence for microscopic diffusion playing a role for F stars and in the helioseismic inversions of the solar sound speed relative to theoretical models (see Guzik & Cox 1993, Richard et al. 1996, Basu et al. 2000). Large enough amounts of mass loss to directly cause lithium depletion are inconsistent with the observed population I pattern (Swenson & Faulkner 1992). As of this time a single theoretical model capable of explaining all of the Population I data has not yet been found (Talon & Charbonnel 1998, PWSN).
### 3.2. Overall Population II Properties
Beginning with the pioneering work of Spite & Spite (1982), there have been a series of progressively more sophisticated observational studies of lithium in halo field stars; the largest sample is that of Thorburn (1994). Ryan et al. (1999, hereafter RNB) obtained a smaller sample with a lower formal error ($`\sigma `$ $``$ 0.036 dex) than the errors in earlier studies ($`\sigma `$ $``$ 0.07-0.09 dex). There have also been preliminary studies of small samples of stars near the turnoff in globular clusters. Different investigators of field stars agree on some general properties:
1. Halo stars hotter than 5800 K exhibit a weaker dependence on $`T_{eff}`$, \[Fe/H\], and a smaller dispersion than seen in Population I stars. There is vigorous debate about the existence and magnitude of any dispersion in the field star case, and there are active controversies about trends with $`T_{eff}`$ and \[Fe/H\].
2. Turnoff globular cluster stars were first studied by Molaro & Pasquini (1994), who reported a Li abundance for a turnoff star in NGC 6397 consistent with the halo plateau abundances. These are technically challenging observations owing to the faintness of the stars and the need for high resolution spectroscopy. A subsequent Keck study of the globular cluster M92 (Deliyannis et al. 1995, Boesgaard et al. 1998) revealed a large scatter, similar in morphology to the old open cluster M67. The sample, however, is small (seven stars, including only three with S/N greater than 40). Pasquini & Molaro (1997) also observed a range of Li abundances in three turnoff stars in the intermediate metal abundance globular cluster 47 Tuc. Any successful theory must explain both the cluster and field star patterns; more data is clearly needed for the globular cluster stars (see also Part III, Charbonnel, Deliyannis & Pinsonneault).
We now turn to a summary of the most recent data on the important global features in field halo stars.
### 3.3. Trends with $`T_{eff}`$
There has been a spirited debate about the existence and magnitude of trends in the halo star data with metallicity and effective temperature. The slope with metal abundance is important for constraining the galactic chemical evolution contribution to the observed lithium abundances, and it may also contribute to the small scatter in the data at fixed effective temperature. Trends with effective temperature are important as a diagnostic of the mass dependence of any physical processes which affect the surface abundances. Thorburn (1994) reported evidence for a positive slope of lithium with respect to $`T_{eff}`$; this conclusion was challenged by Molaro et al. (1995). Subsequently Ryan et al. (1996) reanalyzed their data, claiming confirmation of the original results; see also Bonifacio & Molaro (1997). The existence of a modest rising trend with increased T<sub>eff</sub> is only predicted in the models including atomic diffusion and stellar winds (Vauclair & Charbonnel 1995). However, what is even more important than the existence of a mild mean trend is the thing which is not seen : any evidence for a decline in lithium among the hottest stars.
As discussed above, models which include only microscopic diffusion predict a decline in surface lithium for the hottest halo stars; models with strong depletion from mixing also predict a downwards trend in lithium for the hottest stars (Chaboyer & Demarque 1994). This observational fact therefore constitutes an important limit on lithium depletion in halo stars.
### 3.4. Trends with \[Fe/H\]
The existence of trends of lithium with metallicity is a signature of galactic chemical evolution, and it could also contribute to the dispersion observed in halo stars. The majority of the observational investigations have looked for a correlation between \[Li\] and \[Fe/H\]; in parallel to the controversy over trends with effective temperature, there have been conflicting results on metallicity trends. Both Thorburn (1994) and Ryan et al. (1996) found some evidence for an increase in \[Li\] with \[Fe/H\] with a slope of order 0.1. This was disputed by Molaro et al. (1995) and Bonifacio & Molaro (1997). RNB obtained a sample with a small intrinsic range in effective temperature, but a wider range in metallicity. They could therefore evaluate metallicity but not effective temperature trends, and found a slope of \[Li\] with respect to \[Fe/H\] consistent with the Thorburn (1994) level. Chemical evolution trends should most logically be evaluated in the linear Li - linear Fe/H plane (see Olive and Matteucci, these proceedings).
A general feature of the derived chemical evolution trends is that they are sensitive to the source used for the metallicity, the subset of the data which is used, and the treatment of outliers in the fit. Ryan et al. (1996) also noted that the evidence for trends in the data with $`T_{eff}`$ and \[Fe/H\] is more convincing in a bivariate analysis than when either variable is treated separately. The existence and magnitude of metallicity trends is also important for the interpretation of the dispersion in abundance (see below). The metallicity dependence of any rotational mixing is small (PWSN) and would be difficult to disentangle from chemical evolution effects.
### 3.5. Dispersion
The dispersion in the lithium abundances of halo plateau stars has been the subject of a number of studies. This is largely because the existence or absence of a detectable range in abundance at fixed metallicity and $`T_{eff}`$ is the best direct test and constraint on the transport processes of chemicals in these stars. Lithium abundances can be studied as a function of age in the Population I case, which makes it easier to unambiguously distinguish between different classes of theoretical models (or at least rule bad models out). All of the metal-poor stars that we observe are old, and we therefore cannot directly reconstruct the depletion of lithium by sorting stars of progressively increased age into an evolutionary sequence.
Studies of the dispersion tend to fall into two groups. Some investigators (Deliyannis et al. 1993, Thorburn 1994) found evidence for a dispersion at a low level; others (Spite et al. 1996, Bonifacio & Molaro 1997) placed bounds on the dispersion consistent with their observational errors. The level of dispersion inferred by Deliyannis et al. (1993) is not inconsistent with the latter two studies (greater than 0.04 dex as compared with less than 0.08 and 0.07 dex respectively). In a recent paper, RNB have claimed a more stringent constraint on the overall dispersion. We examine this most recent data set below; a comparison of models including rotational mixing with the Thorburn (1994) data set was performed by PWSN.
The formal dispersion of the RNB data set is 0.053 dex, greater than their observational error of 0.036<sup>2</sup><sup>2</sup>2(not 0.033 as noted in the paper). They attribute this to a correlation between metallicity and lithium abundance, e.g. chemical evolution rather than stellar depletion. There is a substantial overlap between the RNB and Thorburn (1994) data sets, and the markedly lower dispersion inferred by RNB can be traced directly to differences in equivalent width measurements. In Figure 1 we illustrate and compare the properties of the RNB sample (excluding one upper limit) with the stars in common as measured by Thorburn (1994); both have been shifted to the same effective temperature scale. There are both significant zero-point shifts and a marked difference in the overall dispersion of the sample. RNB attribute this to possible scattered light and sky subtraction issues in the Thorburn (1994) data set. We note, however, that similar differences appear in samples in common with other investigators<sup>3</sup><sup>3</sup>3(e.g. compare G64-12 and CD -33 1173 as measured by both Thorburn 1994 and Spite & Spite 1993 to their relative abundances in RNB). The systematic differences between various observational data sets therefore require more scrutiny, especially given the relatively small sample size of the RNB data set and the small number of overdepleted stars expected for modest stellar depletion.
We compare the theoretical distribution of the lowest depletion case of PWSN to the RNB data in Figure 2. The majority of young low mass stars in open clusters are slow rotators, which implies that they should have experienced similar rotational histories and similar degrees of rotational mixing. However about 1/5 of the stars in open clusters are observed to be rapid rotators and these should manifest themselves as overdepleted objects, producing an excess dispersion which is measureable. The existence of a core in the sample with minimal internal dispersion therefore does not by itself rule out more modest stellar destruction (as noted by RNB). The existence and number of outliers is a more stringent test. In the raw RNB sample there are three stars more than 0.1 dex below the median, one of which has an upper limit of 1.36 for its abundance; this simulation would predict 4 depending on the criterion for defining what constitutes an overdepleted star. It is legitimate to question whether the highly overdepleted star is produced by the same mechanism as the other stars, but in any case the sample size is small and it is certainly difficult to make a persuasive case against modest depletion factors based on the data without chemical evolution corrections.
RNB placed more stringent constraints than the above based upon attributing some of their small dispersion to galactic chemical evolution. RNB fitted the data for a trend with \[Fe/H\] and concluded that there was a 10 % probability that as few outliers (one) as observed would be present by chance. This conclusion depends on the usage of a logarithmic, rather than a linear, relationship between lithium abundance and metallicity (Pinsonneault et al. 2000). In conclusion, the RNB data places more severe constraints on the dispersion in abundance than previous studies, and depending on the treatment of trends with metal abundance it may either be consistent with modest stellar depletion factors or places a bound of order 0.1 dex on the absolute depletion from the class of rotational mixing models considered by PWSN.
The existence of stars above the plateau may also provide some important clues; they could either be underdepleted or they could have experienced lithium production. Stars with very low initial angular momentum would experience much less rotational mixing than the norm and would therefore appear as underdepleted. However, there are strong observational selection effects against detecting very slow rotators in open clusters, and it is therefore difficult to estimate the fraction of such objects that would be expected in rotational mixing models. An alternative explanation would be differential lithium production. King et al. (1996) examined the most prominent such star, BD+23:3912, and found no evidence for lithium production in the abundances of other elements that would be affected by the main mechanisms (see King et al. 1996 for a detailed discussion and caveats).
### 3.6. $`Li{}_{}{}^{6}/Li^7`$
Li<sup>6</sup> is more fragile than Li<sup>7</sup> and it is not produced in significant quantities in standard BBN models. The detection of Li<sup>6</sup> in halo stars can therefore be used to set powerful constraints on the absolute depletion of Li<sup>7</sup>, with the caveat that the initial Li<sup>6</sup> abundance must be inferred from chemical evolution models. Smith et al. (1993) first claimed a detection of Li<sup>6</sup> in the halo star HD 84937. This was confirmed in subsequent studies by different investigators who added two more possible detections and a number of upper limits (see Cayrel et al. 1999, Hobbs et al. 1999, Nissen et al. 1999 for recent work on the subject and Nissen in these proceedings). The detected amount of Li<sup>6</sup> is small, but it appears to be secure. The amount is greater than would be expected from the beryllium and boron data, suggesting that alpha-alpha fusion may contribute to the production of Li<sup>6</sup>.
One important uncertainty in the usage of Li<sup>6</sup> data is therefore what the initial abundance of the species could be; for example, Lemoine et al. (1996) and Cayrel et al. (1999) obtained bounds of a factor of four and three respectively on the absolute depletion of Li<sup>6</sup> in HD 84937. PWSN argued that an even higher initial abundance could not be excluded, and considered the extreme limiting case where the halo Li<sup>6</sup> abundance could have been as high as the solar system value.
The second uncertainty is the ratio of Li<sup>6</sup> to Li<sup>7</sup> depletion. Nuclear burning in the convective envelope in standard models would produce strong Li<sup>6</sup> depletion before any Li<sup>7</sup> depletion occurred; this has been used to argue that any detected Li<sup>6</sup> implies negligible $`Li^7`$ depletion (e.g. Lemoine et al. 1996). Both models with microscopic diffusion and models with mild envelope mixing, however, predict simultaneous detectable depletion of both isotopes to varying degrees, with Li<sup>6</sup> being more sensitive but not infinitely so (see PWSN). Nonetheless, the Li<sup>6</sup> data does provide one of the best independent checks on any stellar depletion of Li<sup>7</sup>, and it indicates that large depletion factors are very unlikely to be consistent with the observed detections (see below.)
## 4. Theoretical Constraints on Lithium Depletion
In light of the observational data above, what can we infer about the depletion of lithium in halo stars? The first and most generally agreed-upon conclusion is that a variety of observational tests make a large depletion factor unlikely. There have been three major features of the halo data which have been used to constrain the absolute depletion: the degree of dispersion in the halo plateau, the detection of $`Li^6`$ in some halo stars, and the absence of a decline in surface Li for hotter halo stars.
### 4.1. Bounds from the Dispersion in the Plateau
PWSN inferred a range of 0.2-0.4 dex depletion factors from models including rotational mixing; the lower end of the range was more consistent with the dispersion inferred from the Thorburn (1994) data set, while the upper end of the range permitted the rare highly overdepleted stars to be explained within the framework of rotational mixing. As discussed above, the most recent data set of RNB is marginally consistent with the lower end of the depletion range in PWSN (of order 0.2 dex.) Other investigators (e.g. Bonifacio & Molaro 1997, RNB) have claimed more stringent limits of order 0.1 dex on the absolute depletion based upon the small (or, in their view, nonexistent!) dispersion in the halo Li data. We will return to this claim after reviewing other measures based upon the detection of $`Li^6`$ and the absence of the observed signature of pure microscopic diffusion in halo stars.
### 4.2. Bounds from $`Li^6`$ / $`Li^7`$ Measurements
PWSN set a less severe, but firm, limit of 0.5-0.6 dex $`Li^7`$ depletion from the measured $`Li^6`$ / $`Li^7`$ abundance ratios in halo stars under the assumption that the halo stars did not have a $`Li^6`$ abundance higher than the solar system value. Lemoine et al. (1996) derived a bound of a factor of 4 on the absolute $`Li^6`$ depletion of HD 84937, which in the rotationally mixed models of PWSN would imply a bound of 0.25 dex on the $`Li^7`$ depletion, while Cayrel et al. (1999) used the lithium data in the same star to set a bound of 0.1 dex on its $`Li^7`$ depletion; both of these calculations, however, are dependent on the chemical evolution model which is used. We note that if the Cayrel et al. (1999) bound of a factor of three $`Li^6`$ depletion is used in conjunction with the PWSN models, an absolute $`Li^7`$ depletion of 0.15 dex is inferred for HD 84937 rather than an upper limit of 0.1 dex.
### 4.3. Bounds from the Flatness of the Plateau
Vauclair & Charbonnel (1998) used a different set of properties to infer a stellar depletion factor. Pure microscopic diffusion in halo models would produce a strong decrease in surface lithium with increased $`T_{eff}`$ which is not observed. It is therefore clear that something must be inhibiting microscopic diffusion, especially given the improved agreement with helioseismology from the inclusion of micrscopic diffusion in solar model calculations. They noted that there is a subsurface peak in the $`Li^7`$ abundance of the pure diffusion models which does not vary greatly across the plateau.
Vauclair & Charbonnel (1995) (see also Swenson 1995) argued that mass loss at a rate 10-30 times greater than the solar value could counteract the effects of diffusion if the rate was tuned across the lithium plateau. This would have the effect of exposing a uniform abundance across the plateau, with an absolute depletion of 0.15 dex. Vauclair & Charbonnel (1998) noted that in the presence of sufficiently mild mixing the height of this peak could be preserved, implying that a uniform depletion of order 0.15 dex could apply if the absence of a measurable surface signature of diffusion arose from either the competing effects of mass loss or the interaction of diffusion and mild mixing. As noted by Chaboyer & Demarque (1994), sufficiently strong mixing can cancel the effects of diffusion while not preserving the height of the peak; however, models with the high degree of depletion inferred by that latter paper are difficult to reconcile with the other observational tests of lithium depletion.
### 4.4. Is There Any Depletion?
In light of the remarkable observed properties of the halo lithium plateau, it is reasonable to ask whether there is in fact any depletion at all. In other words, are we directly seeing the primordial lithium abundance (e.g. Bonifacio & Molaro 1997) or is the primordial lithium abundance in fact lower than the observed values because of a significant contribution from galactic chemical evolution (RNB)?
“Standard” stellar models are sometimes invoked as evidence against significant depletion. These classical models, however, achieve this prediction by simply neglecting known physics rather than by demonstrating that it is unimportant. For example, one could construct stellar models which ignore the CNO cycle, and they might even agree with some data, but it does not follow that these should be placed on an equal physical basis with models that include the known nuclear physics. In particular, atomic diffusion cannot be excluded arbitrarily from the computations, on the pretext that it produces unobserved features. This disagreement is a simple signature of some macroscopic motions (mass loss or rotation-induced motions) which counteract the diffusion process.
Any model predicting zero stellar depletion in halo field stars must be reconciled with the apparent scatter in the globular cluster turnoff stars. If both the halo stars and the globular cluster stars are depleted by (say) mild mixing, it is possible to explain the difference in the abundance patterns by a different set of initial angular momenta. Such a difference could arise in the context of the currently popular model for the origin of the range of rotation rates, namely that the lifetime of accretion disks determines the rotation rate, if globular cluster stars experienced more frequent interactions which disrupted their accretion disks early in their lifetimes relative to the lower density systems that the halo stars arose from. It is more challenging, however, to explain why stars with similar thermal structures should experience completely different depletion histories.
The complete absence of depletion in Population II stars would also have to be reconciled with the strong evidence for depletion in Population I stars which is not predicted by standard models.
At the same time, it is also clear that none of the existing theoretical models provide a complete description of the complex lithium abundance pattern seen in stars. Although the most recent classes of rotational models are reasonably successful at reproducing the observed angular momentum evolution of low mass Population I stars, they do not reproduce the solar rotation profile as inferred from helioseismology. They also require an extrapolation of the initial conditions from Population I to Population II stars, which may introduce systematic errors in the calculations. Further theoretical work is clearly needed, and a more refined set of models could potentially alter the inferred degree of stellar depletion.
## 5. Summary
We have compared theoretical models with the observational Population II lithium data to obtain bounds on stellar lithium depletion. From a combination of the dispersion in the data, the detection of $`Li^6`$, and the flatness of the halo plateau interesting bounds can be set on the stellar depletion of lithium in Population II stars. The majority of tests are roughly consistent with depletion at the 0.15-0.2 dex level, with a firm upper bound of 0.5-0.6 dex from a combination of the detected $`Li^6`$ abundance in HD 84937 and an extreme chemical evolution model. The most recent data set of RNB places the most severe observational constraints on the dispersion of lithium in halo stars. It is marginally consistent with the least depleted set of models of PWSN including rotational mixing; a larger statistical sample would permit a more definitive test of the limits on depletion from mixing in halo stars. There are unexplained differences between the equivalent width measurements of Thorburn (1994) and RNB which need to be understood; indeed, the systematic differences between investigators on the absolute level of the plateau are approaching the uncertainty in the inferred degree of stellar depletion.
A new generation of theoretical models is needed to further refine our understanding of light element depletion in stars. The interaction between different physical mechanisms, such as microscopic diffusion, mass loss, and rotational mixing, and a better physical description of each of them, may prove important in this context.
Finally, any stellar lithium depletion has implications for BBN. PWSN discussed the implications of a higher primordial lithium abundance; see Olive (these proceedings) for the case of negligible stellar depletion. If the preliminary data from the BOOMERANG mission is confirmed (Lange et al. 2000), we note that there may be a disagreement between the stellar lithium abundances and the predictions of standard BBN as well as the low deuterium results of Tytler (these proceedings.) This is significantly relaxed if stellar lithium depletion has occured.
## 6. Acknowledgements
M.P. would like to acknowledge support from NASA grant NAG5-7150 and NSF grant AST-9731621. C.C. thanks the Action Spécifique de Physique Stellaire and the Conseil National Français d’Astronomie for support. C.P.D. acknowledges support from the United States National Science Foundation under grant AST-9812735.
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# Charge Current Density from the Scattering Matrix
## Abstract
A method to derive the charge current density and its quantum mechanical correlation from the scattering matrix is discussed for quantum scattering systems described by a time-dependent Hamiltonian operator. The current density and charge density are expressed with the help of functional derivatives with respect to the vector potential and the electric potential. A condition imposed by the requirement that these local quantities are gauge invariant is considered. Our formulas lead to a direct relation between the local density of states and the total current density at a given energy. To illustrate the results we consider, as an example, a chiral ladder model.
$``$ Introduction $``$ The scattering matrix gives an important starting point in descriptions of quantum transport phenomena. It connects the incoming current amplitudes to the outgoing current amplitudes, and is calculated from the Hamiltonian operator by the Møller operator or the Green function method . We can obtain information about global characteristics in systems from the scattering matrix directly. For example, the Landauer formula gives a method to calculate the conductance from transmission amplitudes as components of the scattering matrix . The Friedel sum rule connects the density of states to the scattering matrix , so that it is in principle possible to calculate any equilibrium statistical mechanical quantity from the scattering matrix . We can also derive an expression of the persistent current in open conductors caused by a magnetic flux , or the shot noise from the scattering matrix.
It should be emphasized that the scattering matrix even gives information about local characteristics in systems. We consider one-particle and time-dependent Hamiltonian systems. First, the charge current density $`𝑱_n(𝒙,t)`$ at position $`𝒙`$ and time $`t`$ generated by the particle incident in a state with the quantum number $`n`$ is connected to the scattering matrix $`S=(S_{nn^{}})`$ as
$`𝑱_n(𝒙,t)={\displaystyle \frac{c}{2\pi i}}{\displaystyle \underset{n^{}}{}}S_{n^{}n}^{}{\displaystyle \frac{\delta S_{n^{}n}}{\delta 𝑨(𝒙,t)}}`$ (1)
where $`𝑨(𝒙,t)`$ is the vector potential and $`c`$ is the velocity of light. In this Letter we call this quantity the ”injected current density”. Eq. (1) is the main result of this Letter. Second, the probability density $`\rho _n(𝒙,t)`$ of the state caused by the incident particle with the quantum number $`n`$ is given by
$`\rho _n(𝒙,t)={\displaystyle \frac{1}{2\pi i}}{\displaystyle \underset{n^{}}{}}S_{n^{}n}^{}{\displaystyle \frac{\delta S_{n^{}n}}{\delta U(𝒙,t)}}`$ (2)
with the potential $`U(𝒙,t)`$. The probability density $`\rho _n(𝒙,t)`$ also has been called the ”injectivity” and can be regarded as the density of states with a preselection of the incident channel. Eq. (2) in the time-independent Hamiltonian case has been used in treatments of the electron-electron interaction using a self-consistent potential . In this Letter we give a new derivation of this formula and generalize it to the time-dependent Hamiltonian case.
Our technique to derive Eqs. (1) and (2) also can be used to arrive at quantum mechanical correlation functions of local quantities from the scattering matrix. For example we show that the quantum mechanical correlation $`C_n^{\left(\mu \nu \right)}(𝒙,𝒙^{};t)`$ between the $`\mu `$-th charge current density component at position $`𝒙`$ and the $`\nu `$-th charge current density component at position $`𝒙^{}`$ at the same time $`t`$ is given by
$`C_n^{\left(\mu \nu \right)}(𝒙,𝒙^{};t)={\displaystyle \frac{c^2\mathrm{}}{2\pi }}{\displaystyle \underset{n^{}}{}}\text{Re}\left\{{\displaystyle \frac{\delta S_{n^{}n}^{}}{\delta A^{\left(\mu \right)}(𝒙,t)}}{\displaystyle \frac{\delta S_{n^{}n}}{\delta A^{\left(\nu \right)}(𝒙^{},t)}}\right\}`$ (3)
(4)
where $`A^{\left(\mu \right)}(𝒙,t)`$ is the $`\mu `$-th component of the vector potential.
Our formulas include the vector potential and the potential explicitly, so we have to discuss the gauge invariance. We show Eqs. (1 \- 2) and (4) to be gauge invariant, and derive conditions which should be satisfied by locally gauge invariant quantities.
In the time-dependent Hamiltonian system an energy of an outgoing particle can be different from the energy of the corresponding incoming particle, and the scattering matrix element $`S_{nn^{}}`$ describes a transition to such a different energy. On the other hand, if we consider only time-independent Hamiltonian cases, the problem becomes simpler, because the scattering matrix is decomposed into the scattering matrices restricted to the energy shells. Using this feature we obtain a relation between the local density of states and the total current density defined by the sum of the injected current density $`𝑱_n(𝒙,t)`$ with respect to the suffix $`n`$ satisfying the condition $`E_n=E`$ at energy $`E`$.
As a simple example we investigate a ladder model with a directionality, namely a chiral ladder model, termed a ”quantum rail road” in Ref. . We verify the formula (1) for this model calculating separately the current density and the scattering matrix.
$``$ Injected current density and injectivity $``$ We start from the time-dependent Hamiltonian operator $`\widehat{H}(t)`$ in the quantum scattering system. The dynamics of the system is described by the Schrödinger equation using this Hamiltonian operator. We decompose the total Hamiltonian operator $`\widehat{H}(t)`$ into the asymptotic Hamiltonian operator $`\widehat{H}_0`$ and the scattering operator $`\widehat{H}_1(t)`$; $`\widehat{H}(t)=\widehat{H}_0+\widehat{H}_1(t)`$. The operator $`\widehat{H}_0`$ is chosen to be the Hamiltonian operator which describes incoming and outgoing particles in the asymptotic regions. We assume that $`\widehat{H}_0`$ is a time-independent operator determined uniquely. Besides, the system described by the Hamiltonian operator $`\widehat{H}(t)`$ or $`\widehat{H}_0`$ is assumed to have no bound state. Below we use the coordinate representation for any operator and take $`𝒙`$ as the coordinate of particles. Under these conditions the scattering matrix element $`S_{nn^{}}`$ is given by
$`S_{nn^{}}=\underset{t_2+\mathrm{}}{lim}\underset{t_1\mathrm{}}{lim}{\displaystyle 𝑑𝒙\mathrm{\Phi }_n(𝒙)^{}\widehat{U}(t_2,t_1)\mathrm{\Phi }_n^{}(𝒙)}`$ (5)
where $`\widehat{U}(t_2,t_1)`$ is the time evolution operator $`\mathrm{exp}\{i\widehat{H}_0t_2/\mathrm{}\}`$$`\stackrel{}{T}\mathrm{exp}\{i_{t_1}^{t_2}𝑑t\widehat{H}(t)/\mathrm{}\}`$ $`\mathrm{exp}\{i`$ $`\widehat{H}_0t_1/\mathrm{}\}`$ in the interaction picture with $`\stackrel{}{T}`$ being the positive time-ordering operator, and $`\mathrm{\Phi }_n(𝒙)`$ is the eigenstate of the operator $`\widehat{H}_0`$ corresponding to the energy eigenvalue $`E_n`$. Here the limits $`t_1\mathrm{}`$ and $`t_2+\mathrm{}`$ are defined by $`lim_{t\pm \mathrm{}}X(t)lim_{ϵ+0}(\pm ϵ)_0^\pm \mathrm{}𝑑t^{}e^{ϵt^{}}`$$`X(t^{})`$ for any function $`X(t)`$ of $`t`$ . The set $`\{\mathrm{\Phi }_n(𝒙)\}_n`$ of the eigenstates of the operator $`\widehat{H}_0`$ is chosen to satisfy the orthonormality condition and the completeness relation. The scattering matrix $`S=(S_{nn^{}})`$ is shown to be an unitary matrix; $`S^{}S=SS^{}=I`$.
For simplicity we consider the one-particle system. The total Hamiltonian operator $`\widehat{H}(t)`$ and the operator $`\widehat{H}_0`$ are represented as $`(i\mathrm{}/𝒙q𝑨(𝒙,t)/c)^2/(2m)+U(𝒙,t)+U_0(𝒙)`$ and $`(\mathrm{}^2/(2m))^2/𝒙^2+U_0(𝒙)`$, respectively, where $`m`$ is the mass of the particle, $`q`$ is the charge of the particle, $`𝑨(𝒙,t)`$ is the vector potential, and $`U(𝒙,t)`$ is the external potential (plus the induced potential by the interaction of particles), and $`U_0(𝒙)`$ is the confinement potential. We consider a functional derivative $`\delta ^{\left(t\right)}\widehat{S}_{nn^{}}`$ of the scattering matrix element with respect to the vector potential $`𝑨(𝒙,t)`$ or the potential $`U(𝒙,t)`$ at time $`t`$, using the notation $`\delta ^{\left(t\right)}=\delta /\delta 𝑨(𝒙,t)`$ or $`\delta /\delta U(𝒙,t)`$. Using the expression (5) of the scattering matrix elements we obtain a general relation
$`{\displaystyle \frac{1}{2\pi i}}{\displaystyle \underset{n^{}}{}}S_{n^{}n}^{}(\delta ^{\left(t\right)}S_{n^{}n})={\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t^{}\delta ^{\left(t\right)}\widehat{H}_1(t^{})_n^{\left(t\right)}.`$ (6)
Here the notation $`\mathrm{}_n^{\left(t\right)}`$ means the expectation value taken with the scattering state of the quantum number $`n`$, namely $`\widehat{X}_n^{\left(t\right)}𝑑𝒙\mathrm{\Psi }_n(𝒙,t)^{}\widehat{X}\mathrm{\Psi }_n(𝒙,t)`$ $`/(2\pi \mathrm{})`$ for any operator $`\widehat{X}`$, where $`\mathrm{\Psi }_n(𝒙,t)`$ is a solution of the Schrödinger equation using the total Hamiltonian operator $`\widehat{H}(t)`$ and is defined by $`\mathrm{\Psi }_n(𝒙,t)lim_{t_1\mathrm{}}\stackrel{}{T}\mathrm{exp}\{i_{t_1}^t𝑑t^{}\widehat{H}(t^{})/\mathrm{}\}\mathrm{exp}\{iE_nt_1/\mathrm{}\}\mathrm{\Phi }_n(𝒙)`$. The derivation of Eq. (6) is given by using the completeness relation of the set $`\{\mathrm{\Phi }_n(𝒙)\}_n`$, the property $`\widehat{U}(t_2,t_1)=\widehat{U}(t_2,t)\widehat{U}(t,t_1)`$ of the time evolution operator and the relation $`\delta ^{\left(t\right)}\widehat{U}(t_2,t_1)=\widehat{U}(t_2,t)\mathrm{exp}\{i\widehat{H}_0t/\mathrm{}\}\{(i/\mathrm{})_{\mathrm{}}^{\mathrm{}}𝑑t^{}\delta ^{\left(t\right)}\widehat{H}_1(t^{})\}\mathrm{exp}\{i\widehat{H}_0t/\mathrm{}\}\widehat{U}(t,t_1)`$ in $`t_1<t<t_2`$. Eq. (6) is a key result of this Letter. We introduce the injected current density $`𝑱_n(𝒙,t)`$ and the injectivity $`\rho _n(𝒙,t)`$ as
$`𝑱_n(𝒙,t)\widehat{𝑱}(𝒙,t)_n^{\left(t\right)},\rho _n(𝒙,t)\widehat{\rho }(𝒙)_n^{\left(t\right)}`$ (7)
using the charge current density operator $`\widehat{𝑱}(𝒓,t)q\widehat{𝒗}(t)\delta (𝒙𝒓)`$ with $`\widehat{𝒗}(t)`$ being the velocity operator $`(i\mathrm{}/𝒙q𝑨(𝒙,t)/c)/m`$ and the multiplication $``$ being the symmetrized product, and using the probability density operator $`\widehat{\rho }(𝒓)\delta (𝒙𝒓)`$. Eqs. (1) and (2) are derived from Eq. (6), using the relations $`_{\mathrm{}}^{\mathrm{}}𝑑t^{}\delta \widehat{H}_1(t^{})/\delta 𝑨(𝒙,t)=\widehat{𝑱}(𝒙,t)/c`$ and $`_{\mathrm{}}^{\mathrm{}}𝑑t^{}\delta \widehat{H}_1(t^{})/\delta U(𝒙,t)=\widehat{\rho }(𝒙)`$.
$``$ Current density correlation $``$ We consider functional derivatives $`\delta _j^{\left(t\right)}S_{nn^{}}`$, $`j=1,2`$ of the scattering matrix element with respect to the vector potential $`𝑨(𝒙,t)`$ or the potential $`U(𝒙,t)`$ at time $`t`$. Using Eq. (5) we obtain another general relation
$`{\displaystyle \frac{\mathrm{}}{2\pi }}{\displaystyle \underset{n^{}}{}}\left(\delta _1^{\left(t\right)}S_{n^{}n}^{}\right)\left(\delta _2^{\left(t\right)}S_{n^{}n}\right)`$ (8)
$`=\left({\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t^{}\delta _1^{\left(t\right)}\widehat{H}_1(t^{})\right)\left({\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t^{\prime \prime }\delta _2^{\left(t\right)}\widehat{H}_1(t^{\prime \prime })\right)_n^{\left(t\right)}.`$ (9)
Using Eq. (9) we obtain Eq. (4). Here the quantum mechanical current density correlation $`C_n^{\left(\mu \nu \right)}(𝒙,𝒙^{};t)`$ is defined by
$`C_n^{\left(\mu \nu \right)}(𝒙,𝒙^{};t)\widehat{J}^{\left(\mu \right)}(𝒙,t)\widehat{J}^{\left(\nu \right)}(𝒙^{},t)_n^{\left(t\right)}`$ (10)
where $`\widehat{J}^{\left(\mu \right)}(𝒙,t)`$ is the $`\mu `$-th component of the charge current density operator. It should be noted that the average $`\mathrm{}_n^{\left(t\right)}`$ taken in the correlation (10) includes only the quantum mechanical average, but does not include the thermo-statistical average.
In similar ways we can obtain other quantum mechanical correlation functions $`\widehat{J}^{\left(\nu \right)}(𝒙,t)\widehat{\rho }(𝒙^{})_n^{\left(t\right)}`$ and $`\widehat{\rho }(𝒙)\widehat{\rho }(𝒙^{})_n^{\left(t\right)}`$ from the scattering matrix.
$``$ Gauge invariance $``$ The gauge transformation of the electric potential $`\varphi (𝒙,t)`$ and the vector potential $`𝑨(𝒙,t)`$ is represented as $`\varphi (𝒙,t)\varphi (𝒙,t)+\mathrm{\Delta }\varphi (𝒙,t)`$ and $`𝑨(𝒙,t)𝑨(𝒙,t)+\mathrm{\Delta }𝑨(𝒙,t)`$ with $`\mathrm{\Delta }\varphi (𝒙,t)(1/c)\phi (𝒙,t)/t`$ and $`\mathrm{\Delta }𝑨(𝒙,t)\phi (𝒙,t)/𝒙`$ using a function $`\phi (𝒙,t)`$ of $`𝒙`$. The electric and magnetic fields are invariant under the gauge transformation. Using that the Schrödinger equation is gauge invariant, we can show that the $`\mu `$-th component $`J_n^{\left(\mu \right)}(𝒙,t)`$ of the injected current density is also gauge invariant, meaning that the formula (1) is gauge invariant.
The fact that the injected current density is gauge invariant is also represented as
$`{\displaystyle }dt^{}{\displaystyle }d𝒙^{}\{{\displaystyle \frac{\delta J_n^{\left(\mu \right)}(𝒙,t)}{\delta \varphi (𝒙^{},t^{})}}\mathrm{\Delta }\varphi (𝒙^{},t^{})`$ (11)
$`+{\displaystyle \frac{\delta J_n^{\left(\mu \right)}(𝒙,t)}{\delta 𝑨(𝒙^{},t^{})}}\mathrm{\Delta }𝑨(𝒙^{},t^{})\}=0.`$ (12)
A change $`\delta \varphi (𝒙,t)`$ of the electric potential $`\varphi (𝒙,t)`$ modifies the potential $`U(𝒙,t)`$ as $`\delta U(𝒙,t)=q\delta \varphi (𝒙,t)`$. Moreover Eq. (12) should be satisfied for any function $`\phi (𝒙,t)`$ which is zero at the boundary of the integral of the left-hand side of Eq. (12). Using these facts and the partial integrals in Eq. (12) we obtain
$`{\displaystyle \frac{q}{c}}{\displaystyle \frac{}{t^{}}}{\displaystyle \frac{\delta J_n^{\left(\mu \right)}(𝒙,t)}{\delta U(𝒙^{},t^{})}}+{\displaystyle \frac{}{𝒙^{}}}{\displaystyle \frac{\delta J_n^{\left(\mu \right)}(𝒙,t)}{\delta 𝑨(𝒙^{},t^{})}}=0.`$ (13)
Eq. (13) represents a condition for the injected current density imposed by its gauge invariance. Similarly, the formulas (2) and (4) are gauge invariant, and the injectivity $`\rho _n(𝒙,t)`$ satisfies an equation similar to Eq. (13).
$``$ Total current density and local density of states in the time-independent Hamiltonian system $``$ Now, we consider the time-independent Hamiltonian case. In this case the injected current density and the injectivity are independent of time $`t`$. Moreover the state $`\mathrm{\Psi }_n(𝒙,t)`$ used to define the average $`\mathrm{}_n^{\left(t\right)}`$ can be replaced by the state $`\overline{\mathrm{\Psi }}_n(𝒙)\{1+\widehat{G}(E_n)\widehat{H}_1\}\mathrm{\Phi }_n(𝒙)`$ using the Green operator $`\widehat{G}(E)lim_{ϵ+0}(E\widehat{H}+iϵ)^1`$. The state $`\overline{\mathrm{\Psi }}_n(𝒙)`$ is the eigenstate of the total Hamiltonian operator $`\widehat{H}`$ corresponding to the eigenvalue $`E_n`$ and satisfies the Lippmann-Schwinger equation. One of the important features in the time-independent Hamiltonian system is that the scattering matrix element $`S_{nn^{}}`$ takes a non-zero value only in the case of $`E_n=E_n^{}`$. Therefore, in the formulas (1) and (2) we can replace the full scattering matrix element $`S_{n^{}n}`$ with the matrix element $`𝒮_{l^{}l}(E_n)`$ of the scattering matrix $`𝒮(E_n)`$ restricted to the energy shell of energy $`E_n`$, and exchange the sum of the quantum number $`n^{}`$ with the sum of the channel number $`l^{}`$.
It is important to note a difference between the time-dependent and time-independent Hamiltonian cases in our formulas. In the time-dependent Hamiltonian case the functional derivative $`\delta S_{nn^{}}/\delta 𝑨(𝒙,t)`$ is given by $`\delta S_{nn^{}}=𝑑𝒙𝑑t`$ $`(\delta S_{nn^{}}/\delta 𝑨(𝒙,t))\delta 𝑨(𝒙,t)`$. On the other hand, in the time-independent Hamiltonian case the functional derivative $`\delta 𝒮_{ll^{}}(E)/\delta 𝑨(𝒙)`$ is given by $`\delta 𝒮_{ll^{}}(E)=𝑑𝒙(\delta 𝒮_{ll^{}}(E)/\delta 𝑨(𝒙))\delta 𝑨(𝒙)`$, which does not include the time integral like in the time-dependent Hamiltonian case. This difference appears in the current density correlation formula (4), where a simple exchange of the scattering matrix element $`S_{nn^{}}`$ and the vector potential $`𝑨(𝒙,t)`$ with the on-shell scattering matrix element $`𝒮_{ll^{}}(E)`$ and the vector potential $`𝑨(𝒙)`$, respectively, is not allowed to obtain its time-independent Hamiltonian version. However we can obtain the time-independent Hamiltonian versions of the injected current density formula (1) and the injectivity formula (2) by such a simple exchange. This technical point is better considered in a separate paper.
We proceed to consider the total current density $`𝑱(E,𝒙)`$ and the local density of states $`\rho (E,𝒙)`$ defined by
$`𝑱(E,𝒙){\displaystyle \underset{\begin{array}{c}n\\ \left(E_n=E\right)\end{array}}{}}𝑱_n(𝒙),\rho (E,𝒙){\displaystyle \underset{\begin{array}{c}n\\ \left(E_n=E\right)\end{array}}{}}\rho _n(𝒙).`$ (18)
Eqs. (1) and (2) lead to
$`𝑱(E,𝒙)={\displaystyle \frac{c}{\pi }}{\displaystyle \frac{\delta \theta _f(E)}{\delta 𝑨(𝒙)}},\rho (E,𝒙)={\displaystyle \frac{1}{\pi }}{\displaystyle \frac{\delta \theta _f(E)}{\delta U(𝒙)}}`$ (19)
where $`\theta _f(E)`$ is called ”Friedel phase” and is defined by $`\theta _f(E)(1/(2i))\mathrm{ln}\text{Det}\{𝒮(E)\}`$. It may be noted that using the first equation in Eq. (19) the relation $`𝑱(E,𝒙)=𝑱(E,𝒙)|_{𝑨𝑨}`$ is obtained under the condition $`S_{nn^{}}=S_{n^{}n}|_{𝑨𝑨}`$. The second formula in Eq. (19) have already been used in some works . It follows from Eq. (19) that
$`{\displaystyle \frac{\delta 𝑱(E,𝒙^{})}{\delta U(𝒙)}}=c{\displaystyle \frac{\delta \rho (E,𝒙)}{\delta 𝑨(𝒙^{})}}.`$ (20)
This equation shows a connection of the current density to the local density of states.
The total current density $`𝑱(E,𝒙)`$ takes a non-zero value in some important phenomena, such as edge currents in the quantum hall effect, persistent currents caused by a magnetic flux and so on. Earlier work obtained the total persistent current from the flux derivative of the scattering matrix for a ring connected to a lead . For a purely one-dimensional ring the functional derivative with respect to the vector potential can simply be replaced by the derivative with respect to the flux. In such a special case the first formula in Eq. (19) leads to the same expression as the earlier work.
$``$ Example $``$ As a specific example we consider a chiral ladder model threaded by a weak magnetic field. In this model particles move on the both legs of the ladder only in one direction (See Fig. 1). This may be regarded as a model of two edge channels at one edge of a quantum Hall bar with impurities, by which electrons transfer from an edge channel to another channel.
We assume that the ladder has one-dimensional legs each of which has one channel. We introduce the scattering matrix $`T^{\left(j\right)}`$ which connects the current amplitude to the left of the $`j`$-th rung to the current amplitude to the left of the $`j+1`$-th rung. The dependence of the scattering matrix $`T^{\left(j\right)}`$ on the vector potential in the legs is
$`T^{\left(j\right)}=\left(\begin{array}{cc}t_{11}^{\left(j\right)}\mathrm{exp}\{i\varphi _1^{\left(j\right)}/\varphi _0\}& t_{12}^{\left(j\right)}\mathrm{exp}\{i\varphi _1^{\left(j\right)}/\varphi _0\}\\ t_{21}^{\left(j\right)}\mathrm{exp}\{i\varphi _2^{\left(j\right)}/\varphi _0\}& t_{22}^{\left(j\right)}\mathrm{exp}\{i\varphi _2^{\left(j\right)}/\varphi _0\}\end{array}\right)`$ (23)
where $`\varphi _0\mathrm{}c/q`$ and $`\varphi _1^{\left(j\right)}`$ ($`\varphi _2^{\left(j\right)}`$) is the integral of $`A(x)`$ over the upper (lower) leg between the $`j`$-th and the $`j+1`$-th rungs with $`A(x)`$ being the vector potential element in the direction of the legs at position $`x`$. Here the matrix $`t^{\left(j\right)}(t_{ll^{}}^{\left(j\right)})`$ is the corresponding scattering matrix in the case that the vector potential in the legs is zero. In this model the scattering matrix $`T^{\left(j\right)}`$ is the same as the corresponding transfer matrix, so the scattering matrix $`𝒮^{\left(j\right)}(𝒮_{ll^{}}^{\left(j\right)})`$ of the sub-system consisting of the first $`j`$ number of rungs is given by $`T^{\left(j\right)}T^{\left(j1\right)}\mathrm{}T^{\left(1\right)}`$.
We consider the injected current density $`J_{(\mu ,1)}^{\left(j\right)}`$ ($`J_{(\mu ,2)}^{\left(j\right)}`$) in the upper (lower) leg between the $`j`$-th and the $`j+1`$-th rungs, which is caused by the incident current $`𝒂_\mu `$ shown in Fig. 1. This current is represented as $`J_{(\mu ,\nu )}^{\left(j\right)}=q\rho _0v|𝒮_{\nu \mu }^{\left(j\right)}|^2`$ where $`v`$ is the particle velocity in the upper (or lower) leg and $`\rho _0`$ is the local density of states $`1/(2\pi \mathrm{}|v|)`$ in the one-dimensional perfect wire.
Now we connect the scattering matrix $`𝒮^{\left(N\right)}`$ of the system consisting of $`N`$ number of the rungs ($`N>1`$) to the injected current density $`J_{(\mu ,\nu )}^{\left(j\right)}`$. Noting unitarity of the matrices $`T^{\left(j^{}\right)}`$, $`j^{}=N,N1,\mathrm{},l+1`$ and using the relation $`\delta \varphi _\nu ^{\left(j\right)}/\delta A(x_\nu ^{}^{\left(j^{}\right)})=\delta _{jj^{}}\delta _{\nu \nu ^{}}`$, where $`x=x_1^{\left(j\right)}`$ ($`x_2^{\left(j\right)}`$) is a point in the upper (lower) leg between the $`j`$-th and the $`j+1`$-th rungs, we obtain $`(1/(2\pi i))_{l=1}^2𝒮_{l\mu }^{\left(N\right)}\delta 𝒮_{l\mu }^{\left(N\right)}/\delta A(x_\nu ^{\left(j\right)})=(1/c)J_{(\mu ,\nu )}^{\left(j\right)}`$, which is just the injected current density formula (1) in the time-independent Hamiltonian case.
$``$ Conclusion and remarks $``$ In this Letter we have discussed formulas to derive the charge current density, the charge density and their quantum mechanical current density correlations from the scattering matrix in one-particle and time-dependent Hamiltonian systems. The gauge invariance requires Eq. (13) which has to be satisfied by these local quantities. Using our formulas we obtained a relation between the local density of states and the total current density produced by incident particles at an energy in time-independent Hamiltonian systems. Specifically we verified the current density formula for a chiral ladder model.
In this Letter for simplicity we considered one-particle systems only. However our technique to derive Eqs. (6) and (9) can be used to obtain some generalizations to the formulas including many-particles’ effects. For instance, we can generalize Eq. (1) to the inelastic scattering cases by dynamical scatterers with no charge. Eqs. (6) and (9) also suggest that we can derive formulas for other local physical quantities from functional derivatives of the scattering matrix with respect to other local fields. For example, if the potential $`U(𝒙,t)`$ includes the term $`q\widehat{𝒔}𝑩(𝒙,t)/2`$ as an interaction effect of a spin $`\widehat{𝒔}`$ with a magnetic field $`𝑩(𝒙,t)`$, then we can calculate the local expectation value of a spin component from a functional derivative of the scattering matrix with respect to the magnetic field.
As another approach to the charge current density we can use the linear response theory . The connection of the scattering theoretical approach to the linear response theoretical approach in the description of local quantities is left as a future problem.
$``$
I am very grateful to M. Büttiker for stimulating discussions, encouragements and a careful reading of this Letter. Especially he gave me valuable comments and suggestions relating to the injectivity formula, time-dependent Hamiltonian problems and the gauge invariance, and introduced the chiral ladder model to me.
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# From the Bethe Ansatz to the Gessel-Viennot Theorem
## 1 Introduction
The problem of non-intersecting paths or vicious walkers has been studied by the statistical mechanics community who have been interested in them as simple models of various polymer and other physical systems and independently by the combinatorics community who have been interested in them in connection with binomial determinants. In statistical mechanics they are generally known as vicious walkers, a term coined by Fisher , who studied the continuous version of the model. Various cases of the lattice problem were later studied by Forrester . Independently in the area of combinatorics the problem of non-intersecting paths was solved by a very general theorem of Gessel and Viennot , following the work of Lindström , and Karlin and McGregor . All these studies express the number of configurations as the value of a determinant. Non-intersecting walks arose in yet another context, that of vertex models in statistical mechanics, where it was noticed that if the vertices of the six-vertex model are drawn in a particular way they could be interpreted as lattice paths . If one of the vertices had weight zero, giving a five-vertex model, the resulting paths were non-intersecting. The vertex models are traditionally solved by expressing the partition function (a generating function) in terms of transfer matrices. The partition function is then evaluated by either of two very powerful techniques, that of commuting transfer matrices or by direct diagonalisation of the transfer matrices using the coordinate Bethe Ansatz . In this paper we bring the independent results of the two communities together for the case of $`N`$ non-intersecting paths. We will show that the Bethe Ansatz (from statistical mechanics) and the Gessel-Viennot Theorem (from combinatorics) are essentially equivalent for a fairly general problem on the square lattice. The theorems proved should be easily generalisable to other planar lattices. The connection between the six-vertex model and non-intersecting path problems, from which this correspondence stems, has been recently discussed in general terms . Here we shall consider a model equivalent to a generalised five-vertex model on the square lattice where the (Boltzmann) weights associated with walk edges are inhomogeneous in one direction.
## 2 The model
A lattice path or walk in this paper is a walk on a square lattice rotated $`45^{}`$ which has steps in only the north-east or south-east directions, and with sites labelled $`(m,y)`$ (see figure 1). A set of walks is *non-intersecting* if they have no sites in common. We are concerned with enumerating the number of configurations of $`N`$ non-intersecting walks, starting and ending at given positions, in various geometries: 1) walks in a plane without boundaries; 2) walks which are confined to the upper half plane; and 3) walks which are confined to a strip of a given width, $`L`$. More generally, one may be interested in interacting cases where the walks nearest the boundaries are attracted or repulsed by contact interactions: combinatorially this requires knowledge of the number of walks with particular numbers of contacts with each of the boundaries. In this paper we shall focus on case 3 since the other 2 cases can be easily obtained from this case as limits.
To more easily describe our model we require the following sub-domains of $`^N`$
$`\underset{L}{\overset{o}{𝒮}}`$ $`=\{y|\mathrm{\hspace{0.17em}1}yL,y\text{ and }y\text{ odd}\},`$ (2.1a)
$`\underset{L}{\overset{e}{𝒮}}`$ $`=\{y|\mathrm{\hspace{0.17em}0}yL,y\text{ and }y\text{ even}\},`$ (2.1b)
$`\underset{L}{\overset{}{𝒮}}`$ $`=\{y|\mathrm{\hspace{0.17em}0}yL,y\},`$ (2.1c)
$`\underset{L}{\overset{o}{𝒰}}`$ $`=\{(y_1,\mathrm{},y_N)|\mathrm{\hspace{0.17em}1}y_1<\mathrm{}<y_NL,\text{ }y_i\underset{L}{\overset{o}{𝒮}}\}`$ (2.1d)
$`\underset{L}{\overset{e}{𝒰}}`$ $`=\{(y_1,\mathrm{},y_N)|\mathrm{\hspace{0.17em}0}y_1<\mathrm{}<y_NL,\text{ }y_i\underset{L}{\overset{e}{𝒮}}\}`$ (2.1e)
$`\underset{L}{\overset{}{𝒰}}`$ $`=\{(y_1,\mathrm{},y_N)|\mathrm{\hspace{0.17em}0}y_1<\mathrm{}<y_NL,\text{ }y_i\underset{L}{\overset{}{𝒮}}\}`$ (2.1f)
We will use $`\underset{L}{\overset{p}{𝒰}}`$ to denote $`\underset{L}{\overset{o}{𝒰}}`$ or $`\underset{L}{\overset{e}{𝒰}}`$. Let $`N`$ non-intersecting walks, confined to a strip of width $`L`$, start at $`y`$-coordinates $`𝐲^i=(y_1^i,\mathrm{},y_N^i)\underset{L}{\overset{p}{𝒰}}`$ in column $`m=0`$ of the lattice sites and terminate after $`t`$ steps at $`y`$-coordinates $`𝐲^f=(y_1^f,\mathrm{},y_N^f)\underset{L}{\overset{p^{}}{𝒰}}`$ in the $`t^{th}`$ column. If $`t`$ is even then $`p^{}=p`$ else $`p^{}=\overline{p}`$, where $`\overline{p}`$ is the opposite parity to $`p`$. We will only consider the case that $`L`$ is odd so that $`|\underset{L}{\overset{e}{𝒰}}|=|\underset{L}{\overset{o}{𝒰}}|=\left(\genfrac{}{}{0pt}{}{\frac{1}{2}(L+1)}{N}\right)`$. (If $`L`$ is even a null space enters the subsequent analysis of the transfer matrices leading to a distracting complication.) We are considering paths such that a) if $`(m1,y)`$ is the position of a path in column $`m1`$ the only possible positions for that path in column $`m`$ are $`(m,y^{})`$ with $`y^{}=y\pm 1`$ and $`0y^{}L`$ and b) the non-intersection is defined through the constraint that if there are $`N`$ sites occupied at $`m=0`$ then in each column of sites ($`0mt`$) there are exactly $`N`$ occupied sites. We generalise the walk problem associated with the five-vertex problem by assigning a weight $`w(y,y^{})`$ to the lattice edge from site $`(m1,y)`$ to $`(m,y^{})`$ with $`y^{}=y\pm 1`$ (see figure 1). Notice that, since $`w(y,y^{})`$ is assumed independent of the column index $`m`$, due to the square lattice structure the weights are periodic in the $`m`$ direction with period two: Note if $`y\underset{L}{\overset{p}{𝒮}}`$ then $`y^{}\underset{L}{\overset{\overline{p}}{𝒮}}`$, and in general $`w(y,y^{})w(y^{},y)`$. For the sake of generality we also associate an arbitrary weight $`v(y^i)`$ with each of the sites occupied at $`m=0`$. The weight associated with a given set of walks is the product of $`w`$ weights over all edges occupied by the walks multiplied by the product of the $`v`$ weights for each of the initial sites occupied. The generating function $`\stackrel{=}{Z}_t^𝒩(𝐲^i𝐲^f)`$, of $`N`$ walks of length $`t`$ starting at $`𝐲=𝐲^i`$ in column $`m=0`$ and finishing at $`𝐲=𝐲^f`$ in column $`m=t`$ is the sum of these weights over all sets of walks connecting $`𝐲^i`$ and $`𝐲^f`$:
$$\stackrel{=}{Z}_t^𝒩(𝐲^i𝐲^f)=\underset{𝒴}{}\underset{j=1}{\overset{N}{}}v(y_j(0))\underset{m=1}{\overset{t}{}}w(y_j(m1),y_j(m))$$
(2.2)
where $`y_j(m)`$ is the position of the $`j^{th}`$ walk in column $`m`$ and the set $`𝒴`$ is given by
$`𝒴`$ $`=`$ $`\{y_j(m)|1jN,0mt,1y_1(m)<y_2(m)<\mathrm{}<y_N(m)L,`$ (2.3)
$`y_j(m)=y_j(m1)\pm 1\text{ and }y_j(0)=y_j^i,y_j(t)=y_j^f.\}`$
With homogeneous weights away from the boundaries but extra weights at the boundaries the associated six-vertex model has been considered in . However, we note that in , and in most other studies of the six-vertex model, only such properties of the model are calculated that are averages over all numbers of walks, $`N`$. Here in contrast we are considering the generating function for a fixed number of walks, $`N`$, of a fixed finite length $`t`$. We will use the transfer matrix method from statistical mechanics to find this generating function in terms of a determinant of one-walk generating functions.
Precisely, we will show that the partition or generating function, $`\stackrel{=}{Z}_t^𝒩(𝐲^i𝐲^f)`$, for non-intersecting walk configurations of $`N`$ walks in which the $`j^{th}`$ walk starts at $`y_j^i`$ and arrives at $`y_j^f`$ after $`t`$ steps is given by the following determinant:
$$\stackrel{=}{Z}_t^𝒩(𝐲^i𝐲^f)=\left|\begin{array}{cccc}\stackrel{=}{Z}_t^𝒮(y_1^iy_1^f)& \stackrel{=}{Z}_t^𝒮(y_1^iy_2^f)& \mathrm{}& \stackrel{=}{Z}_t^𝒮(y_1^iy_N^f)\\ \stackrel{=}{Z}_t^𝒮(y_2^iy_1^f)& \stackrel{=}{Z}_t^𝒮(y_2^iy_2^f)& \mathrm{}& \stackrel{=}{Z}_t^𝒮(y_2^iy_N^f)\\ .& .& .& .\\ .& .& .& .\\ .& .& .& .\\ \stackrel{=}{Z}_t^𝒮(y_N^iy_1^f)& \stackrel{=}{Z}_t^𝒮(y_N^iy_2^f)& \mathrm{}& \stackrel{=}{Z}_t^𝒮(y_N^iy_N^f)\end{array}\right|$$
(2.4)
where $`\stackrel{=}{Z}_t^𝒮(y_j^iy_k^f)`$ is the generating function for configurations of a single walk starting at $`y_j^i`$ and ending at $`y_k^f`$ in a strip of width $`L`$. This is then a generalisation of the ‘master formulae’ of Fisher (equation (5.9) of ) and of Forrester (equation (4) of ) where unweighted non-intersecting walks are considered. Importantly, this determinantal result is precisely that given by the very general Gessel-Viennot Theorem for this problem: the walks considered on this lattice with the generalised weights considered satisfy the conditions of the theorem (see for a discussion of a one wall case with homogeneous weights).
## 3 From Bethe Ansatz to determinant
### 3.1 Transfer matrix formulation
The generating function of our walk problem in a strip can be formulated as the matrix element of a product of so-called transfer matrices (see below for their precise definition). This ‘transfer matrix’ contains the weights of all the possible edge configurations of $`N`$ walks of two adjacent columns of sites. As such it is related to an invariant subspace of the full six-vertex transfer matrix . In other words the non-zero elements of the six-vertex transfer matrix occur in diagonal blocks each of which is an “$`N`$-walk transfer matrix” for some $`N`$. There are of course several ways of setting up a strip transfer matrix. The most common are the site-to-site and the edge-to-edge matrices which add either edges or sites respectively to the walks. In the case of the six-vertex model the edge-to-edge matrix is usually chosen. Our choice of vertex-to-vertex (we shall use vertex and site interchangeably) matrix is determined by the requirements that the corresponding eigenvalue problem be as simple as possible while ensuring the non-intersecting constraint can be easily implemented.
The calculation of the generating function using the transfer matrix formulation then hinges on the spectral decomposition of the matrix. The matrices we diagonalise will be matrices that add two-steps to the evaluation of the generating function at a time. However the generating function is initially constructed in terms of one-step transfer matrices as these are useful since the non-intersecting condition is explicit.
###### Definition 1.
Let $`y\underset{L}{\overset{e}{𝒮}}`$ and $`y^{}\underset{L}{\overset{o}{𝒮}}`$. For $`N=1`$ the one-step transfer matrices are defined as
$`\left(\underset{1}{\overset{eo}{𝐓}}\right)_{y,y^{}}`$ $`=\{\begin{array}{cc}0\hfill & \text{if }|yy^{}|>1\hfill \\ w(y,y^{})\hfill & \text{if }|yy^{}|=1\hfill \end{array}`$ (3.5a)
and
$`\left(\underset{1}{\overset{oe}{𝐓}}\right)_{y^{},y}`$ $`=\{\begin{array}{cc}0\hfill & \text{if }|y^{}y|>1\hfill \\ w(y^{},y)\hfill & \text{if }|yy^{}|=1.\hfill \end{array}`$ (3.5b)
The $`N`$-walk transfer matrices for $`N>1`$ are constructed from sub-matrices of a direct product of the above $`N=1`$ matrices:
$`(\underset{N}{\overset{oe}{𝐓}})_{𝐲,𝐲^{}}`$ $`=\left({\displaystyle \underset{i=1}{\overset{N}{}}}\underset{1}{\overset{oe}{𝐓}}\right)_{𝐲,𝐲^{}}𝐲\underset{L}{\overset{o}{𝒰}}\text{ and }𝐲^{}\underset{L}{\overset{e}{𝒰}}`$ (3.6a)
and
$`(\underset{N}{\overset{eo}{𝐓}})_{𝐲^{},𝐲}`$ $`=\left({\displaystyle \underset{i=1}{\overset{N}{}}}\underset{1}{\overset{eo}{𝐓}}\right)_{𝐲^{},𝐲}𝐲^{}\underset{L}{\overset{e}{𝒰}}\text{ and }𝐲\underset{L}{\overset{o}{𝒰}}`$ (3.6b)
Note that since we are only considering $`L`$ being odd we have that $`\underset{N}{\overset{oe}{𝐓}}`$ and $`\underset{N}{\overset{eo}{𝐓}}`$ are square matrices and that in general $`\left(\underset{1}{\overset{oe}{𝐓}}\right)_{y^{},y}\left(\underset{1}{\overset{eo}{𝐓}}\right)_{y,y^{}}`$. More importantly, the restriction of the row and column spaces of the direct product eliminates the possibility of two walks arriving at the same lattice point. Furthermore the condition that the one-walker transfer matrix vanishes for $`|y^{}y|>1`$ prevents the generation of configurations in which pairs of walks “cross” *without* sharing a common lattice site (only nearest neighbour steps are allowed in all cases). This “non-crossing” condition is unnecessarily restrictive in the one-walk case. However, for $`N>1`$, if further neighbour steps are allowed then, in general, it is *not* possible to use the Bethe Ansatz. This condition is the analogue of the “non-crossing condition” of the Gessel-Viennot Theorem.
The generating function $`\stackrel{=}{Z}_t^𝒩(𝐲^i𝐲^f)`$ of $`N`$ non-intersecting walks of length $`t`$ in a strip is related to $`\stackrel{=}{Z}_{t1}^𝒩(𝐲^i𝐲)`$ by recurrence, the coefficients of which are the elements of one of the two one-step transfer matrices defined above. This relationship is given by the following lemma.
###### Lemma 1.
The generating function $`\stackrel{=}{Z}_t^𝒩(𝐲^i𝐲^f)`$ is given for $`t>0`$, depending on whether $`𝐲^f\underset{L}{\overset{e}{𝒰}}`$ or not, by
$$\stackrel{=}{Z}_t^𝒩(𝐲^i𝐲^f)=\{\begin{array}{c}_{𝐲\underset{L}{\overset{e}{𝒰}}}\stackrel{=}{Z}_{t1}^𝒩(𝐲^i𝐲)(\underset{N}{\overset{eo}{𝐓}})_{𝐲,𝐲^f}\text{for}𝐲^f\underset{L}{\overset{o}{𝒰}}\hfill \\ \\ _{𝐲\underset{L}{\overset{o}{𝒰}}}\stackrel{=}{Z}_{t1}^𝒩(𝐲^i𝐲)(\underset{N}{\overset{oe}{𝐓}})_{𝐲,𝐲^f}\text{for}𝐲^f\underset{L}{\overset{e}{𝒰}}.\hfill \end{array}$$
(3.7)
###### Proof.
A simple proof of this Lemma can be constructed using induction on $`t`$. ∎
Together with the initial condition
$$\stackrel{=}{Z}_0^𝒩(𝐲^i𝐲^f)=\delta _{𝐲^i,𝐲^f}V(𝐲^i),$$
(3.8)
where
$$V(𝐲^i)=\underset{\alpha =1}{\overset{N}{}}v(y_\alpha ^i),$$
(3.9)
equation (3.7) determines $`\stackrel{=}{Z}_t^𝒩(𝐲^i𝐲^f)`$.
A simple corollary of this Lemma (again shown by induction) is that the partition function can be written in terms of “two-step” transfer matrices.
###### Corollary.
The generating function $`\stackrel{=}{Z}_t^𝒩(𝐲^i𝐲^f)`$ is given, depending on whether $`t`$ is even or odd, as
$`\stackrel{=}{Z}_{2r}^𝒩(𝐲^i𝐲^f)`$ $`=\{\begin{array}{cc}V(𝐲^i)\left((\underset{N}{\overset{ee}{𝐓}})^r\right)_{𝐲^i,𝐲^f}\hfill & \text{for}𝐲^i\underset{L}{\overset{e}{𝒰}}\text{ and }𝐲^f\underset{L}{\overset{e}{𝒰}}\hfill \\ & \\ V(𝐲^i)\left((\underset{N}{\overset{oo}{𝐓}})^r\right)_{𝐲^i,𝐲^f}\hfill & \text{for}𝐲^i\underset{L}{\overset{o}{𝒰}}\text{ and }𝐲^f\underset{L}{\overset{o}{𝒰}}\hfill \end{array}`$ (3.10d)
or
$`\stackrel{=}{Z}_{2r+1}^𝒩(𝐲^i𝐲^f)`$ $`=\{\begin{array}{cc}V(𝐲^i)_{𝐲\underset{L}{\overset{e}{𝒰}}}\left((\underset{N}{\overset{ee}{𝐓}})^r\right)_{𝐲^i,𝐲}(\underset{N}{\overset{eo}{𝐓}})_{𝐲,𝐲^f}\hfill & \text{for}𝐲^i\underset{L}{\overset{e}{𝒰}}\text{ and }𝐲^f\underset{L}{\overset{o}{𝒰}}\hfill \\ & \\ V(𝐲^i)_{𝐲\underset{L}{\overset{o}{𝒰}}}\left((\underset{N}{\overset{oo}{𝐓}})^r\right)_{𝐲^i,𝐲}(\underset{N}{\overset{oe}{𝐓}})_{𝐲,𝐲^f}\hfill & \text{for}𝐲^i\underset{L}{\overset{o}{𝒰}}\text{ and }𝐲^f\underset{L}{\overset{e}{𝒰}}\hfill \end{array}`$ (3.10h)
for $`r=0,1,2,`$, where the two-step transfer matrices are defined as
$`\underset{N}{\overset{ee}{𝐓}}`$ $`=\underset{N}{\overset{eo}{𝐓}}\underset{N}{\overset{oe}{𝐓}}`$ (3.11a)
$`\underset{N}{\overset{oo}{𝐓}}`$ $`=\underset{N}{\overset{oe}{𝐓}}\underset{N}{\overset{eo}{𝐓}}.`$ (3.11b)
The two-step transfer matrices correspond to adding *two* steps to the paths. We show below that the generating function $`\stackrel{=}{Z}_t^𝒩(𝐲^i𝐲^f)`$ may in fact be expressed in terms of the eigenvalues and eigenvectors of the two-step matrices. In general the two-step matrices are *not* symmetric and hence one has to consider left and right eigenvectors. The $`N`$-walk eigenvectors will be constructed from the one-walk eigenvectors via the Bethe Ansatz. Used in this context the Ansatz expresses the components of the $`N`$-walk eigenvectors as a determinant of the components of one-walk eigenvectors (see (3.17) below).
### 3.2 From transfer matrices to determinants
This section contains our main results in the form of two theorems that state under what conditions the $`N`$-walk generating function can be written as the determinant (2.4). In summary, our theorems proven below show that the *equivalence* of the Bethe Ansatz in the form of equation (3.17) and the result of the Gessel-Viennot Theorem in the form (3.33) rests on showing that for any given problem the Ansatz is sufficiently good to provide a spanning (or complete) set of eigenvectors. This in turn depends only on the completeness of the one-walk eigenvectors (conditions of Lemma 2).
In this paper it is not our purpose to prove that all choices of the functions $`w(y,y^{})`$ which make up the elements of the one-walk transfer matrices allow for the conditions of our theorems to be satisfied. However, these theorems have been subsequently used to analyse some interesting cases in detail where as a consequence we show that an arbitrary boundary weight on one boundary (and special weight on the other) with homogeneous weights otherwise satisfies the conditions of the theorems. Also, we point out that since the Gessel-Viennot determinant holds for any function $`w(y,y^{})`$ it is almost certainly true that our conditions are satisfied always. It is however our main purpose here to demonstrate how the Bethe Ansatz gives the solution of the $`N`$-walk problem *given* the solution of the one-walk problem.
The first theorem states the conditions under which the $`N`$-walk transfer matrices can be diagonalised using a Bethe Ansatz: the major condition is that the one-walk transfer matrix problem can be solved — see Lemma 2. The second theorem basically states that if the Bethe Ansatz gives a complete set of eigenvectors for the $`N`$-walk problem for any $`N`$, then the $`N`$-walk generating function is a determinant of one-walk generating functions.
To begin we set up the conditions that define a solution of the one-walk problem: these will be the conditions our two theorems require.
###### Lemma 2.
Suppose that there exist linearly independent sets of column vectors $`\{\underset{k}{\overset{\text{R}}{\stackrel{o}{\phi }}}\}_{k𝒦_1}`$ and $`\{\underset{k}{\overset{\text{R}}{\stackrel{e}{\phi }}}\}_{k𝒦_1}`$, where $`𝒦_1`$ is some index set, which satisfy
$$\underset{1}{\overset{eo}{𝐓}}\underset{k}{\overset{\text{R}}{\stackrel{o}{\phi }}}=\lambda _k\underset{k}{\overset{\text{R}}{\stackrel{e}{\phi }}}\text{and}\underset{1}{\overset{oe}{𝐓}}\underset{k}{\overset{\text{R}}{\stackrel{e}{\phi }}}=\lambda _k\underset{k}{\overset{\text{R}}{\stackrel{o}{\phi }}}$$
(3.12)
with $`\lambda _k`$, and which span the column spaces of $`\underset{1}{\overset{eo}{𝐓}}`$ and $`\underset{1}{\overset{oe}{𝐓}}`$ respectively (in which case they are said to be complete). Further let $`\underset{1}{\overset{ee}{𝐓}}`$ and $`\underset{1}{\overset{oo}{𝐓}}`$ be defined by (3.11) then
* $`\underset{k}{\overset{\text{R}}{\stackrel{o}{\phi }}}`$ and $`\underset{k}{\overset{\text{R}}{\stackrel{e}{\phi }}}`$ are right eigenvectors of $`\underset{1}{\overset{oo}{𝐓}}`$ and $`\underset{1}{\overset{ee}{𝐓}}`$ respectively with eigenvalue $`\lambda _k^2`$.
* corresponding sets $`\{\underset{k}{\overset{\text{L}}{\stackrel{o}{\phi }}}\}_{k𝒦_1}`$ and $`\{\underset{k}{\overset{\text{L}}{\stackrel{e}{\phi }}}\}_{k𝒦_1}`$ of row vectors may be found such that
$$\underset{k}{\overset{\text{L}}{\stackrel{p}{\phi }}}\underset{k^{}}{\overset{\text{R}}{\stackrel{p}{\phi }}}=\delta _{k,k^{}}\text{and}\underset{k𝒦_1}{}\underset{k}{\overset{\text{R}}{\stackrel{p}{\phi }}}(y)\underset{k}{\overset{\text{L}}{\stackrel{p}{\phi }}}(y^{})=\delta _{y,y^{}}$$
(3.13)
for each $`p\{e,o\}`$, where the $``$ denotes complex conjugation. Note that the vectors of (3.13) have components indexed by $`y`$.
* the row vectors of (ii) satisfy
$$\underset{k}{\overset{\text{L}}{\stackrel{o}{\phi }}}\underset{1}{\overset{oe}{𝐓}}=\lambda _k\underset{k}{\overset{\text{L}}{\stackrel{e}{\phi }}}\text{and}\underset{k}{\overset{\text{L}}{\stackrel{e}{\phi }}}\underset{1}{\overset{eo}{𝐓}}=\lambda _k\underset{k}{\overset{\text{L}}{\stackrel{o}{\phi }}}$$
(3.14)
and also $`\underset{k}{\overset{\text{L}}{\stackrel{o}{\phi }}}`$ and $`\underset{k}{\overset{\text{L}}{\stackrel{e}{\phi }}}`$ are left eigenvectors of $`\underset{1}{\overset{oo}{𝐓}}`$ and $`\underset{1}{\overset{ee}{𝐓}}`$ respectively with eigenvalue $`\lambda _k^2`$.
Note that since the vectors in (3.12) span the space the cardinality of the index set $`𝒦_1`$ is $`(L+1)/2`$. The proof of the lemma is elementary linear algebra and we omit it. Notice that if $`\lambda _k`$ is a solution of (3.12) then so is $`\lambda _k`$ with vector $`\underset{k}{\overset{\text{R}}{\stackrel{e}{\phi }}}`$ replaced by $`\underset{k}{\overset{\text{R}}{\stackrel{e}{\phi }}}`$. These vectors are clearly not independent and normally sufficient independent vectors to form a spanning set are obtained by taking only the positive values of $`\lambda _k`$.
From the above left and right one-walk vectors we now construct the $`N`$-walk vectors and hence eigenvectors of $`\underset{N}{\overset{ee}{𝐓}}`$ and $`\underset{N}{\overset{oo}{𝐓}}`$.
###### Theorem 1.
Let $`\underset{N}{\overset{ee}{𝐓}}`$ and $`\underset{N}{\overset{oo}{𝐓}}`$, $`N>1`$, be given by equations (3.11). By imposing an arbitrary ordering on the elements of $`𝒦_1`$ define
$$𝒦_N=\{𝐤=(k_1,k_2,\mathrm{}k_N)|k_i𝒦_1\text{ and }k_1<k_2<\mathrm{}<k_N\}$$
(3.15)
and
$$\mathrm{\Lambda }_𝐤=\underset{\alpha =1}{\overset{N}{}}\lambda _{k_\alpha }.$$
(3.16)
(a) If for $`C\{L,R\}`$ and $`p\{e,o\}`$, $`\{\underset{k}{\overset{\text{C}}{\stackrel{p}{\phi }}}\}_{k𝒦_1}`$ satisfy the conditions of Lemma 2 then the vectors $`\{\underset{𝐤}{\overset{\text{C}}{\stackrel{p}{\mathrm{\Phi }}}}\}_{𝐤𝒦_N}`$ given by the Bethe Ansatz,
$`\underset{𝐤}{\overset{\text{C}}{\stackrel{p}{\mathrm{\Phi }}}}(𝐲)=`$ $`{\displaystyle \underset{\sigma P_N}{}}ϵ_\sigma {\displaystyle \underset{\alpha =1}{\overset{N}{}}}\underset{k_{\sigma _\alpha }}{\overset{\text{C}}{\stackrel{p}{\phi }}}(y_\alpha )={\displaystyle \underset{\sigma P_N}{}}ϵ_\sigma {\displaystyle \underset{\alpha =1}{\overset{N}{}}}\underset{k_\alpha }{\overset{\text{C}}{\stackrel{p}{\phi }}}(y_{\sigma _\alpha })𝐲\underset{L}{\overset{p}{𝒰}},`$ (3.17)
where $`P_N`$ is the set of $`N!`$ permutations of $`\{1,2,\mathrm{},N\}`$, $`\sigma =(\sigma _1,\sigma _2,\mathrm{},\sigma _N)P_N`$ and $`ϵ_\sigma `$ is the signature of the permutation $`\sigma `$, satisfy
$$\underset{N}{\overset{oe}{𝐓}}\underset{𝐤}{\overset{\text{R}}{\stackrel{e}{\mathrm{\Phi }}}}=\mathrm{\Lambda }_𝐤\underset{𝐤}{\overset{\text{R}}{\stackrel{o}{\mathrm{\Phi }}}}\text{and}\underset{N}{\overset{eo}{𝐓}}\underset{𝐤}{\overset{\text{R}}{\stackrel{o}{\mathrm{\Phi }}}}=\mathrm{\Lambda }_𝐤\underset{𝐤}{\overset{\text{R}}{\stackrel{e}{\mathrm{\Phi }}}}.$$
(3.18)
(b) Moreover the conclusions of parts (i), (ii) and (iii) of Lemma 2 hold with $`\underset{k}{\overset{\text{C}}{\stackrel{p}{\phi }}}`$ replaced by $`\underset{𝐤}{\overset{\text{C}}{\stackrel{p}{\mathrm{\Phi }}}}`$, $`\underset{1}{\overset{eo}{𝐓}}`$ and $`\underset{1}{\overset{oe}{𝐓}}`$ replaced by $`\underset{N}{\overset{eo}{𝐓}}`$ and $`\underset{N}{\overset{oe}{𝐓}}`$, $`𝒦_1`$ replaced by $`𝒦_N`$, and $`\lambda _k`$ replaced by $`\mathrm{\Lambda }_𝐤`$.
The proofs of part of this theorem and Theorem 2 require the following result.
###### Proposition 1.
For $`𝐤𝒦_N`$ and $`𝐲\underset{L}{\overset{p}{𝒰}}`$ let
$$\mathrm{\Phi }_𝐤(𝐲)=\underset{\sigma P_N}{}ϵ_\sigma \underset{\alpha =1}{\overset{N}{}}\varphi _{k_{\sigma _\alpha }}(y_\alpha )\text{and}\mathrm{\Psi }_𝐤(𝐲)=\underset{\sigma P_N}{}ϵ_\sigma \underset{\alpha =1}{\overset{N}{}}\psi _{k_{\sigma _\alpha }}(y_\alpha ).$$
(3.19)
Also let
$$f(𝐤)=\underset{\alpha =1}{\overset{N}{}}f(k_\alpha )$$
(3.20)
then
$$\underset{𝐤𝒦_N}{}f(𝐤)\mathrm{\Phi }_𝐤(𝐲)\mathrm{\Psi }_𝐤(𝐲^{})=\underset{\sigma P_N}{}ϵ_\sigma \underset{\alpha =1}{\overset{N}{}}\left(\underset{k_\alpha 𝒦_1}{}f(k_\alpha )\varphi _{k_\alpha }(y_\alpha )\psi _{k_\alpha }(y_{\sigma _\alpha }^{})\right)$$
(3.21)
and
$$\underset{𝐲\underset{L}{\overset{p}{𝒰}}}{}\mathrm{\Phi }_𝐤(𝐲)\mathrm{\Psi }_𝐤^{}(𝐲)=\underset{\sigma P_N}{}ϵ_\sigma \underset{\alpha =1}{\overset{N}{}}\left(\underset{y_\alpha \underset{L}{\overset{p}{𝒮}}}{}\varphi _{k_\alpha }(y_\alpha )\psi _{k_{\sigma _\alpha }^{}}(y_\alpha )\right)$$
(3.22)
###### Proof.
$`{\displaystyle \underset{𝐤𝒦_N}{}}f(𝐤)\mathrm{\Phi }_𝐤(𝐲)\mathrm{\Psi }_𝐤(𝐲^{})=`$ $`{\displaystyle \underset{\sigma P_N}{}}ϵ_\sigma {\displaystyle \underset{𝐤𝒦_N}{}}{\displaystyle \underset{\tau P_N}{}}ϵ_\tau {\displaystyle \underset{\alpha =1}{\overset{N}{}}}f(k_\alpha )\varphi _{k_{\tau _\alpha }}(y_\alpha )\psi _{k_{\sigma _\alpha }}(y_\alpha ^{})`$ (3.23)
$`=`$ $`{\displaystyle \underset{\sigma ^{}P_N}{}}ϵ_\sigma ^{}{\displaystyle \underset{𝐤𝒦_N}{}}{\displaystyle \underset{\tau P_N}{}}{\displaystyle \underset{\alpha =1}{\overset{N}{}}}f(k_\alpha )\varphi _{k_{\tau _\alpha }}(y_\alpha )\psi _{k_{\tau _\alpha }}(y_{\sigma _\alpha ^{}}^{})`$ (3.24)
The double sum over permutations, $`\tau `$ and $`𝐤𝒦_N`$ is equivalent to summing each $`k_\alpha `$ independently over $`𝒦_1`$ (terms for which two or more components of $`𝐤`$ are equal make zero contribution) and the first result follows. The second result follows in the same way by interchanging the roles of $`k`$ and $`y`$. ∎
###### Proof.
(of Theorem 1) (a) We first obtain the cyclic property (3.18) as follows.
$`\left(\underset{N}{\overset{oe}{𝐓}}\underset{𝐤}{\overset{\text{R}}{\stackrel{e}{\mathrm{\Phi }}}}\right)_𝐲`$ $`={\displaystyle \underset{𝐲^{}\underset{L}{\overset{e}{𝒰}}}{}}{\displaystyle \underset{\sigma P_N}{}}ϵ_\sigma \left(\underset{N}{\overset{oe}{𝐓}}\right)_{𝐲,𝐲^{}}{\displaystyle \underset{\alpha =1}{\overset{N}{}}}\underset{k_{\sigma _\alpha }}{\overset{\text{R}}{\stackrel{e}{\phi }}}(y_\alpha ^{})`$ (3.25a)
$`={\displaystyle \underset{\sigma P_N}{}}ϵ_\sigma {\displaystyle \underset{𝐲^{}\underset{L}{\overset{e}{𝒰}}}{}}{\displaystyle \underset{\alpha =1}{\overset{N}{}}}\left(\underset{1}{\overset{oe}{𝐓}}\right)_{y_\alpha ,y_\alpha ^{}}\underset{k_{\sigma _\alpha }}{\overset{\text{R}}{\stackrel{e}{\phi }}}(y_\alpha ^{})\left(\text{using}(\text{3.6})\right)`$ (3.25b)
$`={\displaystyle \underset{\sigma P_N}{}}ϵ_\sigma \left[{\displaystyle \underset{y_1^{}\underset{L}{\overset{e}{𝒮}}}{}}\left(\underset{1}{\overset{oe}{𝐓}}\right)_{y_1,y_1^{}}\underset{k_{\sigma _1}}{\overset{\text{R}}{\stackrel{e}{\phi }}}(y_1^{})\right]\mathrm{}\left[{\displaystyle \underset{y_N^{}\underset{L}{\overset{e}{𝒮}}}{}}\left(\underset{1}{\overset{oe}{𝐓}}\right)_{y_N,y_N^{}}\underset{k_{\sigma _N}}{\overset{\text{R}}{\stackrel{e}{\phi }}}(y_N^{})\right]`$ (3.25c)
$`={\displaystyle \underset{\sigma P_N}{}}ϵ_\sigma \left[\lambda _{k_{\sigma _1}}\underset{k_{\sigma _1}}{\overset{\text{R}}{\stackrel{o}{\phi }}}(y_1)\right]\mathrm{}\left[\lambda _{k_{\sigma _N}}\underset{k_{\sigma _N}}{\overset{\text{R}}{\stackrel{o}{\phi }}}(y_N)\right]\left(\text{using}(\text{3.12})\right)`$ (3.25d)
$`=\mathrm{\Lambda }_𝐤\underset{𝐤}{\overset{\text{R}}{\stackrel{o}{\mathrm{\Phi }}}}(𝐲)`$ (3.25e)
The critical step, and the whole reason for introducing the Bethe Ansatz, is to enable one to go from the restricted sums of (3.25b) to the unrestricted sums in (3.25c). This is justified for two reasons,
1. since $`\underset{𝐤}{\overset{\text{R}}{\stackrel{e}{\mathrm{\Phi }}}}(𝐲^{})`$ is a determinant, if any of the $`y_\alpha `$’s are equal then $`\underset{𝐤}{\overset{\text{R}}{\stackrel{e}{\mathrm{\Phi }}}}=0`$ – this allows the restriction $`y_1^{}<y_2^{}\mathrm{}<y_N^{}`$ on the sum to be relaxed to $`y_1^{}y_2^{}\mathrm{}y_N^{}`$
2. the $`y_\alpha `$ are in strictly increasing order combined with the fact that the matrix elements of $`\underset{1}{\overset{oe}{𝐓}}`$, are only non-zero if $`|y_\alpha y_\alpha ^{}|1`$ allows the restriction on the sum to be removed altogether.
The second part of (3.18) follows mutatis mutandis.
(b) (i) The vector $`\underset{𝐤}{\overset{\text{R}}{\stackrel{e}{\mathrm{\Phi }}}}`$ is a right eigenvector of $`\underset{N}{\overset{ee}{𝐓}}`$ with eigenvalue $`\mathrm{\Lambda }_𝐤^2`$, since
$$\underset{N}{\overset{ee}{𝐓}}\underset{𝐤}{\overset{\text{R}}{\stackrel{e}{\mathrm{\Phi }}}}=\underset{N}{\overset{eo}{𝐓}}\underset{N}{\overset{oe}{𝐓}}\underset{𝐤}{\overset{\text{R}}{\stackrel{e}{\mathrm{\Phi }}}}=\underset{N}{\overset{eo}{𝐓}}\mathrm{\Lambda }_𝐤\underset{𝐤}{\overset{\text{R}}{\stackrel{o}{\mathrm{\Phi }}}}=\mathrm{\Lambda }_𝐤^2\underset{𝐤}{\overset{\text{R}}{\stackrel{e}{\mathrm{\Phi }}}}$$
which follows from (3.18). Similarly $`\underset{𝐤}{\overset{\text{R}}{\stackrel{o}{\mathrm{\Phi }}}}`$ is a right eigenvector of $`\underset{N}{\overset{oo}{𝐓}}`$ with eigenvalue $`\mathrm{\Lambda }_𝐤^2`$.
(ii) Let us start by deriving the first result of part (ii), namely the orthogonality and normalisation condition. Using (3.22), for $`𝐤,𝐤^{}𝒦_N`$
$`{\displaystyle \underset{𝐲\underset{L}{\overset{p}{𝒰}}}{}}\underset{𝐤}{\overset{\text{L}}{\stackrel{p}{\mathrm{\Phi }}}}(𝐲)\underset{𝐤^{}}{\overset{\text{R}}{\stackrel{p}{\mathrm{\Phi }}}}(𝐲)=`$ $`{\displaystyle \underset{\sigma P_N}{}}ϵ_\sigma {\displaystyle \underset{\alpha =1}{\overset{N}{}}}\left({\displaystyle \underset{y_\alpha \underset{L}{\overset{p}{𝒮}}}{}}\underset{k_\alpha }{\overset{\text{L}}{\stackrel{p}{\phi }}}(y_\alpha )\underset{k_{\sigma _\alpha }^{}}{\overset{\text{R}}{\stackrel{p}{\phi }}}(y_\alpha )\right)`$
$`=`$ $`{\displaystyle \underset{\sigma P_N}{}}ϵ_\sigma {\displaystyle \underset{\alpha =1}{\overset{N}{}}}\delta _{k_\alpha ,k_{\sigma _\alpha }^{}}\left(\text{using}(\text{3.13})\right)`$
$`=`$ $`{\displaystyle \underset{\alpha =1}{\overset{N}{}}}\delta _{k_\alpha ,k_\alpha ^{}}`$
since the components of $`𝐤`$ and $`𝐤^{}`$ are in the same order only the identity permutation gives a non-zero delta function product. Thus
$$\underset{𝐤}{\overset{\text{L}}{\stackrel{p}{\mathrm{\Phi }}}}\underset{𝐤}{\overset{\text{R}}{\stackrel{p}{\mathrm{\Phi }}}}=\delta _{𝐤,𝐤^{}}.$$
(3.26)
Our derivation of the second part of (ii), namely the “completeness condition”, closely parallels the derivation of the orthogonality and normalisation condition. Using (3.21), for $`𝐲,𝐲^{}\underset{L}{\overset{p}{𝒰}}`$
$`{\displaystyle \underset{𝐤𝒦_N}{}}\underset{𝐤}{\overset{\text{R}}{\stackrel{p}{\mathrm{\Phi }}}}(𝐲)\underset{𝐤}{\overset{\text{L}}{\stackrel{p}{\mathrm{\Phi }}}}(𝐲^{})=`$ $`{\displaystyle \underset{\sigma P_N}{}}ϵ_\sigma {\displaystyle \underset{\alpha =1}{\overset{N}{}}}\left({\displaystyle \underset{k_\alpha 𝒦_1}{}}\underset{k_\alpha }{\overset{\text{R}}{\stackrel{p}{\phi }}}(y_\alpha )\underset{k_\alpha }{\overset{\text{L}}{\stackrel{p}{\phi }}}(y_{\sigma _\alpha }^{})\right)`$
$`=`$ $`{\displaystyle \underset{\sigma P_N}{}}ϵ_\sigma {\displaystyle \underset{\alpha =1}{\overset{N}{}}}\delta _{y_\alpha ,y_{\sigma _\alpha }^{}}\left(\text{using}(\text{3.13})\right)`$
$`=`$ $`{\displaystyle \underset{\alpha =1}{\overset{N}{}}}\delta _{y_\alpha ,y_\alpha ^{}}`$
so
$$\underset{𝐤𝒦_N}{}\underset{𝐤}{\overset{\text{R}}{\stackrel{p}{\mathrm{\Phi }}}}(𝐲)\underset{𝐤}{\overset{\text{L}}{\stackrel{p}{\mathrm{\Phi }}}}(𝐲^{})=\delta _{𝐲,𝐲^{}}$$
(3.27)
Notice that $`|𝒦_N|=\left(\genfrac{}{}{0pt}{}{\frac{1}{2}(L+1)}{N}\right)`$ which is the row (and column) space dimension, as it should be for completeness.
(iii) Using basic linear algebra gives
$$\underset{𝐤}{\overset{\text{L}}{\stackrel{o}{\mathrm{\Phi }}}}\underset{N}{\overset{oe}{𝐓}}=\underset{𝐤}{\overset{\text{L}}{\stackrel{e}{\mathrm{\Phi }}}}\mathrm{\Lambda }_𝐤\underset{𝐤}{\overset{\text{L}}{\stackrel{e}{\mathrm{\Phi }}}}\underset{N}{\overset{eo}{𝐓}}=\underset{𝐤}{\overset{\text{L}}{\stackrel{o}{\mathrm{\Phi }}}}\mathrm{\Lambda }_𝐤$$
(3.28)
and
$$\underset{𝐤}{\overset{\text{L}}{\stackrel{e}{\mathrm{\Phi }}}}\underset{N}{\overset{ee}{𝐓}}=\underset{𝐤}{\overset{\text{L}}{\stackrel{e}{\mathrm{\Phi }}}}\mathrm{\Lambda }_𝐤^2\underset{𝐤}{\overset{\text{L}}{\stackrel{o}{\mathrm{\Phi }}}}\underset{N}{\overset{oo}{𝐓}}=\underset{𝐤}{\overset{\text{L}}{\stackrel{o}{\mathrm{\Phi }}}}\mathrm{\Lambda }_𝐤^2$$
(3.29)
###### Lemma 3.
If the conditions of Theorem 1 hold then
$$\stackrel{=}{Z}_t^𝒩(𝐲^i𝐲^f)=V(𝐲^i)\underset{𝐤𝒦_N}{}\underset{𝐤}{\overset{\text{R}}{\stackrel{p^{}}{\mathrm{\Phi }}}}(𝐲^i)\mathrm{\Lambda }_𝐤^t\underset{𝐤}{\overset{\text{L}}{\stackrel{p}{\mathrm{\Phi }}}}(𝐲^f)𝐲^i\underset{L}{\overset{p^{}}{𝒰}}\text{and}𝐲^f\underset{L}{\overset{p}{𝒰}}$$
(3.30)
where if $`t`$ is even, $`p^{}=p`$ but otherwise $`p`$ and $`p^{}`$ are of opposite parity.
###### Proof.
If the conditions of Theorem 1 hold then we have that (3.26), (3.27), (3.28) and (3.29) are valid. Using (3.28) and (3.27) it follows that
$$(\underset{N}{\overset{oe}{𝐓}})_{𝐲,𝐲^{}}=\underset{𝐤𝒦_N}{}\mathrm{\Lambda }_𝐤\underset{𝐤}{\overset{\text{R}}{\stackrel{o}{\mathrm{\Phi }}}}(𝐲)\underset{𝐤}{\overset{\text{L}}{\stackrel{e}{\mathrm{\Phi }}}}(𝐲^{})\text{and}(\underset{N}{\overset{eo}{𝐓}})_{𝐲,𝐲^{}}=\underset{𝐤𝒦_N}{}\mathrm{\Lambda }_𝐤\underset{𝐤}{\overset{\text{R}}{\stackrel{e}{\mathrm{\Phi }}}}(𝐲)\underset{𝐤}{\overset{\text{L}}{\stackrel{o}{\mathrm{\Phi }}}}(𝐲^{}).$$
(3.31)
Also, using (3.29) and (3.27) it follows that
$$(\underset{N}{\overset{ee}{𝐓}})_{𝐲,𝐲^{}}=\underset{𝐤𝒦_N}{}\mathrm{\Lambda }_𝐤^2\underset{𝐤}{\overset{\text{R}}{\stackrel{e}{\mathrm{\Phi }}}}(𝐲)\underset{𝐤}{\overset{\text{L}}{\stackrel{e}{\mathrm{\Phi }}}}(𝐲^{})\text{and}(\underset{N}{\overset{oo}{𝐓}})_{𝐲,𝐲^{}}=\underset{𝐤𝒦_N}{}\mathrm{\Lambda }_𝐤^2\underset{𝐤}{\overset{\text{R}}{\stackrel{o}{\mathrm{\Phi }}}}(𝐲)\underset{𝐤}{\overset{\text{L}}{\stackrel{o}{\mathrm{\Phi }}}}(𝐲^{}).$$
(3.32)
Substituting these into (3.10) and using (3.26) gives the result immediately. ∎
###### Theorem 2.
If the conditions of Theorem 1 hold then
$$\stackrel{=}{Z}_t^𝒩(𝐲^i𝐲^f)=det||\stackrel{=}{Z}_t^𝒮(y_\alpha ^iy_\beta ^f)||_{\alpha ,\beta =1\mathrm{}N}.$$
(3.33)
###### Proof.
Using (3.9), (3.21) and (3.30) (which follows from (3.13) by Lemma 3)
$`\stackrel{=}{Z}_t^𝒩(𝐲^i𝐲^f)`$ $`={\displaystyle \underset{\sigma P_N}{}}ϵ_\sigma {\displaystyle \underset{\alpha =1}{\overset{N}{}}}\left(v(y_\alpha ^i){\displaystyle \underset{k_\alpha 𝒦_1}{}}\lambda _{k_\alpha }^t\underset{k_\alpha }{\overset{\text{R}}{\stackrel{p^{}}{\phi }}}(y_\alpha ^i)\underset{k_\alpha }{\overset{\text{L}}{\stackrel{p}{\phi }}}(y_{\sigma _\alpha }^f)\right)`$
$`={\displaystyle \underset{\sigma P_N}{}}ϵ_\sigma {\displaystyle \underset{\alpha =1}{\overset{N}{}}}\stackrel{=}{Z}_t^𝒮(y_\alpha ^iy_{\sigma _\alpha }^f).`$ (3.34)
which is an expansion of the required determinant. ∎
Finally, we point out that the result of Theorem 2 is also the conclusion of the Gessel-Viennot Theorem.
### 3.3 One wall and no wall geometries
When $`L>(t|y_N^iy_N^f|)/2+\mathrm{max}(y_N^i,y_N^f)`$ the walk closest to the wall at $`y=L`$ cannot touch it since the $`N^{th}`$ walk (of $`t`$ steps) needs at least $`|y_N^iy_N^f|`$ steps to go from $`y_N^i`$ to $`y_N^f`$ and any excursion close to the wall from the end point closest to the wall ($`\mathrm{max}(y_N^i,y_N^f)`$) needs just as many steps to return (hence the factor of $`1/2`$). Hence, when this condition holds, the strip generating function, $`\stackrel{=}{Z}_t^𝒮(y_\alpha ^iy_\beta ^f)`$ becomes equal to the generating function for walks that are affected by only one wall, $`\overline{Z}_t^𝒮(y_\alpha ^iy_\beta ^f)`$. Hence taking the limit $`L\mathrm{}`$ gives the following corollary:
###### Corollary.
For $`𝐲^i\underset{L}{\overset{p}{𝒰}}`$ and $`𝐲^f\underset{L}{\overset{p^{}}{𝒰}}`$, the $`N`$-walk generating function with only one wall at height $`y=0`$ is given by,
$$\overline{Z}_t^𝒩(𝐲^i𝐲^f)=det\overline{Z}_t^𝒮(y_\alpha ^iy_\beta ^f)_{\alpha ,\beta =1\mathrm{}N}.$$
(3.35)
If we also condition the walk closest to the wall at $`y=0`$ so that it cannot touch that wall we will end up with the “no boundary” results.
## Acknowledgements
Financial support from the Australian Research Council is gratefully acknowledged by RB and ALO. JWE is grateful for financial support from the Australian Research Council and for the kind hospitality provided by the University of Melbourne during which time this research was begun.
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# On Effect of Equilibrium Fluctuations on Superfluid Density in Layered Superconductors
## Abstract
We calculate suppression of inter- and intralayer superconducting currents due to equilibrium phase fluctuations and find that, in contrast to a recent prediction, the effect of thermal fluctuations cannot account for linear temperature dependence of the superfluid density in high-T<sub>c</sub> superconductors at low temperatures. Quantum fluctuations are found to dominate over thermal fluctuations at low temperatures due to hardening of their spectrum caused by the Josephson plasma resonance. Near T<sub>c</sub> sizeable thermal fluctuations are found to suppress the critical current in the stack direction stronger, than in the direction along the layers. Fluctuations of quasiparticle branch imbalance make the spectral density of voltage fluctuations at small frequencies non zero, in contrast to what may be expected from a naive interpretation of Nyquist formula.
One of important problems in layered high-T<sub>c</sub> cuprate superconductors that is still under discussion is the origin of the observed linear low-temperature dependence of the superfluid density (phase stiffness). The latter is directly related to magnetic penetration depths, $`\lambda _{}`$ and $`\lambda _{}`$, for a magnetic field screened by currents flowing in directions parallel and perpendicular to the superconducting planes, respectively . This dependence is usually attributed to contribution of quasiparticles near the nodes of the $`d`$-wave gap. According to the alternative explanation suggested in refs. the linear decrease in the temperature dependence of $`1/\lambda ^2`$ is induced entirely by classical thermal phase fluctuations. Recently the role of fluctuations was reconsidered for a $`d`$-wave superconductor by means of a microscopic approach within a functional integral framework, and quantum phase fluctuations were found to lead to a sizeable renormalization of the superfluid density, the effect of thermal fluctuations being small for $`T<T_c`$.
Thermal and quantum fluctuations considered in refs. are equilibrium fluctuations, therefore, the problem can be solved by means of a phenomenological approach based on the fluctuation–dissipation theorem. Such an approach is more simple and physically transparent, being in the same time more general because it is not restricted to a specific model of high-T<sub>c</sub> superconductivity. In the present work we reconsider the role of equilibrium phase fluctuations in layered superconductors applying the fluctuation-dissipation theorem to equations of the linear response, and come to conclusions different from those of refs. . Our calculations demonstrate that, in agreement with the results of ref. , at low temperatures quantum fluctuations dominate over thermal fluctuations due to Coulomb effects. This happens because non-uniform phase fluctuations induce a fluctuating electric field resulting in a finite energy of such fluctuations which is reflected in a finite value of the Josephson plasma frequency. So unlike refs. we assert that a contribution of thermal fluctuations cannot account for temperature dependence of the penetration depth. However, in contrast to the results of ref. we find that renormalization of the superfluid density due to fluctuations at low temperatures is not large. Near T<sub>c</sub>, when quasiparticle density dominates over the superfluid density, we find a different picture. In this case thermally induced fluctuations of the quasiparticle branch imbalance are found to be important, and a reduction of the superfluid density may become sizeable. We find that a suppression of the superconducting critical current due to thermal fluctuations is larger in the stack direction, than in the direction parallel to the conducting layers. This points out a possibility of a destruction of the superconductivity in the perpendicular direction at temperatures lower, than the critical temperature for the in-plane direction.
The fluctuation-dissipation theorem relates correlation functions of fluctuations to dissipative parts of kinetic coefficients describing the linear response of a system. So we start with the linear response of a layered superconductor. We consider expressions for current and charge density written in a form of functions of gauge invariant vector and scalar potentials. These potentials are
$$𝐏_s=(1/2)\chi (1/c)𝐀,\mu =(1/2)_t\chi +\mathrm{\Phi },$$
(1)
where $`𝐀`$ is the vector potential, $`\chi `$ is the phase of the order parameter, and $`\mathrm{\Phi }`$ is the electric potential (we use units with $`e=1`$, $`\mathrm{}=1`$, and $`k_B=1`$). The gauge invariant vector potential plays a role of the superconducting momentum, while the gauge invariant scalar potential is responsible for branch imbalance and charging effects. The electric field is expressed in terms of the potentials as
$$𝐄=_t𝐏_s\mu .$$
(2)
In layered superconductors the variables above are related to specific layers and must be supplied by layer numbers, e. g., the parallel component of $`𝐏_s`$ will be denoted as $`𝐏_n`$. Furthermore, the role of the perpendicular component of $`𝐏_s`$ plays the gauge invariant phase difference between the layers, $`\phi _n=\chi _{n+1}\chi _n2sA_{}/c`$, where $`s`$ is the lattice period in the stack direction. So, we may denote $`P_n=\phi _n/2s`$.
Equations determining current and charge densities read
$`𝐣_n={\displaystyle \frac{c^2}{4\pi \lambda _{}^2}}𝐏_n+\sigma _0_t𝐏_n\sigma _1_{}\mu _n,`$ (3)
$`j_n=j_c\mathrm{sin}\phi _n+{\displaystyle \frac{\sigma _0}{2s}}_t\phi _n\sigma _1_n\mu _n,`$ (4)
$`_t\rho _n=(\gamma _t+\nu _b){\displaystyle \frac{\kappa ^2}{4\pi }}\mu _n(\sigma _2_{}^2+\sigma _2_n^2)\mu _n+`$ (5)
$`_t(\sigma _1_{}𝐏_n+\sigma _1_n\phi _n/2s),`$ (6)
where $`_n\mu _n=(\mu _{n+1}\mu _n)/s`$ is the discrete version of the spatial derivative in the transverse direction, $`\kappa ^1`$ is the Thomas–Fermi screening radius, and $`\nu _b`$ is the branch imbalance relaxation rate.
The first terms in equations (3-4) describe the superconducting current for relative directions, and the rest terms are related for quasiparticle contributions. Equation (6) for charge density can be interpreted as a continuity equation for quasiparticles. It was derived for conventional superconductors in ref. and generalized later for layered superconductors with $`s`$-wave and $`d`$-wave pairing.
Note that the quasiparticle current does not have the form $`j=\sigma E`$ valid for a normal state. Instead, the response of a superconductor to the components of the field $`E`$ expressed in terms of the temporal derivative of $`𝐏_s`$ and of the spatial derivative of $`\mu `$ (cf. (2)) are described by different generalized conductivities $`\sigma _i,i=0,1,2`$. The respective parts of the electric field $`E`$ are induces by time dependent perturbations of current density, and by perturbations of charge density.
Explicit expressions for the coefficients in equations (3-6) depend on temperature, on scattering, and on a mechanism of the superconductivity (see , note that the explicit expressions for coefficients in ref. were presented in the dynamic limit $`\omega >Dq^2`$ where $`D`$ is a diffusion coefficient in quasiclassic approximation). At temperatures $`T0`$ all conductivities are small (e. g. in a $`s`$-wave superconductor they are exponentially small), and $`\gamma 1`$. At $`TT_c`$ the conductivities $`\sigma _i`$ approach the normal state conductivities. Furthermore, the differences between $`\sigma _i`$ for a given direction vanishes as $`\mathrm{\Delta }/T`$, $`\gamma \mathrm{\Delta }/T,\nu _b\mathrm{\Delta }/T`$, so that equation (6) approaches the continuity equation as $`\mathrm{\Delta }0`$. Strictly speaking products in (3-6) assume the convolution with respect to time and coordinates, i. e. in the Fourier transformed form the coefficients in the equations are frequency and wave vector dependent.
In order to derive equations relating fluctuations $`\delta \mu `$, $`\delta \phi `$ and $`\delta 𝐏`$ with fluctuations of charge and current densities, $`\delta \rho `$ and $`\delta 𝐣`$, we insert expression (3-6) into Maxwell equation
$$\times 𝐇=\frac{4\pi }{c}𝐣+\frac{1}{c}_t𝐃$$
and into the Poisson equation. Then making Fourier transformation with respect to time, to in-layer coordinates $`𝐫_{}`$, and to layer numbers $`n`$, we find
$`\left(\begin{array}{c}\delta j_{}\\ \delta j_{}\\ \delta \rho \end{array}\right)=\widehat{A}\left(\begin{array}{c}\delta \phi /2s\\ \delta 𝐏_{}\\ \delta \mu \end{array}\right),`$ (13)
$`\widehat{A}={\displaystyle \frac{4}{\pi }}(\begin{array}{ccc}ϵ(\omega \omega _0\stackrel{~}{\omega }_p^2)& c^2q_{}\widehat{q}_{}& ϵ\widehat{q}_{}\omega _1\\ c^2q_{}\widehat{q}_{}& \omega \mathrm{\Omega }_0\stackrel{~}{\mathrm{\Omega }}_p^2& q_{}\mathrm{\Omega }_1\\ ϵ\widehat{q}_{}\omega _1& q_{}\mathrm{\Omega }_1& (\gamma \frac{\nu _b}{i\omega })\kappa ^2+\frac{ϵ\widehat{q}_{}^2\omega _2+q_{}^2\mathrm{\Omega }_2}{\omega }\end{array}).`$ (17)
Here $`\widehat{q}_{}=(2/s)\mathrm{sin}(q_{}s/2)`$, where $`|q_{}|<\pi /s`$ is the wave number obtained from the discrete Fourier transformation with respect to layer numbers, $`\omega _i=\omega +i\omega _{ir}`$, $`\mathrm{\Omega }_i=\omega +i\mathrm{\Omega }_{ir}`$, $`\omega _{ir}=4\pi \sigma _iϵ`$ and $`\mathrm{\Omega }_{ir}=4\pi \sigma _i`$ are dielectric relaxation frequencies. Furthermore, $`\stackrel{~}{\omega }_p^2=\omega _p^2(1+\lambda _{}^2q_{}^2)`$, $`\stackrel{~}{\mathrm{\Omega }}_p^2=\mathrm{\Omega }_p^2(1+\lambda _{}^2\widehat{q}_{}^2)`$, where $`\mathrm{\Omega }_p=c/\lambda `$ and $`\omega _p=c/\lambda _{}\sqrt{ϵ}`$ are the plasma frequencies for directions parallel and perpendicular to the layers, $`ϵ`$ is a dielectric constant in transverse direction, $`\lambda _{}=c/\sqrt{8\pi sj_c}`$. The in-layer plasma frequency $`\mathrm{\Omega }_p`$, is much larger than typical frequencies of the problem which are of order of the Josephson plasma frequency, $`\omega _p`$.
Note that matrix $`\widehat{A}`$ satisfy the Onsager symmetry relations. Now we apply the fluctuation-dissipation theorem (cf. ) to equations (13).
From the expression for energy dissipation density $`Q`$ which in superconductors is given by
$$Q=\rho _t\mu +𝐣_t𝐏_s,$$
(18)
one can identify potentials $`\mu `$ and $`𝐏_s`$ with generalized forces related to variables $`\rho `$ and $`𝐣`$, respectively. So according to the fluctuation-dissipation theorem correlation functions of $`\delta j_{}`$, $`\delta j_{}`$ and $`\delta \rho `$ are determined by imaginary parts of related coefficients of the matrix $`\widehat{A}`$ in (13). For example, we find
$$\delta j_{}(𝐪,\omega )\delta j_{}(𝐪^{},\omega ^{})=(2\pi )^4(j_{}^2)_\omega \delta (\omega +\omega ^{})\delta (𝐪+𝐪^{})$$
(19)
with $`(j_{}^2)_\omega =2\stackrel{~}{T}\sigma _0^{}/s`$, where $`\stackrel{~}{T}=(\omega /2)\mathrm{coth}(\omega /2T)`$, $`\sigma ^{}=\mathrm{}\sigma `$. Similarly, $`(j_{}\rho )_\omega =2\stackrel{~}{T}\widehat{q}_{}\sigma _1^{}/s\omega `$, $`(\rho ^2)_\omega =2\stackrel{~}{T}(\nu _b/4\pi +q_{}^2\sigma _2^{}+q_{}^2\sigma _2^{})/s\omega ^2`$ and so on. An alternative way to calculate correlation functions based on the Langevin sources in equations of motion, and on the dissipation function of the system gives similar results.
Correlation functions of $`\delta 𝐏_s`$ and $`\delta \mu `$ are related to the inverse matrix $`\widehat{A}^1`$. For example, using the upper diagonal component of $`\widehat{A}^1`$ we can find the spectral density of phase difference fluctuations
$$(\delta \phi ^2)_\omega =\mathrm{}\frac{8\pi \stackrel{~}{T}(\omega \mathrm{\Omega }_0\stackrel{~}{\mathrm{\Omega }}_p^2)[(\gamma \frac{\nu _b}{i\omega })\kappa ^2+\frac{ϵ\widehat{q}_{}^2\omega _2+q_{}^2\mathrm{\Omega }_2}{\omega }]}{D}$$
(20)
where $`D`$ is the determinant of the matrix $`\widehat{A}`$.
Zeros of $`D`$ determine collective modes and penetration of electric and magnetic fields into a superconductor. They are important in calculation of spectral density of voltage noise. The voltage between contacts separated by $`N`$ layers, can be expressed according to eq.(2) as $`V=_t\underset{n=0}{\overset{N1}{}}\phi _n/2(\mu _N\mu _0)`$. Then we make Fourier transformation and calculate mean square value of the voltage using (20) and similar expressions for other correlation functions. We assume that the width of the contacts is large enough, therefore, we need functions at $`q_{}=0`$. Finally, after some algebra we find the expression for the spectral density of voltage fluctuations
$`(\delta V^2)_\omega ={\displaystyle \frac{16\stackrel{~}{T}}{S\epsilon }}{\displaystyle }dq_{}{\displaystyle \frac{\mathrm{sin}^2\frac{q_{}L}{2}}{\widehat{q}_{}^2}}\times `$ (21)
$`\mathrm{}{\displaystyle \frac{\omega (\omega \gamma +i\nu _b)\kappa ^2+ϵ\widehat{q}_{}^2[\omega (\omega _0+\omega _22\omega _1)\omega _p^2]}{(\omega _p^2\omega \omega _0)(\omega \gamma +i\nu _b)\kappa ^2+ϵ\widehat{q}_{}^2[\omega (\omega _1^2\omega _0\omega _2)+\omega _2\omega _p^2]}}`$ (22)
with $`L=Ns`$.
At low temperatures, $`T\mathrm{\Delta }`$, the relaxation frequencies are small in comparison to plasma frequencies $`\omega _p`$, and zero of the denominator gives the underdamped Josephson plasma mode at $`q_{}=0`$. In the opposite limit $`T\mathrm{\Delta }`$, near T<sub>c</sub>, the plasma frequencies are small, $`\mathrm{\Omega }_p^2,\omega _p^2N_s(\mathrm{\Delta }/T)^2`$ where $`N_s`$ is a fraction of the condensed electrons. On the other hand, all conductivities $`\sigma _i`$ and $`\sigma _i`$ at $`TT_c`$ approach the normal state conductivities for parallel and perpendicular directions, respectively. So near T<sub>c</sub> the dielectric relaxation frequencies are larger, than plasma frequencies for respective directions. Furthermore, calculations for $`s`$\- and $`d`$-wave pairing give $`\gamma =\pi \mathrm{\Delta }/4T`$ and $`\gamma =\pi \mathrm{\Delta }_0/2T\sqrt{i\nu /\omega }`$ , respectively. The branch imbalance relaxation rate is determined by energy relaxation for $`s`$-wave pairing and by elastic scattering for $`d`$-wave pairing, and in both cases $`\nu _b\mathrm{\Delta }/T`$. Then zeros of the denominator in equation (22) give the Carlson-Goldman mode for direction perpendicular to the layers, that is underdamped in the case of isotropic pairing in a narrow frequency region, $`\nu _b,\omega _p^2/\omega _r\omega (T/\mathrm{\Delta })\omega _p^2/\omega _r`$. The spectrum of the mode is $`\omega ^2ϵq_{}^2\gamma /\kappa ^2`$. This spectrum in isotropic superconductors was found in clean limit in ref. and in dirty limit in ref. . Evolution of the spectrum from the Josephson plasma mode to the anisotropic Carlson–Goldman mode in layered superconductors was studied in ref. . In the case of $`d`$-wave pairing $`\gamma `$ is not real, and the Carlson–Goldman mode is never underdamped (cf. ).
For $`N1`$ the leading contribution to the voltage fluctuations is
$$(\delta V^2)_\omega ^{(0)}=\frac{2\stackrel{~}{T}L}{S\sigma _0}\frac{\omega ^2\omega _{0r}^2}{(\omega _p^2\omega ^2)^2+\omega ^2\omega _{0r}^2}.$$
(23)
This expression corresponds to the Nyquist formula and exhibits the Josephson plasma resonance. At zero frequency this contribution vanishes, but there is an additional contribution to voltage fluctuations which remains finite at low frequencies $`\omega \nu _b`$. The latter contribution is related to the quasiparticle branch imbalance fluctuations, and is especially pronounced near $`T_c`$:
$$(\delta V^2)_\omega ^{(1)}=\frac{4T}{S\sigma _{}}\frac{l_E^2}{\sqrt{4l_E^2+s^2}},$$
(24)
where $`l_E^2=4\pi \sigma _{}/\nu _b\kappa ^2`$ (we assumed a thick sample, $`Ll_E`$). Equation (24) contradicts to a naive interpretation of the Nyquist theorem, according to which voltage noise at zero frequency is absent because the static resistivity of a superconductor is equal to zero. The non-zero contribution is related to a voltage drop near the superconductor boundary due to penetration of the electric field into layered superconductor over a distance needed for the branch imbalance relaxation. Though $`(\delta V^2)_\omega `$ in (24) is independent on the distance $`L`$ between the contacts, it may give a sizeable contribution which can be measured easily in small mesa structures.
Using equation (20) we can calculate mean square fluctuation of the phase difference.
$$\delta \phi ^2=(\delta \phi ^2)_\omega \frac{d\omega d𝐪}{(2\pi )^4}.$$
(25)
Since $`j_c\mathrm{sin}(\phi +\delta \phi )j_c(1\delta \phi ^2/2)\mathrm{sin}\phi `$, the latter term determines renormalization of the superfluid density. We calculate this renormalization, first, in the limit of low temperature.
At $`T\mathrm{\Delta }`$ the relaxation frequencies are small in comparison to the related plasma frequencies $`\omega _p`$ and $`\mathrm{\Omega }_p`$. We assume a simplifying condition $`\kappa s/ϵ1`$ which holds in high-T<sub>c</sub> superconductors, therefore, perturbations of $`\delta \mu `$ are small and can be neglected. Then the leading terms in the denominator of the matrix $`\widehat{A}`$ are
$$D=\omega \kappa ^2\mathrm{\Omega }_p^2\{ϵ(\omega _p^2\omega ^2)(1+\lambda _{}^2\widehat{q}_{}^2)+c^2q_{}^2$$
$$4\pi i\omega [\sigma _0(1+\lambda _{}^2\widehat{q}_{}^2)+\sigma _0ϵ(\omega _p^2\omega ^2+c^2q_{}^2ϵ)/\mathrm{\Omega }_p^2]\}.$$
A phase volume near zeros of $`D`$ related to the Josephson plasma mode contributes much to the integral (25) resulting in a finite fluctuations at $`T=0`$. Furthermore, the integral over $`q_{}`$ diverges at large $`q_{}`$, and we cut-off it at $`q_{}1/\xi _0`$, where $`\xi _0`$ is the superconducting correlation length. In dimensional units we obtain
$`\delta \phi ^2{\displaystyle \frac{16}{\pi \sqrt{ϵ}}}{\displaystyle \frac{e^2}{\mathrm{}c}}{\displaystyle \frac{\lambda _{}}{\xi _0}}\left(1+{\displaystyle \frac{T^2}{T_0^2}}\right),`$ (26)
$`T_0{\displaystyle \frac{\mathrm{}\mathrm{\Omega }_p}{2\pi k_B}}\sqrt{{\displaystyle \frac{3\mathrm{\Omega }_ps}{2\sigma _0\xi _0\sqrt{ϵ}\mathrm{ln}\lambda _J/\xi _0}}}.`$ (27)
This result differs from that obtained in refs. . For parameter values typical for layered cuprates we find $`T_0T_c`$. The magnitude of $`\delta \phi ^2`$ is not large because the large factor $`\lambda /\xi _0`$ in (26) is multiplied by the small fine structure constant. Thus the renormalization of the penetration depths due to phase fluctuations at low temperatures is practically temperature independent, and is not large.
Calculation of $`\delta P_{}^2\xi _0^2`$ which determines suppression of the superfluid density in the in-plane direction gives value similar to (26).
Now we calculate fluctuations at high temperatures, $`TT_c`$. Since at such temperatures fluctuations are quasi-stationary they can be found using the standard approach based on functional integration of the free energy of the system, which includes energy of the magnetic field and of the superconducting current. We calculate $`\delta \phi ^2`$ and $`\delta P_{}^2\xi _0^2`$, which determine the renormalization of the superfluid stiffness in perpendicular and parallel directions, respectively, using again the cut-off at $`q_{}1/\xi _0`$. Then in the limit $`\xi _0\lambda _J`$, which definitely holds for layered high-T<sub>c</sub> superconductors, we find in dimensional units
$$\delta \phi ^2\frac{16e^2k_BT\lambda _{}^2}{\mathrm{}^2c^2s}\mathrm{ln}\frac{\lambda _{}}{\xi _0}4\mathrm{ln}\frac{\lambda _{}}{\xi _0}\delta P_{}^2\xi _0^2.$$
(28)
Since in layered high-T<sub>c</sub> superconductors $`\lambda _{}`$ is by few orders of magnitude larger, than the correlation length $`\xi _0`$, the logarithm in equation (28) is large as well. Then suppression of the critical current by fluctuations in the direction perpendicular to the layers is larger, than in the parallel direction. Using parameters typical for BSCCO, $`\lambda _{}(0)1.5\times 10^5`$ cm, $`s1.5\times 10^7`$ cm and $`\lambda _{}/\xi _010^4`$, we estimate
$$\delta \phi ^20.5\left[\frac{\lambda _{}(T)}{\lambda _{}(0)}\right]^2\frac{T}{80\text{K}}.$$
Since $`[\lambda _{}(T)/\lambda _{}(0)]^2`$ diverges as $`TT_c`$ we conclude that thermal fluctuations may lead to a sizeable reduction of the critical current. Note that we study fluctuations in the linear approximation and, hence, do not take into account that the phase perturbations due to thermal fluctuations may overcome a finite potential barrier. The latter process would result in a destruction of a superconducting current analogous to Josephson junctions .
Observation of smaller values of T<sub>c</sub> in the stack direction, than in direction parallel to the layers, was reported in YBCO single crystals with low oxygen content . Such crystals are expected to have large $`\lambda _{}`$, which according to (28) is in favor of large fluctuations. However, it is difficult to explain by the fluctuation mechanism so large differences between critical temperatures observed in ref. for different directions.
We are grateful to A. A. Varlamov for attracting our attention to the results of ref. and for useful discussion, and to V. F. Gantmakher, K. E. Nagaev and A. Ya. Shulman for useful discussion. This work was supported by project 98-02-17221 of Russian Foundation for Basic Research, and by project 96053 of Russian program on superconductivity.
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# 1 Introduction
## 1 Introduction
Differentially rotating disks of dust have already been studied by Ansorg and Meinel . They considered the class of hyperelliptic solutions to the Ernst equation introduced by Meinel and Neugebauer , see also -. These hyperelliptic solutions depend on a number of complex parameters and a real potential function. Ansorg and Meinel concentrated on the case in which one complex parameter can be prescribed. They determined the real potential function in order to satisfy a particular boundary condition valid for all disks of dust. To generate their solutions, they used Neugebauer’s and Meinel’s rigorous solution to the boundary value problem of a rigidly rotating disk of dust which also belongs to the hyperelliptic class.
A subclass of Ansorg’s and Meinel’s solutions is made up of Bäcklund transforms of seed solutions of the Weyl class<sup>2</sup><sup>2</sup>2The construction of solutions to the Ernst equation by means of Bäcklund transformations belongs to the powerful analytic methods developed by several authors -. For a detailed introduction see .. Solutions of this type are of particular interest since their mathematical structure is much simpler than that of the more general hyperelliptic solutions.
With this in mind, the following questions arise:
* Is it possible to find solutions corresponding to more general differentially rotating disks of dust by increasing the number of prescribed complex parameters?
* If so, is there a rapidly converging method for approximating arbitrary differentially rotating disks of dust with given boundary conditions (i. e. proper mass density or angular velocity)?
* Is it perhaps possible to construct such a method by restriction to the much simpler solutions of the Bäcklund type?
To answer these questions, the paper is organized as follows. In the first section the metric tensor, Ernst equation, and boundary conditions are introduced and the class of solutions of the Bäcklund type is represented. As will be discussed in the second section, the properties of these solutions can be used to obtain more general solutions by a suitable limiting process. Since these more general solutions depend on two real analytic functions defined on the interval $`[0,1]`$, a rapidly converging numerical scheme to satisfy arbitrary boundary conditions for disks of dust can be created. This is depicted in the third section. Finally, the fourth section contains particular examples of differentially rotating disks of dust, including disks with a realistic profile for the angular velocity and more exotic disks possessing two spatially separated ergoregions.
In what follows, units are used in which the velocity of light as well as Newton’s constant of gravitation are equal to 1.
### 1.1 Metric Tensor, Ernst equation, and boundary conditions
The metric tensor for axisymmetric stationary and asymptotically flat space-times reads as follows in Weyl-Papapetrou-coordinates $`(\rho ,\zeta ,\phi ,t)`$:
$$ds^2=e^{2U}[e^{2k}(d\rho ^2+d\zeta ^2)+\rho ^2d\phi ^2]e^{2U}(dt+ad\phi )^2.$$
For this line element, the vacuum field equations are equivalent to a single complex equation – the so-called Ernst equation
$$(\mathrm{}f)\mathrm{}f=(f)^2,$$
(1)
$$\mathrm{}=\frac{^2}{\rho ^2}+\frac{1}{\rho }\frac{}{\rho }+\frac{^2}{\zeta ^2},=(\frac{}{\rho },\frac{}{\zeta }),$$
where the Ernst potential $`f`$ is given by
$$f=e^{2U}+\text{i}b\text{with}b_{,\zeta }=\frac{e^{4U}}{\rho }a_{,\rho },b_{,\rho }=\frac{e^{4U}}{\rho }a_{,\zeta }.$$
(2)
The remaining function $`k`$ can be calculated from the Ernst potential $`f`$ by a line integration:
$$\frac{k_{,\rho }}{\rho }=(U_{,\rho })^2(U_{,\zeta })^2+\frac{1}{4}e^{4U}[(b_{,\rho })^2(b_{,\zeta })^2]$$
$$\frac{k_{,\zeta }}{\rho }=2U_{,\rho }U_{,\zeta }+\frac{1}{2}e^{4U}b_{,\rho }b_{,\zeta }.$$
To obtain the boundary conditions for differentially rotating disks of dust, one has to consider the field equations for an energy-momentum-tensor
$$T^{ik}=ϵu^iu^k=\sigma _p(\rho )e^{Uk}\delta (\zeta )u^iu^k,$$
where $`ϵ`$ and $`\sigma _p`$ stand for the energy-density and the invariant (proper) surface mass-density, respectively, $`\delta `$ is the usual Dirac delta-distribution, and $`u^i`$ denotes the four-velocity of the dust material<sup>3</sup><sup>3</sup>3$`u^i`$ has only $`\phi `$\- and $`t`$\- components..
Integration of the corresponding field equations from the lower to the upper side of the disk (with coordinate radius $`\rho _0`$) yields the conditions (see , pp. 81-83)
$$2\pi \sigma _p=e^{Uk}(U_{,\zeta }+\frac{1}{2}Q)$$
(3)
$$e^{4U}Q^2+Q(e^{4U})_{,\zeta }+(b_{,\rho })^2=0$$
(4)
for $`\zeta =0^+,\mathrm{\hspace{0.33em}0}\rho \rho _0`$ and
$$Q=\rho e^{4U}[b_{,\rho }b_{,\zeta }+(e^{2U})_{,\rho }(e^{2U})_{,\zeta }].$$
(5)
Note that boundary condition (4) for the Ernst potential $`f`$ does not involve the surface mass-density $`\sigma _p`$. This condition comes from the nature of the material the disk is made of. Therefore, equation (4) will be referred to as the dust-condition.
Instead of prescribing the proper surface mass-density $`\sigma _p`$ \[which leads to the boundary condition (3)\] one can alternatively assume a given angular velocity $`\mathrm{\Omega }=\mathrm{\Omega }(\rho )=u^\phi /u^t`$ of the disk which results in the boundary condition ($`\zeta =0^+`$, $`0\rho \rho _0`$):
$$\mathrm{\Omega }=\frac{Q}{a_{,\zeta }aQ}.$$
(6)
The following requirements due to symmetry conditions and asymptotical flatness complete the set of boundary conditions:
* Regularity at the rotation axis is guaranteed by
$$\frac{f}{\rho }(0,\zeta )=0.$$
* At infinity asymptotical flatness is realized by $`U0`$ and $`a0`$. For the potential $`b`$ this has the consequence $`bb_{\mathrm{}}=\text{const}`$. Without loss of generality, this constant can be set to $`0`$, i.e. $`f1`$ at infinity.
* Finally, reflectional symmetry with respect to the plane $`\zeta =0`$ is assumed, i.e. $`f(\rho ,\zeta )=\overline{f(\rho ,\zeta )}`$ (with a bar denoting complex conjugation).
### 1.2 Solutions of the Bäcklund type
For a given integer $`q1`$, a set $`\{Y_1,\mathrm{},Y_q\}=\{Y_\nu \}_q`$<sup>4</sup><sup>4</sup>4In the following, the notation $`\{Y_1,\mathrm{},Y_q\}`$ will be abbreviated by $`\{Y_\nu \}_q`$. of complex parameters, and a real analytic function $`g`$ defined on the interval $`[0,1],`$ the following expression
$`f=f_0{\displaystyle \frac{\left|\begin{array}{cccccccc}1& 1& 1& 1& 1& \mathrm{}& 1& 1\\ 1& \alpha _1\lambda _1& \alpha _1^{}\lambda _1^{}& \alpha _2\lambda _2& \alpha _2^{}\lambda _2^{}& \mathrm{}& \alpha _q\lambda _q& \alpha _q^{}\lambda _q^{}\\ 1& \lambda _1^2& (\lambda _1^{})^2& \lambda _2^2& (\lambda _2^{})^2& \mathrm{}& \lambda _q^2& (\lambda _q^{})^2\\ 1& \alpha _1\lambda _1^3& \alpha _1^{}(\lambda _1^{})^3& \alpha _2\lambda _2^3& \alpha _2^{}(\lambda _2^{})^3& \mathrm{}& \alpha _q\lambda _q^3& \alpha _q^{}(\lambda _q^{})^3\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 1& \lambda _1^{2q}& (\lambda _1^{})^{2q}& \lambda _2^{2q}& (\lambda _2^{})^{2q}& \mathrm{}& \lambda _q^{2q}& (\lambda _q^{})^{2q}\end{array}\right|}{\left|\begin{array}{cccccccc}1& 1& 1& 1& 1& \mathrm{}& 1& 1\\ 1& \alpha _1\lambda _1& \alpha _1^{}\lambda _1^{}& \alpha _2\lambda _2& \alpha _2^{}\lambda _2^{}& \mathrm{}& \alpha _q\lambda _q& \alpha _q^{}\lambda _q^{}\\ 1& \lambda _1^2& (\lambda _1^{})^2& \lambda _2^2& (\lambda _2^{})^2& \mathrm{}& \lambda _q^2& (\lambda _q^{})^2\\ 1& \alpha _1\lambda _1^3& \alpha _1^{}(\lambda _1^{})^3& \alpha _2\lambda _2^3& \alpha _2^{}(\lambda _2^{})^3& \mathrm{}& \alpha _q\lambda _q^3& \alpha _q^{}(\lambda _q^{})^3\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 1& \lambda _1^{2q}& (\lambda _1^{})^{2q}& \lambda _2^{2q}& (\lambda _2^{})^{2q}& \mathrm{}& \lambda _q^{2q}& (\lambda _q^{})^{2q}\end{array}\right|}}`$ (19)
with (a bar denotes complex conjugation)
$$f_0=\mathrm{exp}(\underset{1}{\overset{1}{}}\frac{(1)^qg(x^2)dx}{Z_D}),Z_D=\sqrt{(\text{i}x\zeta /\rho _0)^2+(\rho /\rho _0)^2},\mathrm{}(Z_D)<0$$
$$\lambda _\nu =\sqrt{\frac{Y_\nu \text{i}\overline{z}}{Y_\nu +\text{i}z}},z=\frac{1}{\rho _0}(\rho +\text{i}\zeta ),\lambda _\nu ^{}\overline{\lambda }_\nu =1$$
$$\alpha _\nu =\frac{1\gamma _\nu }{1+\gamma _\nu },\gamma _\nu =\mathrm{exp}\left(\lambda _\nu (Y_\nu +\text{i}z)\underset{1}{\overset{1}{}}\frac{(1)^qg(x^2)dx}{(\text{i}xY_\nu )Z_D}\right),\alpha _\nu ^{}\overline{\alpha }_\nu =1$$
satisfies the Ernst equation. With the additional requirement that for each parameter $`Y_\nu `$ there is also a parameter $`Y_\mu `$ with $`Y_\nu =\overline{Y}_\mu `$, reflectional symmetry, $`f(\rho ,\zeta )=\overline{f(\rho ,\zeta )}`$, is ensured<sup>5</sup><sup>5</sup>5Hence, the set $`\{\text{i}Y_\nu \}_q`$ consists of real parameters and/or pairs of complex conjugate parameters.. Moreover, the parameters $`Y_\nu `$ are assumed to lie outside the imaginary interval $`[\text{i},\text{i}]`$.
The above Ernst potential $`f=f(\rho /\rho _0,\zeta /\rho _0;\{Y_\nu \}_q;g)`$ is obtained by a Bäcklund transformation applied to the real seed solution $`f_0`$, see . On the other hand, as demonstrated in appendix A, it can be constructed from the hyperelliptic solutions by a suitable limiting process (see also ). The particular ansatz chosen for the seed solution $`f_0`$ guarantees a resulting Ernst potential which corresponds to a disk-like source of the gravitational field (see also section 1.2 of ).
Furthermore, $`f`$ does not possess singularities at $`(\rho ,\zeta )=\rho _0(|\mathrm{}[Y_\nu ]|,\mathrm{}[Y_\nu ])`$. This is due to the fact that $`\alpha _\nu \lambda _\nu `$ is a function of $`\lambda _\nu ^2`$, and this means that $`f`$ does not behave like a square root function near the critical points $`(\rho ,\zeta )=\rho _0(|\mathrm{}[Y_\nu ]|,\mathrm{}[Y_\nu ])`$, but rather like a rational function. Now, in the whole area of physically interesting solutions that will be treated in the subsequent sections, each zero of the denominator is cancelled by a corresponding zero of the numerator in (19) such that the resulting gravitational field is regular outside the disk.
The real function $`g`$ that enters the Ernst potential is assumed to be analytic on $`[0,1]`$ in order to guarantee an analytic behaviour of the angular velocity $`\mathrm{\Omega }`$ for all $`\rho [0,\rho _0]`$. Moreover, the additional requirement
$$g(1)=0$$
leads to a surface mass density $`\sigma _p`$ of the form
$$\sigma _p(\rho )=\sigma _0\psi _p[(\rho /\rho _0)^2]\sqrt{1(\rho /\rho _0)^2}\text{[with }\psi _p\text{ analytic in }[0,1],\psi _p(0)=1\text{}$$
(20)
and therefore ensures that $`\sigma _p`$ vanishes at the rim of the disk.
In this article the question as to whether the above expression for the Ernst potential is sufficiently general to approximate arbitrary differentially rotating disks of dust is investigated. Of particular interest is a rapidly converging method to perform this approximation. To this end, the set $`\{Y_\nu \}_q`$ of complex parameters will be translated into an analytic function
$$\xi :[0,1].$$
Thus the Ernst potential will depend on two real analytic functions defined on $`[0,1]`$:
$$f=f(\rho /\rho _0,\zeta /\rho _0;\xi ;g),$$
which eventually proves to be sufficient to satisfy both the dust condition (4) and the boundary condition (3) \[or alternatively (6)\]. The rapid and accurate approximation can be realized since both $`g`$ and $`\xi `$ are analytic on $`[0,1]`$ and thus permit elegant expansions in terms of Chebyshev polynomials.
## 2 Generalization of the Bäcklund type solutions by a limiting process
As demonstrated in for the Bäcklund type solutions with $`q=1`$, the dust condition (4) can be satisfied by an appropriate choice of the function $`g`$ if the complex parameters $`Y_\nu `$ are prescribed. To fulfil a second boundary condition, (3) or (6), the set $`\{Y_\nu \}_q`$ of these parameters has to be translated into a real analytic function $`\xi `$.
To this end, consider the following equalities for the above solutions $`f=f(\{Y_\nu \}_q;g)`$<sup>6</sup><sup>6</sup>6In the following the Ernst potentials $`f`$ given by (19) are considered as complex functions depending on the set $`\{Y_\nu \}_q`$ of complex parameters and on $`g`$. which are proved in appendix B:
$`f[\{Y_1,\mathrm{},Y_{q2},Y_{q1},Y_q\};g]`$ $`=`$ $`f[\{Y_1,\mathrm{},Y_{q2}\};g]`$
if $`Y_{q1}=Y_q`$
$`f[\{Y_1,\mathrm{},Y_{q2},Y_{q1},Y_q\};g]`$ $`=`$ $`f[\{Y_1,\mathrm{},Y_{q2}\};g]`$
$`\text{if }Y_{q1}=\overline{Y}_q`$
$`\underset{t\mathrm{}}{lim}f[\{Y_1,\mathrm{},Y_{q1},\text{i}t\};g]`$ $`=`$ $`f[\{Y_1,\mathrm{},Y_{q1}\};g]`$
if $`t`$
$`\underset{Y_q\mathrm{}}{lim}f[\{Y_1,\mathrm{},Y_{q2},Y_{q1},Y_q\};g]`$ $`=`$ $`f[\{Y_1,\mathrm{},Y_{q2}\};g]`$
if $`Y_{q1}=\overline{Y}_q.`$
In order to find an approximation scheme, the desired function $`\xi =\xi (\{Y_\nu \}_q)`$ is supposed to be invariant under the modifications (2-2) of the set $`\{Y_\nu \}_q`$ that do not effect the Ernst potential. This property will be necessary to solve the boundary conditions uniquely.
It is realized by the real analytic function
$$\xi (x^2;\{Y_\nu \}_q)=\frac{1}{x}\mathrm{ln}\left[\underset{\nu =1}{\overset{q}{}}\frac{\text{i}Y_\nu x}{\text{i}Y_\nu +x}\right],x[1,1],$$
(25)
which can be proved by considering that for each parameter $`Y_\nu `$ there is also a parameter $`Y_\mu `$ with $`Y_\nu =\overline{Y}_\mu `$, and that, moreover, the parameters $`Y_\nu `$ do not lie on the imaginary interval $`[\text{i},\text{i}]`$.
The set $`𝒳`$ of all functions $`\xi =\xi (x^2;\{Y_\nu \}_q),q`$, which are defined by (25) forms a dense subset of the set $`𝒜`$ of all real analytic functions on $`[0,1]`$. Now, for a given function $`g`$, each $`\xi 𝒳`$ is mapped by (19) onto a uniquely defined Ernst potential $`f`$ <sup>7</sup><sup>7</sup>7Here, $``$ denotes the set of all Ernst potentials corresponding to disk-like sources. :
$$\mathrm{\Phi }_g:𝒳,\mathrm{\Phi }_g(\xi )=f(\{Y_\nu \}_q;g),$$
(26)
where the set $`\{Y_\nu \}_q`$ results from $`\xi `$ by (25).
In the following, it is assumed that this mapping $`\mathrm{\Phi }_g`$ can be extended to form a continuous function defined on $`𝒜`$.<sup>8</sup><sup>8</sup>8The mathematical aspects of this assumption will be discussed in section 5. Then, given the two real functions $`g`$ and $`\xi `$, defined and analytic on the interval $`[0,1]`$, the Ernst potential
$$f(\xi ;g)=\underset{q\mathrm{}}{lim}f(\{Y_\nu ^{(q)}\}_q;g)$$
exists and is independent of the particular choice of the sequence $`\{\{Y_\nu ^{(q)}\}_q\}_{q=q_0}^{\mathrm{}}`$ which serves to represent $`\xi `$ by
$`\xi (x^2)={\displaystyle \frac{1}{x}}\underset{q\mathrm{}}{lim}\mathrm{ln}\left[{\displaystyle \underset{\nu =1}{\overset{q}{}}}{\displaystyle \frac{\text{i}Y_\nu ^{(q)}x}{\text{i}Y_\nu ^{(q)}+x}}\right]\text{ for}x[1,1].`$
This provides the groundwork for the approximation scheme that will be developed in the next section. The treatment additionally assumes that the boundary conditions (3) and (4) \[or (4) and (6)\] interpreted as functions of $`g`$ and $`\xi `$ are invertible. The accurate and rapid convergence of the numerical methods justifies this assumption although a rigorous proof cannot be given.
## 3 An approximation scheme for arbitrary differentially rotating disks of dust
It is now possible to attack general boundary value problems for differentially rotating disks of dust. With the above generalized solutions $`f=f(\xi ;g)`$ the boundary conditions \[see formulas (3-6, 20)\] become a problem of inversion to determine $`g`$ and $`\xi `$ from $`\sigma _p`$ or $`\mathrm{\Omega }`$:
(A) $`S(g;\xi )=\{e^{Uk}[U_{,\zeta }+\frac{1}{2}Q]/[\sigma _0\sqrt{1(\rho /\rho _0)^2}]\}(\xi ;g)2\pi \psi _p\text{or}`$ (A’) $`O(g;\xi )=\{Q/[\mathrm{\Omega }(0)(a_{,\zeta }aQ)]\}(\xi ;g)\mathrm{\Omega }/\mathrm{\Omega }(0)=\mathrm{\Omega }^{}`$ (B) $`D(g;\xi )=\left\{\rho _0^2\left[Q^2e^{4U}+Q(e^{4U})_{,\zeta }+(b_{,\rho })^2\right]\right\}(\xi ;g)0,g(1)0`$
This inversion problem is tackled in the following manner:
1. The only way to treat the complicated system (LABEL:AS1) numerically seems to be by restricting it to a finite, discretized version and solving this by means of a Newton-Raphson method.
2. For this method, a good initial guess for the solution is needed. As shown in appendix C.1, there exists a representation of the functions $`g`$ and $`\xi `$ in terms of $`\sigma _p`$ or $`\mathrm{\Omega }`$ in the Newtonian regime $`\epsilon 1`$ where $`\epsilon =M^2/J`$ and the gravitational mass $`M`$ and the total angular momentum $`J`$ are given by
$$M=2\underset{S}{}(T_{ab}\frac{1}{2}Tg_{ab})n^a\xi ^b𝑑V,J=\underset{S}{}T_{ab}n^a\eta ^b𝑑V,T_{ab}=g_{ab}T^{ab}.$$
(28)
($`S`$ is the spacelike hypersurface $`t=constant`$ with the unit future-pointing normal vector $`n^a`$; the Killingvectors $`\xi ^a`$ and $`\eta ^a`$ correspond to stationarity and axisymmetry, respectively.)
3. This motivates the following finite version which results from expansions of (LABEL:AS1) in terms of Chebyshev-polynomials $`T_j(\tau )=\mathrm{cos}[j\mathrm{arccos}(\tau )]`$:
$`F_j(v_k)0(1j,kN_1+N_21):`$
$`\begin{array}{c}F_j=D_j(1jN_11),F_{N_1}=\epsilon (g_m;\xi _n)\epsilon ,\hfill \\ F_{N_1+j1}=S_j2\pi \psi _j\text{or}F_{N_1+j1}=O_j\mathrm{\Omega }_j^{}(2jN_2),\hfill \\ v_k=g_{k+1}(1kN_11),v_{N_1+k1}=\xi _k(1kN_2)\hfill \\ g(x^2)\underset{j=1}{\overset{N_1}{}}g_jT_{j1}(2x^21)\frac{1}{2}g_1,g(1)0g_1=2\underset{j=2}{\overset{N_1}{}}g_j\hfill \\ \xi (x^2)\underset{j=1}{\overset{N_2}{}}\xi _jT_{j1}(2x^21)\frac{1}{2}\xi _1\hfill \\ \psi _p(x^2)\underset{j=1}{\overset{N_2}{}}\psi _jT_{j1}(2x^21)\frac{1}{2}\psi _1,\hfill \\ \psi _p(0)1\psi _1=2\underset{j=2}{\overset{N_2}{}}(1)^j\psi _j+2\hfill \\ \mathrm{\Omega }^{}[(\rho /\rho _0)^2]=\mathrm{\Omega }(\rho )/\mathrm{\Omega }(0):\hfill \\ \mathrm{\Omega }^{}(x^2)\underset{j=1}{\overset{N_2}{}}\mathrm{\Omega }_j^{}T_{j1}(2x^21)\frac{1}{2}\mathrm{\Omega }_1^{},\hfill \\ \mathrm{\Omega }^{}(0)1\mathrm{\Omega }_1^{}=2\underset{j=2}{\overset{N_2}{}}(1)^j\mathrm{\Omega }_j^{}+2\hfill \\ S(x^2=\rho ^2/\rho _0^2;g;\xi )\underset{j=1}{\overset{N_2}{}}S_j(g_m;\xi _n)T_{j1}(2x^21)\frac{1}{2}S_1(g_m;\xi _n)\hfill \\ O(x^2=\rho ^2/\rho _0^2;g;\xi )\underset{j=1}{\overset{N_2}{}}O_j(g_m;\xi _n)T_{j1}(2x^21)\frac{1}{2}O_1(g_m;\xi _n)\hfill \\ D(x^2=\rho ^2/\rho _0^2;g;\xi )\underset{j=1}{\overset{N_11}{}}D_j(g_m;\xi _n)T_{j1}(2x^21)\frac{1}{2}D_1(g_m;\xi _n)\hfill \end{array}`$ (42)
\[The function $`\epsilon (g_m;\xi _n)=M^2/J`$ is determined using (28) for the above functions $`g`$ and $`\xi `$.\]
4. For the above system, the boundary values are assumed to be given in the form of the $`\psi _k`$’s or $`\mathrm{\Omega }_k^{}`$’s $`(k=2,\mathrm{},N_2)`$. Moreover, some $`\epsilon 1`$ has to be prescribed. Then, good initial $`v_k`$’s come from the Newtonian expansion. The Newton-Raphson method improves the $`v_k`$’s and yields a very accurate solution to (LABEL:AS1) for the chosen small $`\epsilon `$. Now, this solution serves as the initial estimate for the $`v_k`$’s belonging to a marginally increased value for $`\epsilon `$. Again, the Newton-Raphson method improves the solution, and one continues in this manner until this procedure ceases to converge. This occurs for some finite value $`\epsilon _0`$, at the latest for $`\epsilon =1`$. A further discussion of this limit is given below.
5. A rather technical detail is the retranslation of the $`\xi _j`$ into a set $`\{Y_\nu \}_q`$ which then gives a satisfactory approximation of $`\xi `$ in terms of (25). There are many ways to do this. Here, the following one has been chosen.
One rewrites equation (25) in the equivalent form
$$\mathrm{exp}\left[x\xi (x^2;\{Y_\nu \}_q)\right]=\underset{\nu =1}{\overset{q}{}}\frac{\text{i}Y_\nu x}{\text{i}Y_\nu +x}=\frac{P_q(x)}{P_q(x)}\text{with}P_q(x)=\underset{\nu =1}{\overset{q}{}}b_\nu x^\nu .$$
The coefficients $`b_\nu `$ of the polynomial $`P_q`$ can be determined by evaluating the left hand side at $`q`$ arbitrary different points $`x_\mu [0,1]`$ <sup>9</sup><sup>9</sup>9Here, zeros of Chebyshev-polynomials have been used. and solving the following linear system:
$$\mathrm{exp}\left[x_\mu \xi (x_\mu ^2;\{Y_\nu \}_q)\right]\underset{\nu =1}{\overset{q}{}}b_\nu x_\mu ^\nu =\underset{\nu =1}{\overset{q}{}}b_\nu (x_\mu )^\nu $$
The zeros of $`P_q`$ determine the $`Y_\nu `$.
The above scheme has been performed for many different prescribed surface mass densities and angular velocities. This provides strong evidence for the conjecture that, in this manner, all Newtonian disks can be extended into the relativistic regime. It has been found that the value for $`\epsilon _0`$, the limiting parameter for the convergence of this scheme, depends on the chosen profile for $`\psi _p`$ (or equivalently for $`\mathrm{\Omega }^{}`$). It is illustrated in appendix C.2, how the Ernst potential always tends to the extreme Kerr solution as $`\epsilon 1`$. This supports a conjecture by Bardeen and Wagoner . But $`\epsilon _0=1`$ does not hold for all given surface mass densities. Even in the Newtonian regime there are surface mass densities for which a realistic physical disk cannot be found since the corresponding angular velocity would become imaginary. If one chooses a profile for $`\sigma _p`$ not very different from these, then the Newtonian limit still might exist, but some $`\epsilon _0<1`$ turns up, beyond which the method does not converge. In the case of prescribed angular velocity, the situation is similar. Here, for any sequence $`f=f(g_\epsilon ;\xi _\epsilon )`$ the angular velocity $`\mathrm{\Omega }^{}`$ tends for all $`x^2[0,1]`$ to 1 as $`\epsilon 1.`$ So, each nonuniform rotation law will lead to some $`\epsilon _0<1`$ (see section 4 for examples).
The above expansions in terms of Chebyshev-polynomials allow a very accurate representation with only a small number of coefficients. However, the retranslation of $`\xi `$ (see the above point 5) leads to functions that are not especially well suited for an approximation. In particular, if the boundary condition $`\psi _p`$ is chosen to be close to those for which there is no Newtonian disk, then the accuracy cannot be driven particularly high by the computer program used, although the method in priniciple allows arbitrary approximation (see section 4.2).
For $`\psi _p`$’s sufficiently far away from those critical ones, the accuracy obtained was very high. By choosing appropriate values for $`N_1`$ and $`N_2`$ one can always achieve extremly good agreement with the dust condition (4) (12 digits and beyond) which ensures a realistic physical interpretation of the solution. The accuracy to which the second boundary condition, (3) or (6), can be satisfied, depends on the parameter $`\epsilon `$. It is usually around 8 digits in the weak relativistic regime, and falls as $`\epsilon `$ increases, but is still around 4 digits as $`\epsilon `$ tends to $`\epsilon _0`$. These values arose for $`N_1=30,N_2=12,`$ and typical $`\psi _p`$’s (like $`\psi _p`$’s depending linearly on $`x^2`$) and $`\mathrm{\Omega }^{}`$’s (e. g. the realistic one considered in section 4.1). The number $`q`$ of the parameters $`Y_\nu `$ by which $`\xi `$ is represented, was chosen to be between 20 and 30 (independently of $`N_2`$).
What remains to be discussed is the regularity of the Ernst potentials that were obtained. For a few of the solutions, the functions $`e^{2U}`$ and $`b`$ were plotted over the coordinates $`\rho `$ and $`\zeta `$. Moreover, the agreement of the alternative representations of $`M`$ and $`J`$, as given by the behaviour of the Ernst potential at infinity
$$U=\frac{M}{r}+𝒪(r^2),b=2J\frac{\mathrm{cos}\theta }{r^2}+𝒪(r^3),(r=\sqrt{\rho ^2+\zeta ^2},\zeta =r\mathrm{cos}\theta )$$
with the results from formulas (28) yields good confirmation of the regularity. This agreement was checked for all solutions that were calculated.
## 4 Representative examples
From the numerous solutions obtained, three particular sets of differentially rotating disks are discussed in more detail. The first one is an example of disks revolving with a realistic rotation law. The second set illustrates the break down of the numerical method for a specially prescribed surface mass density $`\sigma _p`$ at some $`\epsilon _0<1`$. On the other hand it is demonstrated that, for the same $`\sigma _p`$, regular solutions can be found in the highly relativistic regime. Finally, the third example concerns the occurence of a second ergoregion for a particular series of disks and, moreover, the gradual merging of the two spatially separated ergoregions as $`\epsilon `$ increases.
The deviations between the boundary values obtained for particular numerical solutions and the given boundary conditions are listed in tables. The quantities $`\mathrm{\Delta }_D,\mathrm{\Delta }_\mathrm{\Omega }`$, and $`\mathrm{\Delta }_\sigma `$ therein are defined by
$$\mathrm{\Delta }_D=\underset{x^2[0,1]}{\mathrm{max}}|D_{\text{obt}}(x^2;g;\xi )|$$
$$\mathrm{\Delta }_\mathrm{\Omega }=\underset{x^2[0,1]}{\mathrm{max}}\left|\mathrm{\Omega }_{\text{obt}}^{}(x^2)\mathrm{\Omega }_{\text{giv}}^{}(x^2)\right|$$
$$\mathrm{\Delta }_\sigma =\underset{x^2[0,1]}{\mathrm{max}}\left|\psi _p^{\text{obt}}(x^2)\psi _p^{\text{giv}}(x^2)\right|,$$
where the indices ’obt’ and ’giv’ refer to obtained and given quantities, respectively. Moreover, by letters (a)$`,\mathrm{},`$(e), special examples are marked, for which illustrative graphs have been made. Here, curves drawn in the same line style belong to the same solution. The graphs show the dimensionless quantities $`\rho _0\sigma _p`$ and $`\rho _0\mathrm{\Omega }`$ as well as $`g`$ and $`\xi `$ plotted against the normalized radial coordinate $`\rho /\rho _0`$ and $`x`$, respectively.
### 4.1 Disks possessing a realistic rotation law
As motivated by observations in astrophysics the rotation law of a galaxy is often modelled by an equation of the form (see )
$$\mathrm{\Omega }(\rho )=\frac{\mathrm{\Omega }(0)}{\sqrt{1+\rho ^2/\rho _1^2}}.$$
(43)
Here, the parameter $`\rho _1`$ varies for different galaxies. In the following series of solutions illustrated in figure 1, $`\rho _1=0.7\rho _0`$ has been chosen. As described in section 3, there is a limiting parameter $`\epsilon _00.935`$, for which the numerical method ceases to converge.
Figure 1: Disks possessing the rotation law (43) with $`\rho _1=0.7\rho _0`$ $`(N_1=30,N_2=12)`$.
### 4.2 Disks with a critical surface mass density
For the following sequence of solutions, a surface mass density of the form
$$\sigma _p(\rho )=\sigma _0\left(13\frac{\rho ^2}{\rho _0^2}+\beta \frac{\rho ^4}{\rho _0^4}\right)\sqrt{1\frac{\rho ^2}{\rho _0^2}}$$
(44)
has been assumed.
Figure 2: Disks possessing the surface mass density (44) with $`\beta =6`$ $`(N_1=30,N_2=12)`$.
It turns out that for $`\beta >\beta _N7`$ no Newtonian disks with a real angular velocity can be found. On the other hand, for $`\beta =5.5`$, all relativistic solutions for $`0\epsilon 1`$ exist. The table and graphs of figure 2 refer to the case $`\beta =6`$. Starting here from the Newtonian solution, one soon recognizes a first limiting parameter $`\epsilon _00.60`$ for which the method breaks down. However, by coming from solutions with $`\beta =5.5`$ and $`\epsilon `$ close to 1, it is possible to create highly relativistic solutions with $`\beta =6`$. In fact, there is another limiting parameter, $`\epsilon _10.97`$, above which the solutions with $`\beta =6`$ exist once again. Due to the nearness to the critical surface mass density (for $`\beta =\beta _N`$), the accuracy obtained for the boundary condition (3) is not very high.
### 4.3 Disks possessing spatially separated ergoregions
The particular set of disks depicted in figure 3 demonstrates the occurence of a second ergoregion<sup>10</sup><sup>10</sup>10An ergoregion is a portion of the $`(\rho ,\zeta )`$-space within which the function $`e^{2U}`$ is negative.. These solutions do not satisfy a specially prescribed boundary condition (3) or (6), but have been constructed in the following manner as intermediate solutions.
| | (a) | (b) | (c) | (d) | (e) |
| --- | --- | --- | --- | --- | --- |
| $`\epsilon =M^2/J`$ | 0.84038 | 0.84054 | 0.84079 | 0.84120 | 0.84162 |
| $`\mathrm{\Delta }_D10^{12}`$ | 3 | 2 | 2 | 5 | 2 |
Figure 3: Example for a series of disks possessing spatially separated ergoregions. In the uppermost picture, the rims of the ergoregions in the $`(\rho /\rho _0,\zeta /\rho _0)`$-space are to be seen $`(N_1=40,N_2=9)`$.
If one investigates solutions with surface mass densities similar to those of (44), one recognizes two minima for $`e^{2U}`$ (taken as a function of $`\rho ,\mathrm{\hspace{0.17em}0}\rho \rho _0,\zeta =0`$), say at $`\rho _a`$ and $`\rho _b>\rho _a`$. Now, for a particular choice of $`\sigma _p`$ it is possible to get $`e^{2U}(\rho _a)>0`$ and $`e^{2U}(\rho _b)<0`$, whilst by another choice one can achieve $`e^{2U}(\rho _a)<0`$ and $`e^{2U}(\rho _b)>0`$. This makes clear, that disks with spatially separated ergoregions can be constructed by interpolating between these solutions.
For the chosen example, there is only a narrow interval $`(\epsilon _a,\epsilon _b)`$ for which the two separated ergoregions occur. As can be seen from figure 3, after creation of the second ergoregion at $`\epsilon _a0.8403`$, both ergoregions grow as $`\epsilon `$ increases. Eventually, at $`\epsilon _b0.8415`$, the ergoregions merge into one ergoregion.
## 5 Discussion of mathematical aspects
As already mentioned in section 2, the assumption that the function $`\mathrm{\Phi }_g`$ introduced in (26) can be extended to form a continuous mapping defined on $`𝒜`$, lies at the heart of the above numerical methods. Although this assumption seems to be intuitive, it is not trivial. Consider the following example:
For any analytic function $`\psi :[0,1]`$ one finds the equality<sup>11</sup><sup>11</sup>11To verify this formula one simply expands the logarithms in the form $`\mathrm{ln}(1+ϵ)=ϵ+𝒪(ϵ^2)`$ and notes that the resulting sum tends to the Riemann integral of the right hand side.:
$$\underset{q\mathrm{}}{lim}\underset{\nu =1}{\overset{q}{}}\mathrm{ln}\left[1+\frac{1}{q}\psi \left(\frac{\nu }{q}\right)\right]=\underset{0}{\overset{1}{}}\psi (t)𝑑t.$$
From this it follows that
$$2_0^1\varphi (t)𝑑t=\frac{1}{x}\underset{q\mathrm{}}{lim}\underset{\nu =1}{\overset{q}{}}\mathrm{ln}\frac{q+x\varphi (\nu /q)}{qx\varphi (\nu /q)}\text{with }\psi (t)=\pm x\varphi (t)\text{.}$$
Hence, the function $`\xi (x^2)2`$ can be represented by any sequence of the form
$$Y_\nu ^{(q)}=\text{i}\frac{q}{\varphi (\nu /q)}\text{with}_0^1\varphi (t)𝑑t=1.$$
Since these sequences might be quite different from each other, it is rather surprising that all of them approximate the same Ernst potential given by (19). But this follows from the above assumption.
This already indicates the difficulties which are connected with a rigorous proof of this assumption because the Ernst potential is only given in terms of the set $`\{Y_\nu \}_q`$ and not directly in terms of $`\xi `$.
A further conjecture is strongly confirmed by extensive numerical investigations:
> For the hyperelliptic class of solutions represented by (45) in appendix A, the functions $`\xi `$ and $`g`$ are given by
>
> $`\xi (x^2)`$ $`=`$ $`{\displaystyle \frac{1}{2x}}\mathrm{ln}\left[{\displaystyle \underset{\nu =1}{\overset{p}{}}}{\displaystyle \frac{\text{i}X_\nu x}{\text{i}X_\nu +x}}\right]`$
> $`g(x^2)`$ $`=`$ $`\text{sign}\left({\displaystyle \underset{\nu =1}{\overset{p}{}}}X_\nu \right)A_g(x^2)h(x^2),`$
> $`A_g(x^2)=\sqrt{{\displaystyle \underset{\nu =1}{\overset{p}{}}}(\text{i}xX_\nu )(\text{i}x\overline{X}_\nu )},A_g(x^2)>0.`$
> In particular, in this formulation, the solution for the Neugebauer-Meinel-disk assumes the form $`f=f(\xi ;g)`$ where
>
> $`\xi (x^2)`$ $`=`$ $`{\displaystyle \frac{1}{2x}}\mathrm{ln}{\displaystyle \frac{x^2C_1(\mu )x+C_2(\mu )}{x^2+C_1(\mu )x+C_2(\mu )}},`$
> $`C_1(\mu )=\sqrt{2[1+C_2(\mu )]},C_2(\mu )={\displaystyle \frac{1}{\mu }}\sqrt{1+\mu ^2},`$
> $`g(x^2)`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}\text{arsinh}[\mu (1x^2)],`$
> and the parameter $`\mu ,\mathrm{\hspace{0.17em}0}<\mu <\mu _0=4.62966184..,`$ is related to the angular velocity by
>
> $$\mu =2\mathrm{\Omega }^2\rho _0^2e^{2V_0},V_0=U(\rho =0,\zeta =0).$$
As already mentioned, a direct proof of the above assumptions promises to be very complicated. But there might be an alternative proof which relies on relating a general solution of the Ernst equation to the solution of a so-called Riemann-Hilbert problem, see . In this treatment, an appropriately introduced matrix function, from which the Ernst potential can be extracted, is supposed to be regular on a two-sheeted Riemann surface of genus zero except for some given curve, where it possesses a well-defined jump behaviour. The freedom of two jump functions defined on this curve corresponds to the freedom to choose $`\xi `$ and $`g`$. Now, if one succeeds in finding a particular formulation of a Riemann-Hilbert problem in which $`\xi `$ and $`g`$ are involved, then the final solution for $`f`$ proves to depend only on $`\xi `$ (and $`g`$) and not on a particular global representation in terms of $`\{Y_\nu \}_q`$. This deserves further investigation.
There is very strong numerical evidence for the validity of both assumptions. For various functions $`\xi `$ (and functions $`g`$), different representations $`\{Y_\nu \}_q`$ have been seen to approximate the same Ernst potential. In particular, the approximation of the Neugebauer-Meinel-solution in terms of Bäcklund solutions was carried out to give an agreement up to the 12th digit with the hyperelliptic solution, which confirms both assumptions.
## Appendix A The transition from the hyperelliptic solutions to the Bäcklund type solutions
In this section the Bäcklund type solutions are derived from the hyperelliptic class. The latter is assumed to be given in the form represented in <sup>12</sup><sup>12</sup>12The parameters $`K_\nu `$, the upper integration limits $`K^{(\nu )}`$, and the integration variable $`K`$ have to be replaced by their ’normalized’ values $`X_\nu =K_\nu /\rho _0,X^{(\nu )}=K^{(\nu )}/\rho _0`$, and $`X=K/\rho _0`$, respectively. for an even integer $`p2`$:
$$f=\mathrm{exp}\left(\underset{\nu =1}{\overset{p}{}}\underset{X_\nu }{\overset{X^{(\nu )}}{}}\frac{X^pdX}{V(X)}u_p\right)$$
(45)
$$V(X)=\sqrt{(X+\text{i}z)(X\text{i}\overline{z})\underset{\nu =1}{\overset{p}{}}(XX_\nu )(X\overline{X}_\nu )},z=\frac{1}{\rho _0}(\rho +\text{i}\zeta )$$
$$\underset{\nu =1}{\overset{p}{}}\underset{X_\nu }{\overset{X^{(\nu )}}{}}\frac{X^jdX}{V(X)}=u_j,0j<p$$
(46)
$$u_j=\underset{1}{\overset{1}{}}\frac{(\text{i}x)^jh(x^2)dx}{Z_D},0jp,h:[0,1],\text{analytic},$$
$`Z_D`$ as defined in (19)
The set $`\{\text{i}X_\nu \}_p`$ consists of arbitrary real parameters and/or pairs of complex conjugate parameters (in order to guarantee reflectional symmetry). The ($`z`$-dependent) values for the $`X^{(\nu )}`$ as well as the integration paths on a two-sheeted Riemann surface result from the Jacobian inversion problem (46).
The transition to the Bäcklund type solutions (19) can be obtained in the limit $`\epsilon 0`$ by the following assumptions:
* $`p=2q`$
* $`X_{2\nu 1}=Y_\nu +\epsilon \beta _\nu ,X_{2\nu }=Y_\nu (1\nu q),\{\beta _\nu \}_q\text{ arbitrary}`$
* $`g(x^2)=(1)^qh(x^2)A(\text{i}x),A(X)=\underset{\nu =1}{\overset{q}{}}(XY_\nu )(X\overline{Y}_\nu )`$ .
To this end, the above expression for $`f`$ is rewritten in the equivalent form:
$$f=\mathrm{exp}\left[\underset{\nu =1}{\overset{q}{}}\left(_{X_{2\nu 1}}^{X^{(2\nu 1)}}\frac{A(X)dX}{V(X)}+_{X_{2\nu }}^{X^{(2\nu )}}\frac{A(X)dX}{V(X)}\right)\underset{1}{\overset{1}{}}\frac{(1)^qg(x^2)dx}{Z_D}\right]$$
The Jacobian inversion problem (46) reads as follows in a similarly rewritten form $`(1\mu q)`$:
$$\underset{\nu =1}{\overset{q}{}}(_{X_{2\nu 1}}^{X^{(2\nu 1)}}\frac{A(X)dX}{V(X)(XY_\mu )}+_{X_{2\nu }}^{X^{(2\nu )}}\frac{A(X)dX}{V(X)(XY_\mu )})=\underset{1}{\overset{1}{}}\frac{(1)^qg(x^2)dx}{(\text{i}xY_\mu )Z_D}$$
$$\underset{\nu =1}{\overset{q}{}}(_{X_{2\nu 1}}^{X^{(2\nu 1)}}\frac{A(X)dX}{V(X)(X\overline{Y}_\mu )}+_{X_{2\nu }}^{X^{(2\nu )}}\frac{A(X)dX}{V(X)(X\overline{Y}_\mu )})=\underset{1}{\overset{1}{}}\frac{(1)^qg(x^2)dx}{(\text{i}x\overline{Y}_\mu )Z_D}$$
Furthermore
$`{\displaystyle _{X_{2\nu 1}}^{X^{(2\nu 1)}}}{\displaystyle \frac{A(X)dX}{V(X)(XY)}}+{\displaystyle _{X_{2\nu }}^{X^{(2\nu )}}}{\displaystyle \frac{A(X)dX}{V(X)(XY)}}=`$
$`{\displaystyle _{X_{2\nu }}^{X_{2\nu 1}}}{\displaystyle \frac{A(X)dX}{V(X)(XY)}}+{\displaystyle _{X^{(2\nu )}}^{X^{(2\nu 1)}}}{\displaystyle \frac{A(X)dX}{V(X)(XY)}}`$
with $`X^{(2\nu )}`$ now lying in the other sheet of the Riemann surface.
In the limit $`\epsilon 0`$, one obtains
$$\underset{\epsilon 0}{lim}_{X_{2\nu }}^{X_{2\nu 1}}\frac{A(X)dX}{V(X)(XY)}=\{\begin{array}{c}\pm \pi \text{i}\delta _{\mu \nu }/[\lambda _\mu (Y_\mu +\text{i}z)]\text{for}Y=Y_\mu \hfill \\ 0\text{for}Y=\overline{Y}_\mu \hfill \end{array}$$
with $`\delta _{\mu \nu }`$ being the usual Kronecker symbol and $`\lambda _\mu `$ as defined in (19).
The second term amounts to
$`\underset{\epsilon 0}{lim}{\displaystyle _{X^{(2\nu )}}^{X^{(2\nu 1)}}}{\displaystyle \frac{A(X)dX}{V(X)(XY)}}={\displaystyle _{X^{(2\nu )}}^{X^{(2\nu 1)}}}{\displaystyle \frac{dX}{(XY)\sqrt{(X+\text{i}z)(X\text{i}\overline{z})}}}=`$
$`{\displaystyle \frac{1}{\lambda (Y)(Y+\text{i}z)}}\mathrm{ln}\left({\displaystyle \frac{[\lambda (X^{(2\nu 1)})\lambda (Y)][\lambda (X^{(2\nu )})+\lambda (Y)]}{[\lambda (X^{(2\nu 1)})+\lambda (Y)][\lambda (X^{(2\nu )})\lambda (Y)]}}\right),`$
where for evaluation of the second integral the substitution
$$\lambda =\lambda (X)=\sqrt{\frac{X\text{i}\overline{z}}{X+\text{i}z}}$$
has been used.
Hence, the Jacobian inversion problem reads as follows in the limit $`\epsilon 0`$:
$`{\displaystyle \underset{\nu =1}{\overset{q}{}}}{\displaystyle \frac{[\lambda (X^{(2\nu 1)})\lambda _\mu ][\lambda (X^{(2\nu )})+\lambda _\mu ]}{[\lambda (X^{(2\nu 1)})+\lambda _\mu ][\lambda (X^{(2\nu )})\lambda _\mu ]}}`$ $`=`$ $`\gamma _\mu `$ (47)
$`{\displaystyle \underset{\nu =1}{\overset{q}{}}}{\displaystyle \frac{[\lambda (X^{(2\nu 1)})\lambda _\mu ^{}][\lambda (X^{(2\nu )})+\lambda _\mu ^{}]}{[\lambda (X^{(2\nu 1)})+\lambda _\mu ^{}][\lambda (X^{(2\nu )})\lambda _\mu ^{}]}}`$ $`=`$ $`\overline{\gamma }_\mu `$ (48)
and in an analogous manner
$`f=f_0{\displaystyle \underset{\nu =1}{\overset{q}{}}}{\displaystyle \frac{[\lambda (X^{(2\nu 1)})+1][\lambda (X^{(2\nu )})1]}{[\lambda (X^{(2\nu 1)})1][\lambda (X^{(2\nu )})+1]}}`$ (49)
\[with $`\gamma _\mu ,\lambda _\mu ^{}`$ and $`f_0`$ as defined in (19)\].
Instead of evaluating the quantities $`\lambda (X^{(\nu )}),(1\nu 2q),`$ the coefficients $`b_\nu `$ and $`c_\nu (1\nu q)`$ of the polynomial
$`P(\lambda )`$ $`=`$ $`{\displaystyle \underset{\nu =1}{\overset{q}{}}}[\lambda \lambda (X^{(2\nu 1)})][\lambda +\lambda (X^{(2\nu )})]`$ $`=`$ $`\lambda ^{2q}+\lambda {\displaystyle \underset{\nu =1}{\overset{q}{}}}b_\nu \lambda ^{2\nu 2}+{\displaystyle \underset{\nu =1}{\overset{q}{}}}c_\nu \lambda ^{2\nu 2}`$
are determined. Since
$`{\displaystyle \frac{P(\lambda _\mu )}{P(\lambda _\mu )}}=\gamma _\mu ,{\displaystyle \frac{P(\lambda _\mu ^{})}{P(\lambda _\mu ^{})}}=\overline{\gamma }_\mu ,f=f_0{\displaystyle \frac{P(1)}{P(1)}},`$ (51)
the following system of linear equations for the quantities $`b_\nu ,c_\nu ,P(1),`$ and $`P(1)`$ emerges:
$$\underset{\nu =1}{\overset{q}{}}[b_\nu \alpha _\mu \lambda _\mu ^{2\nu 1}+c_\nu \lambda _\mu ^{2\nu 2}]=\lambda _\mu ^{2q},$$
$$\underset{\nu =1}{\overset{q}{}}[b_\nu \alpha _\mu ^{}(\lambda _\mu ^{})^{2\nu 1}+c_\nu (\lambda _\mu ^{})^{2\nu 2}]=(\lambda _\mu ^{})^{2q}$$
$$\underset{\nu =1}{\overset{q}{}}(b_\nu c_\nu )+P(1)=1$$
$$\underset{\nu =1}{\overset{q}{}}(b_\nu +c_\nu )P(1)=1,$$
with $`\alpha _\mu `$ and $`\alpha _\mu ^{}`$ as defined in (19).
Finally, if the solution of this linear system for $`P(\pm 1)`$ is expresseed by means of Cramer’s rule, the desired form (19) of the Bäcklund type is obtained.
## Appendix B Invariance properties of the Ernst potential
For the proof of the properties (2-2), the Ernst potential (19) is reformulated by
$$f(\{Y_\nu \}_q;g)=f_0\frac{D(1;\{Y_\nu \}_q;g)}{D(1;\{Y_\nu \}_q;g)}$$
(53)
with
$$\begin{array}{c}D(\lambda ;\{Y_\nu \}_q;g)\hfill \\ =\left|\begin{array}{cccccccc}a_1& (a_1x_1)& \mathrm{}& (a_1x_1^{q1})& 1& x_1& \mathrm{}& x_1^q\\ a_2& (a_2x_2)& \mathrm{}& (a_2x_2^{q1})& 1& x_2& \mathrm{}& x_2^q\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ a_{2q+1}& (a_{2q+1}x_{2q+1})& \mathrm{}& (a_{2q+1}x_{2q+1}^{q1})& 1& x_{2q+1}& \mathrm{}& x_{2q+1}^q\end{array}\right|\hfill \\ a_1=\lambda ,a_{2\nu }=\alpha _\nu \lambda _\nu ,a_{2\nu +1}=\alpha _\nu ^{}\lambda _\nu ^{},\hfill \\ x_1=\lambda ^2,x_{2\nu }=\lambda _\nu ^2,x_{2\nu +1}=(\lambda _\nu ^{})^2.\hfill \end{array}$$
The above expression for $`D(\lambda ;\{Y_\nu \}_q;g)`$ is a Vandermonde-like determinant. These determinants have been studied in detail by Steudel, Meinel and Neugebauer . By their reduction formula \[(8) of \], $`D`$ assumes the form:
$`D(\lambda ;\{Y_\nu \}_q;g)`$
$`=𝒱_{q,q+1}(a_r;b_r|x_r)\text{[with }b_r=1\text{ for }r=1\mathrm{}(2q+1)\text{]}`$
$`={\displaystyle \underset{P}{}}\epsilon _P\left({\displaystyle \underset{j=1}{\overset{q}{}}}a_{r(j)}\right)𝒱_q[x_{r(1)},\mathrm{},x_{r(q)}]𝒱_{q+1}[x_{r(q+1)},\mathrm{},x_{r(2q+1)}]`$
where
* the sum runs over all permutations $`P=[r(1),\mathrm{},r(2q+1)]`$ of $`(1,2,\mathrm{},2q+1)`$ with $`r(k)<r(j)`$ for $`k<j<q`$ as well as for $`qk<j`$
* $`\epsilon _P=\{\begin{array}{c}+1\text{for }P\text{ even}\hfill \\ 1\text{for }P\text{ odd}\hfill \end{array}`$
* the Vandermonde determinants are given by
$`𝒱_N[x_1,\mathrm{},x_N]=\underset{k>j}{}(x_kx_j).`$
In this formulation the following properties can be proved:
1. If $`x_{2q+1}=x_{2q}`$ then
$`D(\lambda ;\{Y_\nu \}_q;g)`$
$`=(1)^q(a_{2q}a_{2q+1})\left[{\displaystyle \underset{j=1}{\overset{2q1}{}}}(x_{2q}x_j)\right]D(\lambda ;\{Y_\nu \}_{q1};g)`$
2. If $`x_{2q}=1+\kappa ϵ+𝒪(ϵ^2),x_{2q+1}=1\kappa ϵ+𝒪(ϵ^2),`$ and $`(a_{2q}a_{2q+1})=1+𝒪(ϵ),`$ then
$`D(1;\{Y_\nu \}_q;g)`$
$`=\kappa ϵ\left[{\displaystyle \underset{j=2}{\overset{2q1}{}}}(1x_j)\right](a_{2q}+a_{2q+1}\pm 2)D(1;\{Y_\nu \}_{q1};g)+𝒪(ϵ^2).`$
With (A) the equalities (2) and (2) can be derived whilst (B) serves to confirm (2) and (2). In order to prove (A) consider the following groups of permutations separately:
$$\begin{array}{c}P_1:r(q1)=2q,r(q)=2q+1\hfill \\ P_2:r(2q)=2q,r(2q+1)=2q+1\hfill \\ P_3:r(q)=2q,r(2q+1)=2q+1\hfill \\ P_4:r(q)=2q+1,r(2q+1)=2q\hfill \end{array}$$
For $`x_{2q+1}=x_{2q}`$, all terms belonging to $`P_1`$ and $`P_2`$ vanish while all terms belonging to $`P_3`$ and $`P_4`$ possess a common factor, $`[a_{2q}_{j=1}^{2q1}(x_{2q}x_j)]`$ and $`[a_{2q+1}_{j=1}^{2q1}(x_{2q}x_j)]`$, respectively. After reordering (from which the factor $`(1)^q`$ results), (A) is easily obtained.
The proof for (B) works similarly. Now, eight groups of permutations have to be considered separately:
$$\begin{array}{c}P_{1a}:r(1)=1,r(q1)=2q,r(q)=2q+1\hfill \\ P_{1b}:r(q+1)=1,r(2q)=2q,r(2q+1)=2q+1\hfill \\ P_{2a}:r(q)=2q,r(q+1)=1,r(2q+1)=2q+1\hfill \\ P_{2b}:r(1)=1,r(q)=2q,r(2q+1)=2q+1\hfill \\ P_{3a}:r(q)=2q+1,r(q+1)=1,r(2q+1)=2q\hfill \\ P_{3b}:r(1)=1,r(q)=2q+1,r(2q+1)=2q\hfill \\ P_{4a}:r(q1)=2q,r(q)=2q+1,r(q+1)=1\hfill \\ P_{4b}:r(1)=1,r(2q)=2q,r(2q+1)=2q+1\hfill \end{array}$$
All terms of permutations with a coinciding first index can be combined to give<sup>13</sup><sup>13</sup>13Here the requirements $`a_1=1,x_1=1`$ are necessary. Additionally, for $`P_{4a}`$ and $`P_{4b}`$, the constraint $`a_{2q}a_{2q+1}=1+𝒪(ϵ)`$ is needed. :
$$\begin{array}{c}\{P_{1a},P_{1b}\}𝒪(ϵ^3)\hfill \\ \{P_{2a},P_{2b}\}a_{2q}F+𝒪(ϵ^2)\hfill \\ \{P_{3a},P_{3b}\}a_{2q+1}F+𝒪(ϵ^2)\hfill \\ \{P_{4a},P_{4b}\}\pm 2F+𝒪(ϵ^2)\hfill \end{array}$$
with
$$F=(1)^{q+1}\kappa ϵ\underset{j=2}{\overset{2q1}{}}(1x_j)D(\pm 1;\{Y_\nu \}_{q1};g).$$
## Appendix C Newtonian and ultrarelativistic limits
### C.1 The Newtonian limit
In the limit of small functions $`g`$ and $`\xi `$, i. e.
$$g(x^2)=\epsilon _gg_0(x^2)+𝒪(\epsilon _g^2),\xi (x^2)=\epsilon _\xi \xi _0(x^2)+𝒪(\epsilon _\xi ^2),$$
the Ernst potential $`f=f(\xi ;g)`$ as introduced in section 2 is given by
$`f(\xi ;g)=`$ $`1\epsilon _g{\displaystyle \underset{1}{\overset{1}{}}}{\displaystyle \frac{g_0(x^2)dx}{Z_D}}\text{i}\epsilon _g\epsilon _\xi {\displaystyle \underset{1}{\overset{1}{}}}{\displaystyle \frac{(\text{i}x)g_0(x^2)\xi _0(x^2)dx}{Z_D}}+𝒪(\epsilon _g^2)+𝒪(\epsilon _g\epsilon _\xi ^2).`$
In this section, the above property will be proved and the functions $`g_0`$ and $`\xi _0`$ will be derived as they result from the Newtonian expansion of the boundary conditions.
#### C.1.1 The Ernst potential for small functions $`g`$ and $`\xi `$
Due to the assumption that the function $`\mathrm{\Phi }_g`$ introduced in (26) can be extended to form a continuous mapping defined on $`𝒜`$ (see sections 2 and 5), the representation of $`\xi `$ in terms of $`\{Y_\nu \}_q`$ can be chosen arbitrarily. Here, the following set $`\{Y_\nu \}_q`$ is used:
* $`q=4r`$
* $`\left\{\begin{array}{c}Y_{4\nu 3}=Z_\nu (1+\epsilon _\xi z_\nu ),Y_{4\nu 2}=\overline{Y}_{4\nu 3}\hfill \\ Y_{4\nu 1}=\overline{Z}_\nu (1\epsilon _\xi z_\nu ),Y_{4\nu }=\overline{Y}_{4\nu 1}\hfill \end{array}\right\},\begin{array}{c}\mathrm{}(Z_\nu )0,z_\nu \hfill \\ (\nu =1\mathrm{}r)\hfill \end{array}`$
Then, it follows from (25) that $`\xi (x^2)=\epsilon _\xi \xi _0(x^2)+𝒪(\epsilon _\xi ^2)`$ with
$$\xi _0(x^2)=4\text{i}\underset{\nu =1}{\overset{r}{}}\frac{z_\nu (Z_\nu \overline{Z}_\nu )(x^2Z_\nu \overline{Z}_\nu )}{(x^2+Z_\nu ^2)(x^2+\overline{Z}_\nu ^{\mathrm{\hspace{0.33em}2}})}.$$
To evaluate the Ernst potential in this limit, the formulation (47-LABEL:AG5) in appendix A is used and the following steps are performed:
1. At first, it turns out that in the limit $`\epsilon _\xi 0`$ the coefficients $`b_\nu `$ of the polonomial (LABEL:AG3) vanish. This can be seen by considering the solution to linear system (LABEL:AG5).
$`b_\nu ={\displaystyle \frac{D_\nu }{D}}:`$ (62)
$`\begin{array}{c}D=\left|\begin{array}{ccccccc}a_2& \mathrm{}& (a_2x_2^{q1})& 1& x_2& \mathrm{}& x_2^{q1}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ a_{2q+1}& \mathrm{}& (a_{2q+1}x_{2q+1}^{q1})& 1& x_{2q+1}& \mathrm{}& x_{2q+1}^{q1}\end{array}\right|\hfill \\ a_{2\eta }=\alpha _\eta \lambda _\eta ,a_{2\eta +1}=\alpha _\eta ^{}\lambda _\eta ^{},x_{2\eta }=\lambda _\eta ^2,x_{2\eta +1}=(\lambda _\eta ^{})^2\hfill \\ D_\nu \text{ is derived from }D\text{ by replacing the }\nu \text{-th column by the vector}\hfill \\ \{x_2^q,\mathrm{},x_{2q+1}^q\}.\hfill \end{array}`$
For $`1\nu q`$, $`D_\nu `$ can be expanded in terms of Vandermonde determinants
$$𝒱_{q+1}(x_{r(1)},\mathrm{},x_{r(q+1)}),r(\eta )\{2,\mathrm{},2q+1\},r(\eta )<r(\mu )\text{for}\eta <\mu .$$
In the limit $`\epsilon _\xi 0`$, any set $`\{x_{r(\eta )}\}_{q+1}`$ contains at most $`q`$ different values, and therefore all $`D_\nu `$ vanish. On the other hand, $`D`$ remains finite (here only Vandermonde determinants $`𝒱_q`$ are involved), and hence all $`b_\nu `$ tend to zero.
2. Thus, with any zero $`\stackrel{~}{\lambda }_\nu `$ of the Polynomial (LABEL:AG3), $`(\stackrel{~}{\lambda }_\nu )`$ also becomes a zero as $`\epsilon _\xi 0`$. This set of zeros is ordered in the following way:
$$\{\lambda (X^{(1)}),\lambda (X^{(2)}),\mathrm{},\lambda (X^{(2q1)}),\lambda (X^{(2q)})\}=\{\stackrel{~}{\lambda }_1,\stackrel{~}{\lambda }_1\mathrm{},\stackrel{~}{\lambda }_q,\stackrel{~}{\lambda }_q\},$$
Suppose there is a $`\lambda _\mu `$ different from all zeros :
$$\lambda _\mu \lambda (X^{(2\nu 1)})=\lambda (X^{(2\nu )})\text{and}\lambda _\mu \lambda (X^{(2\nu 1)})\text{for all }\nu =1\mathrm{}q\text{.}$$
Then, since $`\gamma _\mu 1`$ for small $`g`$, (47) cannot be satisfied.
3. This gives rise to the following ansatz $`(\nu =1\mathrm{}q)`$:
$$\lambda ^2(X^{(2\nu 1)})=\lambda _\nu ^2+\epsilon _\xi \kappa _{2\nu 1}+𝒪(\epsilon _\xi ^2),\lambda ^2(X^{(2\nu )})=\lambda _\nu ^2+\epsilon _\xi \kappa _{2\nu }+𝒪(\epsilon _\xi ^2),$$
by which the system (47/48) can easily be solved to get the set $`\{\kappa _\nu \}_{2q}`$.
4. Finally, if $`g(x^2)=\epsilon _gg_0(x^2)+𝒪(\epsilon _g^2)`$ is considered, then (LABEL:AGNL1) follows from (49) by inserting the values obtained for $`\{\lambda (X^{(\nu )})\}_{2q}`$.
#### C.1.2 The functions $`g_0`$ and $`\xi _0`$ as resulting from the boundary conditions
For any family of Ernst potentials $`f=f(g_\epsilon ;\xi _\epsilon )`$ describing a sequence of differentially rotating disks of dust with the parameter $`\epsilon =M^2/J`$ \[$`M`$ and $`J`$ as defined in (28)\], the following expansion is valid (see , pp. 83-89):
$$f=1+e_2(\rho ,\zeta )\epsilon ^2+\text{i}b_3(\rho ,\zeta )\epsilon ^3+𝒪(\epsilon ^4).$$
By comparison with (LABEL:AGNL1) one gets
$$\epsilon _g=\epsilon ^2,\epsilon _\xi =\epsilon ,$$
$$e_2(\rho ,\zeta )=\underset{1}{\overset{1}{}}\frac{g_0(x^2)dx}{Z_D},b_3(\rho ,\zeta )=\underset{1}{\overset{1}{}}\frac{(\text{i}x)g_0(x^2)\xi _0(x^2)dx}{Z_D}.$$
If the boundary conditions,
$`\sigma _p(\rho )`$ $`=`$ $`\sigma _0\psi _2[(\rho /\rho _0)^2]\sqrt{1(\rho /\rho _0)^2}\epsilon ^2+𝒪(\epsilon ^4)\text{(with }\psi _2(0)=1\text{)}\text{or}`$
$`\mathrm{\Omega }(\rho )`$ $`=`$ $`\mathrm{\Omega }_0\mathrm{\Omega }_1[(\rho /\rho _0)^2]\epsilon +𝒪(\epsilon ^3)\text{(with }\mathrm{\Omega }_1(0)=1\text{)},`$
are given, then it follows from equations (3-6) that
$`(e_2)_{,\zeta }`$ $`=`$ $`4\pi \sigma _0\psi _2\sqrt{1(\rho /\rho _0)^2}\text{or}(e_2)_{,\rho }=2\mathrm{\Omega }_0^2\mathrm{\Omega }_1^2\rho \text{and}`$
$`(b_3)_{,\rho }`$ $`=`$ $`2\rho \mathrm{\Omega }_0\mathrm{\Omega }_1(e_2)_{,\zeta }.`$
By expressing $`e_2`$ and $`b_3`$ in terms of $`g_0`$ and $`\xi _0`$ in these equations, one gets Abelian integral equations for $`\xi _0`$ and $`g_0`$. Their solutions read as follows:
$`g_0(x^2)`$ $`=`$ $`4\sigma _0(1x^2){\displaystyle _0^{\pi /2}}(\mathrm{sin}^2\varphi )\psi _2(\mathrm{cos}^2\varphi +x^2\mathrm{sin}^2\varphi )𝑑\varphi `$
$`g_0(x^2)\xi _0(x^2)`$ $`=`$ $`8\sigma _0\mathrm{\Omega }_0(1x^2){\displaystyle _0^{\pi /2}}(\mathrm{sin}^2\varphi )\stackrel{~}{\mathrm{\Omega }}_1(\mathrm{cos}^2\varphi +x^2\mathrm{sin}^2\varphi )𝑑\varphi `$
\[with $`\stackrel{~}{\mathrm{\Omega }}_1(x^2)=\mathrm{\Omega }_1(x^2)\psi _2(x^2)].`$
Note that only one of the functions $`\psi _2`$ and $`\mathrm{\Omega }_1`$ can be prescribed since both represent different boundary conditions of the same Newtonian potential $`e_2`$. Likewise, the constants $`\sigma _0`$ and $`\mathrm{\Omega }_0^2`$ depend on each other. Moreover, these constants in terms of $`\psi _2`$ and $`\mathrm{\Omega }_1`$ are prescribed by the equation $`\epsilon =M^2/J`$.
### C.2 The ultrarelativistic limit
It is difficult to relate the functions $`g`$ and $`\xi `$ of an Ernst potential $`f=f(g;\xi )`$ to its physical properties like $`M`$ and $`J`$. Nevertheless, if a sequence $`f(g_\epsilon ;\xi _\epsilon )`$ can be extended to arbitrary values $`\epsilon <1`$, then, in the limit $`\epsilon 1`$, the universal solution of an extreme Kerr black hole is reached. It is illustrated how this limit results from the form (19) of the Ernst potential.
If the limit $`\rho _00`$ is considered for finite values of $`r=\sqrt{\rho ^2+\zeta ^2}`$, then by using the formulation (53) one gets (with $`\zeta =r\mathrm{cos}\theta `$):
$$f=(1\frac{\rho _0}{r}\underset{1}{\overset{1}{}}(1)^qg(x^2)𝑑x+𝒪(\rho _0^2))\left[\frac{E_1r+\rho _0[E_3\mathrm{cos}\theta (1)^qE_2]}{E_1r+\rho _0[E_3\mathrm{cos}\theta +(1)^qE_2]}+𝒪(\rho _0^2)\right].$$
The $`E_j`$ do not depend on $`\rho `$ and $`\zeta `$ but on $`g`$ and $`\xi `$. In particular:
$$E_1=\left|\begin{array}{cccccccc}b_1& (b_1Z_1)& \mathrm{}& (b_1Z_1^{q1})& 1& Z_1& \mathrm{}& Z_1^{q1}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ b_{2q}& (b_{2q}Z_{2q})& \mathrm{}& (b_{2q}Z_{2q}^{q1})& 1& Z_{2q}& \mathrm{}& Z_{2q}^{q1}\end{array}\right|$$
$$E_2=\left|\begin{array}{cccccccc}b_1& (b_1Z_1)& \mathrm{}& (b_1Z_1^{q2})& 1& Z_1& \mathrm{}& Z_1^q\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ b_{2q}& (b_{2q}Z_{2q})& \mathrm{}& (b_{2q}Z_{2q}^{q2})& 1& Z_{2q}& \mathrm{}& Z_{2q}^q\end{array}\right|$$
$$b_{2\nu 1}=\mathrm{tanh}\left[\frac{1}{2}\underset{1}{\overset{1}{}}\frac{(1)^qg(x^2)dx}{\text{i}xY_\nu }\right],b_{2\nu }\overline{b}_{2\nu 1}=1$$
$$Z_{2\nu 1}=Y_\nu ,Z_{2\nu }=\overline{Y}_\nu .$$
Clearly, if $`E_10`$ then $`lim_{\rho _00}f=1`$. The Ernst potential passes to an ultrarelativistic limit if $`E_1`$ and $`\rho _0`$ tend simultaneously to zero such that<sup>14</sup><sup>14</sup>14It can be shown that $`E_1^2`$. Hence, the ultrarelativistic limit for the family $`f(g_\epsilon ;\xi _\epsilon )`$ is performed when some function $`E_a=E_a(g_\epsilon ;\xi _\epsilon )=E_b(\epsilon )E_1^2(\epsilon )`$, which is independent of the representation $`\{Y_\nu \}_q`$, vanishes.
$$\mathrm{\Omega }_U=\underset{\rho _00}{lim}\frac{(1)^qE_1}{\mathrm{\hspace{0.17em}2}E_2\rho _0}$$
exists. Then one gets
$$f=\frac{2\mathrm{\Omega }_Ur+E_4\mathrm{cos}\theta 1}{2\mathrm{\Omega }_Ur+E_4\mathrm{cos}\theta +1}.$$
The only Ernst potential of this form which is asymptotically flat and regular for $`r>0`$ is the extreme Kerr solution. The constant $`\mathrm{\Omega }_U`$ is then real and describes the ‘angular velocity of the horizon’. Moreover, $`J=1/(4\mathrm{\Omega }_U^2)=M^2`$, and hence $`\epsilon =1`$.
## ACKNOWLEDGEMENTS
The author would like to thank A. Kleinwächter, R. Meinel, and G. Neugebauer for many valuable discussions.
The support from the DFG is gratefully acknowledged.
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# 1 Introduction
## 1 Introduction
There has been considerable interest over the last few years in space-time geometries where gauge and matter degrees of freedom are confined to a four-dimensional submanifold, while gravity is allowed to propagate in the whole of space-time. The motivation for such geometries comes from Hořava-Witten theory, or M-theory, compactified on $`S_1/Z_2`$ , or models built from D-branes. For instance, in Hořava–Witten theory, after compactification on a Calabi-Yau three-fold, the five-dimensional vacuum solution is found to consist of two three-branes, each coinciding with an orbifold fixed plane . The brane tensions have opposite sign, and their magnitude is related to the bulk cosmological constant. The existence of a higher-dimensional bulk is likely to have deep implications for the dynamics of the early universe, in particular for inflation.
More recently, Randall and Sundrum proposed that the setup of two domain walls bounding an orbifold could solve the hierarchy problem. Cosmological solutions for similar geometries were presented in Refs. , but all assuming some potential that stabilizes the orbifold separation. While a stabilizing potential might be necessary to obtain a realistic picture of our universe today, there is no reason to believe that the extra dimension was static during inflation. In Ref. , general inflationary solutions were given for compactified Hořava-Witten theory, where the brane energy density is constant and the bulk cosmological constant is set equal to zero. For the case of negative bulk cosmological constant, a class of solutions was derived in Ref. where the bulk geometry is pure anti-de Sitter (AdS) space. However, the most general solution was not obtained.
In section 2, we find a class of solutions in the case where the bulk geometry is AdS<sub>5</sub>, the stress energy on the walls is given by their tensions, and the walls are located at fixed orbifold coordinates. The coordinate system used to find the solutions is not the most general, but has the nice property of making homogeneity and isotropy along the three spatial directions (i.e., planar symmetry) manifest at each point along the orbifold direction and, in particular, on each brane. While our ansatz coincides with that of Ref. , our class of solutions (described algebraically in section 2) is obtained by explicitly solving Einstein’s equations, and hence is the most general within the set-up described above. It is parameterized by a few constants and an arbitrary periodic function of one variable, and is very similar in form to that found in Ref. for the case of Minkowskian bulk geometry.
Since we are interested in inflationary solutions, we restrict to the case where the tensions of the branes are greater in magnitude than the bulk cosmological constant. With this minor assumption, our solutions all describe the same cosmology on the boundaries, namely that of a de Sitter (dS) phase. Thus, in order for the branes to generate AdS in the bulk, it is necessary that they follow de Sitter trajectories in AdS. It is therefore natural to interpret our solutions as embeddings of two dS<sub>4</sub> surfaces in AdS<sub>5</sub>, and this is the focus of the second half of the paper. As a warm-up, in section 3 we will describe how the solution of Ref. (with zero bulk cosmological constant) describes slices of $`_5`$ (Minkowski space). The generalization to AdS<sub>5</sub> is discussed in section 4. A nice feature of this analysis is that it provides a geometrical interpretation of the various arbitrary parameters describing our class of solutions. It is also a useful tool to describe the time-evolution of the orbifold as well as the causal properties of the solutions (sect. 3 and 4). We find that our solution describes scenarios where the orbifold direction is static, collapses or expands, as seen from an observer living on one of the branes. In solutions with expanding extra dimension, two-way communication between the walls is only possible for a finite amount of time. Afterwards, communication can proceed in one way only: signals emitted from the negative-tension brane do not reach the positive-tension brane in finite affine parameter.
Finally, the geometrical picture allows us to realize that the embeddings of section 2 display an unexpected amount of symmetry, leading us to postulate that our solutions are in fact a subclass of all possible embeddings of two dS<sub>4</sub> in AdS<sub>5</sub>. In section 5, we derive the most general solution for two domain walls in AdS<sub>5</sub>. In particular, it includes novel scenarios where the brane tensions are not restricted to have opposite signs. The stability of the solutions as well as generalizations to any FRW cosmologies on the walls are discussed in section 6.
## 2 A class of solutions in five dimensions
We are interested in two four-dimensional hypersurfaces (branes), $`_4^{(1)}`$ and $`_4^{(2)}`$, embedded in a five-dimensional manifold such that $`Z_2`$ symmetry holds across each brane. The stress-energy in the bulk is given by a negative cosmological constant, $`\lambda <0`$, although, at least for part of this section, the given solution also holds for $`\lambda >0`$. We shall also assume that the energy density on each wall is dominated by a cosmological constant (tension) denoted by $`\rho _i`$, $`i=1,2`$. With the metric signature $`(,+,+,+,+)`$ and in units where the five-dimensional Newton constant equals unity, the action is given by
$$S=__5\sqrt{g_5}[_512\lambda ]\underset{i=1}{\overset{2}{}}_{_4^{(i)}}d^4x\sqrt{g_4^{(i)}}12\rho _i,$$
(1)
where $`g_4^{(i)}`$ is the induced metric on the domain walls. Without loss of generality, we can define a coordinate $`y`$ so that the first domain wall is located at $`y=0`$. To obtain cosmologically relevant solutions, we must impose (spatial) homogeneity and isotropy (i.e., planar symmetry) on this wall. Assuming that planar symmetry is also a symmetry of the bulk (i.e., we can foliate the space-time with planar-symmetric hypersurfaces), a general ansatz for the bulk metric is
$$ds_5^2=e^{2\beta }\left(d\tau ^2+dy^2\right)+e^{2\alpha }d\stackrel{}{x}^2$$
(2)
where $`\alpha =\alpha (\tau ,y)`$ and $`\beta =\beta (\tau ,y)`$. Note that we have used the freedom of coordinate reparameterization in the $`(\tau ,y)`$ plane to choose conformal gauge for that part of the metric. With the assumption of planar symmetry, we shall see that the bulk geometry must be AdS or AdS-Schwarzschild. For such bulk geometries, the set-up is so far general, that is, one can find coordinates in which the first domain wall is located at $`y=0`$. If we want to embed a second domain wall, then, in general, its location will be described by a function of the coordinates, $`y_2=y_2(t,\stackrel{}{x})`$, such that homogeneity and isotropy on that brane might not be manifest. However, for simplicity, we shall assume that the second wall lies at $`y=R`$ (and is thus aligned with the first). Note that by a simple change of coordinates in the $`(\tau ,y)`$ plane, one can show that our set-up also includes cases where the branes move uniformly according to $`y_i=y_i(\tau )`$. Later, we will discuss more general configurations.
Let us first focus our attention on the bulk solutions. One obtains the following (bulk) equations of motion (a dot represents time differentiation while a prime denotes differentiation with respect to $`y`$)
$`(0,0):\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}3}(\dot{\alpha }^2+\dot{\alpha }\dot{\beta }\alpha ^{\prime \prime }+\alpha ^{}\beta ^{}2\alpha ^2)=6\lambda e^{2\beta }`$
$`(5,5):\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}3}(\ddot{\alpha }+\dot{\alpha }\dot{\beta }2\dot{\alpha }^2+\alpha ^2+\alpha ^{}\beta ^{})=6\lambda e^{2\beta }`$
$`(i,i):(2\alpha ^{\prime \prime }+\beta ^{\prime \prime }+3\alpha ^2)(2\ddot{\alpha }+\ddot{\beta }+3\dot{\alpha }^2)=6\lambda e^{2\beta }`$
$`(0,5):\dot{\alpha }^{}+\dot{\alpha }\alpha ^{}\dot{\alpha }\beta ^{}\dot{\beta }\alpha ^{}=0.`$ (3)
By taking linear combinations of the equations, one can rewrite them in terms of light-cone coordinates, $`x^\pm =\tau \pm y`$, as
$$_+_{}\alpha +_+_{}\beta =\frac{\lambda }{2}e^{2\beta }$$
(4)
$$_+_{}\alpha +3_+\alpha _{}\alpha =\lambda e^{2\beta }$$
(5)
$$_+^2\alpha 2_+\alpha _+\beta +(_+\alpha )^2=0$$
(6)
$$_{}^2\alpha 2_{}\alpha _{}\beta +(_{}\alpha )^2=0.$$
(7)
As shown in Ref. , one can derive a first integral of motion by rewriting the $`(0,0)`$ and $`(5,5)`$ equations as
$`F^{}(\tau ,y)=4\lambda \alpha ^{}e^{4\alpha }`$
$`\dot{F}(\tau ,y)=4\lambda \dot{\alpha }e^{4\alpha },`$ (8)
where $`F(\tau ,y)e^{4\alpha }e^{2\beta }(\alpha ^2\dot{\alpha }^2)`$. Equations (8) can then be integrated to yield a single expression
$$\dot{\alpha }^2\alpha ^2=\left[\lambda +Ce^{4\alpha }\right]e^{2\beta },$$
(9)
where $`C`$ is an integration constant.
The assumption of planar symmetry has reduced the allowed bulk geometries to a one-parameter family of solutions depending on $`C`$. From now on we will focus on the the case $`C=0`$. As we will see in the following sections, the geometry of the corresponding bulk space-times is then particularly simple, namely AdS. However, it is also possible to consider the case $`C0`$, which corresponds to a bulk Schwarzschild-AdS geometry .
Taking $`C=0`$, in light-cone coordinates, equation (9) reads
$$_+\alpha _{}\alpha =\frac{\lambda }{4}e^{2\beta }.$$
(10)
We are therefore left with the $`(0,5)`$ equation and equation (10) (since the $`(i,i)`$ equation then follows from the Bianchi identity). Combining equations (4) and (5) with equation (10) yields
$$_+_{}\alpha =_+_{}\beta =\frac{\lambda }{4}e^{2\beta }.$$
(11)
The solution to equations (11) and $`(0,5)`$ can be expressed in terms of two arbitrary functions $`f=f(x^+)`$ and $`g=g(x^{})`$ as
$$ds^2=\frac{1}{(f+g)^2}\left(\frac{4}{\lambda }f^{}g^{}(d\tau ^2+dy^2)+d\stackrel{}{x}^2\right),$$
(12)
where $`f^{}df/dx^+`$ and $`g^{}dg/dx^{}`$. Although we will be interested in the case $`\lambda <0`$, this solution is, in fact, valid for any non-zero value of $`\lambda `$.
Finally, to fix the specific coordinate patch on which we analyze our solutions, let us assume that the range of $`\tau `$ and $`y`$ are such that $`\tau `$ is always time-like, $`y`$ space-like, and further, that $`f+g`$ is positive. Given that $`\lambda <0`$ this means that we will take
$$f(x^+)+g(x^{})>0$$
(13)
$$f^{}(x^+)g^{}(x^{})<0.$$
(14)
for all $`x^+`$ and $`x^{}`$.
Let us now turn to the domain walls. We must supplement the equations of motion with appropriate boundary conditions. Since the branes are localized objects along the transverse direction (delta-function sources), they result in discontinuities in the normal derivative of the metric. Formally, these discontinuities are described by the Israel matching conditions , which yield the following boundary conditions (assuming $`Z_2`$ symmetry across each wall)
$$e^\beta \alpha ^{}|_{y=0}=\rho _1,e^\beta \alpha ^{}|_{y=R}=+\rho _2,$$
(15)
$$e^\beta \beta ^{}|_{y=0}=\rho _1,e^\beta \beta ^{}|_{y=R}=+\rho _2.$$
(16)
where we note that there is a difference in sign between the two walls.
The Israel junction conditions impose restrictions on the functions $`f`$ and $`g`$. From the boundary condition (15) evaluated at $`y=0`$ and given the conditions (13) and (14), we obtain
$`g^{}(\tau )f^{}(\tau )=\rho _1\sqrt{{\displaystyle \frac{4}{\lambda }}f^{}(\tau )g^{}(\tau )}`$
$`\left({\displaystyle \frac{g^{}(\tau )}{f^{}(\tau )}}\right)^{1/2}={\displaystyle \frac{1}{\sqrt{\lambda }}}(\rho _1\pm H_1)\gamma _1,`$ (17)
where we have introduced
$$H_i=\sqrt{\rho _i^2+\lambda }.$$
(18)
Since $`H_i`$ should be real we find that solutions will only exist if we have the condition on the tension,
$$|\rho _i|\sqrt{\lambda }.$$
(19)
Note we are using here the condition (14), which implies that $`g^{}(x^{})/f^{}(x^+)<0`$ for all $`x^+`$ and $`x^{}`$. It turns out that the relation (17) then implies that the second boundary condition (16) is also satisfied at $`y=0`$.
At $`y=R`$, the boundary conditions imply a similar relationship between $`g(\tau )`$ and $`f(\tau +2R)`$. One finds that
$$\left(\frac{g^{}(\tau )}{f^{}(\tau +2R)}\right)^{1/2}=\frac{1}{\sqrt{\lambda }}(\rho _2\pm H_2)\gamma _2$$
(20)
satisfies both boundary conditions.
Together, equations (17) and (20) imply that
$`g(\tau )=\gamma _1^2f(\tau )k_1,`$
$`g(\tau )=\gamma _2^2f(\tau +2R)k_2,`$ (21)
where $`k_i`$ are constants and
$`\gamma _1{\displaystyle \frac{1}{\sqrt{\lambda }}}(\rho _1\pm H_1),`$
$`\gamma _2{\displaystyle \frac{1}{\sqrt{\lambda }}}(\rho _2\pm H_2).`$ (22)
Note that, by (17) and (20) both $`\gamma _i`$ must be positive. Since $`|\rho _i|>H_i`$ this requires
$$\rho _1>0,\rho _2<0,$$
(23)
that is, for solutions of this type, the brane tensions must have opposite sign. As we shall see in section 5, this condition on the brane tensions is a consequence of our particular choice of coordinate system (2). We will later also describe solutions where this condition does not hold.
Together, the relationships (21) imply a periodicity condition on $`f`$ which can be written as
$$f(x+2R)=\left(\gamma _1/\gamma _2\right)^2f(x)+K,$$
(24)
where the constant $`K`$ is given by
$$K=(k_1k_2)/\gamma _2^2$$
(25)
The general solution of this periodicity condition gives
$$f(x)=e^{\xi x}p(x)+\frac{e^{\xi x}1}{e^{2\xi R}1}K,$$
(26)
where
$$\xi \frac{1}{R}\mathrm{ln}(\gamma _1/\gamma _2),$$
(27)
and $`p(x)`$ is a periodic function $`p(x+2R)=p(x)`$. One notes that in the limit $`\xi 0`$ the general expression (26) becomes
$$f(x)=p(x)+\frac{Kx}{2R}.$$
(28)
In the following, we will sometimes find it easier to use $`kk_1`$ and $`K`$ as the independent constants in our solution, and will sometimes stick to $`k_1`$ and $`k_2`$. For completeness let us now give the general expressions for the functions $`\alpha `$ and $`\beta `$ in the metric in terms of $`f`$. We have
$`e^\alpha ={\displaystyle \frac{1}{f(x^+)\gamma _1^2f(x^{})+k}},`$
$`e^\beta =\sqrt{{\displaystyle \frac{4\gamma _1^2f^{}(x^+)f^{}(x^{})}{\lambda }}}e^\alpha .`$ (29)
We see that the general solution is determined by the choice of a periodic function $`p(x)`$ and some constants, which must be chosen so that equations (13) and (14) are satisfied. While it is difficult at this point to get a feel for what these parameters represent, we will find in sections 3 and 4 a geometrical interpretation for them. Note that our solution reduces to the solution presented in Ref. if we choose $`p=\mathrm{constant}`$ (as we shall see later, such a choice amounts to a coordinate transformation). Nevertheless, we shall find evidence in section 4 (and show in section 5) that the solution in this section is not the most general.
Let us now consider the induced geometry on the branes. A simple approach is to substitute the boundary conditions (15) and (16) in equation (9) (with $`C=0`$) to find the Friedmann equation on the domain wall following . Choosing cosmological time on the brane, one finds
$$\dot{\alpha }^2=\rho _i^2+\lambda =H_i^2.$$
(30)
Given the condition on the tensions (19), equation (30) implies that the induced geometry on each brane is precisely that of de Sitter space with cosmological constant $`H_i`$.
We can also see this more explicitly from our final expression (29). Considering the brane at $`y=0`$ for instance, we can perform the following time redefinition
$$e^{H_1t}=(1\gamma _1^2)f(\tau )+k$$
(31)
so that $`t`$ is cosmological time at $`y=0`$. In terms of the new time variable $`t`$, the induced metric on our domain wall is found to be
$$ds_4^2=dt^2+e^{2H_1t}d\stackrel{}{x}^2,$$
(32)
which describes de Sitter space, in agreement with equation (30). A similar analysis holds for the brane at $`y=R`$.
To conclude this section, let us discuss an important subclass of solutions, namely those with a static orbifold in the coordinates of (2). Starting with the case $`\rho _1=\rho _2`$ (i.e., $`\xi =0`$), we see from equations (29) that a static bulk is obtained only if $`\gamma _1^2=1`$ and $`p=\mathrm{constant}`$. We will show in section 4 that this solution corresponds to the Randall–Sundrum scenario .
In the case $`\rho _1\rho _2`$ (i.e., $`\xi 0`$), equations (29) tell us that static solutions require $`p=\mathrm{constant}`$ and $`k=\frac{(1\gamma _1^2)K}{e^{2\xi R}1}`$. For such choices, one finds the following constraint relating the physical distance between the walls, $`R_{static}`$, the bulk cosmological constant and the brane tensions
$$\mathrm{tanh}(\sqrt{\lambda }R_{static})=\frac{\sqrt{\lambda }(\rho _1+\rho _2)}{\lambda +\rho _1\rho _2},$$
(33)
in agreement with Refs. and . Since $`\rho _i^2+\lambda 0`$ and $`\rho _1\rho _2<0`$, the denominator in equation (33) is negative. In turn, this implies that $`\rho _2>\rho _1>0`$.
## 3 Geometrical analysis of $`\lambda =0`$ solution
In this section, we shall discuss the limit $`\lambda 0^{}`$ of our solution from a geometrical perspective. In the process, we will find geometrical interpretations for the various parameters describing the general solution (such as $`k`$, $`K`$, etc.).
Let us first determine the bulk geometry in this limit. We note that general solutions with a metric in the form (2) with $`\lambda =0`$ were first discussed (albeit in four dimensions) by Taub (see also ). In that paper it was shown that all solutions with $`C=0`$ were simply different parametrizations of flat space, while all solutions with $`C0`$ were simply different parametrizations of the Kasner “rolling radii” solutions . Here we will rederive the $`C=0`$ case as a limit of our general solution. From equation (12), we note that taking $`\lambda 0`$ requires letting either $`f^{}`$ or $`g^{}`$ go to zero at the same rate. Suppose we let $`g^{}0`$. Defining
$$g(x^{})=\frac{\lambda }{4M^2}\stackrel{~}{g}(x^{})+\stackrel{~}{k},$$
(34)
where $`\stackrel{~}{k}`$ is a constant and $`M`$ some energy scale, equation (12) becomes
$$ds^2=\frac{1}{(f(x^+)+\stackrel{~}{k})^2}\left(\frac{1}{M^2}df(x^+)d\stackrel{~}{g}(x^{})+d\stackrel{}{x}^2\right).$$
(35)
To express this metric in a more familiar form, we can perform the coordinate transformation to $`(T,Y,\stackrel{}{X})`$
$`Y+T={\displaystyle \frac{1}{M(f(x^+)+\stackrel{~}{k})}}`$
$`YT={\displaystyle \frac{1}{M}}\stackrel{~}{g}(x^{}){\displaystyle \frac{M\stackrel{}{x}^2}{f(x^+)+\stackrel{~}{k}}}`$
$`\stackrel{}{X}={\displaystyle \frac{\stackrel{}{x}}{f(x^+)+\stackrel{~}{k}}},`$ (36)
which gives
$$ds^2=dT^2+dY^2+d\stackrel{}{X}^2.$$
(37)
Thus we have shown that the case $`C=0`$ case indeed corresponds to five-dimensional Minkowski space ($`_5`$) in the limit $`\lambda 0`$. If we now consider the branes at the orbifold fixed points, this imposes boundary conditions on $`f`$ and $`g`$ (see eqns. (21)). For these to be consistent with equation (34), we need to choose the signs in $`\gamma _i`$ (see eq. (22)) such that $`\gamma _i0`$ as $`\lambda 0`$ (note also that this yields $`\gamma _1/\gamma _2|\rho _2/\rho _1|`$). Furthermore, we must identify $`\stackrel{~}{k}`$ with the constant $`k`$ introduced in section 2, $`M`$ with the tension $`\rho _1`$ of the $`y=0`$ brane, and $`\stackrel{~}{g}(x^{})=f(x^{})`$ (from eq. (21)). Similarly, the above derivation can be repeated for the choice $`f^{}0`$. Note that our solution in the limit $`\lambda 0`$ agrees with the results presented in Ref. . Finally, combining equation (13) with the first of equations (36), we note that our original coordinates are mapped to the range $`Y+T>0`$.
We saw in the previous section that the induced geometry on each domain wall is that of de Sitter space. We can therefore interpret the solutions obtained in section 2 (for $`\lambda 0`$) as possible embeddings of two dS<sub>4</sub> surfaces in $`_5`$. As in well-known (e.g., see Ref. ), a de Sitter surface can be described as a hyperboloid embedded in $`_5`$. Indeed, evaluating equations (36) at $`y=0`$ and $`y=R`$, one finds that the branes are given by
$`\left(T{\displaystyle \frac{k}{2\rho _1}}\right)^2+\left(Y+{\displaystyle \frac{k}{2\rho _1}}\right)^2+\stackrel{}{X}^2={\displaystyle \frac{1}{\rho _1^2}}\mathrm{for}y=0,`$
$`\left(T{\displaystyle \frac{\rho _1(k+K)}{2\rho _2^2}}\right)^2+\left(Y+{\displaystyle \frac{\rho _1(k+K)}{2\rho _2^2}}\right)^2+\stackrel{}{X}^2={\displaystyle \frac{1}{\rho _2^2}}\mathrm{for}y=R.`$ (38)
From equations (38), we note that:
* The curvature of the hyperboloids are given by $`|\rho _i|`$.
* The origins of the light cones asymptotic to the hyperboloids are separated by the null vector
$$D=\left[(\rho _1/2\rho _2^2)(K+k)k/2\rho _1\right](1,1,\stackrel{}{0}).$$
(39)
* For fixed null separation $`D`$, the quantity $`k`$ simply translates the two hyperboloids by the same vector. It therefore describes the invariance of the solution upon translations in the embedding space and corresponds to a change of coordinates in (37).
While the intrinsic curvature of each hyperboloid is given by $`|\rho _i|`$, the sign of $`\rho _i`$ determines on which side of the hyperboloid the bulk extends. The bulk lies outside (inside) the hyperboloid if $`\rho _i<0`$ ($`>0`$). As an example, figure 1 shows the solution with $`\rho _1=\rho _2`$ and parameters $`k=0`$ and $`K>0`$, (where we have suppressed two spatial dimensions). In accordance with the above discussion, the portion of space-time consistent with $`\rho _1>0`$ and $`\rho _2<0`$ is indicated by the shaded region. From (38) we see that the two hyperboloids intersect in the cylinder $`Y+T=0`$, $`\stackrel{}{X}^2=1/\rho _1^2`$. Since $`Y+T=0`$ is a null plane and given our initial coordinate patch, it is easy to show that no observer in the shaded region can ever reach this intersection. Due to the symmetry in the $`\stackrel{}{X}`$ directions, it is sufficient for our purposes to focus on the two-dimensional projection onto the $`\stackrel{}{X}=0`$ plane (shown in fig. 2 for the above example). An important point to notice is that the hyperbolas are separated by a null vector (as mentioned in the second point above), and therefore share a common asymptote.
We can also consider examples where the curvatures of the hyperboloids are not equal, once again focusing on the subspace $`\stackrel{}{X}=0`$. For simplicity, let us set $`k=0`$. The case $`K=0`$ with $`\rho _2>\rho _1>0`$ is shown in figure 3, with the bulk indicated by the shaded region. Note that the branes never intersect for $`K=0`$. Furthermore, it is crucial that $`\rho _2>\rho _1`$ in order to obtain a bulk region that is simultaneously inside one brane and outside the other, which is equivalent to the condition stated below equation (33).
For the case $`K0`$, we see from (38) that the hyperboloids now do intersect in the paraboloid
$`Y+T=B{\displaystyle \frac{\rho _2^2\rho _1^2}{(\rho _1/\rho _2^2)(K+k)k/\rho _1}}`$
$`YT={\displaystyle \frac{1}{B}}\stackrel{}{X}^2+{\displaystyle \frac{1}{B\rho _1^2}}{\displaystyle \frac{k}{\rho _1}}.`$ (40)
Note that, as above, the intersection lies in a null plane. Furthermore, since our coordinates are restricted to $`Y+T>0`$, the intersection lies within our coordinate patch only if $`B>0`$. The solution for $`\rho _2>\rho _1>0`$ with $`K<0`$ and $`K>0`$ are illustrated in figures 4 and 5, respectively. These two figures are related by reflections about the $`T`$ and $`Y`$ axes, but they should be viewed as distinct solutions in our analysis due to the constraint $`Y+T>0`$. In particular, the intersection lies in the physically accessible region in figure 4, where $`B>0`$, but not in figure 5, where $`B<0`$. Finally, if $`\rho _1>\rho _2`$, the only choice consistent with there being a bulk between the two branes in the region $`Y+T>0`$ is $`K>0`$, in which case there is also an intersection in our coordinate patch, as illustrated in figure 6. (For the case $`K<0`$, the signs of the brane tensions do not allow for a bulk region between them in the region $`Y+T>0`$.)
We have shown that a broad class of configurations can be embedded in flat space, namely those consisting of null-separated de Sitter surfaces. The fact that we have only encountered cases with null separation is surprising, and leads us to question the generality of the solutions obtained in section 2. Intuitively, one would also expect cases with time-like and space-like separated hyperboloids. We shall confirm this intuition in section 5.
### 3.1 Dynamics of the orbifold
We next want to investigate the evolution of the orbifold as seen by an observer on the brane at $`y=0`$. First, let the cosmological time for this observer be denoted by $`t^{(1)}`$ (i.e., $`t^{(1)}`$ puts the induced metric at $`y=0`$ in the form given in eq. (32)). The range $`\mathrm{}<t^{(1)}<\mathrm{}`$ is mapped to the range $`\mathrm{}<T<\mathrm{}`$, with the rule that $`t^{(1)}`$ increases with $`T`$ in the region $`Y+T>0`$. Note that requiring the time coordinate to be cosmological time at $`y=0`$ does not uniquely fix the surfaces of constant time away from the brane. In this sense, the distance between the branes is not a coordinate-independent notion. However, a natural prescription is the following. First, choose coordinate such that the metric has the form (2). Then recall that we have seen that that neither the function $`p(x)`$ nor $`k`$ in the solution encode physical information. Thus it is natural to take the case $`p=\mathrm{const}`$ and $`k=0`$. One then denotes the separation as the distance in the $`y`$-direction with fixed $`t`$ and $`\stackrel{}{x}`$ in this metric.
As seen in figures 2-6, we have five physically distinct solutions to consider, which we classify into three cases.
* $`\rho _2>\rho _1`$:
+ $`K=0`$: This solution is shown in figure 3. As mentioned at the end of section 2, it has a static orbifold due to the choice $`p=const`$. The physical length of the orbifold is given by the $`\lambda 0`$ limit of equation (33), namely
$$R_{static}\frac{1}{\rho _1}+\frac{1}{\rho _2},$$
(41)
and agrees with the expression found in Ref. . Finally, the curves of constant $`t^{(1)}`$ are straight lines trough the origin.
+ $`K<0`$: This solution corresponds to the shaded region in fig. 4. If we denote the physical distance between the branes by $`R_{phys}`$ in this case, one finds that $`R_{phys}R_{static}^{}`$ as $`t^{(1)}\mathrm{}`$. As $`t^{(1)}`$ increases, the distance between the domain walls decreases until it reaches zero size (corresponding to the point where the hyperboloids intersect). Hence, if we define the onset of inflation to occur at at some $`t_0^{(1)}`$, where $`\mathrm{}<t_0^{(1)}<t_{collapse}^{(1)}`$, this solution describes an initial condition where the branes are within $`R_{static}`$ apart, and subsequently move toward each other.
+ $`K>0`$: In this case (shown in fig. 5), $`R_{phys}R_{static}^+`$ as $`t^{(1)}\mathrm{}`$. Subsequently, $`R_{phys}`$ increases forever.
Hence, for the case $`\rho _2>\rho _1`$, the orbifold remains static if the initial distance is $`R_{static}`$, collapses if it is less than $`R_{static}`$, and greater if it is larger than $`R_{static}`$. This conclusion agrees with the analysis of Ref. .
* $`\rho _2<\rho _1`$: This solution was not discussed in Ref. . From figure 6, we see that it describes an orbifold starting from zero size and expanding to infinity. Hence, in this case, any initial distance leads to expansion.
* $`\rho _2=\rho _1`$: By comparing figures 2 and 6, we see that the evolution is the same as in the case $`\rho _2<\rho _1`$.
As mentioned earlier, different choices of $`p`$ correspond to different observers in the sense that they are associated with diffeomorphisms which leave the bulk metric in the form (2) and the induced metric at $`y=0`$ in the form (32). The notion of the separation of the branes is coordinate dependent. To see this explicitly suppose, for instance, we allow for a generic choice of $`p`$. It turns out that an observer at $`y=0`$ will then see, on top of the average behavior of the extra dimension, an oscillatory component. For instance, in the case $`\rho _2<\rho _1`$, the orbifold would oscillate as it expands. Note that, given the restrictions imposed by equations (13) and (14), the function $`p`$ must be chosen such that its oscillations will not lead to $`R_{phys}`$ reaching zero size where it would not in the case $`p=const`$.
Nonetheless, there is also an observer-independent statement one can make about the various solutions described above. We noted that, in general, the branes intersect in a paraboloid (3). This intersection lies in the coordinate patch $`Y+T>0`$ if $`B>0`$. Furthermore, it lies a finite distance in the future of any observer in the bulk (corresponding to an orbifold collapsing to a final singularity) if $`\rho _2>\rho _1>0`$, or in the past (corresponding to an orbifold expanding from an initial singularity) if $`\rho _1>\rho _2>0`$. Finally, let us note that the surface of intersection is formally a singular surface and, thus, our solution breaks down in its neighborhood. A full treatment of brane collisions would require including effects due to the finite thickness of the branes (usually of string size).
### 3.2 Causal properties
Let us now consider the causal properties of the solution. More precisely, we want to see if the walls are causally-connected during the inflationary period, and whether two-way or one-way communication between the walls is possible.
In the cases where the orbifold collapses or remains static (i.e. $`K0`$), two-way communication is possible for all times, as expected.
The remaining solutions involve expanding orbifolds. Since the extra dimension is in fact inflating, one would expect that a signal sent from one brane could never reach the other. However, we must keep in mind that the length of the orbifold is computed along a space-like path, while gravitons propagate along null geodesics. Consequently, the causal properties are not as trivial as one would have guessed. The analysis yields qualitatively the same result for all solutions which describe expanding orbifolds, so, without loss of generality, we can focus on the case $`\rho _1=\rho _2`$ for concreteness. Suppose observer 2 at $`y=R`$ sends signals to observer 1 at $`y=0`$. The two dotted lines in figure 7 correspond to two such signals. We see that the first signal reaches $`y=0`$ in finite affine parameter while the other does not. Hence, from the point of view of observer 2, there exists some time on his clock where his signals stop reaching the $`y=0`$ brane. We shall say that a causal boundary has appeared for observer 2. Now, consider signals sent from $`y=0`$ towards $`y=R`$. It is clear from figure 7 that such signals always reach $`y=R`$ in finite affine parameter. However, if $`t^{(1)}`$ is physical time for observer 1, then after a while observer 1 sees his signals taking infinite time $`t^{(1)}`$ to reach $`y=R`$. From the point of view of observer 1, we shall say that there is a horizon somewhere in the bulk. Note that this is not a genuine horizon in the five-dimensional space-time, but corresponds to a surface with the following property: once a signal sent by observer 1 has passed through the horizon, it cannot return to observer 1. Note that causality implies that the appearance of a horizon requires the appearance of a causal boundary, and vice versa. To see this, suppose the contrary is true; that is, suppose that there is a horizon for observer 1 but no causal boundary for observer 2. If observer 1 sends a graviton towards $`y=R`$, because of the horizon he will see his graviton frozen somewhere in the bulk. On the other hand, observer 2 will receive this signal within finite time on his clock. If observer 2 replies to the signal, his reply will be received at $`y=0`$ in finite time on observer 1’s clock (because we have assumed that there were no causal boundary). From the point of view of observer 1, he received a reply from observer 2 before his initial signal made it to $`y=R`$, an obvious violation of causality. It is easily seen that the horizon and causal boundary describe the same surface in the bulk, namely the future light cone asymptotic to the $`y=0`$ hyperboloid. That is, once observer 2 is within this surface, his signals will not reach $`y=0`$ in finite affine parameter.
To summarize our conclusions for expanding extra dimensions, our solution predicts that two-way communication is only possible for a finite amount of time. Afterwards, only the brane with positive tension can send signals to (and, hence, have influence on) the negative-tension brane.
## 4 Geometrical analysis of $`\lambda <0`$ solution
In this section, we generalize the analysis of the previous section to the case of negative bulk cosmological constant. As there, we will find that all the solutions presented in section 2 correspond to the same bulk space-time. In this case it is AdS<sub>5</sub>. We can therefore think of our solution as a class of possible ways to embed two dS<sub>4</sub> surfaces in AdS<sub>5</sub>. As in the case of flat space discussed in the previous section, we shall find evidence that our solution does not describe the full spectrum of such embeddings.
Recall that the general bulk solution has the form (12)
$$ds^2=\frac{1}{(f+g)^2}\left(\frac{4}{\lambda }f^{}g^{}(d\tau ^2+dy^2)+d\stackrel{}{x}^2\right),$$
(42)
With the line element in this form, it is not clear that, in fact, for $`\lambda <0`$, in all cases the bulk geometry is that of AdS<sub>5</sub>. However, upon the coordinate transformation
$`z+t={\displaystyle \frac{2}{\sqrt{\lambda }}}f`$
$`zt={\displaystyle \frac{2}{\sqrt{\lambda }}}g,`$ (43)
we have
$$ds^2=\frac{1}{(\lambda )z^2}(dt^2+d\stackrel{}{x}^2+dz^2),$$
(44)
which is a more familiar form for the line element of AdS<sub>5</sub>. We see that the general functions $`f`$ and $`g`$ simply represented different choices of conformal coordinates $`\tau `$ and $`y`$ in the AdS<sub>5</sub> space.
As is well known, this coordinate system does not cover the whole of AdS space. More generally one can describe AdS as a hyperboloid
$$X^2T_1^2+T_2^2Y^2\stackrel{}{X}^2=\frac{1}{\lambda }$$
(45)
embedded in a flat six-dimensional space $`X=(T_1,T_2,Y,\stackrel{}{X})`$ with line element
$$ds^2=dXdXdT_1^2dT_2^2+dY^2+d\stackrel{}{X}^2.$$
(46)
We should note here that the space described contains closed time-like curves. The full AdS space is really the universal covering space. This will not enter our analysis here, since in all our solutions the domain walls will lie in the same sheet of the universal cover. In terms of the coordinates in (44) the embedding is given by
$`T_1+Y={\displaystyle \frac{1}{(\lambda )z}}`$
$`T_1Y={\displaystyle \frac{\stackrel{}{x}^2t^2+z^2}{z}}`$
$`T_2={\displaystyle \frac{t}{\sqrt{\lambda }z}}`$
$`\stackrel{}{X}={\displaystyle \frac{\stackrel{}{x}}{\sqrt{\lambda }z}}.`$ (47)
Note that since $`z=(f+g)/\sqrt{\lambda }>0`$, our original coordinates are mapped to the range $`T_1+Y>0`$. (The boundary of this region, $`z=0`$, is shown in figure 8.)
Let us now turn to the description of the domain walls. Recall that in our original coordinates the walls were fixed at $`y=0`$ and $`y=R`$, or equivalently $`x^+=x^{}`$ and $`x^+=x^{}+2R`$. After the conformal transformation (43) the walls will now be moving in the $`z`$ direction. Explicitly, from the relations (21), the equations of the walls in terms of $`z`$ and $`t`$ are the linear relations
$`tz=\gamma _1^2(t+z)+{\displaystyle \frac{2k_1}{\sqrt{\lambda }}}`$ $`\text{for }y=0;`$ (48)
$`tz=\gamma _2^2(t+z)+{\displaystyle \frac{2k_2}{\sqrt{\lambda }}}`$ $`\text{for }y=R.`$
It is easy to see that in these coordinates, the domain walls move in the $`z`$ direction with constant velocity $`H_i/|\rho _i|`$.
Finally we can transform these equations into the flat six-dimensional embedding space. One finds that the walls correspond to planes
$`\pm T_2+k_1\left({\displaystyle \frac{\rho _1}{H_1}}1\right)(T_1+Y)={\displaystyle \frac{\rho _1}{\sqrt{\lambda }H_1}}`$ $`\text{for }y=0;`$ (49)
$`\pm T_2+k_2\left({\displaystyle \frac{|\rho _2|}{H_2}}1\right)(T_1+Y)={\displaystyle \frac{|\rho _2|}{\sqrt{\lambda }H_2}}`$ $`\text{for }y=R.`$
The choice of signs correspond to the choice of sign in the expressions (22) for $`\gamma _1`$ and $`\gamma _2`$.
Note that both these equations have the form, for $`i=1,2`$,
$$n_iX=c_i$$
(50)
where $`n_i`$ is a time-like unit vector and $`c_i>1/\sqrt{\lambda }`$. Intersecting with the AdS<sub>5</sub> hyperboloid (45), one finds that the domain walls are dS<sub>4</sub> submanifolds of the embedding space of curvature
$$\left(c_i^2+\frac{1}{\lambda ^2}\right)^{1/2}=H_i,$$
(51)
where we have substituted the particular form of $`c_i`$ from equation (49). (This can be seen explicitly by using the $`SO(4,2)`$ symmetry of the embedding flat space to put $`\stackrel{}{n}=(1,0,\mathrm{},0)`$ and substituting in eqn. (45)). Thus we see again how the domain walls are indeed dS<sub>4</sub> surfaces in the AdS<sub>5</sub> space with curvatures $`H_i`$.
To get a sense of the global structure of the solution consider the case where $`k_1=k_2=0`$ and choose the upper signs in the solution (49). The planes are then simply given by $`T_2=|\rho _i|/(H_i\sqrt{\lambda })`$. The intersection is sketched in figure 8. Note that in this case, the branes never intersect. The figure allows to understand how the curves describing the domain walls would turn into the hyperbolas shown in figure 3 if we let $`\lambda 0`$. Intuitively, we can generate the solutions described above by taking the hyperbolas of the previous section and “pasting” them on the hyperboloid as shown in figure 8. We can use this intuition to realize that the various parameters (e.g., $`k`$, $`p`$, etc.) describing the solution play a similar role as in the case $`\lambda =0`$. However, we should also note that, in general, there is a second class of solutions where we make the opposite choice of signs for the two walls in (49). (Actually this choice is not possible for the specific case $`k_1=k_2=0`$.) Now the planes are on “opposite sides” of the AdS hyperboloid. This has no natural $`\lambda 0`$ limit. It corresponds to the result of section 3 that we were forced to choose a correlated signs in $`\gamma _i`$ to get the flat-space solution.
As a second example, we can briefly mention that the domain walls in the Randall–Sundrum scenario (taking the $`|\rho _i|\sqrt{\lambda }`$ limit in eqns. (49)) are described by the null planes
$`T_1+Y={\displaystyle \frac{1}{k_1\sqrt{\lambda }}}`$ $`\text{for }y=0,`$ (52)
$`T_1+Y={\displaystyle \frac{1}{k_2\sqrt{\lambda }}}`$ $`\text{for }y=R,`$
Again these branes will never intersect.
To cast the above discussion in a more general set-up, recall that in the previous section, the relative position of domain walls was characterized in flat space by their Minkowskian distance (more formally, by the distance between their asymptotic light cones). Furthermore, in general this separation was null for the solutions of section 2. Similarly here, the asymptote of each dS in the six-dimensional embedding space is a light cone, whose origin lies at
$$a_i=c_in_i\text{for }i=1,2.$$
(53)
Note that, in general, the vector $`a_i`$ will not lie on the bulk AdS hyperboloid. We can now use the distance between the light cones, given by
$$Da_1a_2,$$
(54)
to characterize the relative location of the dS surfaces in a coordinate-independent way. Furthermore, it is easily seen that the quantity $`D`$ does reduce to the distance between hyperboloids in the limit $`\lambda 0`$, as defined in section 3. Just as in the case $`\lambda =0`$ where our solutions all described null separated hyperboloids, similarly, it can be easily verified from equations (49) that all our $`\lambda <0`$ solutions with a flat limit yield $`D^2=0`$, corresponding to null separation. The size of the separation is controlled by the parameters $`k_1`$ and $`k_2`$.
Once again, we expect that solutions with time-like and space-like separated dS surfaces are also allowed. We shall discuss these more general configurations in section 5.
Finally, we note that the dynamics of the orbifold as well as the causal properties of the solution are qualitatively the same as for the case $`\lambda =0`$. The nature of the intersection between the walls is controlled by $`n_i`$ and $`\rho _i`$. For general $`n_i`$, the intersection is a null paraboloid. As before there are two cases, one where the space-time has either expands from an initial singularity and one which contracts to a final singularity. For the special case where the $`n_i`$ are parallel or anti-parallel, there is no intersection and the distance between the walls is fixed. For example, in the case $`\rho _2>\rho _1>0`$, we find that the orbifold remains static if the initial distance between the branes equals $`R_{static}`$ (see eqn. (33)) and expands (collapses) if it is larger (smaller) than $`R_{static}`$. Furthermore, two-way communication is only possible for a finite amount of time if the extra dimension is expanding.
## 5 General embedded solutions
In this section we will step back a little and reconsider the solutions given thus far. In doing so we will see that they are in fact special cases of a more general configuration of a pair of branes in bulk AdS space. Note that in this section our solutions will also no longer be confined to a particular coordinate patch of AdS space.
From the discussion of previous section, the solutions we found correspond to a pair of dS surfaces embedded in AdS space. The dS surfaces are not in general position but are null separated in the sense discussed below eqn. (54). Recall that the solutions were found by solving the bulk Einstein equation with negative cosmological constant together with the Israel matching conditions describing the discontinuity in the normal derivative of the metric at the brane.
The important point to note is that once we fix the bulk solution, in this case to be AdS, the solution of the Israel conditions for the two walls are completely independent. Each set of conditions are local equations relating the shape of the brane embedded in the bulk space to the stress energy on that brane, independent of the second brane. From the analysis of the previous two sections, it appears that for pure cosmological constant $`\rho `$ on the brane, the solution to the Israel conditions is that the brane describe a dS surface embedded in AdS. If this is correct, it is then clear that the general solution corresponds to a pair of dS brane embedded in AdS with arbitrary separation.
To see that this is indeed the case, we can consider the Israel conditions in a coordinate independent form and show that, when the bulk space is AdS<sub>5</sub>, they imply that the brane is a dS<sub>4</sub> surface. If we let $`t^\mu `$ be the normal vector to the brane, the induced metric $`g_{\mu \nu }^4`$ is then given by
$$g_{\mu \nu }=g_{\mu \nu }^4+t_\mu t_\nu ,$$
(55)
where $`g_{\mu \nu }`$ is the bulk metric. As has been noted in various papers (see e.g., ), assuming $`Z_2`$ reflection invariance at the brane, the Israel conditions relate the extrinsic curvature $`K_{\mu \nu }`$ of the brane to the stress energy $`T_{\mu \nu }^B=6\rho g_{\mu \nu }^4`$ on the brane. One has
$$K_{\mu \nu }=\frac{1}{2}\left(T_{\mu \nu }^B\frac{1}{3}T^Bg_{\mu \nu }^4\right)=\rho g_{\mu \nu }^4,$$
(56)
where the extrinsic curvature is given by
$$K_{\mu \nu }=g_\mu ^{4\rho }g_\nu ^{4\kappa }_\rho t_\kappa .$$
(57)
Since $`g_{\mu \nu }^4`$ and $`t_\mu `$ are functions of the embedding, equation (56) is a local differential equation for the functions describing the embedding of the brane in AdS<sub>5</sub>, completely independent of the presence of a second brane.
To show that a dS<sub>4</sub> surface satisfies the Israel conditions, we could evaluate (56) in a particular set of coordinates. This is essentially what was done in section 2. Alternatively, we can argue simply by symmetry that if the brane is dS the extrinsic curvature must be proportional $`g_{\mu \nu }^4`$, since this is the only symmetric tensor on the brane with the correct symmetries. The only question is then what curvature of the dS space must we choose to make the constant of proportionality exactly that in (56). To answer this we recall that there is an expression for the intrinsic curvature $`R_4`$ of $`g_{\mu \nu }^4`$ in terms of the bulk curvature $`R`$ and $`K_{\mu \nu }`$. In particular, we have the general expression
$$R_{\kappa \lambda \mu \nu }^4=g_\kappa ^{4\kappa ^{}}g_\lambda ^{4\lambda ^{}}g_\mu ^{4\mu ^{}}g_\nu ^{4\nu ^{}}R_{\kappa ^{}\lambda ^{}\mu ^{}\nu ^{}}+K_{\kappa \mu }K_{\lambda \nu }K_{\kappa \nu }K_{\lambda \mu }.$$
(58)
Substituting the form of $`K_{\mu \nu }`$ and the bulk AdS space curvature $`R_{\kappa \lambda \mu \nu }=\lambda \left(g_{\kappa \mu }g_{\lambda \nu }g_{\kappa \nu }g_{\lambda \mu }\right)`$ gives the intrinsic scalar curvature
$$R_4=12\left(\rho ^2+\lambda \right).$$
(59)
We have reproduced the result we derived in section 2. The curvature of the brane dS space is such that the square of the Hubble constant (see eq. (30)) is $`\rho ^2+\lambda `$.
One notes that the curvature (59) of the brane is independent of the sign of the brane tension $`\rho `$. What then distinguishes the $`\rho >0`$ case from $`\rho <0`$? Since the brane is a codimension-one boundary in AdS space, we can either take the bulk space-time to be the space “inside” the dS boundary or “outside” the dS boundary. Consider Figure 9, which shows the intersection of two planes with the AdS hyperboloid. Let us focus on the left-hand plane, corresponding to a single dS submanifold. By “inside” we mean that the bulk space-time includes the throat of the AdS hyperboloid. By “outside” we mean that the bulk space-time is one of the two disconnected regions which do not include the throat. For example, the solid circle in figure 9 lies ”inside” the dS boundary, while the open circle lies ”outside” it. These two regions are distinguished by the direction of the normal vector $`t_\mu `$. Furthermore they have opposite extrinsic curvature $`K_{\mu \nu }`$. It is then easy to show that one has the following conditions
$$\begin{array}{c}\text{if }\rho >0\text{ then the spacetime is “inside” the dS boundary}\\ \text{if }\rho <0\text{ then the spacetime is “outside” the dS boundary}.\end{array}$$
(60)
We can now give a geometrical description of the general embedding solution. As we argued above, since the Israel conditions are local we can choose the branes to lie on any pair of dS surfaces. As noted in the previous section, in general, these are described by the intersection of an arbitrary pair of planes with the AdS hyperboloid. This is shown in Figure 9. For generic choices of $`n_1`$ and $`n_2`$ the dS spaces will always intersect transversally. This means there are configurations for all values of the signs of $`\rho _1`$ and $`\rho _2`$. However, if $`n_1`$ and $`n_2`$ are either parallel or anti-parallel this is no longer true. In these cases the branes never intersect. As a result, if the vectors are parallel, one only has a solution with a bulk space-time bounded by a pair of branes if $`\rho _1>\rho _2>0`$ or $`\rho _2>\rho _1>0`$. In the case where they are anti-parallel one requires $`\rho _1>0`$ and $`\rho _2>0`$.
The analysis of sections 3 and 4 allowed us to realize that the solutions obtained in section 2 all described de Sitter surfaces separated by a null vector. The above discussion has shown that time-like as well as space-like separation vectors are also allowed. Furthermore, given that two time-like vectors of equal magnitude are related by a boost, the degrees of freedom describing the general solution are: the brane tensions, the magnitude of the separation vector, and whether it is null, space-like, or time-like. In addition, we note that, except for the special case where the planes describing the dS branes are parallel or anti-parallel, there are in general no conditions on the signs of $`\rho _1`$ and $`\rho _2`$ for a solution to exist. The condition we found in section 2 is an artifact of using a particular coordinate system.
The use of a specific coordinate system to analyze the null-separated solutions was useful to describe the dynamics of the extra dimension as viewed by an observer living on either brane (sec. 3.1) as well as the appearance of horizons in the bulk (sec. 3.2). However, we noted that these questions could also be addressed purely geometrically. The same is true for the general embeddings described here. Of course, if one wished, it is possible to describe the general solutions using some global coordinate system adapted to one of the branes. As mentioned in section 2, one way to proceed would be to start with the same ansatz for the metric as equation (2) but allow for a more general location of the second brane.
We end this section by noting that although discussed in the context of negative bulk cosmological constant $`\lambda `$ and dS branes ($`\rho _i`$ satisfying $`\rho _i^2+\lambda >0`$), the derivation of the embedding conditions is completely general. The construction naturally goes over to cases of positive or zero $`\lambda `$ and arbitrary $`\rho _i`$. For instance, in bulk dS or flat space the intrinsic brane curvature (59) is always positive and we are considering embedded dS branes. In AdS space, the intrinsic curvature can also be zero (the Randall–Sundrum case) or negative. In the latter case we are embedding AdS branes in the AdS bulk. Geometrically these arise from intersecting with planes where $`n`$ is null (for flat branes) or space-like (for AdS branes).
## 6 Discussion
All the solutions presented in this work were obtained by assuming (albeit implicitly) that the bulk was AdS. It was then found that a domain wall of uniform energy density can be embedded in this background provided it follows a de Sitter trajectory, and we described the most general configuration with two such trajectories in AdS. Rather than fixing the bulk geometry for all times, a more general approach is to treat the problem as an initial-value problem. Let us assume that the only stress-energy in the problem is a negative cosmological constant $`\lambda `$ in the bulk and brane tensions $`\rho _i`$ such that $`|\rho _i|>\sqrt{\lambda }`$. Suppose we choose some space-like hypersurface on which we specify the initial spatial bulk metric and the boundary branes. In general, one could imagine complicated initial conditions, but a natural configuration to consider is one with a space-like slice of anti-de Sitter spatial bulk metric and two (spatially) flat surfaces with arbitrary separation, velocity, and orientation. If one were to solve for the time evolution of this system, one would generically find that the bulk does not remain AdS but that a non-vanishing Weyl tensor is generated. This would mean that our solutions correspond to initial conditions which are tuned so that the bulk remains AdS. Coming back to the general problem, it is not clear whether all bulk evolutions would yield homogeneous and isotropic cosmologies on the branes. It would be essential to know what subclass of initial conditions would be consistent with the usual assumptions of cosmology.
A related issue has to do with the stability of the solutions. Suppose our tuned initial conditions are perturbed, then one should investigate whether the path is only slightly or greatly disturbed by the variation. The brane motion could be perturbed either by moving the brane as a whole away from its initial trajectory without altering its energy density, by moving a region of the brane off the trajectory while keeping the energy density constant, or by perturbing the energy density in a spatially homogeneous or inhomogeneous way. In the first case, we know of at least one example of instability, namely solutions with static orbifolds in the case $`\rho _1\rho _2`$. Indeed, we saw in section 3 that the branes will eventually collide if brought infinitesimally closer than the static distance, or end up infinitely far apart if pulled away from each other. Note that, even though the path is unstable, the cosmological evolution remains de Sitter and the bulk is still AdS. Nevertheless, more general homogeneous perturbations may drive the bulk away from AdS and, thus, drive the cosmological evolution on each brane away from dS. As for inhomogeneous energy density perturbations, it was argued in Ref. that any spatial inhomogeneities on the brane stress-energy will modify the AdS bulk through gravitational radiation. Cosmological energy density perturbations in brane-worlds have recently been investigated in Refs. .
To conclude, let us briefly point out that our solutions can be straightforwardly generalized to any type of energy density on the brane (e.g., radiation, matter). Fixing the bulk metric to be AdS for simplicity, one can use the Israel junction condition at the location of the brane to solve for its motion in the bulk , which in turn determines its cosmological evolution. If we want to embed two domain walls, each one should be allowed to travel along any trajectory consistent with its energy density.
Acknowledgments
We would like to thank Burt Ovrut, Alexander Polyakov, Andy Strominger, and Herman Verlinde for helpful discussions. J.K. is supported by the Natural Sciences and Engineering Research Council of Canada. P.J.S. is supported in part by US Department of Energy grant DE-FG02-91ER40671. D.W. would like to thank The Rockefeller University and The University of Chicago for hospitality during the completion of this work.
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# One-loop Shift in Noncommutative Chern-Simons Coupling
## I Introduction
Field theory (especially gauge field theory) on a noncommutative space (or spacetime) has attracted much interest recently . That such theory arises from string/M(atrix) theory suggests that space or spacetime noncommutativity should be a general feature of quantum gravity for generic points deep inside the moduli space of M-theory. Moreover, being a natural deformation of usual quantum field theory, noncommutative field theory is of interests in its own right. A charge in the lowest Landau level in a strong magnetic field can be viewed as living in a noncommutative space, because the guiding-center coordinates of the charge are known not to commute. In this paper, we study noncommutative Chern-Simons (NCCS) field theory in 3 dimensions, which is a deformation of ordinary Chern-Simons (CS) theory and may have applications in planar condensed matter systems, especially in the quantum Hall systems. For simplicity, in this paper we mainly consider pure NCCS theory, with gauge group $`U(N)`$ and with no matter fields coupled to it.
Three-dimensional noncommutative spacetime has coordinates satisfying
$$[x^\mu ,x^\nu ]=i\theta ^{\mu \nu }\mu ,\nu =0,1,2,$$
(1)
where $`\theta ^{\mu \nu }`$ are antisymmetric and real parameters of dimension length squared. The action for a pure $`U(N)`$ CS theory on this space reads
$$I_{CS}=\frac{i\kappa }{4\pi }d^3x\epsilon ^{\mu \nu \lambda }\text{Tr}(A_\mu _\nu A_\lambda +\frac{2}{3}A_\mu A_\nu A_\lambda ).$$
(2)
Here the dynamical field is the gauge potential gauge potential $`A^\mu =A_\mu ^aT^a`$, $`T^a`$ the generators of the gauge group $`G=U(N)`$, normalized to $`\text{Tr}(T^aT^b)=\delta ^{ab}/2`$ with $`T^0=i/\sqrt{2N}`$ for the $`U(1)`$ sector. $`\kappa `$ is the CS coupling, $`\epsilon ^{\mu \nu \lambda }`$ the totally antisymmetric tensor with $`\epsilon ^{012}=1`$. In the action (2) we are using a representation, in which the coordinates $`x^\mu `$ are the same as usual, but the product of any two functions of $`x^\mu `$ is deformed to the Moyal star-product:
$$fg(x)=e^{(i/2)\theta ^{\mu \nu }_\mu ^x_\nu ^y}f(x)g(y)|_{y=x}.$$
(3)
The commutator in Eq. (1) is understood as the Moyal bracket with respect to the star product:
$$[x^\mu ,x^\nu ]x^\mu x^\nu x^\nu x^\mu .$$
(4)
For applications to a system in the lowest Landau level, one considers only the spatial noncommutativity: $`\theta ^{01}=\theta ^{02}=0`$ and $`[x^1,x^2]=i\theta `$.
It is obvious that if $`\theta ^{\mu \nu }=0`$, the action (2) reduces to that of ordinary pure CS theory in 3 dimensions , which is known to be a topological quantum field theory , with the partition function and the correlation functions of Wilson loops being topological invariants, independent of spacetime metric. Diagrammatically the ordinary CS theory is renormalizable . Many topological features can be probed in perturbation theory . One interesting result is the one-loop quantum shift of the non-Abelian CS coupling . However, ordinary pure Abelian CS theory has no such shift, though additional matter coupling does at the two-loop level. In this paper we will show that there is a non-vanishing one-loop shift in noncommutative CS coupling even if the gauge group is $`U(1)`$. This shift turns out to be a constant proportional to the integer $`N`$, independent of the noncommutativity parameters $`\theta ^{\mu \nu }`$, and identical to the one-loop shift in ordinary $`SU(N)`$ CS theory when $`N2`$. Possible physical and mathematical implications of our results will be discussed.
## II Regularized Feynman Rules
The action (2) is invariant under the following infinitesimal gauge transformations:
$$\delta A_\mu =D_\mu \lambda \lambda +[A_\mu ,\lambda ].$$
(5)
To do perturbation theory, we follow the standard procedure of path integral quantization to establish the Feynman rules. The full, regularized action after gauge fixing in Euclidean spacetime reads
$$I_{tot}=I_{CS}+I_{YM}+I_{gf}+I_{gh}.$$
(6)
Here we have added the noncommutative Yang-Mills (YM) term
$$I_{YM}=\frac{1}{2e^2}d^3x\text{Tr}(F_{\mu \nu }F^{\mu \nu }),$$
(7)
with the field strength $`F_{\mu \nu }`$ defined by
$$F_{\mu \nu }=_\mu A_\nu _\nu A_\mu +[A_\mu ,A_\nu ].$$
(8)
This gauge invariant term in the action provides a higher-derivative regularization for the CS theory, since the YM coupling $`e^2`$ is of dimension of mass, which is used as a cut-off that is sent to infinity at the end of calculations. The third term in Eq. (6) is the gauge fixing term:
$$I_{gf}=\frac{1}{\alpha e^2}d^3x\text{Tr}(^\mu A_\mu )^2,$$
(9)
a linear, covariant gauge condition convenient for perturbation theory. In the following we are going to take the Landau gauge $`\alpha =0`$, which was known to have computational advantages in the infrared in ordinary CS theory . The last term $`I_{gh}`$ is the ghost action corresponding to the above gauge fixing:
$$I_{gh}=d^3x\text{Tr}(^\mu \overline{c}D_\mu c),$$
(10)
where $`c`$ and $`\overline{c}`$ are the ghost and anti-ghost respectively.
In ordinary Yang-Mills theory, to write down the Feynman rules, in addition to full action (6), one more thing one needs to know is the representation of the gauge group, since it determines the normalization of the group factors. In the following, we will mainly concentrate on $`U(1)`$ theory, and we will come to $`U(N)`$ case naturally after the explicit calculation for $`U(1)`$ case. In noncommutative $`U(1)`$ gauge theory, though the group factor is trivial, what is nontrivial is the noncommutativity of the kernel in Fourier transform. Suppose we have two kernels of Fourier transform $`_k=e^{ikx}`$ and $`_p=e^{ipx}`$, by using Eq. (3), one can easily check the following commutator
$$[_k,_p]=2i\mathrm{sin}(\frac{\theta ^{ij}}{2}p_ik_j)_{k+p}=2i\mathrm{sin}(\frac{\theta }{2}kp)_{k+p}.$$
(11)
where $`kpk_\mu \theta ^{\mu \nu }p_\nu `$, and we have used Eq. (1). After making Fourier transform of the action, one can immediately find out that this commutator plays exactly the same role as that of Lie commutators of the gauge group in ordinary non-Abelian gauge theory. Therefore, we can establish the Feynman rules by following the same procedure as that of ordinary Yang-Mills theory, with the group structure constants, $`f^{abc}`$, being replaced by a momentum-dependent factor , namely,
$$f^{abc}f^{k,p,k+p}=\sqrt{2}i\mathrm{sin}(\frac{\theta }{2}kp),$$
(12)
where $`\sqrt{2}`$ is due to the normalization of $`T^0`$. With the help of this correspondence, we establish the following Feynman rules for $`U(1)`$ NCCS theory:
(i) The gluon propagator:
$$\mathrm{\Delta }_{\mu \nu }(p)=\frac{4\pi }{\kappa }\frac{m}{p^2(p^2+m^2)}(m\epsilon _{\mu \nu \rho }p^\rho +\delta _{\mu \nu }p^2p_\mu p_\nu ).$$
(13)
where $`m=e^2\kappa /4\pi `$. At the end of computations, we remove the cut-off by taking $`e^2\mathrm{}`$ or $`m\mathrm{}`$.
(ii) The ghost propagator:
$$\frac{1}{p^2}.$$
(14)
(iii) The ghost-ghost-gluon vertex:
$$\sqrt{2}q_\nu \mathrm{sin}\left[\frac{\theta }{2}qp\right].$$
(15)
(iv) The three gluon vertex:
$$\frac{\kappa }{4\pi }\frac{\sqrt{2}}{m}\mathrm{sin}\left[\frac{\theta }{2}pq\right][m\epsilon _{\mu \nu \rho }(rq)_\mu \delta _{\nu \rho }(qp)_\rho \delta _{\mu \nu }(pr)_\nu \delta _{\rho \mu }].$$
(16)
(v)The four-gluon vertex:
$`{\displaystyle \frac{\kappa }{4\pi }}{\displaystyle \frac{1}{m}}[f^{p,q,t}f^{r,s,t}(\delta _{\mu \rho }\delta _{\nu \sigma }\delta _{\mu \sigma }\delta _{\nu \rho })`$ $`+`$ $`f^{r,q,t}f^{p,s,t}(\delta _{\mu \rho }\delta _{\nu \sigma }\delta _{\mu \nu }\delta _{\sigma \rho })`$ (17)
$`+`$ $`f^{s,q,t}f^{r,p,t}(\delta _{\mu \nu }\delta _{\rho \sigma }\delta _{\mu \sigma }\delta _{\nu \rho })].`$ (18)
where $`p,q,r,s,t`$ are incoming momenta and $`f^{x,y,z}`$ is given by Eq. (12).
## III Ward-Slavnov-Taylor Identities
In ordinary gauge theories, Ward-Slavnov-Taylor (WST) identities play a very important role in renormalized perturbation theory. For renormalizable gauge theories, they are essentially manifestation of gauge invariance for the regularized and renormalized action (with counter terms included). Conversely, checking WST identities is essentially checking renormalizability and gauge invariance of the renormalized gauge theory. The same is true for noncommutative gauge theories. In the following we are going to check part of the Ward identities to assure gauge invariance, and to use part of them to simplify the calculations.
Renormalizability of the theory requires that the full inverse $`A`$-propagator and the full $`AAA`$-vertex are of the following form as the external momenta tend to zero
$$\mathrm{\Delta }_{\mu \nu }^1(k)\frac{\kappa }{4\pi }Z_Aϵ_{\mu \nu \lambda }k_\lambda +Z_A^{}(k^2\delta _{\mu \nu }k_\mu k_\nu ),(k0),$$
(19)
$$\mathrm{\Gamma }_{\mu \nu \lambda }(p,q,r)Z_gϵ_{\mu \nu \lambda },(p,q,r0).$$
(20)
These equations defines the relevant renormalization constants $`Z_A`$, $`Z_A^{}`$ and $`Z_g`$. In next section, we will confirm the validity of these equations for one-loop two-point and three-point functions, to verify renormalizability of the theory at the one-loop level. Similarly one can defined $`Z_{gh}`$ and $`\stackrel{~}{Z_g}`$, the renormalization constants<sup>*</sup><sup>*</sup>* Here we would like to remind that all the renormalization constants we defined here are consistent with the conventions used in the refs. , while being the inverse of the standard ones used in many textbooks on quantum field theory. for the ghost wave function and the $`\overline{c}Ac`$-vertex respectively, through the full ghost propagator and the full $`\overline{c}Ac`$-vertex
$$\stackrel{~}{\mathrm{\Delta }}(p)\frac{1}{Z_{gh}p^2},(p0),$$
(21)
$$i\mathrm{\Gamma }_\mu (p,q,r)i\stackrel{~}{Z}_gp_\lambda ,(p,q,r0).$$
(22)
In the next section we will see that at least at one loop, these renormalization constants are in fact finite, i.e. they are independent of the cut-off.
Assuming renormalizability and introducing the renormalized fields and renormalization constants, one can write the renormalized action as
$$I_{ren}=I_{CS}^{}+I_{gf}^{}+I_{gh}^{}$$
(23)
with
$$I_{CS}^{}=\frac{i\kappa _r}{4\pi }d^3x\epsilon ^{\mu \nu \lambda }\text{Tr}[Z_A^1A_\mu ^{(r)}_\nu A_\lambda ^{(r)}+\frac{2}{3}Z_g^1A_\mu ^{(r)}A_\nu ^{(r)}A_\lambda ^{(r)}],$$
(24)
$$I_{gf}^{}=\frac{1}{\alpha _r}d^3x\text{Tr}(^\mu A_\mu ^{(r)})^2,$$
(25)
$$I_{gh}^{}=d^3x\text{Tr}(Z_{gh}^1^\mu \overline{c}^{(r)}_\mu c^{(r)}+\stackrel{~}{Z}_g^1^\mu \overline{c}^{(r)}[A_\mu ^{(r)},c^{(r)}]).$$
(26)
The terms (24), (25) and (26) in the renormalized action (23) should be equal to, respectively, the corresponding terms (2), (9) and (10) in the original action (6), if they are expressed in terms of the bare fields through The relation between $`A_\mu `$ and $`A_\mu ^{(r)}`$, though looks unusual, is appropriate for the CS theory. See e.g. Refs. .
$$A_\mu =\frac{Z_A}{Z_g}A_\mu ^{(r)},c=Z_{gh}^{1/2}c^{(r)},\overline{c}=Z_{gh}^{1/2}\overline{c}^{(r)}.$$
(27)
This requires that the renormalized CS coupling be related to the bare one by
$$\kappa _r=\frac{Z_A^3}{Z_g^2}\kappa ,$$
(28)
and that the following WST identity be true for the renormalization constants:
$$\frac{Z_A}{Z_{gh}}=\frac{Z_g}{\stackrel{~}{Z}_g}.$$
(29)
As easy to check, the WST identities guarantee that the renormalized actions (24) and (26) are gauge invariant.
## IV One-loop renormaliztion in $`U(1)`$ theory
In this section, we first study at the one-loop level the renormalization of $`U(1)`$ NCCS theory, in particular the shift in the Chern-Simons coupling $`\kappa `$.
Let us start with the ghost self-energy. It contains a planar diagram contribution
$$\stackrel{~}{\mathrm{\Pi }}_p^{(1)}=\frac{4\pi }{\kappa }\frac{m}{p^2}\frac{d^3k}{(2\pi )^3}\frac{k^2p^2(kp)^2}{k^2(k^2+m^2)(k+p)^2},$$
(30)
and a non-planar diagram contribution
$$\stackrel{~}{\mathrm{\Pi }}_{np}^{(1)}=\frac{4\pi }{\kappa }\frac{m}{p^2}\frac{d^3k}{(2\pi )^3}\frac{k^2p^2(kp)^2}{k^2(k^2+m^2)(k+p)^2}e^{i\theta kp}.$$
(31)
The integral (30) is finite, with the leading term proportional to $`1/|m|`$. Therefore, in the large $`|m|`$ limit, we get the contribution to the ghost self-energy as
$$\stackrel{~}{\mathrm{\Pi }}_p^{(1)}=\frac{2}{3\kappa }\text{sgn}(\kappa ).$$
(32)
On the other hand, the non-planar diagram contribution is evaluated to be
$$\stackrel{~}{\mathrm{\Pi }}_{np}^{(1)}=\frac{2}{3\kappa }\text{sgn}(\kappa )f(p\theta |m|),$$
(33)
where the function $`f(x)`$ is defined by
$$f(x)=\frac{1}{x}_0^x𝑑y(3y^{1/2}y^{3/2})K_{1/2}(y)=\sqrt{2\pi }\frac{(1+x)e^x1}{2x}.$$
(34)
Since $`f(x)0`$ as $`x\mathrm{}`$, the non-planar diagram does not contribute to the ghost self-energy in the limit $`|m|\mathrm{}`$. Therefore, the ghost self-energy correction at the one loop level is finite and independent of external momentum $`p`$. Correspondingly, we get the one-loop ghost wave function renormalization constant
$$Z_{gh}^{(1)}=1\frac{2}{3\kappa }\text{sgn}(\kappa ).$$
(35)
To calculate the gluon self energy $`\mathrm{\Pi }_{\mu \nu }^{(1)}(p)`$, we decompose it into the following structure
$$\mathrm{\Pi }_{\mu \nu }^{(1)}(p)=\frac{1}{m}\mathrm{\Pi }_e^{(1)}(\delta _{\mu \nu }p^2p_\mu p_\nu )+\frac{\kappa }{4\pi }\mathrm{\Pi }_o^{(1)}(p)\epsilon _{\mu \nu \lambda }p^\lambda $$
(36)
Since only the gluon loop diagram has an odd number of $`\epsilon `$ tensors, $`\mathrm{\Pi }_o^{(1)}(p)`$ receives a nonzero contribution only from the gluon loop diagram Fig. 2a. In contrast, $`\mathrm{\Pi }_e^{(1)}`$ picks up contributions from the tadpole diagram (Fig. 2c), the ghost loop diagram Fig. 2b, as well as from the gluon diagram Fig. 2a with an even number of $`\epsilon `$ tensors. Contracting $`\mathrm{\Pi }_{\mu \nu }^{(1)}`$ with $`\kappa /4\pi (\epsilon _{\mu \nu \lambda }p^\lambda /2p^2)`$ and $`m\delta _{\mu \nu }/2p^2`$, we obtain $`\mathrm{\Pi }_o^{(1)}(p)`$ and $`\mathrm{\Pi }_e^{(1)}(p)`$:
$$\mathrm{\Pi }_o^{(1)}=\frac{4\pi }{\kappa }\frac{2m}{p^2}\frac{d^3k}{(2\pi )^3}[\mathrm{sin}^2(\theta /2kp)]\frac{[k^2p^2(kp)^2][5k^2+5(kp)+4p^2+2m^2]}{k^2(k+p)^2(k^2+m^2)[(k+p)^2+m^2]},$$
(37)
and
$$\mathrm{\Pi }_e^{(1)}=\frac{m}{2p^2}\left[\frac{d^3k}{(2\pi )^3}[\mathrm{sin}^2(\theta /2kp)]\frac{N_e(p,k)}{k^2(k+p)^2(k^2+m^2)[(k+p)^2+m^2]}+\frac{5m}{3\pi }\right].$$
(38)
where
$`N_e(p,k)=6k^6+18k^4(kp)+20k^4p^2+22k^2(kp)^2p^212(kp)^3`$ (39)
$`+9k^2p^47(kp)^2p^2+m^2[2k^4+4k^2(kp)+k^2p^2+(kp)^2].`$ (40)
At this point we would like to comment that the structures shown in the above results are similar to those in ordinary non-Abelian Chern-Simons gauge theory in $`2+1`$ dimensions, although here we are dealing with the $`U(1)`$ case. Still they are different in the following two aspects. The first is that we have non-planar diagram contributions due to the oscillating factor $`4\mathrm{sin}^2(\theta /2kp)`$. The second is that the value of the tadpole contribution changes (see the second term in (38)), the reason being that one of the terms in the four-gluon vertex vanishes due to the fact $`\mathrm{sin}(\theta /2pp)=0`$ by using the Feynman rule (17).
The integral in Eq. (37) is finite. To calculate it, again we separate it into planar and non-planar contributions. The calculation of the planar contribution is standard. Taking $`|m|\mathrm{}`$, we obtain
$$\mathrm{\Pi }_{o,p}^{(1)}=\frac{7}{3\kappa }\text{sgn}(\kappa ).$$
(41)
By using Feynman parameterization, we can rewrite the corresponding non-planar contribution as follows:
$$\mathrm{\Pi }_{o,np}^{(1)}=2p^2_0^1𝑑x_0^x𝑑y_0^y𝑑z\{5I_2(\nu )+[5p^2(1+yxz)(yxz)+4p^2+2m^2]I_1(\nu )\},$$
(42)
where the argument $`\nu `$ in the functions $`I_1`$ and $`I_2`$ is defined as
$$\nu =\sqrt{p^2(1+yxz)+m^2(1y)}.$$
(43)
The functions $`I_1`$ and $`I_2`$ are defined as
$$I_1=\frac{\pi ^{3/2}}{2}\left(\frac{\theta p}{2}\right)^{3/2}[\nu ^{3/2}K_{3/2}(\nu \theta p)\frac{\theta p}{3}\nu ^{1/2}K_{1/2}(\nu \theta p)],$$
(44)
and
$$I_2=\frac{\pi ^{3/2}}{12}\left(\frac{\theta p}{2}\right)^{1/2}[15\nu ^{1/2}K_{1/2}(\nu \theta p)20\left(\frac{\theta p}{2}\right)^{3/2}\nu ^{1/2}K_{1/2}(\nu \theta p)+4\left(\frac{\theta p}{2}\right)^{5/2}\nu ^{3/2}K_{3/2}(\nu \theta p)],$$
(45)
where $`K_\beta (x)`$ is the modified Bessel function. It has an exponentially decay profile. In the limit $`|m|\mathrm{}`$, we see that the integrand in Eq. (42) vanishes, therefore the non-planar diagram does not contribute to $`\mathrm{\Pi }_o^{(1)}`$. Namely,
$$\mathrm{\Pi }_{o,np}^{(1)}=0.$$
(46)
Thus, we get
$$\mathrm{\Pi }_o^{(1)}=\frac{7}{3\kappa }\text{sgn}(\kappa ).$$
(47)
Similar analysis can be applied to the integral (38). It turns out that the integral is finite as $`|m|\mathrm{}`$. Therefore, the photon wave function renormalization constant is
$$Z_A^{(1)}=1+\mathrm{\Pi }_o^{(1)}=1+\frac{7}{3\kappa }\text{sgn}(\kappa ).$$
(48)
Furthermore, we study the one loop corrections to the vertex $`\overline{c}Ac`$, we show that the one loop correction vanishes as $`|m|\mathrm{}`$. The reason is that the one-$`\epsilon `$ term of Fig. 2(i) cancels against the three-$`\epsilon `$ term of Fig. 2(h); the one-$`\epsilon `$ term of Fig. 2(h) goes to zero; the non-$`\epsilon `$ terms in the two diagrams cancel each other as well. Finally, we get the renormalization for the $`\overline{c}Ac`$ vertex
$$\stackrel{~}{Z}_g^{(1)}=1.$$
(49)
After we extract $`Z_{gh}^{(1)}`$, $`Z_A^{(1)}`$, and $`\stackrel{~}{Z}_g^{(1)}`$, we can employ the WST identity (29) established in the previous section to get the three-gluon vertex renormalization constant:
$$Z_g^{(1)}=\frac{Z_A^{(1)}}{Z_{gh}^{(1)}}\stackrel{~}{Z}_g^{(1)}=1+\frac{3}{\kappa }\text{sgn}(\kappa ),$$
(50)
where we have worked up to the first order in $`1/\kappa `$, consistent with one-loop perturbative theory.
Now we have shown that all renormalization constants at the one loop level are finite, so the one-loop beta function vanishes, as in ordinary CS theory. Substituting the renormalization constants $`Z_A^{(1)}`$ and $`Z_g^{(1)}`$ in the definition of the renormalized CS coupling $`\kappa _r^{(1)}`$, we have
$$\kappa _r^{(1)}=\kappa +\text{sgn}(\kappa ).$$
(51)
This is the main result of the present paper. The second term is the desired one-loop shift in the $`U(1)`$ Chern-Simons coupling. Note that the shift is just the unity in our normalization for the coupling. Note that it is independent both of the noncommutativity parameters $`\theta ^{\mu \nu }`$ and of the value of the bare coupling $`\kappa `$ except for its sign. Also recall that for ordinary $`U(1)`$ CS theory, the one-loop shift vanishes in the same $`F^2`$ regularization.
## V Generalization to $`U(N)`$
In this section, we generalize the expression (51) we obtained in last section for the one-loop shift of the CS coupling from $`U(1)`$ to $`U(N)`$.
It is known that for a noncommutative gauge theory, the gauge group is restricted to be only $`U(N)`$; even $`SU(N)`$ is not allowed, because the closure of the Moyal commutator is violated . Another way to see this is that the three-gluon coupling in the action (2) mixes the $`U(1)`$ gluon with the $`SU(N)`$ gluons. Thus, unlike ordinary gauge theory which allows the $`U(1)`$ and $`SU(N)`$ sectors to have independent coupling constants, in noncommutative theory $`U(N)`$ gauge invariance enforces the $`U(1)`$ and the $`SU(N)`$ gluons to share the same coupling constant.
For $`U(N)`$ noncommutative Yang-Mills (NCYM) theory, it is this distinct feature that makes the coupling in the $`U(1)`$ sector runs in the same way as that in the $`SU(N)`$ sector, as verified in a recent explicit calculation . Before this calculation was done, a beautiful proof without doing any new calculation had been given in ref. for the statement that the one-loop beta-function of $`U(N)`$ NCYM can be simply read off from the known value of the ordinary $`SU(N)`$ Yang-Mills theory. This is because in NCYM the nonplanar one-loop $`U(N)`$ diagram contributes only to the $`U(1)`$ part of the theory. In the following we will apply the same trick to NCCS, and derive the one-loop shift in the CS coupling in the $`U(N)`$ theory without doing any new calculations.
Following ref. , let us consider the quadratic one-loop 1PI effective action of the ordinary $`U(N)`$ CS theory, in which the $`U(1)`$ and $`SU(N)`$ sectors share the same coupling constant, in the $`F^2`$ regularization. From the results in refs. one infers that after removing the cut-off,
$$\mathrm{\Gamma }_2^{(1)}(\theta =0)=\frac{i}{4\pi }d^3x\epsilon ^{\mu \nu \lambda }[(\kappa +N\text{sgn}(\kappa ))\text{Tr}(A_\mu _\nu A_\lambda )\text{sgn}(\kappa )(\text{Tr}A_\mu )_\nu (\text{Tr}A_\lambda )].$$
(52)
The coefficient $`\kappa +N\text{sgn}(\kappa )`$ in the first term is read off from the known one-loop shift in ordinary $`SU(N)`$ CS theory , while the existence of the second term is due to the necessity for cancelling the $`U(1)`$ part in the first term, since we know there is no one-loop shift in the $`U(1)`$ coupling constant. Also it is easy to check that the second term is the only nonplanar contribution to $`\mathrm{\Gamma }_2^{(1)}`$, coming from the nonplanar part of diagrams like Fig. 2 (a)-(c).
Now let us turn on nonzero $`\theta _{\mu \nu }`$. The planar contributions to $`\mathrm{\Gamma }_2^{(1)}`$ are known to be the same as in ordinary theory , while the nonplanar diagrams are suppressed to zero by an extra rapidly oscillating phase factor. Our result (46) in last section verifies the latter by explicit computation. Thus, with the second term put to zero, one reads from eq. (52) that in $`U(N)`$ NCCS,
$$\mathrm{\Gamma }_2^{(1)}(\theta 0)=\frac{i}{4\pi }d^3x\epsilon ^{\mu \nu \lambda }[\kappa +N\text{sgn}(\kappa )]\text{Tr}(A_\mu _\nu A_\lambda ).$$
(53)
Therefore, without any new calculation, we infer that the one-loop shift in the coupling for $`U(N)`$ NCCS is
$`\kappa _r`$ $`=`$ $`\kappa +N\text{sgn}(\kappa ).`$ (54)
For the $`U(1)`$ case, we plug $`N=1`$ in Eq. (54), reproducing exactly Eq. (51) that we have obtained by explicit calculation in last section.
In summary, the one-loop shift of the CS coupling in $`U(N)`$ NCCS is the same as that in ordinary $`SU(N)`$ CS theory for $`N2`$, while for the $`U(1)`$ case it gives a non-vanishing value in contrast to ordinary CS theory.
## VI Conclusions and Discussions
The induced CS coupling by fermionic fields in noncommutative quantum electrodynamics in 3 dimensions has been studied in ref. . In this paper we have considered instead a pure CS theory without matter, and have studied the one-loop quantum correction to the coupling constant due to self-interactions of the gauge bosons that arise from spacetime coordinate noncommutativity.
First of all, the renormalization constants we have calculated at the one loop level are finite, showing that the beta function at this level vanishes. Moreover all renormalization constants, including the one-loop shift in the CS coupling are shown to be independent of the noncommutativity parameters. This is a bit surprising, since the spacetime noncommutativity parameters appear in the Lagrangian of pure CS theory explicitly. This adds explicit evidence to a theorem proved in Ref. that one-loop results in noncommutative CS theory are all independent of spacetime noncommutativity. It would be interesting to see whether the same independence remains true at higher orders. Our conjectured answer is ”yes”. In other words, we conjecture that noncommutative CS theory is a “deformed” topological field theory, in the sense that the partition function and correlation functions are topological invariants, independent of both metric and noncommutativity parameters.
Furthermore, we have shown that the one-loop shift in the coupling constant does not vanish in a pure $`U(1)`$ NCCS theory. This arises as a consequence of spacetime coordinate noncommutativity, since in ordinary $`U(1)`$ CS theory there is no quantum shift in the coupling constant at all. We notice two features of our result (51): 1) it is independent of the bare coupling $`\kappa `$, except for its sign; 2) it is the simplest integer, the unity, independent of the noncommutatvity parameters. Therefore this result is not smooth in the limit $`\theta _{\mu \nu }0`$.
Finally, we have shown that the one-loop shift of the $`U(N)`$ NCCS coupling is the integer $`N`$ that characterizes the gauge group $`U(N)`$, exactly the same as that in ordinary $`SU(N)`$ CS theory for $`N2`$. In ordinary CS theory the quantization of the shift was interpreted as being consistent with the topological quantization of the non-abelian CS coupling. The latter is known to result from the topological fact that a large gauge transformation changes the CS action by a value proportional to the integer winding number of the gauge transformation, viewed as a map from the (compactified) spacetime to the gauge group. Whether the topological quantization of the CS coupling remains true in the noncommutative theory is not clear at all at this moment. However, our result (54) shows that the one-loop shift is still quantized in the noncommutative case. If one reverses the logic in the above reasoning for the ordinary CS theory, this result seems to indicate that possibly in noncommutative geometry there should be a counterpart of the concept of the usual winding number that remains integer-valued, and that the NCCS coupling should satisfy a similar topological quantization.
Here we would like to mention that the one-loop shift in the $`U(1)`$ NCCS coupling depends on the regularization used. Our result (51) was obtained in the $`F^2`$ regularization, in which the Yang-Mills term was added to the action and its coefficient was taken as cut-off. If we had used dimensional regularization, the shift would be zero. This situation is not surprising, completely similar to the well-known situations for ordinary non-Abelian CS theory; see e.g. Refs. . In our opinion, the $`F^2`$ regularization is more physical, in the sense that the $`F^2`$ term may naturally appear in realistic planar systems.
In field theory on ordinary spacetime, the $`U(1)`$ CS coupling (CS coefficient) is known to have several interesting physical meanings, when the CS gauge field couples to various fields. For example, when there is a Maxwell term in the action, $`\kappa `$ is related to the topological mass of the CS photon . When there are matter fields coupled to the CS field, the CS coefficient $`\kappa `$ will give rise to fractional (exchange) statistics for matter field quanta . Finally a well-known folklore in the community is that an effective CS coupling for the electromagnetic field indicates the Hall effect, with the effective CS coefficient directly related to the Hall conductance. We expect all these physical interpretations should be generalizable to nocommutative spacetime. A systematic inverstigation of NCCS coupled to matter fields (both fermionic and bosonic) will be published elsewhere .
###### Acknowledgements.
One of us, GHC, thanks the Institute for Theoretical Physics, University of California at Santa Barbara, for an ITP Graduate Fellowship, and for the warm hospitality he receives during his stay. This research was supported in part by the National Science Foundation under Grants No. PHY-9407194(GHC) and PHY-9907701(YSW). Notice added: After the paper was put on the e-print archive, we are informed by e-mail from Soo-jong Rey that he has obtained the same result (in agreement with ours) on $`U(1)`$ NCCS theory.
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# A potential inversion study of liquid CuBr
## I Introduction
Copper-(I)-bromide, like several other copper and silver halides, attains a rather large ionic conductivity before melting, and belongs to a class of ionic systems called fast-ion conductors. A representative of this class of substances, the CuCl, belongs to the first substances investigated by the neutron diffraction with isotopic substitution (NDIS) technique . In this early work, a featureless CuCu partial pair distribution function (PPDF) $`g_{CuCu}`$ has been found, and this structural detail attracted much attention - including a redetermination and confirmation of the PPDF – as it provides an easy explanation for the fast-ion conductivity: The small copper ions move more or less gas-like through the network of bromine atoms. According to a hypotheses of Ginoza et al a structureless cation-cation PPDF should be common to all fast (cat)ion conductors.
The structure of liquid CuBr has been investigated by Allen et al. by NDIS making use of the different neutron coherent scattering lengths of <sup>65</sup>Cu(6.72 $``$ 10<sup>-15</sup> m) and <sup>63</sup>Cu(11.09 $``$ 10<sup>-15</sup> m) to determine the PPDF. Again, a flat, gas-like $`g_{CuCu}`$ was found in this study. The structure determination by Allen et al. has been challenged by DiCicco et al. , especially concerning the $`g_{CuBr}`$ PPDF. This group used extended x-ray absorption fine structure (EXAFS) spectroscopy at both the copper and the bromine edge to investigate the short range order in CuBr(l), and found the first shell bromine coordination around copper significantly sharper than determined with neutron diffraction. This is taken by DiCicco et al. as an indication that apart from inter-ionic repulsion there is a significant contribution of covalency to the inter-atomic potential. Further, reverse Monte Carlo (RMC) simulations gave evidence that the flat $`g_{CuCu}`$ is likely to be an artifact of the maximum entropy analysis used in ref. .
This paper will address the question about the sharpness of the first CuBr coordination shell. A set of partial structure factors (PSF) is determined consistent with all available diffraction information. Using the potential inversion scheme of Levesque, Weis, Reatto an empirical two body potential is developed. This allows to check the existence of a three-dimensional arrangement of atoms consistent with this set of PSF and to interprete the one-dimensional pair distribution functions in terms of three-dimensional structures. The structure of both CuBr and ZnCl<sub>2</sub> melt is characterized by tetrahedral structural units - Zn(II) is isoelectronic with Cu(I). The angular correlation between neighboring atoms is hence discussed comparing these two systems.
## II Experimental
Liquid diffraction experiments have been carried out at the high energy beam-line BW5 at the DORIS-III storage ring at HASYLAB in Hamburg, in the set-up for liquid and amorphous substances . The photon energy was set to 121.0 keV. The covered momentum transfer range was 0.43 Å$`{}_{}{}^{1}<Q<`$ 26.5 Å<sup>-1</sup>, with $`Q=4\pi /\lambda \mathrm{sin}(\theta )`$, $`\lambda `$ is the wave-length and $`\theta `$ half the scattering angle.
The sample (Sigma Chemical Co., quoted purity $`>`$99.8%) was sealed under vacuum in a quartz-glass tube of 5 mm inner diameter. The sample was heated to 803 K in a standard neutron diffraction furnace . This furnace – capable to reach a maximum temperature of 2100K – contains several niobium heat shields the inner most with a radius of 50 mm. As in ref. a restrictive collimation has been used to reduce the background scattering contribution of these niobium shields. Fig. 1 shows the scattering contribution of background and sample. The main contribution to the background is the scattering from the silica container. At low momentum transfers, some powder lines from the niobium shields are remaining in spite of the collimation system. The data points at the positions of these niobium lines are excluded from the further analysis. The usual corrections for detector dead-time, absorption and polarization are made and the intensities are converted into a differential cross section normalizing to the known cross-section (the sum of elastic self scattering and Compton cross section) at large momentum transfers. The data correction procedure is described in more detail in .
## III Theoretical summary
The differential cross section of a liquid in a neutron or electromagnetic radiation scattering experiment can be expressed in terms of a total scattering function:
$`S^{(n)}(Q)`$ $`=`$ $`{\displaystyle \frac{\left(\frac{d\sigma }{d\mathrm{\Omega }}\right)^{(n)}_i^{N_{uc}}\nu _ib_i^2}{(_i^{N_{uc}}\nu _ib_i)^2}}+1`$ (1)
$`S^{(x)}(Q)=i(Q)+1`$ $`=`$ $`{\displaystyle \frac{\left(\frac{d\sigma }{d\mathrm{\Omega }}\right)^{(x)}/\sigma _{el}_i^{N_{uc}}\nu _if_i^2}{(_i^{N_{uc}}\nu _if_i)^2}}+1`$ (2)
where $`\left(\frac{d\sigma }{d\mathrm{\Omega }}\right)`$ is the coherent differential cross section, $`Q=4\pi /\lambda \mathrm{sin}(\theta )`$,$`\lambda `$ the wavelength of the radiation and $`\theta `$ the diffraction angle, $`b_i`$ the coherent scattering lengths , $`f_i`$ the X-ray form factors in the independent atom approximation , $`\sigma _{el}`$ the scattering cross section of the free electron, $`\nu _i`$ the stoichiometric coefficient of the atom $`i`$, and where the sums are extending over the number of distinct atoms $`N_{uc}`$ in the unit of composition, CuBr. These total structure functions are composed of partial structure factors $`s_{ij}`$:
$$S^{(n/x)}=\underset{ij}{}w_{ij}(Q)s_{ij}$$
(3)
where
$$w_{ij}(Q)=\frac{\nu _i\nu _jf_i(Q)f_j(Q)}{(\nu _if_i(Q)^2}$$
(4)
and
$$\underset{ij}{}w_{ij}(Q)=1.$$
(5)
For the neutron case $`f(Q)`$ has to be replaced by $`b`$ and the weighting factors become independent of $`Q`$.
The total pair distribution functions are related to the structure factors of Eqs. 1 and 2 via:
$$r(^{(n/x)}g1)=\frac{1}{2\pi ^2\rho _{uc}}Q(^{(n/x)}S(Q)1)\mathrm{sin}(Qr)dQ$$
(6)
with $`\rho _{uc}`$ the density per unit of composition. In this work a value of $`\rho _{uc}=0.0175`$ Å<sup>-1</sup> for CuBr at 800 K has been used to keep consistency with the work of Allen and Howe , although an experimental value of $`\rho _{uc}=0.0181`$ Å<sup>-1</sup> has been reported later by Saito et al . Likewise, the partial structure factors are related to the partial pair distribution functions (PPDF) via :
$$r(g_{ij}1)=\frac{1}{2\pi ^2\rho _{uc}}Q(s_{ij}1)\mathrm{sin}(Qr)𝑑Q$$
(7)
and, hence, the total pair distribution:
$${}_{}{}^{n/x}g(r)=\underset{ij}{}\text{FT}[w_{ij}(Q)]g_{ij}$$
(8)
where FT is the Fourier sine transformation and $``$ the convolution operation. The convolution reduces to a simple product in the case of neutron scattering.
An alternative definition of the neutron total structure factor has been used by Allen and Howe :
$$F(Q)=c_a^2b_a^2[S_{aa}(Q)1]+2c_ac_bb_ab_b[S_{ab}(Q)1]+c_b^2b_b^2[S_{bb}(Q)1]$$
(9)
where $`c_a`$, $`c_b`$ are the concentration of the atomic species, the $`S_{aa}`$, $`S_{bb}`$, $`S_{bb}`$ are partial structure factors differing from the $`s_{ij}`$ in Eq. 3 by a factor of $`\nu `$. It is noted, that in this definition the total neutron structure factor has the dimension of a cross section, while $`S^{(n)}(Q)`$ in Eq. 1 is dimensionless.
In order to generate three dimensional structures from the pair distribution functions, the potential inversion scheme of Levesque, Weiss and Reatto (LWR-scheme) has been applied. The idea of this method is based on the equation:
$$g(r)=\mathrm{exp}\left[\frac{v(r)}{kT}+g(r)1c(r)+B(r,v)\right]$$
(10)
relating the pair distribution function and the pair potential, where $`v(r)`$ is the pair potential, $`c(r)`$ the direct correlation function and B(r,v) the bridge function. Starting with a crude guess of the bridge function, e.g. neglecting $`B(r,v)`$ completely in the hypernetted chain approximation, a first guess of the potential $`v^{(1)}`$ can be calculated. With a Monte Carlo (MC) or Molecular Dynamics (MD) simulation the pair correlation function $`g(r)^{(1)}`$ and the direct correlation function $`c(r)^{(1)}`$ for a system of particles interacting via $`v^{(1)}`$ can be determined. Thus, the bridge function $`B(r,v^{(1)})`$ is determined exactly via Eq. 10. $`B(r,v^{(1)})`$ is assumed to be a better approximation for $`B(r,v)`$ than the complete neglect in the first approximation. Inserting $`B(r,v^{(1)})`$ in equation 10 gives the LWR iteration formula:
$`[v^{(n)}v^{(n1)}]/kT=\mathrm{ln}(g^{(n1)}/g^{(exp)})`$ (11)
$`+c^{(n1)}c^{(exp)}g^{(n1)}+g^{(exp)}`$ (12)
The empirical potential Monte-Carlo (EPMC) scheme is closely related to equation 12, but considers the upper line of Eq. 12, the logarithmic term only. The complete form of Eq. 12 has been preferred here, as this scheme shows faster convergence . The applicability of the LWR scheme to polyatomic systems has been shown by Kahl et al . An example of the application of this technique has been given recently in .
## IV Deduction of the PPDF
In Fig. 2 the total pair distribution function and the total structure factor determined from the hard X-ray experiment are shown. The contribution of $`g_{CuCu}`$ and $`g_{BrBr}`$ to $`{}_{}{}^{x}g(r)`$ $`{}_{}{}^{(CuBr,x)}g(r)`$ is also shown in Fig. 2. $`{}_{}{}^{(CuBr,x)}g(r)`$ is calculated according to Eq. 8 as the sum of the PPDF convoluted with $`\mathrm{FT}(w_{ij}(Q))`$. It can be seen, that the first peak in r-space is entirely due to the CuBr-PPDF and, hence, the height of this peak is independent of possible errors in the estimation of the BrBr- and CuCu- PPDF. The corresponding CuBr-PSF is shown in the lower part of Fig. 2. The height of the peak in $`g_{CuBr}`$ (Fig. 2) is 5.2, clearly higher than than the 3.7 which can be read from Fig. 4 of ref. .
Starting point for the determination of the CuCu- and the BrBr-PPDF are the total neutron pair distribution functions for <sup>63</sup>CuBr, <sup>63</sup>Cu$`{}_{}{}^{65}{}_{0.5}{}^{}`$Cu<sub>0.5</sub>Br and <sup>65</sup>CuBr. The CuCu- and the BrBr-PPDF used in the preceding paragraph are the result of a direct inversion of the equation system 8. Two additional constraints to the PPDF have been applied, one in Q-space and one in r-space: The CuCu- and BrBr-PPDF show only broad features compared to the CuBr-PPDF, the corresponding PSF oscillate only up to a momentum transfer $`Q10`$ Å<sup>-1</sup> and are set to zero for $`Q>10`$ Å<sup>-1</sup>. A minimum allowed distance has been set in r-space to be 2.2 Å for the CuCu-PPDF and 3.1 Å for the BrBr-PPDF. The CuCu- and BrBr-PPDF determined in this way is more structured than the maximum entropy (ME) solution given by Allen et al , as it is to be expected. Using this ME-solution could only reduce the structure in the contribution of these PPDF to the total X-ray pair distribution function in Fig. 2.
Fig. 3 compares the shape of the first CuBr coordination shell determined with EXAFS , neutron diffraction and in this work. It is evident, that the ME solution by Allen et al underestimates the sharpness of the CuBr peak. DiCicco et al , who were the first to deduce a sharper CuBr peak, forwarded the view that this improved result was due to ”the exceptional short-range sensitivity” of the EXAFS experiment. It is in fact the result of the use of additional information. Two EXAFS experiments at the bromine and copper edge are in principle incapable to extract three PPDF. Only with the assignment of the first peak in real space to the CuBr-PPDF and the assumption of a Gaussian distribution of the first neighbor shell the problem becomes tractable. This assumption is of course very reasonable for a system, which melts from a fast ion conductor. Given that, it would be really astonishing, if the like atoms would be the next neighbors in this system.
Introducing the constraints mentioned in the preceding paragraph is largely equivalent to the assumption of next neighborhood of unlike atoms. Thus diffraction gives the information about a sharp CuBr-coordination and, in fact the information about the entire pair distribution function. The very good agreement of the CuBr-PPDF derived form neutron and X-ray diffraction is noted (Fig. 3), if the same constraints are imposed to the analysis of the neutron data. The PPDF determined by the maximum entropy technique are not the real PPDF, nor do they pretend to be the real PPDF. These PPDF have to be interpreted as a lower limit estimation of the structure based on no other information but the neutron diffraction data. They are, in a sense, a prove that the unlike atoms are next neighbors.
Fig. 4 shows a close-up of the first CuBr coordination shell. The peak is strongly asymmetric and clearly cannot be fitted with a single Gaussian distribution. It has been fitted with two Gaussians. The coordination number integrated over both peaks is four.
## V Simulation
In order to interpret the PPDF in terms of a three dimensional structure the Levesque, Weis, Reatto potential inversion scheme has been applied to the PPDF shown in Fig. 5. It is evident, that the converged LWR-potentials produce a three dimensional structure in agreement with these PPDF: According to Pusztai and McGreevy the same is not true for the maximum entropy PPDF of Allen and Howe. Contrary to the hypothesis of Ginoza et al, that the cation-cation PPDF is structureless in all melts of fast cation conductors, the CuCu PPDF is found similarly structured as the BrBr-PPDF. This finding is in agreement with the RMC study . The agreement of the RMC-PPDF with the PPDF determined here is good for the first CuBr coordination shell and the long distance correlations. Some discrepancies appear at intermediate distances ($`36`$Å). Three facts are certain about the PPDF determined here: They are consistent with all available diffraction information including the new X-ray result. The set of PPDF can be mapped to a three-dimensional arrangement of atoms of the correct density. An effective two-body potential exists with the property, that an ensemble of particles interacting via this potential samples the configuration space consistent with these PPDF.
Some of the angular correlations of near neighbor atoms have been determined from a MC simulation with the converged LWR-potential. These are compared in Fig. 6 to the same functions obtained by RMC . Although there are some differences in the underlying PPDF (cf. Fig. 5), the angular correlations are found to be similar. An interesting detail is perhaps, that the maximum in the Br-Cu-Br distribution function is shifted form $`\mathrm{cos}(\mathrm{}_{\mathrm{BrCuBr}})1/3`$ to $`\mathrm{cos}(\mathrm{}_{\mathrm{BrCuBr}})0`$, but it is also evident that quoting the maximum of the distribution alone would be misleading.
Fig. 7 compares the angular correlations between neighboring atoms in ZnCl<sub>2</sub> and CuBr, both obtained using the LWR formalism. Both melts are based on a tetrahedral structural motive and in fact Cu(I) and Zn(II) are isoelectronic species. So some of the distribution functions, the $`+`$ and the $``$ distribution function, show similarities. However, there is a pronounced difference in the distribution function of the $`++`$ angle, the angle joining adjacent tetrahedra. While in ZnCl<sub>2</sub> there is a peak at $`\mathrm{cos}(\mathrm{}_{++})1./3.`$ the corresponding angular correlations in CuBr are very weak. As a consequence also the $`+++`$ correlation is much less structured in CuBr than in ZnCl<sub>2</sub>. Thus a continuous random network is a useful parameterization for the ZnCl<sub>2</sub>, but not the CuBr structure. The fact, that ZnCl<sub>2</sub> but not CuBr can be easily supercooled into a glass is probably intimately related to this difference.
## VI Conclusion
The short range order of the melt of the fast-ion conductor Cu(I)Br has been discussed recently. By diffraction experiments with hard electromagnetic radiation it can be conclusively shown, that the first CuBr coordination shell is definitely sharper than obtained earlier in a maximum entropy analysis of NDIS data by Allen et al. This is in agreement with the result obtained by DiCicco et al using EXAFS spectroscopy. The analysis presented in this work indicates, that the featureless CuCu-PPDF obtained by Allen et al is an artifact of the maximum entropy analysis, contrary to the view that a structureless cation-cation PPDF is a common feature in all fast-ion conductors. The same result has been obtained by reverse Monte Carlo by Pusztai et al.
The LWR potential inversion scheme has been used to derive an effective two-body potential from the PPDF. This method can be used to interpret the PPDF in terms of a three dimensional structure. A MC simulation with the converged potential has been used to determine angular correlations between near neighbor atoms. The angular correlations in ZnCl<sub>2</sub> melt shows similarities in the $`+`$ distribution functions reflecting the similarity in the local environment of the cation. However, while in ZnCl<sub>2</sub> the angle between adjacent ZnCl<sub>4</sub> tetrahedra is — like in silica glass — well defined and the $`++`$ angle distribution function is peaked, the corresponding function in CuBr is, as the $`+++`$ function, broad and featureless. This behavior is probably the key to understand, why ZnCl<sub>2</sub> but not CuBr can be easily supercooled into a glassy state.
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# Blow-up for solutions of hyperbolic PDE and spacetime singularities.
## 1 Introduction
Anyone who has worked with nonlinear hyperbolic equations is familiar with the phenomenon that smooth solutions often cease to exist after a finite time. There are two common ways in which this happens. In the book of Alinhac these have been called the ‘ODE blow-up mechanism’ and the ‘geometric blow-up mechanism’. In the first case the solution itself is unbounded as the blow-up time is approached. In the second case the solution itself remains bounded while the first derivative is unbounded. Although the Einstein equations, which are the subject of the following, are arguably the most geometric of hyperbolic differential equations the blow-up which is typically encountered occurs by the ODE mechanism in the terminology of . At the same time, the question of the sense in which the Einstein equations can be called hyperbolic is delicate. These points are related, as will now be explained
The basic unknown in the Einstein equations is a pseudoriemannian metric $`g`$. The equations are diffeomorphism invariant, in the sense that if a metric $`g`$ on a manifold $`M`$ satisfies the Einstein equations and if $`\varphi `$ is a diffeomorphism of $`M`$ onto itself, then the pull-back $`\varphi ^{}g`$ is also a solution. In order to make contact with PDE theory as it is usually formulated, it is necessary to choose a local coordinate system $`x^\alpha `$ on $`M`$ and treat the components $`g_{\alpha \beta }`$ of the metric in this coordinate system as the basic unknowns. The Einstein equations imply a system of second order partial differential equations for these components. This system cannot have a well-posed Cauchy problem because of diffeomorphism invariance. For suppose that $`g_{\alpha \beta }`$ is a solution of these coordinate equations with given initial data. The functions $`g_{\alpha \beta }`$ are the components of a metric $`g`$. Now let $`\varphi `$ be a diffeomorphism which leaves a neighbourhood of the Cauchy surface invariant and let $`\overline{g}=\varphi ^{}g`$. Then $`\overline{g}_{\alpha \beta }`$ is a solution with the same initial data as $`g_{\alpha \beta }`$. However it is in general not equal to $`g_{\alpha \beta }`$. Thus uniqueness in the Cauchy problem in the usual sense fails for this system. On the other hand, there is a different uniqueness statement which does hold. It says that if $`g`$ and $`\overline{g}`$ are metrics with the same initial data on a suitable Cauchy surface then there exists a diffeomorphism $`\varphi `$ such that $`\overline{g}=\varphi ^{}g`$ which is the identity on the Cauchy surface. In the following the precise mathematical formulation of this statement will not be required, but the intuitive idea is important. It should be noted that this uniqueness property, often called geometric uniqueness, is exactly what is desired from the point of view of physics since metrics related by a diffeomorphism should be regarded as physically indistinguishable.
It has now been indicated that diffeomorphism invariance presents a difficulty when studying the Cauchy problem for the Einstein equations. How can this difficulty be overcome? In principle, one could attempt to construct a diffeomorphism-invariant version of PDE theory but this has apparently never been done. In practice, the procedure is to require some condition relating the coordinates to the metric. This breaks the diffeomorphism invariance and opens up the possibility that with this extra condition imposed the Einstein equations imply a system of hyperbolic equations called the reduced equations, and that conversely suitable solutions of the reduced equations give rise to solutions of the Einstein equations. Details on this procedure of hyperbolic reduction can be found in . There are many different kinds of hyperbolic reduction and the task is to find reductions which are well-adapted to given problems.
The Einstein equations are the basic equations of general relativity, which is the best existing theory of the gravitational field. Within this theory spacetime singularities, such as the big bang, correspond to singularities of $`g_{\alpha \beta }`$ which cannot be removed by a diffeomorphism, even one which itself becomes singular where $`g_{\alpha \beta }`$ does.
## 2 The Gowdy equations
A useful laboratory for studying questions concerning the global behaviour of solutions of the Einstein equations is the class of Gowdy solutions. These have high symmetry and after a suitable hyperbolic reduction give the following system of semilinear wave equations:
$`_t^2Xt^1_tX+_x^2X`$ $`=`$ $`2(_tX_tZ_xX_xZ)`$ (1)
$`_t^2Zt^1_tZ+_x^2Z`$ $`=`$ $`e^{2Z}((_tX)^2(_xX)^2)`$
Here $`(t,x)(0,\mathrm{})\times 𝐑`$ and $`X`$ and $`Z`$ are real-valued functions. Because of the symmetry only one space dimension remains. It will be useful to compare this with the equation
$$_t^2u+_x^2u=e^u$$
(2)
The equation (2) and its higher dimensional analogues have been studied in a series of papers by Kichenassamy and Littman. Partial results of a similar type for (1) have been obtained in and . An obvious difference between the equations (1) and (2) is that while the singularities in solutions of (2) develop dynamically in the $`(t,x)`$ plane there is a singularity in the coefficients of (1). Moreover it was proved by Moncrief that smooth Cauchy data for (1) at $`t=t_0>0`$ develop into a smooth solution on the whole of the interval $`(0,\mathrm{})`$. In fact this difference is more apparent than real. When proving theorems about (2) it is convenient to introduce a function which is constant on the singular hypersurface and rewrite the equations in terms of that. In reality, a corresponding procedure has been carried out when introducing the coordinates used in the Gowdy equation. This is so natural in the context of the Einstein equations, where a priori all time coordinates are equally valid, that it might pass unnoticed.
The nature of all these results is that they demonstrate the existence of large families of solutions of the equations whose singularities can be described precisely. This is achieved by using certain singular hyperbolic equations, whose singularity is of the type known as Fuchsian. The description of the singularity is given by an asymptotic expansion, which contains a number of free functions. These free functions may be called ‘data on the singularity’. If these data are analytic then the expansion actually converges. The case of analytic data is easiest to handle. This was done for (2) in and and for (1) in . The results for (2) were generalized to the case of smooth data (in fact data with sufficiently high order Sobolev regularity) in . The theory developed there does not suffice to cover (1) with data which are smooth but not analytic. This difficulty was overcome in , where there is also a general discussion of the issue of extending results using Fuchsian techniques from the analytic to the smooth case and a survey of the literature where these techniques have been applied to problems in general relativity.
Now the asymptotic expansions for solutions of the Gowdy equations will be presented for illustration.
$`X(t,x)`$ $`=`$ $`X_0(x)+t^{2k(x)}(\psi (x)+v(t,x))`$ (3)
$`Z(t,x)`$ $`=`$ $`k(x)\mathrm{log}t+\varphi (x)+u(t,x)`$ (4)
Here the functions $`u(t,x)`$ and $`v(t,x)`$ are remainder terms which are $`o(t)`$ as $`t0`$. Note the occurrence of powers of $`t`$ with exponents depending on $`x`$ in the expansion for $`X`$. It is these which cause additional technical difficulties in this case as compared to the case of (2). Solutions of the above form are obtained for smooth data $`X_0`$, $`k`$, $`\varphi `$ and $`\psi `$ which are arbitrary except for the restrictions that $`0<k(x)<1`$. The significance of the condition $`k>0`$ will not be discussed here. The condition $`k<1`$, known as the low velocity condition, will play a significant role in the following discussion. It turns out that it can be removed at the expense of requiring that $`X_0`$ should be constant. In this way a class of high velocity solutions can be obtained. They depend, however, on one less free function than the low velocity solutions.
For the physical interpretation of the results on singularities in Gowdy spacetimes it is important to know that the singularity in the metric coefficients at $`t=0`$ cannot be removed by a diffeomorphism. One way of doing this is to consider scalar polynomial quantities in the curvature of the metric. These are real-valued functions which are defined in a coordinate invariant way. One of the simplest is the Kretschmann scalar $`K`$, which is the squared length (with respect to the given metric) of the Riemann curvature tensor. It is a function of $`X`$, $`Z`$ and their first and second order partial derivatives whose explicit form is very complicated. Fortunately, for a solution with an asymptotic expansion of the type given above it is not too difficult to calculate the leading term in the expansion of $`K`$ about $`t=0`$. The result, discussed in , is
$$KQ(k)t^{3(k^2+1)/4}$$
(5)
for a certain quartic polynomial $`Q`$ whose only positive root is at $`k=1`$. Thus, provided $`k1`$ everywhere, $`K`$ blows up uniformly as $`t0`$.
An important issue is that of the generality of the solutions constructed by the above techniques. They are general in the crude sense that they depend on the same number of free functions as data for the same equations on a regular Cauchy surface. In the case of (2) Kichenassamy proved in using the Nash-Moser theorem that a non-empty open set of data on a regular Cauchy surface develop into solutions of the type constructed by Fuchsian techniques. Up to now no corresponding results have been proved in problems arising in general relativity, such as the Gowdy problem. By a more direct approach, Chruściel was able to show that a large open set of data for the Gowdy equations lead to solutions where the Kretschmann scalar blows up as $`t0`$. Note that while the methods used in were special to one space dimension, Fuchsian techniques do not depend essentially on the dimension.
## 3 The Einstein-scalar field equations
As mentioned above, the unknown of central importance in the Einstein equations is a metric $`g`$. Physically, it contains the information about the gravitational field and hence about the mutual gravitational attraction of material bodies. Thus it is natural that in any model involving the Einstein equations a description of the matter present is also necessary. In fact, in general relativity the gravitational field has degrees of freedom which can propagate in the absence of matter. This leads to the phenomenon of gravitational waves. The equations in the absence of matter are known as the vacuum Einstein equations and have the simple geometric form $`\mathrm{Ric}(g)=0`$, where $`\mathrm{Ric}(g)`$ is the Ricci curvature of $`g`$. In the presence of matter these equations acquire a non-vanishing right hand side. An example which will be important in the following is that of the Einstein-scalar field equations. In that case the Einstein equations take the form $`\mathrm{Ric}(g)=\varphi \varphi `$. In general the matter fields must satisfy some equations of motion and in the case of the scalar field the equation of motion is that $`\varphi `$ should satisfy the linear wave equation defined by the metric $`g`$. In coordinate components this takes the form
$$g^{\alpha \beta }(_\alpha _\beta \varphi \mathrm{\Gamma }_{\alpha \beta }^\gamma _\gamma \varphi )=0$$
(6)
where
$$\mathrm{\Gamma }_{\alpha \beta }^\gamma =\frac{1}{2}g^{\gamma \delta }(_\alpha g_{\beta \delta }+_\beta g_{\alpha \delta }_\delta g_{\alpha \beta })$$
(7)
are the Christoffel symbols of $`g`$. Here we use the summation convention and $`g^{\alpha \beta }`$ is the inverse matrix of $`g_{\alpha \beta }`$.
It is known that in some cases the singularities of solutions of the Einstein equations display complicated oscillatory behaviour and it is suspected that this occurs under very general circumstances. The best known example of complicated behaviour is the Mixmaster solution. This is a solution of the vacuum Einstein equations which is symmetric under the group $`SU(2)`$. The Einstein equations reduce to ordinary differential equations but the ODE solutions are very complicated. Until recently our knowledge of the Mixmaster solution was based on numerical calculations and heuristic arguments. Now a number of features of this picture have been proved rigorously by Ringström. The generality of Mixmaster behaviour is still out of reach of present analytical results. There have, however, been interesting advances in numerical approaches. (See .)
General predictions about the nature of singularities in solutions of the Einstein equations come from heuristic work of Belinskii, Khalatnikov and Lifshitz. According to their picture the singularities in solutions of the vacuum Einstein equations should look like a different Mixmaster solution at each spatial point near the singularity. In particular the evolution at different spatial points decouples and the solution of the PDE system can be approximated by solutions of an ODE system. (Here we are reminded of the ODE blow-up mechanism in the terminology of Alinhac.) The same picture suggests that for the Einstein-scalar field system the decoupling should still take place but the ODE solutions should now be monotone rather than oscillatory near the singularity. This gives rise to a hope that the Einstein-scalar field system should be easier to handle analytically than the Einstein vacuum equations. Following this lead, it was shown in that a Fuchsian analysis could be carried out for the Einstein-scalar field system without any need for symmetry assumptions. Solutions are obtained which depend on the same number of free functions as there are in the Cauchy data on a regular Cauchy surface. In this crude sense the solutions are general. Only the analytic case was treated up to now.
The asymptotic expansions obtained will now be presented. This is only done for a certain diagonalizable case. This does not restrict the number of free functions although it is a restriction on the solutions and a large part of the work of was devoted to showing that the restriction is unnecessary. On the other hand, the diagonalizable case is best suited to explain the results. For physical reasons we are interested in solutions on a four-dimensional manifold. The time coordinate is denoted by $`x^0`$ and the spatial coordinates by $`x^a`$ with $`a=1,2,3`$. The Gaussian coordinate conditions
$`g_{00}`$ $`=`$ $`1`$ (8)
$`g_{0a}`$ $`=`$ $`0`$ (9)
are imposed. The asymptotic expansions for the remaining unknowns are as follows:
$`g_{ab}`$ $`=`$ $`t^{2p_1}l_al_b+t^{2p_2}m_am_b+t^{2p_3}n_an_b+o(t^{q_{ab}})`$ (10)
$`\varphi `$ $`=`$ $`A\mathrm{log}t+B+o(1)`$ (11)
Here $`p_1`$, $`p_2`$, $`p_3`$, $`l_a`$, $`m_a`$, $`n_a`$, $`A`$ and $`B`$ depend on the spatial coordinates $`x^a`$ and $`p_1+p_2+p_3=1`$. The $`q_{ab}`$ are suitably chosen real numbers, which will not be discussed in detail here. The important thing is that they represent higher order terms. The functions occurring in the expansion must satisfy certain algebraic and differential equations which will also not be discussed, except to say that there is a good understanding of the solution set of these equations. There is an equivalent of the condition $`k<1`$ in the Gowdy case, namely that all $`p_a`$ should be positive. This implies that all $`p_a`$ are smaller than one.
The conclusion of this work is that it provides strong support for the ideas of Belinskii, Khalatnikov and Lifshitz in the case that a scalar field is present. The conceptual, geometric and calculational aspects of the proof are a lot more complicated than in the Gowdy case, but the analytic input is almost identical.
## 4 The nonlinear scalar field
From the foregoing discussion of the great difference in the nature of singularities in solutions of the Einstein equations caused by the presence of a scalar field the impression might arise that the whole question depends in a delicate way on the exact nature of the matter fields present. In fact this impression is probably false: there should probably only be a small number of universality classes which cover most types of matter. This will be illustrated by a class of matter models which will be referred to collectively as the nonlinear scalar field. It turns out that a wide variety of choices of the parameters which go into the definition all lead to the same asymptotic behaviour, which is that of the (linear) scalar field discussed in the previous section. This will also allow some of the aspects of the Fuchsian theory to be shown in action.
The main interest for physicists of the equations obtained by coupling the Einstein equations to a nonlinear scalar field has been the the phenomenon known as inflation. The mathematics of relevance to inflation concerns the behaviour of the solutions in the direction away from the singularity and so has no direct connection with the discussion below.
Consider first the following equation
$$_t^2\varphi t^1_t\varphi +\mathrm{\Delta }\varphi =m^2\varphi +V^{}(\varphi )$$
(12)
Here $`\varphi `$ is a real-valued function, $`\mathrm{\Delta }`$ is the Laplacian on $`𝐑^n`$ (the case of $`n=3`$ is of most interest for the following) and $`m`$ is a constant (mass) which we may as well take to be non-negative. The function $`V`$ (potential) is smooth and $`V^{}`$ vanishes at least quadratically at the origin. The separation of the right hand side into the sum of a mass term and a potential term has no essential significance in this problem. It is convenient for making contact with the notation used in the physics literature. The aim is to determine whether this equation has solutions which near $`t=0`$ have the leading order asymptotics $`A(x)\mathrm{log}t+B(x)`$ already encountered in the previous section. To this end, introduce $`\psi =t^ϵ(\varphi A\mathrm{log}tB)`$, where $`ϵ`$ is a positive constant, whose purpose will be seen later on. When reexpressed in terms of $`\psi `$ the equation (12) becomes:
$`t^2_t^2\psi (2ϵ+1)t_t\psi `$ $`ϵ^2\psi +t^2\mathrm{\Delta }\psi =t^{2ϵ}\mathrm{log}t\mathrm{\Delta }At^{2ϵ}\mathrm{\Delta }B`$ (13)
$`+t^{2ϵ}m^2(A\mathrm{log}t+B+t^ϵ\psi )+W(t,x,\psi )`$
where $`W(t,x,\psi )=t^{2ϵ}V^{}(A\mathrm{log}t+B+t^ϵ\psi )`$. The equation has been multiplied by $`t^2`$ so as to make the combination $`td/dt`$ appear. Some choices of $`V`$ to be found in the physics literature are $`V_1(\varphi )=V_0`$, a constant, $`V_2(\varphi )=\varphi ^\alpha `$, $`V_3(\varphi )=e^{\beta \varphi }`$ for positive constants $`\alpha `$ and $`\beta `$. The value of the constant $`V_0`$ has no effect on the equation (12) but does have an effect when the coupling to the Einstein equations is considered. The potential $`V_1`$ combined with a non-zero value of $`m`$ occurs in a model known as chaotic inflation. The potential $`V_3`$ occurs in what is called power-law inflation. (It is not the potential itself which has power-law behaviour, but the long-time behaviour of the solution.) The functions $`W`$ corresponding to $`V_1`$, $`V_2`$ and $`V_3`$ are $`W_1(t,x,\psi )=0`$, $`W_2(t,x,\psi )=\alpha t^{2ϵ}(A\mathrm{log}t+B+t^ϵ\psi )^{\alpha 1}`$ and $`W_3(t,x,\psi )=\beta e^{\beta B}t^{2ϵ+\beta A}\mathrm{exp}(\beta t^ϵ\psi )`$. The functions $`W_1`$, $`W_2`$ and $`W_3`$ are smooth for $`t>0`$. Suppose that the function $`A`$ is nowhere vanishing. Then $`W_1`$ and $`W_2`$ extend continuously to $`t=0`$ by defining them to be zero there. This is also true of their derivatives with respect to the arguments $`x`$ and $`\psi `$. This will be abbreviated by saying that the functions $`W_1`$ and $`W_2`$ are regular. As for $`W_3`$, can be made regular in this sense by choosing $`ϵ`$ sufficiently small provided $`A>2\beta ^1`$. This is a restriction on the data on the singularity. In particular, if $`\beta `$ is negative then it bounds $`A`$ away from zero. If $`V`$ had been chosen to grow even faster at infinity, e.g. $`V(\varphi )=e^{\varphi ^2}`$ then the corresponding $`W`$ could not have been made regular by imposing inequalities on $`A`$. This type of potential with faster then exponential growth does not seem to occur in the physics literature and the Fuchsian theory does not seem to apply to it. It will not be considered further here.
The Fuchsian theory we will use concerns first order systems and thus the first step in the analysis is to reduce (13) to first order in a suitable way. Let $`v=t_t\psi +ϵ\psi `$ and $`w=t_x\psi `$. Then the the following system is obtained:
$`t_t\psi +ϵ\psi v`$ $`=`$ $`0`$
$`t_tw+(ϵ1)w`$ $`=`$ $`t_xv`$ (14)
$`t_tv+ϵv`$ $`=`$ $`t_xw+\mathrm{\Delta }At^{2ϵ}\mathrm{log}t+t^{2ϵ}\mathrm{\Delta }B`$
$`t^{2ϵ}m^2(A\mathrm{log}t+B+\psi )+W(t,x,\psi )`$
This system is of the following general form for Fuchsian equations:
$$t_tu+N(x)u=t^\eta f(t,x,u,_xu)$$
(15)
for some $`\eta >0`$. It is also symmetric hyperbolic, a luxury which is usually not easy to obtain in problems in this area. On the other hand, it does have a disadvantage. The theorem proved in required the assumption that $`t^{N(x)}`$ be bounded for $`t`$ near zero. To obtain this the eigenvalues of $`N(x)`$ should have non-negative real parts. Since $`ϵ`$ has to be chosen small in general, depending on the data, the eigenvalue $`ϵ1`$ will be negative and this condition is not satisfied for the equation (4). An alternative form is obtained by substituting the definition of $`v`$ back into (4) in one place. The result is:
$`t_t\psi +ϵ\psi v`$ $`=`$ $`0`$
$`t_tw+ϵw`$ $`=`$ $`t_xv+t_x\psi `$ (16)
$`t_tv+ϵv`$ $`=`$ $`t_xw+\mathrm{\Delta }At^{2ϵ}\mathrm{log}t+t^{2ϵ}\mathrm{\Delta }B`$
$`t^{2ϵ}m^2(A\mathrm{log}t+B+\psi )+W(t,x,\psi )`$
The equation (4) is not symmetric hyperbolic. It does, however satisfy the conditions of the existence theorem for the case of analytic data proved in . In order to ensure this it is important that $`ϵ`$ is positive. If $`ϵ`$ were chosen to be zero then all the eigenvalues of the matrix $`N(x)`$ would vanish but the matrix would not be diagonalizable. This would violate one of the conditions assumed in , namely that zero eigenvalues, if they occur at all, only occur through $`N(x)`$ being a direct sum of a matrix whose eigenvalues have positive real parts with a zero matrix.
The result of the previous discussion is that given analytic functions $`A`$ and $`B`$ with $`A`$ nowhere vanishing then if the restrictions on $`V`$ and $`A`$ already obtained are satisfied there exists a unique solution, analytic in $`x`$ and continuous in $`t`$ of (12) with the desired asymptotics. This shows the existence of a family of solutions depending on two free analytic functions, which is the same as the number of functions which can be given as data on a regular Cauchy surface.
In a method for going beyond the analytic case in Fuchsian problems was applied to the Gowdy equations. Since the matrix $`N(x)`$ occurring in (4) does not depend on $`x`$ it should be possible to prove an existence theorem for that equation with smooth data given on the singularity by the methods of . Here we would like to note that the method of generalizes easily to the case of smooth data for (12). The main idea is to use both the systems (4) and (4). If the smooth data is approximated by a sequence of analytic data then a corresponding sequence of analytic solutions can be produced using (4). These define a sequence of analytic solutions of (4). Then the fact that the latter equation is symmetric hyperbolic allows the use of energy estimates to prove that the sequence of analytic solutions converges to the desired smooth solution. Details of the procedure can be found in .
## 5 Coupling the nonlinear scalar field to the Einstein equations
In this section it will be shown that the Einstein equations coupled to the nonlinear scalar field can be handled by the same method as that used for the case with $`m=0`$ and $`V=0`$ in . The arguments will only be sketched since most of the steps are identical. The potential will be assumed to be one of the functions $`V_1`$, $`V_2`$ or $`V_3`$ considered in the last section. In addition it is assumed that $`A0`$ and that in the case of $`V_3`$ the inequality $`A>2\beta ^1`$ holds. The Einstein equations take the form
$$\mathrm{Ric}(g)=\varphi \varphi +[\frac{1}{2}m^2\varphi ^2+V(\varphi )]g$$
(17)
Evidently this equation does change if $`V`$ is changed by an additive constant. The equation for the scalar field is
$$g^{\alpha \beta }(_\alpha _\beta \varphi \mathrm{\Gamma }_{\alpha \beta }^\gamma _\gamma \varphi )=m^2\varphi +V^{}(\varphi )$$
(18)
Applying the Gaussian coordinate conditions $`g_{00}=1`$ and $`g_{0a}=0`$ gives
$$g^{\alpha \beta }_\alpha _\beta \varphi =_t^2\varphi +g^{ab}_a_b\varphi $$
(19)
and the close relation of equation (18) to equation (12) begins to become apparent. We are looking for solutions to the Einstein equations coupled to a nonlinear scalar field where the metric has the asymptotic form found in for the case of a linear scalar field. This form implies that the quantities $`t^2g^{ab}`$ all vanish like some positive power of $`t`$ as $`t0`$ and thus do not disturb the Fuchsian form of the wave equation. The same holds for the quantities $`t\mathrm{\Gamma }^a`$ and the quantity $`t\mathrm{\Gamma }^0`$ is equal to one up to higher order terms. These facts cannot be easily seen from the equations given in this paper. To check them it is necessary to use more detailed results from . The conclusion is that up to higher order terms the equation (18) agrees with the model equation (12).
What remains to be done in order to check that the Einstein equations coupled to a nonlinear scalar field have solutions which behave asymptotically like those for a linear scalar field is to see that certain source terms in the Einstein equations are equal to those in the case with a linear massless scalar field (which has already been solved in ), up to terms of higher order in $`t`$. The essential thing is to estimate the quantity $`_t^2\varphi +\frac{1}{2}m^2\varphi ^2+V(\varphi )`$. More specifically it is necessary to estimate the deviation of this quantity from the corresponding quantity with $`\psi =0`$ and show that it is $`O(t^{2+\alpha _0})`$ where $`\alpha _0`$ is a positive parameter introduced in . In the course of the proofs in that paper it is permissible to reduce the size of the number $`\alpha _0`$ if required. It is easy to check that, under the assumptions already made on $`V`$ and $`A`$, the parameter $`\alpha _0`$ can be chosen so as to ensure this estimate. This is the same kind of procedure as choosing $`ϵ`$ small in the case of the nonlinear scalar field. The information obtained in this way must be used in two places. It must be used in the Einstein constraints in estimating the quantity $`\overline{C}`$ of . It must also be used in evolution equations for $`\kappa ^a_b`$ in , which constitute the essential part of the Einstein evolution equations. The additional contribution which arises when passing from the massless linear scalar field to the nonlinear scalar field does not change the equations for $`ab`$. For $`a=b`$ the coefficient $`\alpha _0`$ appears. It can be concluded from all this that solutions of the Einstein equations coupled to a nonlinear scalar field which have the same asymptotics as in the case of a massless linear scalar field and which are of the same degree of generality can be constructed by Fuchsian techniques. This has been shown under the assumption that the potential is of one of the forms $`V_1`$, $`V_2`$ or $`V_3`$ and that in the third case the initial datum $`A`$ satisfies an additional restriction. In this proof the velocity-dominated system, which is the system of equations satisfied by the leading order approximation to the solution, is the same as in the case of the massless linear scalar field.
## 6 Formation of localized structure
It has already been indicated that the use of Fuchsian theory to construct spacetimes with singularities of a well-understood type is not likely to be sufficient to handle general singularities of solutions of the Einstein equations, due to the appearance of oscillatory behaviour in time as the singularity is approached. There is, however, another effect which complicates the situation in an essential way. Note that the solutions constructed by Fuchsian techniques are such that the solution blows up like $`\mathrm{log}t`$ as $`t=0`$ is approached and the spatial derivatives of the solutions of all orders are also $`O(\mathrm{log}t)`$. Taking spatial derivatives does not increase the rate of blow-up. It will now be shown that there are solutions for which this property does not hold. Consider the transformation:
$`\stackrel{~}{X}`$ $`=`$ $`X/(X^2+e^{2Z})`$ (20)
$`\stackrel{~}{Z}`$ $`=`$ $`\mathrm{log}[e^Z/(X^2+e^{2Z})]`$ (21)
This transformation is equal to its own inverse and transforms solutions of the Gowdy equations into solutions. The idea is now to take a solution $`(\stackrel{~}{X},\stackrel{~}{Z})`$ with the low velocity asymptotic behaviour and transform it to a new solution $`(X,Z)`$.
The case will be considered that $`\stackrel{~}{X}_0(0)=0`$, $`\stackrel{~}{X}_0(x)0`$ for $`x0`$ and $`_x\stackrel{~}{X}_0(0)0`$. For the transformed solution $`X^2+e^{2Z}`$ is of the form
$$(\stackrel{~}{X}_0)^2+(e^{2\stackrel{~}{\varphi }}+2\stackrel{~}{X}_0\stackrel{~}{\psi })t^{2\stackrel{~}{k}}+o(t^{2\stackrel{~}{k}})$$
(22)
From this it is not hard to see that $`X`$ is bounded on an interval about $`x=0`$ uniformly in $`t`$. Similarly $`Z=O(\mathrm{log}t)`$. Thus the transformed solution $`(X,Z)`$ does not blow up any faster than the original solution $`(\stackrel{~}{X},\stackrel{~}{Z})`$. However spatial derivatives blow up faster in general. To see this it suffices to evaluate $`Z(t,0)`$ and $`Z(t,t^{\stackrel{~}{k}/2})`$. The second of these is bounded for values of $`x`$ close to $`x=0`$ as $`t0`$ while the first behaves like $`\stackrel{~}{k}\mathrm{log}t`$. It follows from the mean value theorem that the maximum value of $`_xZ`$ in a small region about $`x=0`$ must blow up at least as fast as $`t^{\stackrel{~}{k}/2}`$. Of course precise rates of blow-up could be calculated for this and higher order derivatives if desired. It is seen that these solutions of the Gowdy equations, unlike those constructed by the Fuchsian method, have spatial derivatives which blow up faster than the function itself. There is formation of localized structure.
In the explicit example given above the localized structures formed do not have an invariant geometrical meaning and are due to a bad choice of variables. On the other hand, there is convincing numerical evidence due to Berger, Moncrief and collaborators for a second kind of localized structure. In the latter case there is no reason to doubt that there is a real geometrical effect involved. For a survey of this work and some pictures of the structures, see . There is also a lot of interesting work on the subject using heuristic methods. A rigorous treatment is still lacking.
Although this is still in the domain of speculation, there are arguments that in more general situations things will get yet more complicated. The basic idea will now be sketched. The oscillatory behaviour in the Mixmaster model is composed of a series of bounces where expansion turns into contraction and vice versa. In the inhomogeneous case bounces of this kind are supposed to happen independently at different spatial points. However this does not always happen coherently. There are spatial regions of coherent behaviour but these split into smaller regions as the singularity is approached. On the boundaries between regions spatial derivatives become large. The localized features observed in the Gowdy solutions are simple examples of boundaries of this kind. The heuristic models available suggest that in the Gowdy case this kind of formation of structure will only happen a finite number of times, at least for generic solutions. In the general case, on the other hand, it is expected to happen infinitely many times before the singularity is reached, creating structures on arbitrarily small spatial scales. This behaviour has been called ‘spacetime turbulence’. Information on this can be found in . It has also been suggested that this process which is predicted by the picture of Belinskii, Khalatnikov and Lifshitz actually shows that the picture must eventually break down in the generic case before the singularity is reached.
It seems clear that it will take a long time before the ideas mentioned in the last paragraph can be brought into the fold of well-understood mathematical phenomena. Nevertheless, the recent progress described in this paper is a ground for optimism that there will be significant mathematical developments in the area of singularities of solutions of the Einstein equations in the next few years.
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# Introduction
## Introduction
Recently, there have been various proposals entertaining the idea that our $`d=3+1`$ dimensional world is located on a brane in a higher dimensional space-time. In such scenarios, the standard model fields are confined to the brane, whereas gravity permeates the bulk. The idea of compactification by localization of zero modes on defects of a higher dimensional field theory was first advocated in Ref. .
In Refs., the hierarchy problem is addressed through the notion of $`n`$ compact (large) extra dimensions of size $`R1\mathrm{mm}`$. In this scenario the Planck scale $`M_4`$ of our world is not fundamental but can be derived from a higher dimensional scale $`M_{4+n}`$ of the order of the electroweak scale $`M_W`$ by means of the following power-law conversion:
$$M_4^2M_{4+n}^{2+n}R^n.$$
(1)
However, there is still a large hierarchy of the fundamental scales $`M_{4+n},1/R`$.
In contrast, the authors of Ref. suggest to solve the hierarchy problem in a rather different way. The set-up is a five-dimensional universe consisting of one three-brane constituting our world and another, hidden three-brane, which are separated by a (small) distance $`\pi r_c`$. Assuming a special, non-factorizable metric for the five dimensional space-time, there is practically no conversion between (fundamental) Planck scale $`M_5`$ and the Planck scale $`M_4`$ of our world. Scales $`M`$ belonging to the electroweak regime are given in terms of $`M_5`$ as
$$M=\text{e}^{2\pi kr_c}M_5,$$
(2)
where the fundamental scales $`k,1/r_c`$ have no large hierarchy ($`kr_c50`$).
In addition, in Ref. it was proposed that the fermion hierarchy in the standard model may be generated by localizing different families on spatially separated walls along a compact dimension. Thereby, the localized light modes obtain their masses by coupling to a Higgs condensate which falls off exponentially away from the wall on which the Higgs field is localized.
The purpose of this paper is to study the idea of “brane worlds” in a very simple supersymmetric context, namely Wess-Zumino models. One and the same field is used to create domain walls (which play the role of the branes) and the fermionic and bosonic light modes localized on them. The latter play the role of matter in the lower dimensional universes on the walls. We will assume the geometry of space-time as given, and we do not investigate the issue of radius stabilization, although the model we discuss in fact provides a realization of “brane-lattice-crystallization”, proposed in Ref. as a solution to this problem. The models we consider feature several distinct topological sectors. We do not address the question why a particular sector is selected. Obviously, such a limited approach poses severe restrictions: there is no gravity, no gauge fields, and we have to work in lower dimensions. The advantages are transparency, the ability to look at different issues in isolation, and the potential to perform explicit calculations. We will be able to address issues such as spontaneous supersymmetry breaking, the chirality of the localized fermionic modes, and the spectrum of the localized light modes.
For concreteness, we choose to work with a specific model that has the property that we desire; static configurations of well separated domain walls in a compact space dimension. We stress that it is just a toy model; the conclusions we draw should not, and do not, depend on its specific details.
### The model
To be specific, we consider a generalized Wess-Zumino model in $`d=3+1`$ dimensions with one chiral superfield $`\varphi `$.<sup>1</sup><sup>1</sup>1We denote the superfield and its lowest component by the same symbol $`\varphi `$. The superpotential of this model is given by
$$W=\mathrm{ln}(\varphi +1)+\mathrm{ln}(\varphi 1),$$
(3)
and the Kähler potential takes the canonical form. We also consider this model reduced to $`d=2+1`$ and $`d=1+1`$ dimensions.
Since the coordinate $`z`$ is compact, the field $`\varphi `$ obeys periodic boundary conditions $`\varphi (x)=\varphi (x+L)`$, where $`L`$ is the circumference of the compact dimension. In general, static, stable solutions to the second-order equations of motion may or may not be BPS saturated. If a BPS saturated solution is selected as the vacuum of the model, only 1/2 of the SUSY generators is spontaneously broken. We will focus mostly on the BPS saturated solutions, but static, stable, non-BPS solutions also exist, and we will comment on them in passing.
Since the presence of a non-vanishing $`(1,0)`$ central charge $`Z`$ in the SUSY algebra is necessary for the existence of BPS domain walls, it was argued in Ref. that, in generalized Wess-Zumino models with a compact dimension, the superpotential must be a multi-branch function and that the manifold $`M`$ on which the fields are defined must admit non-contractable cycles. The model with the superpotential $`W`$ of Eq. (3) fulfills these criteria. The manifold $`M`$ in this case is a plane with two punctures a the branch points $`\varphi =\pm 1`$.
As we also have occasion to consider the model in less than $`d=3+1`$ dimensions, and as the scalar components of the $`𝒩=1`$ super multiplets in lower dimensions are real, it is convenient to decompose the field $`\varphi `$ according to $`\varphi =1/\sqrt{2}(\varphi _1+ı\varphi _2)`$ with real fields $`\varphi _1`$ and $`\varphi _2`$. By definition, the form of the scalar potential is not altered by the procedure of dimensional reduction. The BPS equation and its solutions are the same in each number of space-time dimensions that we consider. Therefore, we will first characterize all periodic solutions to the BPS equation; armed with this knowledge, we will subsequently investigate various implications in different dimensions. Of course, a solution along a compact dimension is constrained by the requirement that its period $`\lambda `$ fits an integer $`N`$ times on the circumference $`L`$.
The scalar potential of the model is given by
$$V=\frac{dW}{d\varphi }\frac{d\overline{W}}{d\overline{\varphi }}=\left|\frac{1}{\varphi +1}+\frac{1}{\varphi 1}\right|^2.$$
(4)
This potential is invariant under the transformations $`\varphi _1\varphi _1`$ and $`\varphi _2\varphi _2`$. It has a supersymmetric vacuum, $`V=0`$ at $`\varphi =0`$, and also a run-away vacuum at $`|\varphi |\mathrm{}`$. In the vacuum at the origin, the mass of the field is $`m=2`$. In addition, there are two saddle points at $`\varphi =\pm ı`$, and two poles at $`\varphi =\pm 1`$. As we will discuss in detail below, the solutions to the BPS equation either wind around one of these poles, or around both.
At this point we briefly pause to discuss radiative corrections to the classical BPS solutions and the spectrum of the localized modes. Obviously, the model with the superpotential Eq.(3) is not renormalizable in $`d=3+1`$ dimensions. Clearly, it should be regarded as an effective theory, describing the low energy dynamics of a fundamental theory below an energy scale $`M`$. The field $`\varphi `$ is supposedly the only relevant degree of freedom below this scale.
In fact, in order for an effective theory to be useful, it should be in the weakly coupled regime. To that end, the superpotential in Eq.(3) can be slightly modified by the introduction of a coupling constant $`g`$,
$$W=\mathrm{ln}(g\varphi +1)+\mathrm{ln}(g\varphi 1),$$
(5)
so that for $`g1`$ the model is certainly weakly coupled.
The tension of the domain walls is $`T=2\pi `$, independent of $`g`$. The scale $`T`$ should be much smaller than the scale $`M`$ for consistency of the effective theory approach. Moreover, the mass in the supersymmetric vacuum at the origin now becomes $`m=2g^2`$, and the mass scale $`m_{pzm}`$ of the pseudo-zero modes localized on the wall is exponentially suppressed by the ratio of the distance between the walls and their width. We thus find the folowing hierarchy of scales
$$MTmm_{pzm}.$$
(6)
It is the underlying fundamental theory that must create this hierarchy. As a consequence, it is possible to describe the dynamics in different energy ranges by various effective theories.
Radiative corrections to local quantities such as the energy density profile of the domain wall can therefore be calculated in perturbation theory as described in Refs.. The scale $`M`$ forms a physical ultra-violet cut-off in this calculation, and the size of the corrections is controlled by $`g`$.
Moreover, the dynamics of the zero modes and pseudo zero modes that are localized on the walls can be described by a weakly coupled, $`d=2+1`$ dimensional, effective theory. In this case, the scale $`m`$ forms the physical cut-off. Above this scale, more massive modes become accessible and should be incorporated in the theory. Radiative corrections can be calculated in a perturbative expansion and are controlled by powers of the coupling constant $`g`$.
As radiative corrections are not the issue we investigate in this paper, we will just take $`g=1`$ in what follows. Our results can be trivially generalized for other values of $`g`$. In addition, we will not speculate about the nature of an underlying fundamental theory that could give rise to the superpotential in Eq.(5).
The paper can be outlined as follows: In Section 1 we study and classify all periodic solutions to the BPS equation. We show that there are BPS saturated configurations that contain well separated domain walls. In Section 2 we investigate such configurations with a sequence of equidistant walls in $`(3+1)`$-dimensions. Apart from the bosonic and fermionic zero modes associated with the spontaneous breaking of translational invariance and half of supersymmetry, there are other light modes with exponentially suppressed masses. The sequence of domain walls can be considered as a crystal, and the spectrum of the pseudo-zero modes is analyzed using tight binding methods. We then study to what extend the domain walls can be considered as parallel universes. After integrating out the heavy modes, the $`(2+1)`$-dimensional theory of the light modes has $`𝒩=1`$ supersymmetry. In the multi-wall background, light states that are localized on one particular wall are no longer mass eigenstates. This leads to oscillation phenomena. Section 3 investigates the same model in $`(2+1)`$ dimensions. The interesting new feature is that we are now able to study the “chirality” (whether they are left or right moving) of the fermionic zero modes localized on the domain walls (there are no chiral fermions in $`d=2+1`$ dimensions). We show that in the BPS saturated background there are always one left moving and one right moving fermion. We argue that this continues to be the case even if the $`𝒩=2`$ supersymmetry is explicitly broken to $`𝒩=1`$, in which case only one SUSY generator is spontaneously broken. In Section 4 we study the model in $`d=1+1`$ dimensions. What is different here is that the domain walls with infinite mass in the higher dimensions now become solitons with a finite mass. We show that in particular topological sectors, and for sufficiently large circumference of the compact dimension, the classical ground state is BPS saturated, but non-perturbatively the ground state energy is lifted above the BPS bound. This phenomenon was first observed in Ref. . The ground states are connected by instanton tunneling. The tunneling probability and the energy shift of the ground state are obtained as a function of the circumference of the compact dimension. The Appendix contains technicalities concerning the fermionic zero modes.
## 1 Periodic solutions to the BPS equation
In this section we study periodic solution to the BPS equation,
$$\frac{d\varphi }{dx}=e^{ı\delta }\frac{d\overline{W}}{d\overline{\varphi }},$$
(7)
where the phase $`e^{ı\delta }`$ takes the value $`\pm ı`$ for periodic solutions. The two possible choices for the phase correspond to whether the solutions wrap around the poles clockwise or counter clockwise in the complex $`\varphi `$ plane. Solutions for one choice can be obtained from solutions with the opposite choice of the sign by the transformation $`xx`$, mapping walls into anti-walls and vice versa. For definiteness, we will from here on consider the BPS equation with $`e^{i\delta }=+i`$. It is well known that solutions of the BPS equation have a “constant of the motion” $`I=\mathrm{Im}[e^{i\delta }W]`$. This constant of the motion can in principle be used to mark the individual solutions. Here, instead of $`I`$, we use the related quantity
$$k=e^{2I}=\frac{1}{4}(\varphi _1^2+2+\varphi _2^2)^22\varphi _1^2,$$
(8)
to label the solutions, which takes a simpler form. The “constant of the motion” $`k`$ is positive semi-definite. It only vanishes at the poles $`\varphi =\pm 1`$. Therefore, there are solutions for each value $`k>0`$. Since the BPS equation is a first order differential equation, solutions, when viewed as trajectories in the complex $`\varphi `$ plane, do not intersect each other or themselves, except possibly at the supersymmetric vacuum $`\varphi =0`$.
The target space is a plane with two punctures at the location of the poles of the scalar potential. The trajectories of solutions to the BPS equation fall into three homotopy classes of non-contractable cycles in the complex $`\varphi `$ plane.<sup>2</sup><sup>2</sup>2Here, we are actually classifying the functions $`f_k(s)=\varphi _k(s\lambda (k))`$ ($`s[0,1]`$), where $`\lambda (k)`$ is the wavelength of the solution labeled by $`k`$. For any given value of $`k`$ in the range $`0<k<1`$ there are two distinct solutions to the BPS equation. These two solutions are mapped onto each other by the transformation $`\varphi _1\varphi _1`$, which is a symmetry of the scalar potential. One solutions winds around the pole at $`\varphi =1`$, whereas the other winds around the pole at $`\varphi =1`$. We will denote these two homotopy classes as Class IA and Class IB, respectively. Solutions in a third homotopy class, which we refer to as Class II, are obtained for $`k>1`$. For each value of $`k>1`$ there is only one solution, which winds around both poles of the scalar potential and is invariant under the transformation $`\varphi _1\varphi _1`$.
There is a critical solution for $`k=k_c=1`$ which separates solutions in the three homotopy classes in the complex $`\varphi `$ plane. This critical solution contains the point $`\varphi =0`$, where the scalar potential has a minimum. The wavelength of this solution is therefore infinite. It represents a domain wall that interpolates between the same supersymmetric vacuum at $`x\mathrm{}`$ and $`x+\mathrm{}`$, winding either around the pole at $`\varphi =1`$ or the pole at $`\varphi =+1`$. We will refer to these walls as “left” and “right” walls respectively. Examples of solutions in each of the homotopy classes are shown in Fig.(1).
At $`k=4`$ the cycle passes through the two saddle points at $`\varphi \pm 1`$. For the cycles in Class II, the wavelength tends to infinity both in the limit $`k1`$ and $`k\mathrm{}`$. In between, the wavelength reaches a minimum which we numerically determined to take the value $`\lambda _m=4.96`$ for $`k_m=1.42`$.
For very small values of $`k`$, $`k1`$, the cycles are small circles centered around the poles in the complex $`\varphi `$ plane. The wavelength of such solutions is approximately $`\lambda =\pi k/2`$, and the energy density is approximately constant. For large $`k`$, $`k1`$, the cycles are approximately large circles centered around the origin. In this limit the wavelength is approximately given by $`\lambda =\pi \sqrt{k}`$, and again, the energy density is approximately constant.
### 1.1 Domain wall solutions
For $`k`$ close to the critical value $`k_c=1`$ the solutions exhibit well separated wall-like structures. To illustrate this, we show in Fig.(3) a numerical solution $`\varphi (x)`$ of Class IA subject to the initial condition $`\mathrm{Re}\left[\varphi (0)\right]=0.0001`$ and $`\mathrm{Im}\left[\varphi (0)\right]=0`$. The value of the “constant of the motion” $`k`$ that is associated with this solution is $`k110^8`$. The corresponding energy-density $`\epsilon (x)`$ is also indicated. It is clear that if the trajectory of the solution in the complex $`\varphi `$ plane passes the supersymmetric vacuum at the origin, the ratio of the width of the walls to distance between them is small since the force driving the trajectory away from the minimum is small. In Fig.(4) a numerical solution, belonging to Class II and subject to the initial condition $`\mathrm{Re}\left[\varphi (0)\right]=0`$ and $`\mathrm{Im}\left[\varphi (0)\right]=0.0001`$, is shown. The value of the “constant of the motion” $`k`$ that is associated with this solution is $`k1+10^8`$. For $`k1`$ the solution consists of equally spaced “right” walls (Class IA) or equally spaced “left” walls (Class IB), and for $`k1`$ the solution consists of equally spaced, alternating “left” and “right” walls. As the critical value $`k_c=1`$ is approached, the distance between the walls increases, but the width and the shape of the walls converge.
### 1.2 Approximate solutions in between the walls
In between the walls, both the real and the imaginary part of $`\varphi `$ are very small. In this region the trajectory is approximately governed by the BPS equation linearized about $`\varphi =0`$
$`{\displaystyle \frac{d\varphi _1}{dx}}`$ $`=`$ $`2\varphi _2,`$
$`{\displaystyle \frac{d\varphi _2}{dx}}`$ $`=`$ $`2\varphi _1.`$ (9)
The general solution of Eq. (1.2) is
$$\varphi _1(x)=C_+e^{2x}+C_{}e^{2x},\varphi _2(x)=C_+e^{2x}+C_{}e^{2x}.$$
(10)
For Class I cycles, with the initial condition $`\varphi _1(0)=u`$ and $`\varphi _2(0)=0`$, the solution is
$`\varphi _1(x)=u\mathrm{cosh}2x,`$
$`\varphi _2(x)=u\mathrm{sinh}2x.`$ (11)
The constant $`u`$, which is very small for well separated walls, is related to the constant of motion $`k`$ of the solution by $`k=(1u^2/2)^2`$. The approximation breaks down when $`x`$ is so large that the magnitude of the fields $`\varphi _1`$ and $`\varphi _2`$ becomes of the order one. This happens at the location of the walls. The wavelength of the solution is therefore related to the constant $`k`$ by $`\lambda =1/2\mathrm{ln}(1k)`$ for $`k1`$.
Similarly, the approximate solution for Class II solutions in between a “left” wall and a “right” wall is given by
$`\varphi _1(x)`$ $`=`$ $`v\mathrm{sinh}2x,`$
$`\varphi _2(x)`$ $`=`$ $`v\mathrm{cosh}2x.`$ (12)
The constant $`v`$, which also has to be very small for well separated walls, is related to the constant of motion by $`k=(1+v^2/2)^2`$. From this approximate solution it follows that the wavelength of the Class II solution is related to the constant $`k`$ by $`\lambda =\mathrm{ln}(k1)`$ in the limit $`k1`$.
### 1.3 BPS domain walls on a circle
So far we have discussed generic periodic solutions to the BPS equation. Now we want to determine what BPS saturated configurations are possible on a compact space with given circumference $`L`$. It is clear that we have to select only those solutions in the total solution space for which the wavelength $`\lambda `$ fits an integer $`N`$ times on the circumference $`L`$. The allowed values of $`k`$ are thus the solutions of the equation $`L=N\lambda (k)`$, where the integer $`N`$ is a winding number that specifies how many times the solution winds around the poles when $`x`$ varies between $`0`$ and $`L`$. In Fig.(2) we show the period function $`\lambda (k)`$, which gives the wavelength of a solution as a function of $`k`$.
The simplest BPS configurations have winding number $`N=1`$; for any value of $`L`$ there is such a configuration of Class IA and Class IB, but only if $`L>\lambda _m`$ is there such a configuration of Class II. The tension $`T2\mathrm{\Delta }W`$ of the Class I configurations is equal to $`T=4\pi `$ while the tension of the Class II configuration is equal to $`T=8\pi `$. When $`L`$ is very large, the Class IA and IB configuration describe just one “right” wall and one “left” wall, respectively. The Class II configuration contains one “right” wall and one “left” wall at a distance of $`L/2`$.
Similar considerations apply to configurations with higher winding number $`N`$. The tension of the BPS configurations is proportional to $`N`$. As long as $`L`$ is much larger than $`N`$ times the width of the walls, the configurations contain well separated equidistant domain walls, $`N`$ “right” (“left”) walls in Class IA (Class IB) configurations, and $`N`$ “left” walls alternating with $`N`$ “right” walls in Class II configurations.
### 1.4 Non-BPS domain walls
As explained in the previous section there is no Class II BPS saturated configuration if the circumference $`L`$ of the compact dimension is smaller than $`\lambda _m`$. However, there is a stable, static solution to the second order equations of motion with the same topology. Such a solution can not be BPS saturated. It is a non-contractable cycle, so it can not be shrunk to a point. At the same time it has to remain at finite distance of the origin in the complex $`\varphi `$ plane because for fixed circumference $`L`$ the kinetic energy grows if the length of the trajectory in the complex $`\varphi `$ plane increases.
Following the same argument, it is easy to see that there are stable, static configurations that are not BPS saturated in all other topological sectors. The simplest example is the homotopy class of non-contractable cycles which wind first clockwise around one pole and then counterclockwise around the other pole (”race-track” topology). With obvious notation, we will refer to these topologies as $`RL^1`$ and $`LR^1`$.
Summarizing, in any non-trivial topological sector and for any value of $`L`$ there is at least one static, stable solution to the second order equation of motion with minimal energy in that sector, a classical ground state. In certain topological sectors, those with the topologies $`R^N`$, $`L^N`$, $`(R^1)^N`$ and $`(L^1)^N`$, the stable solution is BPS saturated for any value of $`L`$, and only half of supersymmetry is broken in the classical ground state. In other topological sectors, those with the topologies $`(RL)^N`$ and $`(R^1L^1)^N`$, there are two stable solutions that are BPS saturated if the circumference $`L`$ of the compact dimension is larger than a minimum value, $`N\lambda _m`$. If $`L`$ is smaller than the minimum value, then there is still a static solution to the equation of motion, but it is not BPS saturated.
## 2 d=3+1; localized light modes
In this section we study the massless and light modes in a BPS saturated background that takes the form of a sequence of equidistant domain walls. We have discussed configurations of this type in Section 1. As only half of supersymmetry is spontaneously broken, the theory of the massless and light modes in $`d=2+1`$ dimensions has $`𝒩=1`$ supersymmetry. We will therefore focus on the study of bosonic modes; clearly, each bosonic mode is accompanied by a fermionic mode that completes the supersymmetry multiplet.
There is an exact bosonic zero mode originating from the invariance of the action under translations along the compact dimension of the BPS saturated background as a whole. We will now study the spectrum of the bosonic pseudo zero modes corresponding to independent translations of the domain walls<sup>3</sup><sup>3</sup>3For large separation between the walls the action is nearly invariant under independent translations of individual walls..
### 2.1 Tight binding approach
In oder to get a more quantitative picture of the bosonic pseudo zero modes, we follow Ref. and expand the action in fluctuations of the fields around the classical, BPS saturated background configuration up to quadratic order. In contrast to the case of a real wall background as it was investigated in Ref. , we have to keep track of the fact that here the BPS saturated wall $`\varphi _w`$ is complex. Writing $`\varphi =\varphi _w+\varphi ^{}`$ and $`\varphi ^{}=\frac{1}{\sqrt{2}}(\varphi _1^{}+ı\varphi _2^{})`$ with real fields $`\varphi _1^{}`$ and $`\varphi _2^{}`$, the following equations of motion are obtained
$`\left(_\mu ^\mu +\overline{W}^{\prime \prime }W^{\prime \prime }+\text{Re}\left[\overline{W}^{}W^{\prime \prime \prime }\right]\right)\varphi _1^{}\text{Im}\left[\overline{W}^{}W^{\prime \prime \prime }\right]\varphi _2^{}`$ $`=`$ $`0,`$
$`\left(_\mu ^\mu +\overline{W}^{\prime \prime }W^{\prime \prime }\text{Re}\left[\overline{W}^{}W^{\prime \prime \prime }\right]\right)\varphi _2^{}\text{Im}\left[\overline{W}^{}W^{\prime \prime \prime }\right]\varphi _1^{}`$ $`=`$ $`0.`$ (13)
The primes denote derivatives with respect to $`\varphi `$ ($`\overline{\varphi }`$) in the case of $`W`$ ($`\overline{W}`$) evaluated at $`\varphi =\varphi _w`$ ($`\overline{\varphi }=\overline{\varphi }_w`$). For localized massless and massive modes we make the following decomposition
$$\varphi _1^{}=\chi _1(z)A(t,x,y),\varphi _2^{}=\chi _2(z)A(t,x,y),$$
(14)
where $`\chi _i`$ and $`A`$ are real fields. From the equation of motion, Eq.(2.1), we obtain the following equations of motion for these fields
$`{\displaystyle \frac{_n^nA}{A}}={\displaystyle \frac{\left(_z^2+\overline{W}^{\prime \prime }W^{\prime \prime }+\text{Re}\left[\overline{W}^{}W^{\prime \prime \prime }\right]\right)\chi _1\text{Im}\left[\overline{W}^{}W^{\prime \prime \prime }\right]\chi _2}{\chi _1}}m^2,`$
$`{\displaystyle \frac{_n^nA}{A}}={\displaystyle \frac{\left(_z^2+\overline{W}^{\prime \prime }W^{\prime \prime }\text{Re}\left[\overline{W}^{}W^{\prime \prime \prime }\right]\right)\chi _2\text{Im}\left[\overline{W}^{}W^{\prime \prime \prime }\right]\chi _1}{\chi _2}}m^2,`$ (15)
where $`n=0,1,2`$ is a Lorentz index in $`d=2+1`$ dimensions. Hence, the field $`A(t,x,y)`$ has a mass $`m`$ in $`d=2+1`$ dimensions. The vector
$$\stackrel{}{\chi }\left(\begin{array}{c}\chi _1\\ \chi _2\end{array}\right),$$
(16)
satisfies the eigenvalue equation
$$H\stackrel{}{\chi }=m^2\stackrel{}{\chi },$$
(17)
where the hermitian operator $`H`$ is defined as
$$H=\left(\begin{array}{cc}_z^2+\overline{W}^{\prime \prime }W^{\prime \prime }+\text{Re}\left[\overline{W}^{}W^{\prime \prime \prime }\right]& \text{Im}\left[\overline{W}^{}W^{\prime \prime \prime }\right]\\ \text{Im}\left[\overline{W}^{}W^{\prime \prime \prime }\right]& _z^2+\overline{W}^{\prime \prime }W^{\prime \prime }\text{Re}\left[\overline{W}^{}W^{\prime \prime \prime }\right]\end{array}\right).$$
(18)
This equation has the form of a one-dimensional Schrödinger equation, and, in analogy, we will refer to $`H`$ as the Hamiltonian and to $`\stackrel{}{\chi }`$ as the wave function. Using the BPS equation, Eq. (7), it is not difficult to explicitly show that $`H`$ annihilates $`\chi =_z\varphi _w`$, which corresponds to the Goldstone boson of spontaneously broken overall translational invariance along the $`\widehat{z}`$ direction.
Since we are not able to solve the eigenvalue problem in Eq.(17) exactly, we resort to the tight binding approximation to study the masses of the pseudo zero modes corresponding to independent translations of the individual walls. For definiteness, we take the $`R^N`$ topological sector, with the background of $`N`$ walls separated by a distance $`\lambda =L/N`$. Considerations for the $`(RL)^N`$ sectors are similar, but slightly more complicated due to the sequence of alternating walls and anti-walls. We will uncover the light bosonic modes of the wall lattice by first considering a massless mode bound to a single wall in isolation, just as one can find the electronic levels of a one dimensional crystal from the energy levels of the individual atoms when they are far apart and overlaps are small.
Thus, we consider a single wall, $`\varphi _c^w`$, along an infinite dimension. In this limit $`\chi ^0\frac{1}{\sqrt{2\pi }}_z\varphi _c^w`$ is a normalized eigen function with $`m^2=0`$. In the case of $`N`$ such walls separated by an infinite distance, there are $`N`$ localized (bound) wave functions with vanishing mass. However, when the walls are at a large, but finite, distance from each other, the $`N`$-fold degeneracy is lifted due to the small overlap of potentials and wave functions centered at different wall sites, and a band structure emerges. We will now elucidate this band structure, and we will estimate the width of the band.
To this end we decompose the operator $`H^c=H^0+V^c`$, where
$$H^0=\left(\begin{array}{cc}_z^2+4& 0\\ 0& _z^2+4\end{array}\right),$$
(19)
and
$$V^c=\left(\begin{array}{cc}\overline{W}_c^{\prime \prime }W_c^{\prime \prime }4+\text{Re}\left[\overline{W}_c^{}W_c^{\prime \prime \prime }\right]& \text{Im}\left[\overline{W}_c^{}W_c^{\prime \prime \prime }\right]\\ \text{Im}\left[\overline{W}_c^{}W_c^{\prime \prime \prime }\right]& \overline{W}_c^{\prime \prime }W_c^{\prime \prime }4\text{Re}\left[\overline{W}_c^{}W_c^{\prime \prime \prime }\right]\end{array}\right).$$
(20)
Here the superscript $`c`$ indicates that the quantity is evaluated at $`\varphi =\varphi _c^w`$. With this definition, $`V^c`$ vanishes exponentially fast away from the wall. We approximate $`H`$ in the case of $`N`$ walls separated by a distance $`\lambda `$ as
$$H=H^0+\underset{n}{}V^c(zn\lambda ).$$
(21)
This Hamiltonian commutes with the generator of the $`R_N`$ symmetry, the translations over distances $`n\lambda `$. The (pseudo zero mode) wave functions, the equivalent of Bloch waves, are linear combinations of $`\stackrel{}{\chi ^0}`$ centered at each of the $`N`$ wall sites. These linear combinations transform under irreducible representations of $`R_N`$. They take the form
$$\stackrel{}{\chi }_k=\underset{n}{}a_k\left[n\right]\stackrel{}{\chi }^0(zn\lambda ),$$
(22)
where $`\stackrel{}{a}_k`$, with $`k=1..N`$, are the $`N`$ real, orthonormal eigenvectors of the $`N\times N`$ matrix
$$\left(\begin{array}{ccccccccc}2& 1& 0& .& .& .& .& 0& 1\\ 1& 2& 1& 0& .& .& .& .& 0\\ 0& 1& 2& 1& 0& .& .& .& 0\\ .& .& .& .& .& .& .& .& .\\ 0& .& .& .& .& 0& 1& 2& 1\\ 1& 0& .& .& .& .& 0& 1& 2\end{array}\right).$$
(23)
This matrix has a non-degenerate eigenvalue equal to zero for any value of $`N`$, with eigenvector $`\stackrel{}{a}_1=\frac{1}{\sqrt{N}}(1,1,..,1)`$. For even values of $`N`$ there is a second non-degenerate eigenvalue equal to $`4`$ with eigen vector $`\stackrel{}{a}_2=\frac{1}{\sqrt{N}}(1,1,1,..,1,1)`$. All other eigenvalues are between $`0`$ and $`4`$ and are doubly degenerate. Therefore, the $`N`$ dimensional representation of $`R_N`$ is reduced to $`2`$ one-dimensional and $`N/21`$ two-dimensional irreducible representations in case $`N`$ is even, and to $`1`$ one-dimensional and $`(N1)/2`$ two-dimensional irreducible representations in case $`N`$ is odd.
As an aside, the eigenvectors and eigenvalues of the matrix in Eq.(23) also yield the normal modes and frequencies of a closed linear chain of $`N`$ identical masses and springs. In $`d=3+1`$ dimensions, the walls extend into two infinite dimensions and have therefore infinite mass. However, for solitons in $`d=1+1`$ dimensions, which have finite energy, these vibrations are physical. If the distance between the solitons is large as compared to their width, these vibrations are very soft.
Given the wave functions in Eq.(22), the approximate mass eigenvalues are then given by
$$m_{TB}^2=\frac{_0^L𝑑z\stackrel{}{\chi }^TH\stackrel{}{\chi }}{_0^L𝑑z\stackrel{}{\chi }^T\stackrel{}{\chi }}.$$
(24)
In fact, for the wave functions that transform under the two dimensional representation of $`R_N`$, degenerate perturbation theory should be used. Including only overlap integrals from neighboring sites, and using the asymptotic form for $`\stackrel{}{\chi }^c`$ far away from the wall, we estimate that the masses are exponentially suppressed,
$$m_{TB}^2=Ce^{2\lambda },$$
(25)
where $`C`$ is of the order one. We thus see that the $`N`$ degenerate vanishing masses in the limit where the walls are infinitely far apart evolve into a (mostly double degenerate) band of masses when the distance between the walls becomes finite.
### 2.2 Oscillations
For an observer who can not resolve the circumference $`L`$ of the compact dimension, the physics of the light modes is described by an $`𝒩=1`$ supersymmetric theory in $`d=2+1`$ dimensions with a spectrum as we discussed in the previous section. In such a world there is a hierarchy between the heavy modes, with masses comparable to the wall tension, and the light modes, with masses that are exponentially suppressed by the ratio of the distance between the walls and their width. It is interesting to note that in topological sectors with multi-wall backgrounds that are not BPS saturated the physics of the light modes is described by a $`𝒩=1`$ model with supersymmetry breaking. The scale of the breaking is equal to the scale of the light masses. It seems that under such conditions there are no hierarchy or naturalness problems.
Alternatively, an observer, who is confined to a particular wall and can not resolve its width, experiences a different physical reality. Such an observer does not have the means to do experiments in which the massive bulk modes are excited. Still, due to the BPS saturation of the background, his world has $`𝒩=1`$ supersymmetry. But he will be perplexed by the appearance and disappearance of particles in his world. From the ansatz of Eq. (22) it is clear that a state localized on a particular wall must be a superposition of the wave functions that represent the light modes. This in turn means that such a state is not a mass eigenstate. The implication is oscillation of localized particles between the various walls with frequencies corresponding to the mass differences. In the analogy of a linear crystal, this would correspond to electrons hopping from atom to atom. For example, in the case of the $`R^2`$ sector with two walls, the approximate massless mode is
$$\stackrel{}{\chi }_1(z)=\stackrel{}{\chi }^0(z)+\stackrel{}{\chi }^0(z\lambda ),$$
(26)
and there is one light mode, approximately given by
$$\stackrel{}{\chi }_2(z)=\stackrel{}{\chi }^0(z)\stackrel{}{\chi }^0(z\lambda ).$$
(27)
A state localized one of the walls is constructed by the superposition of the mass eigenstates,
$$\stackrel{}{\chi }_0(z)\pm \stackrel{}{\chi }_1(z).$$
(28)
## 3 d=2+1; chiral fermions
When the model is dimensionally reduced to $`d=2+1`$ dimensions, it has $`𝒩=2`$ supersymmetry. It has inherited four supersymmetry generators, whereas minimal supersymmetry in $`d=2+1`$ requires only two generators. The scalar potential may be written as
$$V=\underset{i=1}{\overset{2}{}}\left(\frac{𝒲}{\varphi _i}\right)^2.$$
(29)
The object $`𝒲`$ plays the role of the superpotential in the dimensionally reduced model. It is related to the superpotential $`W`$ in $`d=3+1`$ dimensions by
$$𝒲=2\mathrm{I}\mathrm{m}\left[W\left(\frac{1}{\sqrt{2}}(\varphi _1+ı\varphi _2)\right)\right]=2\mathrm{arctan}\frac{\varphi _2}{\varphi _1+\sqrt{2}}+2\mathrm{arctan}\frac{\varphi _2}{\varphi _1\sqrt{2}}.$$
(30)
The extended supersymmetry manifests itself in the fact that the superpotential $`𝒲`$ is harmonic. In terms of the scalar components of the $`𝒩=1`$ superfields in $`d=2+1`$ dimensions, $`\varphi _1`$ and $`\varphi _2`$, the BPS equations are
$`{\displaystyle \frac{d\varphi _1}{dx}}`$ $`=`$ $`\alpha {\displaystyle \frac{𝒲}{\varphi _1}}=2\alpha \left\{{\displaystyle \frac{\varphi _2}{(\varphi _1+\sqrt{2})^2+\varphi _2^2}}{\displaystyle \frac{\varphi _2}{(\varphi _1\sqrt{2})^2+\varphi _2^2}}\right\},`$
$`{\displaystyle \frac{d\varphi _2}{dx}}`$ $`=`$ $`\alpha {\displaystyle \frac{𝒲}{\varphi _2}}=2\alpha \left\{{\displaystyle \frac{\varphi _1+\sqrt{2}}{(\varphi _1+\sqrt{2})^2+\varphi _2^2}}+{\displaystyle \frac{\varphi _1\sqrt{2}}{(\varphi _1\sqrt{2})^2+\varphi _2^2}}\right\},`$ (31)
where $`\alpha =\pm 1`$. The interesting new feature in $`d=2+1`$ dimensions is that now “chiral” fermions exist on the $`d=1+1`$ dimensional worlds of the domain walls. In $`d=1+1`$ dimensions there is no spin, but massless fermions can be either left or right moving; this is what we will allude to as “chirality”. As outlined in the Appendix, in a BPS saturated background configuration there is always one left moving and one right moving fermionic zero mode, corresponding to the two spontaneously broken supersymmetry generators. Whether a mode that is associated with a particular broken generator is left or right moving depends on the sign $`\alpha `$ in the BPS equation that provides the background configuration. In other words, it is not important around which pole the background configuration winds in the complex $`\varphi `$ plane, but whether it winds clockwise or counter clockwise.
### 3.1 Wess-Zumino model
In order to further investigate the issue of the “chirality” of massless fermions on the wall, we discuss in this section the regular Wess-Zumino model reduced to $`d=2+1`$ dimensions. It is well known that this model allows for real, BPS saturated domain walls (kink or anti-kink). However, these solutions do not exist in a compact dimension, and there are no BPS saturated multi-wall solutions (although almost static configurations with walls at large distances do exist). For the issues we want to discuss here, these differences are irrelevant.
Our analysis is similar to the discussion of fermionic zero modes in Ref., where the same model was investigated in $`d=1+1`$ dimensions. The superpotential of the dimensionally reduced, regular Wess-Zumino model is
$$𝒲=\frac{1}{\sqrt{2}}\left(2\varphi _1\frac{1}{3}\varphi _1^3+\varphi _1\varphi _2^2\right).$$
(32)
There are two vacua, $`\varphi _1=\pm \sqrt{2}`$ with $`\varphi _2=0`$. The BPS saturated domain walls connecting these vacua are determined by
$`{\displaystyle \frac{d\varphi _1}{dy}}`$ $`=`$ $`\alpha {\displaystyle \frac{1}{\sqrt{2}}}\left(2\varphi _1^2+\varphi _2^2\right),`$
$`{\displaystyle \frac{d\varphi _2}{dy}}`$ $`=`$ $`\alpha {\displaystyle \frac{1}{\sqrt{2}}}\left(2\varphi _1\varphi _2\right),`$ (33)
where $`\alpha =\pm 1`$. The solutions that interpolate between the two vacua (the well-known kink or anti-kink, depending on the sign $`\alpha `$), take the form
$`\varphi _1^w`$ $`=`$ $`\alpha \sqrt{2}\mathrm{tanh}(xx_0)`$
$`\varphi _2^w`$ $`=`$ $`0.`$ (34)
The equation of motion for the fermions in the background is
$$ı\gamma ^m_m\psi _i𝐖_{ij}\psi _j=0,$$
(35)
where the “mass” matrix $`𝐖`$ is given as
$$𝐖=\sqrt{2}\left(\begin{array}{cc}\varphi _1^w& \varphi _2^w\\ \varphi _2^w& \varphi _1^w\end{array}\right).$$
(36)
There are two massless fermions on the wall corresponding to Rebbi–Jackiw zero modes,
$`\psi _1(x,y,t)`$ $`=`$ $`A(y)\zeta (x,t),`$
$`\psi _2(x,y,t)`$ $`=`$ $`0,`$ (37)
with
$`A(y)`$ $`=`$ $`e^{\alpha _{y_0}^y\mathrm{W}_{11}(y^{})𝑑y^{}},`$
$`ı\gamma ^2\zeta `$ $`=`$ $`\alpha \zeta ,`$
and
$`\psi _1(x,y,t)`$ $`=`$ $`0,`$
$`\psi _2(x,y,t)`$ $`=`$ $`B(y)\eta (x,t),`$ (38)
where
$`B(y)`$ $`=`$ $`e^{\alpha _{y_0}^y\mathrm{W}_{22}(y^{})𝑑y^{}},`$
$`ı\gamma ^2\eta `$ $`=`$ $`\alpha \eta .`$
Both $`\zeta `$ and $`\eta `$ satisfy the massless Dirac equation in $`d=1+1`$ dimensions, $`ı\gamma ^a_a\zeta =0`$, where $`a=0,1`$. Note that $`\gamma ^2`$ is the chiral projection matrix in $`d=1+1`$ dimensions, so that one of the massless fermions is right moving and the other one is left moving. These massless fermions are the goldstino fermions associated with the breaking of two out of four generators of the $`𝒩=2`$ supersymmetry, as discussed in the Appendix.
We proceed by addressing the fate of these massless fermions when the $`𝒩=2`$ supersymmetry is explicitly broken to $`𝒩=1`$. In this case, the wall breaks one of the remaining two supersymmetry generators, and therefore only one Goldstone fermion exists. What happens to the second massless fermion when the explicit supersymmetry breaking is switched on? To answer this question we introduce explicit $`𝒩=2`$ breaking terms in the superpotential.
First, we insert a coupling constant $`\lambda `$ in front of the $`\varphi _1\varphi _2^2`$ term in the superpotential, so that the $`𝒩=2`$ supersymmetry is explicitly broken unless $`\lambda =1`$. The vacuum expectation values and the wall solution do not depend on $`\lambda `$, but $`\mathrm{W}_{22}`$ now is equal to $`\lambda \varphi _1^w`$. Even though there is only one Goldstino, from the perspective of the Rebbi-Jackiw zero modes in Eqs. (37,38), it is clear that the second massless fermion continues to exist, even if $`\lambda 1`$. When $`\lambda `$ changes sign, the second massless fermion changes from right to left moving (or vice versa, depending on the sign $`\alpha `$), and for $`\lambda =0`$ the second massless fermionic mode is not normalizable.
Next, we add a term $`\frac{1}{2}\mu \varphi _2^2`$ to the superpotential, so that the $`𝒩=2`$ supersymmetry is explicitly broken to $`𝒩=1`$ for $`\mu 0`$. Again, the vacuum expectation values and the wall do not depend on $`\mu `$, but now $`\mathrm{W}_{22}=\varphi _1^w+\mu `$. In this case the second massless fermion persists as long as $`|\mu |<\sqrt{2}`$. For $`|\mu |>\sqrt{2}`$ the second zero mode becomes non-normalizable.
It is clear that since the fermionic modes on the wall are chiral, they can only become massive in pairs, combining one left moving and with one right moving fermion. Individual massless modes, either left or right moving, can only disappear when they become non-normalizable.
### 3.2 Complex, multi-wall backgrounds
In the case of a BPS saturated background with one domain wall along a compact dimension, there is one massless scalar localized on the wall corresponding to the spontaneously broken translational symmetry. Moreover, there is one left moving and one right moving massless fermion, corresponding to the two spontaneously broken supersymmetry generators (see the Appendix). These massless modes combine to form a massless superfield in $`d=1+1`$ dimensions.
The supersymmetry multiplet in $`d=1+1`$ dimensions contains a real scalar and a Majorana fermion. Therefore, in the case of $`N`$ walls, there will be a total of $`N1`$ light superfields ($`d=1+1`$, $`𝒩=1`$ supersymmetry) with exponentially suppressed masses, and one massless superfield. In order to make states localized on a particular wall, states with different masses have to be superimposed, just as in the $`d=3+1`$ dimensional case.
The explicit analysis of the Rebbi–Jackiw modes becomes impossible if the fermion “mass” matrix in the background is not diagonal, as is the case with the type of models we consider in this paper. But the results of the previous section must still be valid. It is therefore expected that if the $`𝒩=2`$ supersymmetry is explicitly broken to $`𝒩=1`$, there will still be a left and a right moving fermion on the wall, at least for a finite range of the breaking parameters. Moreover, in a compact dimension, it seems unlikely that fermionic zero modes can disappear by becoming non-normalizable.
## 4 d=1+1; solitons
The main new feature in $`d=1+1`$ dimensions is that the domain walls, which have infinite energy in higher dimensions, now reduce to solitons with a finite mass. Such solitons have a particle interpretation, whereas the domain wall backgrounds in the higher dimensions are interpreted as vacua.
All topological sectors in the $`d=1+1`$ dimensional theory, whether they have a BPS saturated ground state or not, can be assigned a pair of winding numbers $`(N_l,N_r)`$, where $`N_l𝐙`$ is the winding number around the “left” pole at $`\varphi =1`$ in Fig.(1), and $`N_r𝐙`$ is the winding number around the “right” pole at $`\varphi =+1`$. According to our conventions the winding number increases by one if a pole is encircled in the counter clockwise direction. This classification of the sectors according to the pair of winding numbers is not unique; for example, the sectors $`RLRL^1`$ and $`R^2`$ have the same winding numbers, $`(0,2)`$, but they are topologically distinct.
The sectors $`L^{N_l}`$ and $`R^{N_r}`$ have a classical ground state in which only half of supersymmetry is broken for any value of $`L`$. In the $`L^{N_l}`$ sector this ground state corresponds to the solution of the BPS equation which in the large $`L`$ limit contains $`N_l`$ “left” solitons if $`N_l`$ is positive, or $`|N_l|`$ anti “left” solitons if $`N_l`$ is negative, separated by a distance $`L/|N_l|`$. The situation is analogous for the $`R^{N_r}`$ sectors.
For fixed $`L>\lambda _m`$, the sectors $`(RL)^N`$ have two degenerate classical ground states in which only half of supersymmetry is broken. In the large $`L`$ limit, and if $`N`$ is positive, one of these ground states corresponds to the solution of the BPS equation containing $`N`$ “left” solitons which alternate with $`N`$ “right” solitons, separated by a distance $`L/2N`$. Similarly, if $`N`$ is negative we have the BPS solution with $`|N|`$ “anti-left” solitons alternating with $`|N|`$ “anti-right” solitons separated by a distance $`L/2N`$. The second, degenerated ground state corresponds to the BPS solution that winds around both poles at large distance, $`N`$ times in the counterclockwise direction if $`N`$ is positive, and $`|N|`$ times in the clockwise direction if $`N`$ is negative. This solution to the BPS equation has approximately constant energy density in the large $`L`$ limit. There is a tunneling transition with finite action that connects the two ground states and the remaining $`1/2`$ of supersymmetry is broken non-perturbatively. For $`L<N\lambda _m`$, there is no BPS saturated ground state in the sector $`(RL)^N`$ already at the classical level. There are static, stable solutions to the second order equation of motion which break supersymmetry completely at the classical level.
For the same reason, supersymmetry is broken completely in the groundstate at the classical level in all sectors other than the ones we discussed above. In all sectors, the ground state energy tends to $`2\pi (|N_l|+|N_r|)`$ in the large $`L`$ limit. The theory has two types of solitons, “left” and “right”, and two topologically conserved quantum numbers, $`N_l`$ and $`N_r`$. The solitons interact through a short range force, which decreases exponentially with distance. “Left” and “right” solitons repel each other, and so do “anti-left” and “anti-right” solitons. “Left” and “anti-left” solitons on the other hand attract each other.
### 4.1 Semi-classical calculation of the tunneling amplitude
In this section we calculate the amount by which the ground state energy is lifted above the BPS bound by the tunneling transition in the sector $`(RL)^N`$, for $`N=1`$ and $`L>\lambda _m`$. We choose $`N=1`$ for definiteness; the tunneling action scales linearly with $`|N|`$, and therefore the result can be trivially generalized for any value of $`N`$. The two ground states correspond to two solutions of the BPS equation in Class II, one with $`k<k_m`$ and one with $`k>k_m`$, as shown in Fig.(5). At the non-perturbative level the corresponding solitons are not longer BPS saturated because a tunneling process mixes them and lifts their mass above the BPS bound. This phenomenon was first observed in Ref.. In Ref. it was noted that the calculation of the semi-classical instanton action is equivalent to the calculation of the energy of domain wall junctions in $`d=2+1`$ dimensions. In the same work, the general calculational framework of the action was established and illustrated for a specific model.
We assume that the instanton in the imaginary time saturates the BPS bound in the model we consider; there is no obvious reason why this would not be the case. If the instanton is BPS saturated, it satisfies the equation
$$\frac{\varphi }{\zeta }=\frac{1}{2}e^{i\delta }\frac{d\overline{W}}{d\overline{\varphi }},$$
(39)
where $`\zeta =x+it`$, and $`t`$ is the euclidean time. The field $`\varphi `$ depends both on $`\zeta `$ and $`\overline{\zeta }`$. The phase $`\delta `$ has to be chosen so that it is consistent with the boundary conditions that are imposed on the solution (see below). The tunneling action takes the form
$`A`$ $`=`$ $`2{\displaystyle 𝑑x𝑑t\left\{\frac{\varphi }{\zeta }\frac{\overline{\varphi }}{\overline{\zeta }}+\frac{\varphi }{\overline{\zeta }}\frac{\overline{\varphi }}{\zeta }\right\}}`$ (40)
$`+{\displaystyle 𝑑x𝑑t\frac{dW}{d\varphi }\frac{d\overline{W}}{d\overline{\varphi }}}.`$
In Figs.(5, 6) we have indicated the boundary conditions that are imposed on the instanton configuration. At $`t=T/2`$ the instanton field $`\varphi (x,t)`$ takes the form of the “outer” configuration (see Fig.(5)) and at $`t=T/2`$ it takes the form of the “inner” configuration. We always have in mind the limit $`T\mathrm{}`$. Periodic boundary conditions apply to the vertical edges in Fig.(6).
We have chosen the solutions to wind counter-clockwise in the $`\varphi _1`$-$`\varphi _2`$ plane when followed in the positive $`\widehat{x}`$ direction, which corresponds to $`\delta =\pi /2`$ in the BPS equation. The initial and final configuration are positioned such that $`\varphi (L/2,T/4)=iv_0`$ and $`\varphi (L/2,T/4)=iv_i`$. Note that in order for the instanton configuration to satisfy the BPS equation, it has to connect the solution with the larger value of $`k`$ as the initial configuration with the solution with the smaller value of $`k`$ as the final configuration. The anti-instanton that interpolates in the reverse situation satisfies an anti-BPS equation. The action can now be written as
$$A_{BPS}=𝑑x𝑑t\left[2\stackrel{}{}\stackrel{}{S}(\stackrel{}{}\times \stackrel{}{a})_z\right],$$
(41)
where
$$\stackrel{}{S}=\left(\begin{array}{c}\mathrm{cos}\delta \mathrm{Re}\left[W\right]+\mathrm{sin}\delta \mathrm{Im}\left[W\right]\\ \mathrm{cos}\delta \mathrm{Im}\left[W\right]\mathrm{sin}\delta \mathrm{Re}\left[W\right]\end{array}\right),$$
(42)
and
$$\stackrel{}{a}=\left(\begin{array}{c}\mathrm{Im}\left[\varphi \frac{\overline{\varphi }}{x}\right]\\ \mathrm{Im}\left[\varphi \frac{\overline{\varphi }}{t}\right]\end{array}\right).$$
(43)
With the application of Gauss’ and Stoke’s theorems the surface integrals can be converted into line integrals over the boundaries of the surface,
$$A_{BPS}=2\stackrel{}{S}𝑑\stackrel{}{n}\stackrel{}{a}𝑑\stackrel{}{x},$$
(44)
where the boundaries are the edges of the rectangle in Fig.(6). Due to the periodic boundary conditions only the integrals of $`\stackrel{}{S}`$ contribute over the vertical edges of the contour, as the superpotential is multi-valued. Turning now to the explicit calculation of the tunneling action, there are three non-vanishing contributions. First, consider
$`A_{BPS}^0`$ $`=`$ $`2{\displaystyle _{T/2}^{T/2}}\left\{\mathrm{Im}\left[W\right]_{x=\frac{L}{2}}\mathrm{Im}\left[W\right]_{x=\frac{L}{2}}\right\}𝑑t,`$ (45)
$`=`$ $`8\pi T.`$
This contribution is equal to the classical ground state energy times the elapsed time $`T`$. It is equal to the action in case there is no tunneling, when the initial and final configuration are the same. We define the tunneling probability by subtracting this trivial contribution in the exponent. Next, consider
$`A_{BPS}^1`$ $`=`$ $`2{\displaystyle _{L/2}^{L/2}}\left\{\mathrm{Re}\left[W\right]_{t=\frac{T}{2}}\mathrm{Re}\left[W\right]_{t=\frac{T}{2}}\right\}𝑑x`$ (46)
$`=`$ $`\mathrm{ln}\left({\displaystyle \frac{k_0}{k_i}}\right)L,`$
where $`k_o`$ and $`k_i`$ are the two solutions to the equation $`\lambda (k)=L`$ with $`k>k_c`$. We define $`k_o`$ to be the larger of the two. Note that the integrands are “constants of motion” for the solutions to the BPS equation, and therefore the integrals can be expressed in $`k_o`$ and $`k_i`$, respectively. In the limit $`\lambda =L\lambda _m`$, $`k_i`$ approaches $`1`$ from above, and $`k_o`$ is related to $`L`$ by $`L=\pi \sqrt{k}`$. In this limit, the second contribution to the tunneling action is therefore $`A_{BPS}^1=2L\mathrm{ln}L/\pi `$. Finally, we calculate
$`A_{BPS}^2`$ $`=`$ $`{\displaystyle \left\{\mathrm{Im}\left[\varphi \frac{\overline{\varphi }}{x}\right]_{t=\frac{T}{2}}\mathrm{Im}\left[\varphi \frac{\overline{\varphi }}{x}\right]_{t=\frac{T}{2}}\right\}𝑑x}`$ (47)
$`=`$ $`S(k_o)+S(k_i),`$
where $`S(k)`$ is the area enclosed in the $`\varphi _1`$-$`\varphi _2`$ plane by the solution to the BPS equation marked by $`k`$. $`S(k)`$ is a monotonically increasing function of $`k`$. In the limit $`kk_c=1`$, $`S(k)`$ approaches $`S(1)=4`$ from above. For large values of $`k`$ on the other hand $`S(k)=2\pi \sqrt{k}`$. Therefore, in the limit $`\lambda =L\lambda _m`$, the third contribution to the tunneling action is approximately $`A_{BPS}^2=2\pi \sqrt{k_o}+1`$, or, in terms of $`L`$, $`A_{BPS}^2=2L+1`$.
We have calculated the tunneling action solely in terms of quantities related to the boundary configurations. An explicit interpolating instanton configuration is not needed, just the existence of a BPS saturated instanton configuration is sufficient. In the large $`L`$ limit, the semi-classical approximation is valid, and the ground state energy is lifted above the BPS bound by an amount
$$\mathrm{\Delta }E8\pi e^{2A_{BPS}}=8\pi e^{4L\mathrm{ln}L},$$
(48)
to leading order in $`L`$. The solution to the BPS equation that connects the vacuum at $`\varphi =0`$ to the runaway vacuum at $`\varphi i\mathrm{}`$ (this is a solution that is not periodic, obtained with $`\delta =0`$) has infinite tension. The circumference $`L`$ in a sense acts like an infra-red regulator. For this reason, the tunneling action is not linear in $`L`$, but proportional to $`L\mathrm{ln}L`$. We also note that in this model the instanton mixes a state containing two localized solitons with a state that corresponds to a flat energy density. This is in contrast with the model considered in Ref., where the mixing is between two states with different solitons.
## Summary
We have constructed a toy model to study some issues related to the idea of “brane worlds” in a field theoretical context. This toy model is a simple Wess-Zumino model with only one chiral superfield in $`d=3+1`$ dimensions. Different aspects were studied in various dimensions.
We showed that it is possible to have BPS saturated configurations with multiple, well separated domain walls in a compact space dimension. Generically, in such a background there are light modes with masses exponentially suppressed by the ratio of the distance between the walls and their width. These modes form the “matter” of the “brane world”, wheras the domain walls play the role of branes in our simple picture. We also described how particles can oscillate between walls. We illustrated these issues in our specific model.
We also found that the “chiral” fermions localized on $`d=1+1`$ dimensional BPS saturated walls embedded in a $`d=2+1`$ dimensional space-time always come in pairs; one is left moving, and one is right moving.
Finally, we discussed the toy model in $`d=1+1`$ dimensions. The domain walls then reduce to solitons with a finite mass. We showed that some classically BPS saturated states are lifted non-perturbatively by instanton tunneling.
## Acknowledgements
The authors would like to thank M. Shifman for valuable discussions. T.t.V also acknowledges useful conversations with A. Losev and A. Vainshtein. The work of R.H. was supported by a postdoctoral fellowship of Deutscher Akademischer Austauschdienst (DAAD), the work of T.t.V. by the Department of Energy under Grant No. DE-FG02-94ER40823.
## Appendix A
In this Appendix we establish our notation and discuss the zero modes in a BPS saturated background.
### d=3+1 dimensions
We use the conventions outlined in . The signature of the metric $`g^{\mu \nu }`$ is $`(+,,,)`$. The component Lagrangian for an $`𝒩=1`$ Wess-Zumino model in $`d=3+1`$ dimensions with canonical Kähler potential and superpotential $`W`$ is
$``$ $`=`$ $`_\mu \overline{\varphi }^\mu \varphi +ı\psi ^\alpha \sigma _{\alpha \dot{\alpha }}^\mu _\mu \overline{\psi }^{\dot{\alpha }}+F^{}F`$ (49)
$`+F{\displaystyle \frac{dW}{d\varphi }}{\displaystyle \frac{1}{2}}{\displaystyle \frac{d^2W}{d\varphi ^2}}\psi ^\alpha \psi _\alpha `$
$`+\overline{F}{\displaystyle \frac{d\overline{W}}{d\overline{\varphi }}}{\displaystyle \frac{1}{2}}{\displaystyle \frac{d^2\overline{W}}{d^2\overline{\varphi }}}\overline{\psi }_{\dot{\alpha }}\overline{\psi }^{\dot{\alpha }}.`$
The component fields transform under the supersymmetry transformations as
$`\delta \varphi `$ $`=`$ $`\zeta ^\alpha \psi _\alpha ,`$
$`\delta \psi _\alpha `$ $`=`$ $`ı\sigma _{\alpha \dot{\alpha }}^\mu \overline{\zeta }^{\dot{\alpha }}_\mu \varphi +\zeta _\alpha F,`$
$`\delta F`$ $`=`$ $`ı\overline{\zeta }_{\dot{\alpha }}\overline{\sigma }^{\mu \dot{\alpha }\alpha }_\mu \psi _\alpha ,`$ (50)
where the transformation parameter $`\zeta `$ is the two component complex Grassmann variable. The auxiliary field $`F`$ can be eliminated by the use of its equation of motion $`\overline{F}=\frac{d\overline{W}}{d\overline{\varphi }}`$ to give the Lagrangian density
$``$ $`=`$ $`_\mu \overline{\varphi }^\mu \varphi +ı\psi ^\alpha \sigma _{\alpha \dot{\alpha }}^\mu _\mu \overline{\psi }^{\dot{\alpha }}{\displaystyle \frac{dW}{d\varphi }}{\displaystyle \frac{d\overline{W}}{d\overline{\varphi }}}`$ (51)
$`{\displaystyle \frac{1}{2}}{\displaystyle \frac{d^2W}{d\varphi ^2}}\psi ^\alpha \psi _\alpha {\displaystyle \frac{1}{2}}{\displaystyle \frac{d^2\overline{W}}{d\overline{\varphi }^2}}\overline{\psi }_{\dot{\alpha }}\overline{\psi }^{\dot{\alpha }}.`$
The BPS equation for domain walls reads
$$\frac{d\varphi }{dz}=e^{ı\delta }\frac{d\overline{W}}{d\overline{\varphi }},$$
(52)
where in a non-compact direction the phase $`\delta `$ is determined from the equation
$$\mathrm{Im}[e^{ı\delta }W(\varphi (z=\mathrm{})]=\mathrm{Im}[e^{ı\delta }W(\varphi (z=+\mathrm{})].$$
(53)
BPS domain walls in a compact dimension are only possible if the superpotential is not single valued. In the models we consider $`e^{ı\delta }\pm i`$.
Let us now expand the Lagrangian of the model about the non-trivial solution $`\varphi _w`$ of the BPS equation. Setting $`\varphi =\varphi _w+\varphi ^{}`$ and expanding up to quadratic terms in the fluctuations $`\varphi ^{}`$, we obtain
$``$ $`=`$ $`_\mu \overline{\varphi }^\mu \varphi +ı\psi ^\alpha \sigma _{\alpha \dot{\alpha }}^\mu _\mu \overline{\psi }^{\dot{\alpha }}`$ (54)
$`{\displaystyle \frac{1}{2}}\left\{W^{\prime \prime \prime }\overline{W}^{}\varphi ^2+W^{}\overline{W}^{\prime \prime \prime }\overline{\varphi }^2+2W^{\prime \prime }\overline{W}^{\prime \prime }\varphi ^{}\overline{\varphi }^{}+W^{\prime \prime }\psi ^\alpha \psi _\alpha +\overline{W}^{\prime \prime }\overline{\psi }_{\dot{\alpha }}\overline{\psi }^{\dot{\alpha }}\right\},`$
where the primes on $`W`$ ($`\overline{W}`$) indicate derivatives with respect to $`\varphi `$ ($`\overline{\varphi }`$) evaluated at $`\varphi =\varphi _w`$ ($`\overline{\varphi }=\overline{\varphi }_w`$). The equation of motion for the scalar in the wall background is
$$^2\varphi ^{}+\left\{W^{}\overline{W}^{\prime \prime \prime }\overline{\varphi }^{}+W^{\prime \prime }\overline{W}^{\prime \prime }\varphi ^{}\right\}=0,$$
(55)
and the equation of motion for the fermion in the wall background is
$$ı\sigma _{\alpha \dot{\alpha }}^\mu _\mu \overline{\psi }^{\dot{\alpha }}W^{\prime \prime }\psi _\alpha =0.$$
(56)
Half of the supersymmetry generators is broken by the wall solution. In order to determine which generators are unbroken, we act with the SUSY generators on the wall and find
$$\delta \psi _\alpha =\frac{1}{\alpha }\frac{d\varphi _w}{dz}\left(\zeta _\alpha +ı\alpha \sigma _{\alpha \dot{\alpha }}^3\overline{\zeta }^{\dot{\alpha }}\right),$$
(57)
where we have used the BPS equation, and $`\alpha =\pm 1`$. Transformations, for which the parameter satisfies $`\zeta _\alpha =ı\alpha \sigma _{\alpha \dot{\alpha }}^3\overline{\zeta }^{\dot{\alpha }}`$, correspond to unbroken symmetries. The massless fermions localized on the wall which are related to the broken SUSY generators take the form
$$\psi _\alpha (x,y,z,t)=\frac{d\varphi _w}{dz}(z)\zeta _\alpha (x,y,t),$$
(58)
where $`\zeta _\alpha (x,y,t)`$ is constrained by
$$\zeta _\alpha =+ı\alpha \sigma _{\alpha \dot{\alpha }}^3\overline{\zeta }^{\dot{\alpha }}.$$
(59)
It can then be verified that $`\psi _\alpha `$ satisfies the fermionic equation of motion in the domain wall background if $`\zeta _\alpha `$ satisfies the equation
$$ı\sigma _{\alpha \dot{\alpha }}^m\overline{\sigma }^{3\dot{\alpha }\beta }_m\zeta _\beta =0,$$
(60)
where $`m`$ takes the values $`0,1,2`$. This is exactly the Dirac equation in $`d=2+1`$ dimensions,
$$ı\gamma ^m_m\zeta =0,$$
(61)
with the following representation of the $`\gamma `$ matrices; $`\gamma ^0=\sigma ^3`$, $`\gamma ^1=ı\sigma ^2`$ and $`\gamma ^2=ı\sigma _1`$. In this basis the charge conjugation matrix, defined by $`C^T=C`$ and $`C\gamma ^\mu C^1=(\gamma ^\mu )^T`$, takes the form $`C=\pm ı\sigma ^2`$. The Majorana constraint reads
$$\zeta ^TC=ı\zeta ^{}\gamma ^0.$$
(62)
This is identical to the constraint in Eq.(59). We thus find a massless Majorana fermion in the $`d=2+1`$ dimensional world on the wall. Together with the massless real scalar, this fermion forms a SUSY multiplet in the lower dimensional theory.
### d=2+1 dimensions
#### $`𝒩=1`$ SUSY
The Lagrangian of the model with $`𝒩=1`$ supersymmetry and one superfield with canonical kinetic terms and superpotential $`𝒲`$ takes the form
$``$ $`=`$ $`{\displaystyle \frac{1}{2}}\{_m\varphi ^m\varphi +ı\overline{\psi }\gamma ^m_m\psi +F^2`$ (63)
$`+2F{\displaystyle \frac{𝒲}{\varphi }}{\displaystyle \frac{^2𝒲}{\varphi ^2}}\overline{\psi }\psi \}.`$
We choose the Majorana basis for the $`\gamma `$ matrices, $`\gamma ^0=\sigma ^2`$, $`\gamma ^1=ı\sigma ^3`$ and $`\gamma ^2=ı\sigma ^1`$. Here $`\varphi `$ is a real scalar field and $`\psi `$ is a real two component spinor, with the Dirac conjugate defined by $`\overline{\psi }=\psi ^T\gamma ^0`$. The supersymmetry transformations take the form
$$\begin{array}{c}\delta \varphi =\overline{\zeta }\psi ,\hfill \\ \delta \psi =ı_m\varphi \gamma ^m\zeta +F\zeta ,\hfill \\ \delta F=ı\overline{\zeta }\gamma ^m_m\psi ,\hfill \end{array}$$
(64)
where the transformation parameter $`\zeta `$ is a two component real Grassmann variable. The BPS equation is
$$\frac{d\varphi }{dy}=\alpha \frac{d𝒲}{d\varphi },$$
(65)
where $`\alpha =\pm 1`$.
Expanding the Lagrangian about the non-trivial solution $`\varphi _w`$ to Eq. (65) up to quadratic terms in the fluctuations, we obtain
$$=\frac{1}{2}\left\{_m\varphi ^{}^m\varphi ^{}+ı\overline{\psi }\gamma ^\mu _\mu \psi \left(𝒲^{\prime \prime 2}+𝒲^{}𝒲^{\prime \prime \prime }\right)\varphi ^2𝒲^{\prime \prime }\overline{\psi }\psi \right\}.$$
(66)
The equations of motion for the scalar and the fermion are
$$^2\varphi ^{}+\left\{𝒲^{\prime \prime 2}+𝒲^{}𝒲^{\prime \prime \prime }\right\}\varphi ^{}=0,$$
(67)
and
$$ı\gamma ^m_m\psi 𝒲^{\prime \prime }\psi =0.$$
(68)
The scalar zero mode related to the broken translational symmetry in the $`\widehat{y}`$ direction is
$$\varphi ^{}(x,y,t)=\frac{d\varphi _w}{dy}(y)A(x,t),$$
(69)
where $`A(x,t)`$ satisfies the equation of motion of a massless real scalar, $`^2A=0`$, in $`d=1+1`$ dimensions. By acting with a supersymmetry generator on $`\varphi _w`$ we find
$$\delta \psi =\alpha \frac{d\varphi _w}{dy}(1\alpha ı\gamma ^2)\zeta .$$
(70)
For the unbroken generator $`\zeta =1/2(1+\alpha ı\gamma ^2)\zeta `$, whereas for the broken generator $`\zeta =1/2(1\alpha ı\gamma ^2)\zeta `$. The fermionic zero mode related to the broken supersymmetry generator is
$$\psi (x,y,t)=\frac{d\varphi _w}{dy}(y)\zeta (x,t),$$
(71)
where $`\zeta `$ satisfies the massless Dirac equation in $`d=1+1`$ dimensions $`ı\gamma ^a_a\zeta =0`$, with $`a=0,1`$, satisfies the Majorana constraint and in addition the constraint $`\zeta =1/2(1\alpha ı\gamma ^2)\zeta `$, that is, $`\zeta `$ is chiral. Consequently, the massless fermion on the wall is either left or right moving depending on the sign of $`\alpha `$.
#### $`𝒩=2`$ SUSY
The minimal model with $`𝒩=2`$ supersymmetry in $`d=2+1`$ dimensions contains two $`𝒩=1`$ supermultiplets. The Lagrangian in components is
$``$ $`=`$ $`{\displaystyle \frac{1}{2}}\{_m\varphi _i^m\varphi _i+ı\overline{\psi }_i\gamma ^m_m\psi _i+F_iF_i`$ (72)
$`+2F_i{\displaystyle \frac{𝒲}{\varphi _i}}{\displaystyle \frac{^2𝒲}{\varphi _i\varphi _j}}\overline{\psi }_i\psi _j\},`$
where the indices $`i`$, $`j`$ are $`1,2`$ and the superpotential is harmonic, so that
$$\frac{^2𝒲}{\varphi _1^2}+\frac{^2𝒲}{\varphi _2^2}=0.$$
(73)
An $`𝒩=2`$ model in $`d=2+1`$ dimensions can be obtained from an $`𝒩=1`$ model with one chiral superfield in $`d=3+1`$ dimensions. In this procedure, the fields are assumed to be independent of one of the spatial coordinates. The superpotential in the lower dimensional theory is obtained from the superpotential in the higher dimensional theory by
$$𝒲(\varphi _1,\varphi _2)=2\mathrm{R}\mathrm{e}\left[W(\frac{1}{\sqrt{2}}(\varphi _1+ı\varphi _2))\right].$$
(74)
The lower dimensional superpotential obtained in this way is automatically harmonic, because the higher dimensional superpotential is holomorphic. The $`N=2`$ model in $`d=2+1`$ dimensions has, apart from the regular supersymmetry transformations of $`N=1`$ models,
$$\begin{array}{cc}\delta \varphi _1=\overline{\zeta }\psi _1,\hfill & \delta \varphi _2=\overline{\zeta }\psi _2,\hfill \\ \delta \psi _1=ı_m\varphi _1\gamma ^m\zeta +F_1\zeta ,\hfill & \delta \psi _2=ı_m\varphi _2\gamma ^m\zeta +F_2\zeta ,\hfill \\ \delta F_1=ı\overline{\zeta }\gamma ^m_m\psi _1,\hfill & \delta F_2=ı\overline{\zeta }\gamma ^m_m\psi _2,\hfill \end{array}$$
(75)
additional supersymmetry transformations that take the form:
$$\begin{array}{cc}\delta \varphi _1=\overline{\eta }\psi _2,\hfill & \delta \varphi _2=\overline{\eta }\psi _1,\hfill \\ \delta \psi _1=ı_\mu \varphi _2\gamma ^\mu \eta F_2\eta ,\hfill & \delta \psi _2=ı_\mu \varphi _1\gamma ^\mu \eta +F_1\eta ,\hfill \\ \delta F_1=ı\overline{\eta }\gamma ^\mu _\mu \psi _2,\hfill & \delta F_2=ı\overline{\eta }\gamma ^\mu _\mu \psi _1.\hfill \end{array}$$
(76)
From the perspective of dimensional reduction, these extra transformations arise because the underlying $`d=3+1`$ dimensional theory has four supersymmetry generators.
The auxiliary fields $`F_i`$ can be eliminated to give the Lagrangian
$$=\frac{1}{2}\left\{_m\varphi _i^m\varphi _i+ı\overline{\psi }_i\gamma ^m_m\psi _i\frac{𝒲}{\varphi _i}\frac{𝒲}{\varphi _i}\frac{^2𝒲}{\varphi _i\varphi _j}\overline{\psi }_i\psi _j\right\}.$$
(77)
The BPS equation is
$$\frac{d\varphi _i}{dy}=\alpha \frac{𝒲}{\varphi _i},$$
(78)
with $`\alpha =\pm 1`$.
Expanding the Lagrangian of Eq. (77) about the solution $`\varphi _i^w`$ up to quadratic fluctuations, we obtain
$$=\frac{1}{2}\left\{_m\varphi _i^{}^m\varphi _i^{}+ı\overline{\psi }_i\gamma ^m_m\psi _i𝒲_i𝒲_{ikl}\varphi _k^{}\varphi _l^{}𝒲_{ik}𝒲_{il}\varphi _k^{}\varphi _l^{}𝒲_{ij}\overline{\psi }_i\psi _j\right\},$$
(79)
where the subscripts on $`𝒲`$ indicate derivatives with respect to $`\varphi _i`$ evaluated at $`\varphi _i=\varphi _i^w`$.
In the domain wall background, the scalar equation of motion is
$$_m^m\varphi _n^{}+\left\{𝒲_i𝒲_{inl}+𝒲_{in}𝒲_{il}\right\}\varphi _l^{}=0,$$
(80)
and fermions satisfy
$$ı\gamma ^m_m\psi _i𝒲_{ij}\psi _j=0.$$
(81)
There is a massless scalar, $`A(x,t)`$, on the wall which is related to the broken translational symmetry,
$$\varphi _i^{}(x,y,t)=\frac{d\varphi _i^w}{dy}(y)A(x,t),$$
(82)
and a massless fermion, $`\zeta (x,t)`$, (either right or left moving, depending on the sign of $`\alpha `$), due to the broken regular supersymmetry generator,
$$\psi _i=\frac{d\varphi _i^w}{dy}\zeta (x,t),$$
(83)
where $`\zeta `$ is a (real) two component Majorana spinor satisfying the $`1+1`$ dimensional massless Dirac equation and the constraint $`\zeta =1/2(1ı\alpha \gamma ^2)\zeta `$. In addition, one of the extra SUSY generators is broken, which results in another massless fermion on the wall,
$`\psi _1(x,y,t)`$ $`=`$ $`{\displaystyle \frac{d\varphi _2^w}{dy}}\eta (x,t),`$
$`\psi _2(x,y,t)`$ $`=`$ $`{\displaystyle \frac{d\varphi _1^w}{dy}}\eta (x,t),`$ (84)
where $`\eta `$ is another (real) two component Majorana spinor satisfying the $`1+1`$ dimensional Dirac equation and the constraint $`\zeta =1/2(1+ı\alpha \gamma ^2)\zeta `$. We therefore find in the case of the BPS wall in the $`d=2+1`$, $`𝒩=2`$ case a massless scalar, a left moving massless fermion and a right moving massless fermion, which together form an $`𝒩=1`$ multiplet in the lower dimensional theory.
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